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# Theory of coherent photoassociation of a Bose-Einstein condensate
## I Introduction
Laser cooling and its spin-offs now routinely produce gaseous samples in which thermal energies, when expressed as frequencies, are smaller than the typical linewidth of an optical dipole transition. Photoassociating transitions, in which two thermal atoms combine in the presence of light to make a molecule, may therefore exhibit linewidths every bit as narrow as the transitions one encounters in nonlinear laser spectroscopy. As a result, photoassociation spectroscopy PATHEO has become the source of the most accurate molecular structure data available. Bose-Einstein condensation BEC is another recent triumph in the quest toward low temperatures in atomic physics. The connection between photoassociation and Bose-Einstein condensation has long been close, though somewhat incidental; photoassociation spectroscopy has provided key numerical data for condensation experiments SCALEN .
There have been early discussions of photoassociation of a condensate itself BUR97 ; JUL98 . However, to us, the true scope of the connection between photoassociation and condensation was only revealed by our explicit observation JAV98 that in a thermal gas it is the same phase space density that governs both the onset of Bose-Einstein condensation and the efficiency of photoassociation.
At the heart of our photoassociation work lies the quasicontinuum (QC) approach JAV98 ; MAC99 . The idea is to enclose two colliding atoms in a box, which has the effect of discretizing the dissociation continuum of the corresponding diatomic molecule. At the end of the calculations, the quantization volume is taken to infinity. Aside from resolving certain mathematical difficulties, this method turns out to have the unexpected benefit that analysis of photoassociation is reverted to analysis of few-level systems, as in quantum optics or laser spectroscopy. For instance, studies of two-color photoassociation schemes may draw from decades of experience in quantum optics and laser spectroscopy MAC99 .
However, in its initial form our QC approach does not apply to a quantum degenerate sample. We sought to rectify this shortcoming by introducing a phenomenological second-quantized Hamiltonian for photoassociation JAV99 . This idea was developed at the same time independently by Drummond et alDRU98 , and mathematically closely related approaches to the Feshbach resonance are also under active study TIM98 ; TIM99 ; ABE99 . Comparison with the QC approach gives the transition matrix elements to insert into our Hamiltonian.
We first considered a two-mode model that only takes into account one C.M. wave function for atoms and one for molecules, the modes containing the atomic and molecular condensates JAV99 . The main finding was coherent photoassociation analogous to coherent transients in few-level systems. For instance, the system may exhibit a form of Rabi flopping between atoms and molecules. Moreover, by properly sweeping the frequency of the photoassociating laser, in a process akin to rapid adiabatic passage the atomic condensate may be turned into a molecular condensate JAV99 . In another development in this direction, we have argued that two-color free-bound-bound stimulated Raman adiabatic passage, STIRAP, is feasible starting from an atomic condensate MAC00 . We have also gone beyond the two- and three-mode approximations, allowing for an arbitrary position dependence of the atomic and molecular condensates, albeit in a classical approximation similar to the one underlying the Gross-Pitaevskii equation JAV99a . It then turns out that an equilibrium with both atomic and molecular condensates present together with the photoassociating light is unstable. The sample tends to collapse spontaneously into clumps whose densities increase with time JAV99a .
The primary purposes of the present paper are to document the numerous technical and physical details of our second-quantized approach to photoassociation that could not be accommodated by the letter format of Refs. JAV99 , MAC00 , and JAV99a , and to extend our discussion in several directions that support those references. To offer a comprehensive account of the field theory version of our approach to coherent photoassociation, we have found it necessary to analyze the dipole matrix element for photoassociation in detail. This endeavor in effect constitutes an alternative derivation (c.f. Ref. MAC99 ) of our entire QC methodology. Second, we add an analysis of two-color photoassociation of a quantum-degenerate sample in a three-mode approximation that is to some extent complementary to the one in Ref. MAC00 . Much as expected, the two-color scheme provides a reprieve from spontaneous-emission losses from the primary photoassociated state. Third, we present a quantitative analysis of “rogue” photodissociation from a molecular condensate to atomic modes outside the condensate. Our suggestion JAV99 that with increasing light intensity the unwanted photodissociation may overtake coherent condensate-condensate transitions is corroborated. We find a minimum usable time scale proportional to the the inverse of the recoil frequency of laser cooling.
Probably the most prominent qualitative finding emerging from our work is the observation that it is Bose enhancement that ultimately facilitates coherent transients such as Rabi flopping, adiabatic following, and STIRAP in photoassociation of a condensate MAC00 ; MAC00a . Throughout this paper, we continue to demonstrate how coherent optical transients come about in a condensate, and argue why they should be absent in a nondegenerate gas.
There has recently been a remarkable experiment on two-color photoassociation of a condensate WYN00 . Accordingly, we include a detailed discussions on the values of experimental parameters in alkalis in general, and the parameters of Ref. WYN00 in particular. The analysis of the actual experiment demonstrates that there is still some way to go before genuinely coherent photoassociation is reached.
In Sec. II we give a walk-through of our second-quantized Hamiltonian, including a detailed discussion of the dipole moment matrix element and both the momentum and position representation of the Hamiltonian. The special case with only one spatial mode for both atoms and molecules is the subject of Sec. III. The classical version of the field theory, but including all spatial modes, is the subject of Sec. IV. The numerous complications to our one color scheme that one is liable to encounter in real experiments, as well as the experimental parameter values, are the subject of Sec. V. The brief remarks in Sec. VI conclude the paper. There are also two appendices, A on the details of the relation of the dipole matrix element between second-quantized and quasicontinuum approaches, and B on the role of atom-atom collisions in our development.
## II Field theory for atoms and molecules
The task of the present section is to develop in detail the second-quantized approach governing photoassociation of atoms into molecules in the prototype case of one laser color only. Simple heuristic arguments based on Refs. JAV98 ; MAC99 could, and did, achieve most of our aims in Ref. JAV99 where we dealt primarily with the momentum representation, but the field theory of Ref. JAV99a calls for a few additional angles. To support them, we present here a partially new ab-initio discussion of our QC method.
We take the photoassociating atoms to be in precisely one internal state, and similarly we assume that photoassociation leads to molecules with precisely one internal state. These assumptions could be relaxed, but then one has to follow the fate of the internal states as well. We do not go into this, but in essence assume that (i) the atoms are polarized and that (ii) the photoassociation resonance in itself selects a unique final state for the molecule.
### II.1 Two atoms
We begin with a pair of atoms, assumedly in the dissociation continuum of a given potential energy curve of a diatomic molecule. As the atoms interact, their relative momentum need not be a constant of the motion. Nonetheless, given a finite range for atom-atom interactions, in free space the wave functions of the relative motion $`\varphi _𝐤(𝐫)`$ could still be characterized by the asymptotic ($`r\mathrm{}`$) wave vector $`𝐤`$.
On the other hand, in the spirit of the QC method JAV98 ; MAC99 , we assume that the relative motion of the atoms is confined to a finite volume $`V`$. There are two basic questions in our two-atom analysis that must be considered. First, our phenomenological many-particle Hamiltonian (Ref. JAV99 and Eq. (12) below) is written down in terms of plane waves, yet it is more common to analyze photoassociation in terms of angular-momentum partial waves. How do we make the connection? Second, in our quasicontinuum method we resort to a finite quantization volume $`V`$, which tends to infinity only at the end of the calculations. How should we handle the finite quantization volume?
We quantize the relative motion of the two atoms in a spherical box of radius $`R`$ and volume $`V=\frac{4}{3}\pi R^3`$ using reflecting boundary conditions. While angular momentum is still a constant of the motion, the usual eigenstates MAC99 of the relative motion cannot be characterized by a momentum vector, even asymptotically. Nonetheless, it is evidently possible to construct orthonormal superpositions $`\overline{\varphi }_𝐤(𝐫)`$ of the eigenstates of the spherical box that in the limit of a large box turn into the states $`\varphi _𝐤(𝐫)`$, i.e., states that behave like plane waves at large distances.
Let us first take the inner product of a true plane wave, normalized to the volume $`V`$, with a spherically symmetric (real) test function $`f(r)`$ with a finite range $`R`$. In the limit $`𝐤0`$ we have
$$I_1=\underset{𝐤0}{lim}\frac{1}{\sqrt{V}}_0^Rd^3re^{i𝐤𝐫}f(r)\frac{1}{\sqrt{V}}_0^{\mathrm{}}4\pi r^2𝑑rf(r).$$
(1)
Next we take the same inner product for the $`l=0`$ partial wave of the plane wave $`e^{i𝐤𝐫}`$, basically the spherical Bessel function $`j_0(kr)\mathrm{sin}kr/r`$. We normalize this partial wave also in the spherical box, or equivalently, in the radial coordinate $`r`$ with respect to the measure $`4\pi r^2dr`$. In the limit of small $`k`$ we have the integral for the inner product
$`I_2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi R}}}{\displaystyle _0^R}4\pi r^2𝑑r{\displaystyle \frac{\mathrm{sin}kr}{r}}f(r)`$ (2)
$``$ $`{\displaystyle \frac{k}{\sqrt{2\pi R}}}{\displaystyle _0^{\mathrm{}}}4\pi r^2𝑑rf(r).`$
Obviously, the ratio of the two,
$$\alpha (k)=\frac{I_1}{I_2}=\frac{1}{k}\sqrt{\frac{2\pi R}{V}},$$
(3)
is the expansion coefficient of the $`s`$-wave in the plane wave, given that both the plane wave and the $`s`$-wave are normalized to the volume $`V`$. Actually, the $`k`$ values of the spherical eigenmodes are quantized, and the smallest value is $`k=\pi /R`$. We therefore have to be careful with the limit $`k0`$ when applying Eq. (3).
When atom-atom interactions are taken into account, the radial eigenstates for a given angular momentum are not just spherical Bessel functions, but also reflect the scattering phase shifts. Henceforth we assume that only $`s`$-wave scattering needs to be considered, as is usually the case for bosons at sufficiently low temperatures. Obviously, even though the partial waves have phase shifts, by tweaking the phases of the partial waves it is still possible to make an asymptotic near-plane wave $`\overline{\varphi }_k(𝐫)`$ out of the eigenstates of the relative motion of the atoms in the sphere. The adjustment of phases has no effect on the weight of the $`s`$-wave component in $`\overline{\varphi }_k(𝐫)`$, which may still be inferred from Eq. (3).
Let us denote the wave function of the particular molecular state we are aiming for by $`\overline{\psi }(𝐫)`$, and make the standard (albeit crude) approximation that the relevant electronic dipole matrix element $`𝐝`$ is a constant independent of the relative coordinate $`𝐫`$ of the atoms comprising the molecule. Then the QC dipole matrix element for photoassociation and photodissociation characterizing a process in which the relative momentum of the colliding atoms is $`\mathrm{}𝐤`$ reads simply
$$𝐝(𝐤)=𝐝d^3r\overline{\psi }^{}(𝐫)\overline{\varphi }_𝐤(𝐫).$$
(4)
We only consider $`s`$-wave collisions, and correspondingly set $`\overline{\psi }(𝐫)=\overline{\psi }(r)`$ as appropriate for a nonrotating $`J=0`$ molecule. Of course, only the $`l=0`$ component of the wave function $`\overline{\varphi }_𝐤(𝐫)`$ counts. Taking into account the weight from Eq. (3), we have
$$𝐝(𝐤)=𝐝\frac{1}{k}\sqrt{\frac{2\pi R}{V}}_0^R4\pi r^2𝑑r\overline{\psi }^{}(r)\overline{\varphi }_k^0(r).$$
(5)
The notation $`\overline{\varphi }_k^0(r)`$ stands for the radial $`l=0`$ wave function of the relative motion of the two atoms corresponding to the wave number $`k`$, and normalized to volume $`V`$ as usual.
A convenient qualitative model is provided by the limiting form
$$\overline{\varphi }_k^0(r)=\frac{1}{\sqrt{2\pi R}}\frac{\mathrm{sin}k(ra)}{r},$$
(6)
where $`a`$ is the $`s`$-wave scattering length. Combination of (5) and (6) gives
$$𝐝(𝐤)=𝐝\frac{4\pi }{k\sqrt{V}}_0^Rr𝑑r\mathrm{sin}k(ra)\overline{\psi }^{}(r).$$
(7)
The form (6) is only valid outside the range of the atom-atom interaction potential, and for small enough $`k`$ so that the scattering phase shift may be written as $`ka`$. Equation (7) is therefore quantitatively reliable only if the vibrational bound-state wave function of the molecule $`\overline{\psi }(r)`$ happens to reside outside the range of the atom-atom interactions of the photoassociating atoms. In fact, this is the case in some of the current photoassociation experiments COT95 ; COT95a .
The matrix element (7) is explicitly proportional to $`V^{1/2}`$. In fact, the $`V^{1/2}`$ scaling of $`𝐝(𝐤)`$ is a generic property of our QC approach, and does not depend on the specific assumptions of Eq. (7). First, the almost-plane wave $`\overline{\varphi }_𝐤(𝐫)`$ is normalized to the volume $`V`$. No matter what kind of (square integrable) structure occurs around $`𝐫0`$, the plane-wave form at large distances makes the normalization constant proportional to $`V^{1/2}`$. Second, because of the finite range of the bounded molecular wave function $`\overline{\psi }(r)`$, its normalization coefficient does not depend on $`V`$. Third, thanks to the same finite range, the integral (4) effectively extends only over a finite, fixed volume. The net result is that the matrix element indeed is proportional to the normalization constant of $`\overline{\varphi }_𝐤(𝐫)`$, $`V^{1/2}`$.
Another important property of the matrix element (7), that it tends to a constant as $`k0`$, is also a generic feature of the physics. Namely, for small enough $`k`$, the form of the $`s`$-wave $`\overline{\varphi }_k^0(r)`$ indeed is $`\mathrm{sin}k(ra)/r`$, except within the range of atom-atom interactions. But within this range, the shape of $`\overline{\varphi }_k^0(r)`$ is independent of $`k`$ at small enough $`k`$. In order that the inner and the outer forms join smoothly, the amplitude of the inner wave function must be $`k`$ for small $`k`$, just as is the amplitude of the outer wave function. So, for small enough $`k`$, the wave function $`\overline{\varphi }_k^0(r)`$ is $`k`$ for all of those $`r`$ for which the bound-state wave function $`\overline{\psi }(r)`$ is effectively nonzero. The matrix elements (4), (5) and (7) therefore all tend to a constant as $`k0`$.
In order to prepare for the comparison between photoassociation and photodissociation, let us assume that a laser field with the amplitude
$$𝐄(t)=\frac{1}{2}𝐄e^{i\omega t}+\mathrm{c}.\mathrm{c}.$$
(8)
is incident on a molecule. We denote the detuning of the light above the photodissociation threshold by $`\delta `$. Given the reduced mass of the colliding atoms $`\mu `$, the corresponding resonant wave number $`k_0`$ and velocity $`v_0`$ are such that $`\mathrm{}k_0^2/2\mu =\delta `$ and $`v_0=\mathrm{}k_0/\mu `$. Using the standard dipole and rotating-wave approximations as well as Fermi’s golden rule, we have the photodissociation rate
$`\mathrm{\Gamma }_0`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}^2}}{\displaystyle 𝑑kD(k)\left|\frac{𝐝(𝐤)𝐄}{2}\right|^2\delta \left(\frac{\mathrm{}k^2}{2\mu }\delta \right)}`$ (9)
$`=`$ $`{\displaystyle \frac{\mu ^2v_0V}{\pi \mathrm{}^2}}\left|{\displaystyle \frac{𝐝(k_0)𝐄}{2\mathrm{}}}\right|^2,`$
where
$$D(k)=\frac{k^2V}{2\pi ^2}$$
(10)
is the density of $`k`$ states, whether in a cubic or in a spherical box MAC99 ; VOLREF . Equation (9) is fully compatible with the development in Refs. JAV98 and MAC99 , as it should.
Since the dipole matrix element tends to a constant in the limit $`k0`$, it is easy to see from (9) that the Wigner threshold law holds; namely, that $`lim_{v_00}\mathrm{\Gamma }_0/v_0`$ is finite, and nonzero except for an unlikely accident. As a matter of fact, our argument about the $`k0`$ limit of the dipole matrix element was nothing but a recital of a standard argument for the Wigner threshold law. Nonetheless, it will furnish a relevant piece of the puzzle when we are to discuss atom-molecule field theory below. To this end, we note from Eqs. (5) and (9) the equality
$`\underset{v_00}{lim}{\displaystyle \frac{\mathrm{\Gamma }_0}{v_0}}`$
$`=`$ $`{\displaystyle \frac{\mu ^2}{\pi \mathrm{}^2}}\left|{\displaystyle \frac{𝐝𝐄}{\mathrm{}}}\underset{k_00}{lim}{\displaystyle \frac{\sqrt{2\pi R}}{k_0}}{\displaystyle 4\pi r^2𝑑r\overline{\psi }^{}(r)\overline{\varphi }_{k_0}^0(r)}\right|^2.`$
### II.2 Many atoms in momentum representation
#### II.2.1 Basic Hamiltonian
In Ref. JAV99 we introduced a phenomenological second-quantized Hamiltonian for photoassociation of bosonic atoms to (obviously) bosonic molecules,
$`{\displaystyle \frac{H}{\mathrm{}}}`$ $`=`$ $`{\displaystyle \underset{𝐤}{}}\left[{\displaystyle \frac{\mathrm{}𝐤^2}{4m}}b_𝐤^{}b_𝐤+\left({\displaystyle \frac{\delta _0}{2}}+{\displaystyle \frac{\mathrm{}𝐤^2}{2m}}\right)a_𝐤^{}a_𝐤\right]`$ (12)
$`{\displaystyle \underset{\mathrm{𝐤𝐤}^{}𝐪}{}}[{\displaystyle \frac{𝐝_{\mathrm{𝐤𝐤}^{}}𝐄_𝐪}{4\mathrm{}}}b_{𝐤+k^{}+𝐪}^{}a_𝐤a_𝐤^{}+\mathrm{H}.\mathrm{c}.].`$
Here $`m`$ stands for the mass of the atom, $`m=2\mu `$. The operators $`a_𝐤`$ and $`b_𝐤`$ are boson annihilation operators for atoms and molecules in the plane wave mode $`𝐤`$. The Hamiltonian trivially includes the kinetic energy of the atoms and molecules. Here photoassociation takes place with an absorption (as opposed to induced emission) of a photon. By momentum conservation, atoms with wave vectors $`𝐤`$ and $`𝐤^{}`$ plus a photon with wave vector $`𝐪`$ must then make a molecule with wave vector $`𝐤+𝐤^{}+𝐪`$. This explains the form of the cubic operator product; annihilate the atoms and a photon, create the molecule.
As is usual in quantum optics, we use a classical field to represent the photons. Specifically, the positive-frequency part of the electric field reads
$$𝐄^+(𝐫)=\frac{1}{2}\underset{𝐪}{}𝐄_𝐪e^{i𝐪𝐫}.$$
(13)
Both $`𝐄^+(𝐫)`$ and the coefficients $`𝐄_𝐪`$ may in principle be slowly varying functions of time, but the leading time dependence of the electric field $`e^{i\omega t}`$ has been absorbed into the detuning $`\delta _0`$ in a transformation to a rotating frame. The seemingly unexpected factor $`\frac{1}{2}`$ in the detuning term correspond to the fact that upon photodissociation one molecule produces two atoms, both of which generally take away kinetic energy with them. The sign of the detuning is chosen in such a way that $`\delta _0>0`$ corresponds to tuning of the laser by the energy $`\mathrm{}\delta _0`$ above the photodissociation threshold. A quick way to verify this is to consider the potential resonance when the (quasi) energy (in the rotating frame) for a system of one molecule and zero atoms would be the same as the energy for a system with zero molecules and two atoms, much like in the Appendix A. The total energy of the atoms must equal the energy of the molecule plus $`\mathrm{}\delta _0`$, which is compatible with Eq. (12).
We assume that the optical transition responsible for photoassociation is a dipole transition. In the dipole approximation, the electronic photodissociation and photoassociation transitions in a molecule do not (cannot!) depend on the propagation direction of light, hence the dipole matrix elements $`𝐝_{\mathrm{𝐤𝐤}^{}}`$ do not depend on photon momenta. Besides, by translational invariance, the dipole matrix elements must be functions of the difference $`𝐤𝐤^{}`$ only. Because of the exchange symmetry of the boson operators, the matrix element $`𝐝_{\mathrm{𝐤𝐤}^{}}`$ may be chosen to be symmetric in the exchange of the momentum indices, hence, an even function of $`𝐤𝐤^{}`$. And finally, since we consider $`s`$-wave photoassociation only, the matrix element is a function of $`|𝐤𝐤^{}|`$.
To pin down the matrix elements $`𝐝_{\mathrm{𝐤𝐤}^{}}`$, in Ref. JAV99 we took the limit of a dilute thermal gas. The known QC results are recovered if one expresses the thus far undefined coefficients $`𝐝_{\mathrm{𝐤𝐤}^{}}`$ in terms of the matrix elements (4) as
$$𝐝_{\mathrm{𝐤𝐤}^{}}\sqrt{2}𝐝\left(\frac{1}{2}(𝐤𝐤^{})\right).$$
(14)
The $`\sqrt{2}`$ is a consequence of the Bose-Einstein statistics, and the factor $`\frac{1}{2}`$ follows from the way that the relative momentum must be defined to make it the conjugate of the conventional relative position. We have recently noticed that there is a subtlety associated with this identification having to do with the the statistics of the atoms. We elaborate in Appendix A, but meanwhile continue according to Eq. (14).
In sum, if one may treat the atoms and the molecules as bosons in their own right, then both the form of, and even the numerical coefficients in, the photoassociation Hamiltonian are unambiguously determined by simple physical considerations. Admittedly, we do not know of any conclusive ab initio argument for the bosonic nature of atoms and molecules in photoassociation. But neither do we know of any for alkali atoms, which nonetheless seem to behave like good bosons in current BEC experiments.
#### II.2.2 Two-mode approximation
We now revisit the situation that was the focus of Ref. JAV99 : an infinite, homogeneous condensate and one plane wave of light with photon momentum $`\mathrm{}𝐪`$. By conservation of momentum, the molecules one may create by photoassociation all have the wave vector equal to $`𝐪`$. On the other hand, as far as momentum conservation is concerned, a molecule with the wave vector $`𝐪`$ may photodissociate into two atoms, neither of which has zero momentum. We will discuss such “rogue photodissociation” in more detail below, Sec. V.1, and so far ignore it.
All told, we only retain the two modes of the model with the annihilation operators $`aa_0`$ and $`bb_𝐪`$. The Hamiltonian reads
$`{\displaystyle \frac{H}{\mathrm{}}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}𝐪^2}{4m}}b^{}b\frac{1}{2}\delta _0a^{}a\frac{1}{2}\kappa (b^{}aa+ba^{}a^{})`$ (15)
$``$ $`{\displaystyle \frac{\mathrm{}𝐪^2}{4m}}b^{}b\frac{1}{2}\delta _0a^{}a\frac{1}{2}\kappa (b^{}aa+ba^{}a^{})`$
$`{\displaystyle \frac{\mathrm{}𝐪^2}{8m}}(2b^{}b+a^{}a)`$
$``$ $`\frac{1}{2}\delta a^{}a\frac{1}{2}\kappa (b^{}aa+ba^{}a^{}).`$
In the “approximate” equality we have added a constant of the motion to the Hamiltonian, a step that has no effect on the ensuing dynamics. The QC Rabi frequency reads
$`\kappa `$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{2\mathrm{}}}𝐄[\underset{𝐤0}{lim}𝐝(𝐤)]`$ (16)
$`=`$ $`\underset{v_00}{lim}\sqrt{{\displaystyle \frac{2\pi \mathrm{}^2\mathrm{\Gamma }_0}{\mu ^2Vv_0}}},`$
where we have used Eq. (9). Without any loss of generality, we have chosen $`\kappa `$ to be real and nonnegative. Seemingly alarmingly, the Rabi frequency still depends on the quantization volume. We will return to this point in Sec. III. At present, our aim is just to set up the second-quantized Hamiltonian in the two-mode approximation. The task is completed by noting that
$$\delta =\delta _0+\frac{\mathrm{}𝐪^2}{4m}$$
(17)
is the detuning corrected for the photon recoil energy of the molecule. From now on we will keep track of this distinction, so that $`\delta `$ always includes the appropriate recoil. In fact, that was already implicitly the case in Eq. (9).
### II.3 Many-atom field theory
Our phenomenological Hamiltonian (12) was written down originally in momentum representation. Nevertheless, as in Refs. JAV99 and JAV99a , it is fairly straightforward to convert it into position representation, i.e., into a quantum field theory.
Let us introduce atomic and molecular fields in a quantization volume $`V`$. A priori, this volume does not have to be the cavity of radius $`R`$, as in Sec. II.1. Often it is actually more convenient to use a cubic box with periodic boundary conditions. On the other hand, when one deals with two atomic fields, it is often expedient to write the integrals in the theory in terms of center-of-mass and relative coordinates. If the integral over the relative coordinates cuts off because the integrand tends to zero at large distances, it is immaterial what volume is used, as long as it is large enough. Thus, when convenient, in such relative-coordinate integrals we may still imagine the spherical potential well. The short of the story is that the quantization volume $`V`$ refers to the geometry that is expedient in the particular context.
We write the atomic and molecular fields in terms of plane wave states as
$$\varphi (𝐫)=\frac{1}{\sqrt{V}}\underset{𝐤}{}e^{i𝐤𝐫}a_k,\psi (𝐫)=\frac{1}{\sqrt{V}}\underset{𝐤}{}e^{i𝐤𝐫}b_k.$$
(18)
The Hamiltonian (12) may be cast as a Hamiltonian density of a field theory for these fields. Using the standard continuum limit
$$\underset{𝐤}{}f(𝐤)\frac{V}{(2\pi )^3}d^3kf(𝐤),$$
(19)
Eq. (13), and the properties of Fourier integrals, the part of the Hamiltonian density depending on the dipole interaction becomes
$$\frac{(𝐫)}{\mathrm{}}=\frac{𝐄^+(𝐫)}{2\mathrm{}}\psi ^{}(𝐫)d^3r^{}\varphi (𝐫+\frac{1}{2}𝐫^{})𝐝(𝐫^{})\varphi (𝐫\frac{1}{2}𝐫^{})+\mathrm{}.$$
(20)
The dipole kernel of photoassociation $`𝐝(𝐫)`$ is given in terms of the dipole matrix element of Eq. (4) as
$$𝐝(𝐫)=\frac{\sqrt{2V}}{(2\pi )^3}d^3ke^{i𝐤𝐫}𝐝(𝐤).$$
(21)
It should be noted that our argument contains a subtle trick, which is exposed in Appendix B.
To make further progress, we study the kernel using the qualitative model for the dipole matrix element (7). Some more juggling with Fourier transforms gives
$`𝐝(𝐫)`$
$`=`$ $`\sqrt{2}𝐝{\displaystyle \frac{\theta (a+r)(a+r)\overline{\psi }^{}(a+r)\theta (ar)(ar)\overline{\psi }^{}(ar)}{r}},`$
where $`\theta `$ is the Heaviside unit step function. This result is only as good as the assumptions of Eq. (7). In particular, as Eq. (7) is not valid for large $`k`$, Eq. (LABEL:DR1) is dubious at short distances. Also, as we already pointed out with Eq. (7), even for small $`k`$ the model is quantitatively accurate only if the vibrational wave function $`\overline{\psi }(r)`$ resides outside the range of the atom-atom interactions. In spite of these caveats, Eq. (LABEL:DR1) gives an exceedingly plausible idea of the range $`\mathrm{\Delta }r`$ of the dipole interaction kernel; $`\mathrm{\Delta }r`$ is of the order of the larger of the spatial extent of the vibrational state $`\overline{\psi }(r)`$ and the absolute value of the scattering length $`a`$.
The relevant length scale for the atomic field is larger than $`\mathrm{\Delta }r`$ if the energies of all of the atoms relevant for the time evolution are smaller than $`\mathrm{}^2/m(\mathrm{\Delta }r)^2`$. Henceforth we assume this to be the case. Then we may simplify the field-theory form of the integral as follows,
$`{\displaystyle \frac{(𝐫)}{\mathrm{}}}`$ $`=`$ $`{\displaystyle \frac{𝐄^+(𝐫)}{2\mathrm{}}}\psi ^{}(𝐫){\displaystyle d^3r^{}\varphi (𝐫+\frac{1}{2}𝐫^{})𝐝(𝐫^{})\varphi (𝐫\frac{1}{2}𝐫^{})}+\mathrm{}{\displaystyle \frac{𝐄^+(𝐫)}{2\mathrm{}}}\psi ^{}(𝐫)\varphi (𝐫)\varphi (𝐫){\displaystyle d^3r^{}𝐝(𝐫^{})}+\mathrm{}`$ (23)
$`=`$ $`\left[{\displaystyle \frac{𝐄^+(𝐫)}{2\mathrm{}}}\sqrt{2V}\underset{𝐤0}{lim}𝐝(𝐤)\right]\psi ^{}(𝐫)\varphi (𝐫)\varphi (𝐫)+\mathrm{}`$
$`=`$ $`\left[\sqrt{2}{\displaystyle \frac{𝐝𝐄^+(𝐫)}{2\mathrm{}}}\underset{k0}{lim}{\displaystyle \frac{\sqrt{2\pi R}}{k}}{\displaystyle 4\pi r^2𝑑r\overline{\psi }^{}(r)\overline{\varphi }_k^0(r)}\right]\psi ^{}(𝐫)\varphi (𝐫)\varphi (𝐫)+\mathrm{}.`$
The second equality follows from the assumedly slow dependence of the atomic field $`\varphi (𝐫)`$ on position $`𝐫`$ compared to $`𝐝(𝐫)`$, the third from Eq. (21) and the properties of Fourier transformations, and in the last equality we have substituted Eq. (5). Comparison with Eq. (LABEL:GC) then gives the contact interaction form of the photoassociation dipole interaction,
$$\frac{(𝐫)}{\mathrm{}}=𝒟(𝐫)\psi ^{}(𝐫)\varphi (𝐫)\varphi (𝐫)+\mathrm{},$$
(24)
with
$$𝒟(𝐫)=e^{i\mathrm{arg}[𝐝𝐄^+(𝐫)]}\underset{v_00}{lim}\sqrt{\frac{2\pi \mathrm{}^2\mathrm{\Gamma }_0(𝐫)}{v_0\mu ^2}}.$$
(25)
Here $`\mathrm{\Gamma }_0(𝐫)`$ is the photodissociation rate as per the prevailing light intensity at position $`𝐫`$, and the exponential records the phase of the dipole interaction term.
For reference, we summarize the contact-interaction version of the Hamiltonian for photoassociation in terms of atomic and molecular fields. While at that, by adding a suitable multiple of the conserved particle number completely analogously to the operation that we did in the chain of equations (15), we move the detuning term to the molecules. The result is
$`H`$ $`=`$ $`{\displaystyle d^3r(𝐫)};`$
$`{\displaystyle \frac{(𝐫)}{\mathrm{}}}`$ $`=`$ $`\varphi ^{}(𝐫)\left({\displaystyle \frac{\mathrm{}^2}{2m}}\right)\varphi (𝐫)+\psi ^{}(𝐫)({\displaystyle \frac{\mathrm{}^2}{4m}}+\delta _0)\psi (𝐫)`$ (26)
$`[𝒟(𝐫)\psi ^{}(𝐫)\varphi (𝐫)\varphi (𝐫)+\mathrm{H}.\mathrm{c}.].`$
## III Coherence in two-mode model
We consider the two-mode model, whose Hamiltonian we reproduce here from Eq. (15),
$$\frac{H}{\mathrm{}}=\frac{1}{2}\delta a^{}a\frac{1}{2}\kappa (b^{}aa+ba^{}a^{}).$$
This is the age-old Hamiltonian for second-harmonic generation, with atoms replacing the fundamental-frequency light and molecules the second harmonic. Needless to say, the literature on the model is extensive. We will cite a few relevant examples below.
Let us take the conserved particle number $`a^{}a+2b^{}b`$ to have the value $`N`$. The Hamiltonian (15) may then be restricted to the space spanned by the states $`|n|n_M|N2n_A`$ with $`n=0,\mathrm{},N/2`$ molecules and $`N2n`$ atoms, and is tridiagonal in that basis. It is easy to find both the eigenstates of the Hamiltonian and the time evolution of any specified initial state numerically. We have done so using inverse iteration, and a variation of the Crank-Nicholson method. While the problem considered in Ref. JAV99b is different, the numerical methods described therein work equally well here.
In this way we have firstly JAV99 ; WAL72 found the fraction of atoms converted into molecules, given that the system starts out at time $`t=0`$ with everything in atoms and is driven by a resonant laser, $`\delta =0`$. We found a nonlinear analog of Rabi flopping, the system oscillating between atoms and molecules. Given that we have a quantum system where the spacing between the successive energy eigenstates is not constant, it is not a surprise that the oscillations collapse, and even revive DRO92 .
Interesting as these features are as a matter of principle, the drawback remains that Rabi oscillations of occupation probabilities are generally not robust, not even in quantum optics or laser spectroscopy. We foresee little experimental utility for nonlinear Rabi oscillations. Our main message of the analysis of Rabi oscillations in Ref. JAV99 rather is that, for $`N1`$, the characteristic frequency scale of the system is not the QC Rabi frequency $`\kappa `$ but
$$\mathrm{\Omega }=\sqrt{N}\kappa =\underset{v_00}{lim}\sqrt{\frac{2N\pi \mathrm{}^2\mathrm{\Gamma }_0}{\mu ^2Vv_0}}=\underset{v_00}{lim}\sqrt{\frac{2\pi \mathrm{}^2\mathrm{\Gamma }_0\rho }{\mu ^2v_0}}.$$
(27)
The $`\sqrt{N}`$ is evidently a Bose enhancement factor. The implicit quantization volume has disappeared from the result. Instead we have the density of atoms if all molecules were dissociated,
$$\rho =\frac{N}{V},$$
(28)
which is a true physical parameter for the system.
Bose enhancement highlights the role of statistics in our analysis. Suppose we were to deal with a nondegenerate gas. Then the occupation probability of each QC state is much less than unity by definition, and there is no Bose enhancement. The frequency of Rabi flopping between a given QC state and a bound molecular states would just be $`\kappa V^{1/2}`$, and vanishes in the limit of large volume. It is not possible to have Rabi flopping in photoassociation of an infinite and homogeneous nondegenerate gas, even in principle MAC00a .
The photoassociation rate for a pair of atoms starting in a given QC state is $`\kappa ^21/V`$, and also vanishes in the limit of a large volume. The saving grace JAV98 ; MAC99 is that, to keep the density constant, the atom number $`N`$ and at the same time the number of candidate atoms to photoassociate with any given atom tends to infinity. Adding the probabilities for photoassociation due to all colliders, the total rate of photoassociation for any given atom is $`N\kappa ^2\rho `$. This remains a constant in the continuum limit, which leads to a finite photoassociation rate $`\rho `$.
In contrast, in a BEC the atom number already ends up in the transition matrix element, giving the effective Rabi frequency $`\mathrm{\Omega }=\sqrt{N}\kappa \sqrt{\rho }`$. One does not add rates $`\kappa ^2`$ due to different atoms, but all the condensate atoms act as a single state that has a finite matrix element for photoassociation even in the limit of infinite quantization volume. This ultimately facilitates coherence effects like Rabi flopping.
The Hamiltonian (15) acts in a rotating frame, and there is no manifest physical significance to its eigenvalues. Nonetheless, let us call them energies, and the state with the lowest energy the ground state. Qualitatively, when the laser is tuned far below the photodissociation threshold, $`\delta \mathrm{\Omega }`$, the $`a^{}a`$ term in the Hamiltonian becomes costly in energy. The ground state of the Hamiltonian should then tend to have the system mostly in molecules. Conversely, for $`\delta \mathrm{\Omega }`$, the ground state favors atoms. We confirm these surmises explicitly in Fig. 1 by plotting the fraction of atoms converted into molecules, the expectation value
$$f=\frac{2b^{}b}{N},$$
(29)
for the numerically obtained ground state. The difference between the dashed curves is the atom number, $`N=10`$ and $`N=100`$.
The idea is near that, when the laser is swept from a large above-threshold detuning to a large (in absolute value) below-threshold value, the system will follow adiabatically, and an atomic condensate is converted to a molecular condensates JAV99 . Moreover, if $`\mathrm{\Omega }`$ were the relevant frequency scale, adiabaticity means that the detuning must change by $`\mathrm{\Omega }`$ in a time of the order or longer than $`\mathrm{\Omega }^1`$. In our model we assume that the detuning is swept as
$$\delta (t)=\xi \mathrm{\Omega }^2t.$$
(30)
We expect adiabatic following when $`\xi \stackrel{<}{}1`$.
This piece of intuition is correct. In Fig. 2 we reproduce the relevant figure from Ref. JAV99 . We fix $`N=100`$. At an initial time giving $`\delta =20\mathrm{\Omega }`$, we start the system in its ground state, almost purely atoms, and integrate the time dependent Schrödinger equation numerically while sweeping the detuning at two rates $`\xi =1`$ and $`\xi =0.1`$. The plot gives the fraction of atoms converted to molecules as a function of the instantaneous detuning. For $`\xi =1`$ we expect to be somewhere at the borderline of adiabaticity, and actually find a conversion efficiency of 0.8. When the laser tuning is swept ten times more slowly, $`\xi =0.1`$, the conversion efficiency has reached 0.97.
We believe that we have discovered a feasible method for preparing a molecular condensate: start with an atomic condensate and a photoassociating laser, then sweep the frequency of the laser in a proper way, and obtain a molecular BEC. In the general manner of adiabatic methods, this “rapid adiabatic passage” should be robust.
Besides its potential utility, rapid adiabatic passage is another example of coherent phenomena that occur in a condensate but not in a nondegenerate gas. Once more, in a nondegenerate gas the frequency scale for adiabaticity is $`\kappa `$, and the time scale is proportional to $`\kappa ^1V^{1/2}`$. Even if nothing else went wrong, in the limit of a large sample the time scale for adiabatic atom-molecule conversion would tend to infinity. Once more, Bose enhancement saves the day by turning the volume dependence of the relevant frequency scale into a density dependence MAC00a .
## IV Field theory for all modes
If only one spatial mode has to be considered for both the atomic and the molecular condensate, the two-mode Hamiltonian (15) is all there is to it. Nonetheless, even if somehow the infinite homogeneous condensate could be approximated in practice, for whatever reason there would always be atoms and molecules present with momenta that are not included in the two-mode picture. As the full Hamiltonian (12) mixes momenta nonlinearly, the possibility of instability arises and should be investigated.
We will study deviations from the two-mode picture, or the possibility that the (infinite) atomic and molecular condensates are not spatially homogeneous. Here we use the quantum field version of our photoassociation theory. But solving nonlinear quantum field theories is generally not a simple matter, and one must approximate. We will resort to classical field theory, the same way one proceeds when an alkali condensate is analyzed using the Gross-Pitaevskii equation.
### IV.1 Gross-Pitaevskii equations
To derive the Gross-Pitaevskii equations (GPE) for this problem, in the usual way we first amend the Hamiltonian by adding to it a multiple of the conserved particle number
$$N=\underset{𝐤}{}(a_𝐤^{}a_𝐤+2b_𝐤^{}b_𝐤)=d^3r[\varphi ^{}(𝐫)\varphi (𝐫)+2\psi ^{}(𝐫)\psi (𝐫)].$$
(31)
Instead of the Hamiltonian density, we then use the “Kamiltonian” density
$$𝒦(𝐫)=(𝐫)\mu [\varphi ^{}(𝐫)\varphi (𝐫)+2\psi ^{}(𝐫)\psi (𝐫)]$$
(32)
in our calculations. This substitution has no effect on the dynamics. At this point the real scalar $`\mu `$ is arbitrary. Though this piece of knowledge has no bearing on our analysis, in the end the constant $`\mu `$ will be analogous to chemical potential in thermodynamics. Hence, $`\mu `$ is called the chemical potential.
Using the standard commutators for boson fields, the Heisenberg equation of motion for the molecular field becomes
$`i\dot{\psi }(𝐫)`$ $`=`$ $`[\psi (𝐫),{\displaystyle d^3r^{}\frac{𝒦(𝐫^{})}{\mathrm{}}}]`$
$`=`$ $`\left[{\displaystyle \frac{\mathrm{}^2}{4m}}+\delta _02{\displaystyle \frac{\mu }{\mathrm{}}}\right]\psi (𝐫)\frac{1}{2}𝒟(𝐫)\varphi (𝐫)\varphi (𝐫).`$
The transformation to the classical field theory is effected by positing that the fields in the equations of motion are no longer boson fields, but commuting $`c`$-number fields. The interpretation is that $`\varphi (𝐫)`$ and $`\psi (𝐫)`$ are the macroscopic wave functions for atomic and molecular condensates.
From this point onward we again take the driving field to be a simple plane wave with the positive frequency part $`\frac{1}{2}𝐄e^{i𝐪𝐫}`$. Second, we scale the atomic and molecular fields with the square root of density, $`\sqrt{N/V}`$. Third, we incorporate the spatial variation of the electric field into the definition of the molecular field. Altogether, we define new atomic and molecular fields $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ as
$$\varphi =\sqrt{\frac{N}{V}}\mathrm{\Phi },\psi =\sqrt{\frac{N}{V}}e^{i𝐪𝐫}\mathrm{\Psi }.$$
(34)
The net results are four. First, the normalization of the fields now reads
$$\frac{1}{V}d^3r(|\mathrm{\Phi }|^2+2|\mathrm{\Psi }|^2)=1.$$
(35)
Second, the coefficient $`𝒟`$ gets multiplied by the square root of the atom density that would prevail if all molecules were to dissociate. Third, all explicit position dependence disappears from the equations of motion. The GPE for the rescaled atomic and molecular wave functions are
$`i\dot{\mathrm{\Phi }}(𝐫)`$ $`=`$ $`\left[{\displaystyle \frac{\mathrm{}^2}{2m}}\mu \right]\mathrm{\Phi }(𝐫)e^{i\phi }\mathrm{\Omega }\mathrm{\Phi }^{}(𝐫)\mathrm{\Psi }(𝐫),`$ (36)
$`i\dot{\mathrm{\Psi }}(𝐫)`$ $`=`$ $`\left[{\displaystyle \frac{\mathrm{}}{4m}}^2+i{\displaystyle \frac{\mathrm{}}{2m}}𝐪+\delta 2{\displaystyle \frac{\mu }{\mathrm{}}}\right]\mathrm{\Psi }(𝐫)`$ (37)
$`\frac{1}{2}e^{i\phi }\mathrm{\Omega }[\mathrm{\Phi }(𝐫)]^2.`$
None other than our characteristic frequency scale for photoassociation, $`\mathrm{\Omega }`$ of Eq. (27), is now explicitly the frequency scale in the field equations as well. The phase factor $`e^{i\phi }`$ accounts for the phases of the electric field and dipole moment, and will soon prove inconsequential. Fourth, a photon recoil term got added to the kinetic energy of the molecules.
While still using full dimensional units, we pause to discuss the mathematical symmetries of the GPE. First, a trivial phase change of one of the fields, e.g., $`\mathrm{\Phi }e^{i\phi /2}\mathrm{\Phi }`$, converts Eqs. (IV.1) to the same equations, except that the phase factors vanish from atom-molecule interaction terms. Therefore, we drop the phase in the interaction term. Second, the GPE have a global gauge invariance. If the fields $`\mathrm{\Phi }`$, $`\mathrm{\Psi }`$ are a solution, then so are the fields $`\mathrm{\Phi }e^{i\phi }`$, $`\mathrm{\Psi }e^{2i\phi }`$ for an arbitrary fixed phase $`\phi `$. In particular, putting $`\phi =2\pi `$, one may see that the equations are invariant under the change of the sign of the field $`\mathrm{\Phi }`$. Third, as a time dependent generalization of the gauge invariance, the GPE is invariant under the replacement
$`\mathrm{\Phi }`$ $``$ $`\mathrm{\Phi }e^{i\mathrm{\Delta }\mu t},`$ (38)
$`\mathrm{\Psi }`$ $``$ $`\mathrm{\Psi }e^{2i\mathrm{\Delta }\mu t},`$ (39)
$`\mu `$ $``$ $`\mu +\mathrm{\Delta }\mu .`$ (40)
Fourth, the GPE is Galilei invariant, in that the replacements
$`\mathrm{\Phi }(𝐫,t)`$ $``$ $`e^{i𝐤𝐫}\mathrm{\Phi }(𝐫{\displaystyle \frac{\mathrm{}𝐤}{m}}t,t),`$ (41)
$`\mathrm{\Psi }(𝐫,t)`$ $``$ $`e^{2i𝐤𝐫}\mathrm{\Psi }(𝐫{\displaystyle \frac{\mathrm{}𝐤}{m}}t,t),`$ (42)
$`\mu `$ $``$ $`\mu +{\displaystyle \frac{\mathrm{}^2𝐤^2}{2m}},`$ (43)
$`\delta `$ $``$ $`\delta {\displaystyle \frac{\mathrm{}𝐤}{m}}𝐪`$ (44)
convert a solution into another solution that corresponds to an added momentum $`\mathrm{}𝐤`$ per atom. The nontrivial part of the transformation, Eq. (44) is that, due to the Doppler shift, the effective detuning changes depending on the overall motion of the atom-molecule system.
Finally, we express the GPE (IV.1) in a specific system of units,
$$t_0=\frac{1}{\mathrm{\Omega }},r_0=\sqrt{\frac{\mathrm{}}{m\mathrm{\Omega }}},m_0=m,$$
(45)
for time, length and mass respectively. Technically, we should rename the scaled variables and parameter; say, $`t=\overline{t}t_0`$, $`\delta =\overline{\delta }/t_0`$. However, we eschew such a heavy notation, and continue to use, e.g., the symbol $`\delta `$ for what now stands for a dimensionless number and should be properly denoted by $`\overline{\delta }=\delta /\mathrm{\Omega }`$. The result is
$`i\dot{\mathrm{\Phi }}`$ $`=`$ $`\left[\frac{1}{2}^2\mu \right]\mathrm{\Phi }\mathrm{\Phi }^{}\mathrm{\Psi },`$ (46)
$`i\dot{\mathrm{\Psi }}`$ $`=`$ $`\left[\frac{1}{4}^2+\frac{1}{2}i𝐪+\delta 2\mu \right]\mathrm{\Psi }\frac{1}{2}\mathrm{\Phi }^2.`$ (47)
The GPE equations are as stated in Ref. JAV99a . The plan now is to analyze them. Of course, equations of the type (IV.1) are standard fare in studies of second-harmonic generation. The literature once more is extensive, and we will only give a few specific pointers as we go along.
### IV.2 Steady state of GPE
Atomic and molecular fields $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ represent a stationary state, one in which the physics does not change with time, if and only if their time evolution is solely in global (position independent) phase factors, possibly different ones for $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$. Now consider fields of the form $`\mathrm{\Phi }(𝐫,t)=\mathrm{\Phi }(𝐫)e^{i\omega t}`$ and $`\mathrm{\Psi }(𝐫,t)=\mathrm{\Psi }(𝐫)e^{2i\omega t}`$ for any real $`\omega `$. One sees right away that for this type of evolution, at least the exponential time dependence properly matches on both sides of Eqs. (IV.1). Although we have not been able to prove it mathematically, we conjecture that, assuming time independent parameters in Eqs. (IV.1), this type of time dependence is also the only possible case in which the physics is independent of time.
But then, by virtue of the transformation (IV.1), by readjusting the chemical potential one can always remove the time dependence of the fields entirely in any stationary state. The chemical potential started its life as an arbitrary parameter. From now on we make use of the arbitrariness and choose the value in such a way that in steady state the atomic and molecular fields are literally constants in time.
We shall see shortly that, for any time independent detuning $`\delta `$, one may pick a suitable value for the chemical potential $`\mu `$ so that the GPE (IV.1) have solutions $`\mathrm{\Phi }_0`$, $`\mathrm{\Psi }_0`$ that are constants in both space and time. But by virtue of the Galilean transformation (IV.1), we may then construct from $`\mathrm{\Phi }_0`$, $`\mathrm{\Psi }_0`$ a stationary solution corresponding to an arbitrary overall flow of atoms and molecules, solutions of the form $`\mathrm{\Phi }e^{i𝐤𝐫}`$ and $`\mathrm{\Psi }e^{2i𝐤𝐫}`$. Conversely, we believe that all spatially homogeneous stationary solutions are such Galilean boosts of the once-and-for-all constant solutions $`\mathrm{\Phi }_0`$, $`\mathrm{\Psi }_0`$. As far as it comes to spatially homogeneous steady states, we therefore may, and will, without any restriction on generality consider fields that are constants in space and time.
In the context of second-harmonic generation it is well known that the GPE may have spatially inhomogeneous solitary-wave type stationary solutions HE96 . By applying the Galilean transformation, one may find traveling solitary waves as well. Here we will make no effort to find, let alone classify, solitary-wave solutions to our GPE, but proceed as if the homogeneous steady states were all there is to it.
Within the scope of the present paper, the stationary solutions thus satisfy
$`0`$ $`=`$ $`\mu \mathrm{\Phi }_0\mathrm{\Phi }_0^{}\mathrm{\Psi }_0,`$ (48)
$`0`$ $`=`$ $`(\delta 2\mu )\mathrm{\Psi }_0\frac{1}{2}\mathrm{\Phi }_0^2`$ (49)
$`1`$ $`=`$ $`|\mathrm{\Phi }_0|^2+2|\mathrm{\Psi }_0|^2,`$ (50)
where the final equation originates from the normalization (35). Moreover, by virtue of the global gauge invariance, one may choose $`\mathrm{\Phi }_0`$ real. Next, because of the invariance of the GPE with respect to the sign change of $`\mathrm{\Phi }`$, we may always choose $`\mathrm{\Phi }_0`$ to be nonnegative. Furthermore, if $`\mathrm{\Phi }_0>0`$, then (48) shows that $`\mathrm{\Psi }`$ (like $`\mu `$) must be real. On the other hand, if $`\mathrm{\Phi }_0=0`$, $`\mathrm{\Psi }_0`$ comes with an arbitrary phase, and may be chosen real. All told, we only need to find the real solutions $`\mathrm{\Phi }_0`$, $`\mathrm{\Psi }_0`$, $`\mu `$ to Eqs. (IV.2), and besides only solutions with $`\mathrm{\Phi }_00`$ need be retained.
There are three distinct solutions. The trivial one reads
$$\mathrm{\Phi }_0^0=0,\mathrm{\Psi }_0^0=\frac{1}{\sqrt{2}},\mu ^0=\frac{1}{2}\delta ;$$
(51)
for $`\delta \sqrt{2}`$ we have
$`\mathrm{\Phi }_0^+`$ $`=`$ $`{\displaystyle \frac{\sqrt{6\delta ^2\delta \sqrt{6+\delta ^2}}}{3}},`$
$`\mathrm{\Psi }_0^+`$ $`=`$ $`{\displaystyle \frac{\delta \sqrt{6+\delta ^2}}{6}},`$
$`\mu ^+`$ $`=`$ $`\mathrm{\Psi }_0^+;`$ (52)
and for $`\delta \sqrt{2}`$ we find
$`\mathrm{\Phi }_0^{}`$ $`=`$ $`{\displaystyle \frac{\sqrt{6\delta ^2+\delta \sqrt{6+\delta ^2}}}{3}},`$
$`\mathrm{\Psi }_0^{}`$ $`=`$ $`{\displaystyle \frac{\delta +\sqrt{6+\delta ^2}}{6}},`$
$`\mu ^{}`$ $`=`$ $`\mathrm{\Psi }_0^{}.`$ (53)
At least two steady states are found for each $`\delta `$, and three in the interval $`\sqrt{2}<\delta <\sqrt{2}`$. The question is, which one represents the desired physical steady state. Here we attempt to mimic the ground state of the quantum-mechanical two-mode model (15), and choose the stationary solution so that JAV99a
$`\delta <\sqrt{2}`$ $`:`$ $`\mathrm{\Phi }_0=\mathrm{\Phi }_0^0,\mathrm{\Psi }_0=\psi _0^0,\mu =\mu ^0,`$
$`\delta \sqrt{2}`$ $`:`$ $`\mathrm{\Phi }_0=\mathrm{\Phi }_0^{},\mathrm{\Psi }_0=\psi _0^{},\mu =\mu ^{}.`$ (54)
The success is evident in Fig. 1, where we plot side by side the fraction of atoms converted to molecules from the quantum-mechanical two-mode model and the corresponding GPE approximation $`f=2\mathrm{\Psi }_0^2`$ (solid line) as a function of detuning. By comparing with the $`N=10`$ and $`N=100`$ quantum results, it may be seen that the agreement between the GPE and the quantum approach gets better as the number of atoms is increased. Even for an atom number as small as 100 and around the nonanalytic point $`\delta =\sqrt{2}`$ of the GPE approximation (54), the difference is only on the order of one per cent.
We conclude by noting that Eqs. (51) and (52) together give a stationary solution that is the exact mirror image of our choice (54) with the substitution $`\delta \delta `$. This corresponds to the state with maximum energy for the two-mode quantum system. It is a stationary state just as is the minimum, and may be used for atom-molecule conversion. The difference is that, in the case of the maximum, the detuning would be swept in the opposite direction, from below to above the dissociation threshold. Otherwise rapid adiabatic passage should work essentially as before.
### IV.3 Small excitations of the system
So far, while analyzing the GPE, we have achieved nothing more than in our studies of the two-mode model; rather less, because in our classical field theory we lose quantum fluctuations. Nevertheless, we now have the tools to analyze the stability of the steady state, and see how spatial inhomogeneities evolve in time.
As in Ref. JAV99a , we linearize the GPE around a stationary solution, and then attempt to find eigenmodes for small deviations from the stationary case. Both of these steps are achieved at once by inserting the Ansatz
$$\mathrm{\Phi }(𝐫,t)=\mathrm{\Phi }_0+u_\mathrm{\Phi }e^{i(𝐩𝐫\omega t)}+v_\mathrm{\Phi }^{}e^{i(𝐩𝐫\omega ^{}t)}$$
(55)
into the GPE, and only retaining the first-order terms in the “small” coefficients $`u_\mathrm{\Phi }`$, $`v_\mathrm{\Phi }`$. Of course, the field $`\mathrm{\Psi }`$ is treated similarly. We need to mix plane waves $`𝐩`$ and $`𝐩`$ because the GPE mixes fields and their complex conjugates. With our Ansatz we also retain the possibility that the evolution frequency of an eigenmode could be complex.
The Ansatz succeeds if the as of yet unknown evolution frequency $`\omega `$ satisfies the eigenvalue equations
$`\left[\frac{1}{2}𝐩^2+\mu +\omega \right]u_\mathrm{\Phi }+\psi _0v_\mathrm{\Phi }+\mathrm{\Phi }_0u_\mathrm{\Psi }=0,`$ (56)
$`\left[\frac{1}{2}𝐩^2+\mu \omega \right]v_\mathrm{\Phi }+\psi _0u_\mathrm{\Phi }+\mathrm{\Phi }_0v_\mathrm{\Psi }=0,`$ (57)
$`\left[\frac{1}{4}(𝐩^2+2𝐩𝐪)+2\mu \overline{\delta }+\omega \right]u_\mathrm{\Psi }+\mathrm{\Phi }_0u_\mathrm{\Phi }=0,`$ (58)
$`\left[\frac{1}{4}(𝐩^22𝐩𝐪)+2\mu \overline{\delta }\omega \right]v_\mathrm{\Psi }+\mathrm{\Phi }_0v_\mathrm{\Phi }=0.`$ (59)
The characteristic equation is fourth order in $`\omega `$, so in principle the solutions can always be written down analytically. However, here we use Mathematica MATHEMATICA to produce results directly numerically, and occasionally to extract analytical forms for special and limiting cases.
The dispersion relations of the excitation modes, $`\omega =\omega (𝐩)`$, depend on the relative propagation directions of the excitation and of light, and also on the size of the photon recoil kick. We encompass these dependences into a dimensionless parameter, which in terms of the original dimensional parameters reads
$$\eta =\sqrt{\frac{\mathrm{}(𝐪𝐩)^2}{m\mathrm{\Omega }𝐩𝐩}}.$$
(60)
Given the parameter $`\eta `$ and the absolute value $`p`$ of the (dimensionless) propagation vector $`𝐩`$, the terms that depend on photon recoil in Eqs. (IV.3) are replaced as follows,
$$\frac{1}{4}(𝐩^2\pm 2𝐩𝐪)\frac{1}{4}(p^2\pm 2\eta p).$$
(61)
For fixed values of the parameters $`\delta `$, $`\eta `$ and $`p`$, there are four excitation modes with four (in general) different evolution frequencies $`\omega `$. If a positive imaginary part is encountered in any one of the four frequencies, the corresponding mode grows exponentially and the steady state is unstable.
Let us take the wave vector characterizing the small-excitation mode to be perpendicular to the propagation direction of light, $`\eta =0`$. We begin with $`p=0`$ and assume $`|\mathrm{\Psi }_0|^2<\frac{1}{2}`$, a nontrivial mix of atoms and molecules. We find the solutions to the characteristic equation for the eigenvalue problem (IV.3)
$`\omega _1^2(0)`$ $`=`$ $`0.`$ (62)
$`\omega _2^2(0)`$ $`=`$ $`\left(3+{\displaystyle \frac{1}{2|\mathrm{\Psi }_0|^2}}\right)\left({\displaystyle \frac{1}{2}}|\mathrm{\Psi }_0|^2\right).`$ (63)
When the wave number of the excitation moves away from $`p=0`$, of the four frequencies the two evolving continuously from $`\pm \omega _2(0)`$ remain real. Their dispersion relations for small $`p`$ are of the type $`\omega _2(p)\omega _2(0)+\frac{1}{2m^{}}p^2`$, with $`m^{}0`$, so that these excitations are akin to optical phonons. On the other hand, the remaining two excitation frequencies are either purely real or purely imaginary, and for small enough $`p`$ they are always imaginary. Such modes do not propagate at all, but grow or shrink in place exponentially.
In fact, in Fig. 3, which is transcribed from Ref. JAV99a , we plot the largest imaginary part among the four frequencies $`\omega `$, $`|\mathrm{}[\omega _1]|`$, as a function of the detuning $`\delta `$ and the wave number $`p`$ of the excitation mode. As above, we set $`\eta =0`$. It may be seen that, for any detuning $`\delta >\sqrt{2}`$, an unstable excitation mode is always found. The largest imaginary part of an evolution frequency, i.e, the largest growth rate of an instability, is encountered at $`\delta =0.154496`$ and $`p=\pm 0.771324`$, and equals $`\mathrm{}(\omega )=0.24256`$.
The analogous instability is naturally known in second-harmonic generation, and goes under the rubric “modulational instability” TRI95 ; HE96a .
We next come to the effect of the direction of the wave vector of the mode on instability. In Fig. 4 we plot the largest positive imaginary part of a mode frequency found for any $`p`$ as a function of the parameter $`\eta `$. The curves are labeled with their corresponding fixed detunings $`\delta `$. In our studies of this kind, the largest growth rate of the instability was always found for $`\eta =0`$. For a given detuning, the small excitations whose momentum direction is perpendicular to light propagation always present the most unstable scenario.
Finally consider the largest growth rate of the instability (among all $`𝐩`$) as a function of the detuning. As may be inferred from Fig. 3, with increasing detuning $`\delta \stackrel{>}{}1`$ it decreases and occurs at smaller momenta, i.e., at larger size scales. In fact, for $`\delta 1`$, the maximum growth rate of an instability $`R(\delta )=\mathrm{max}\{\mathrm{}[\omega (p,\delta )]\}_p`$ behaves as
$$R(\delta )\left[\frac{1}{2\delta }\frac{5}{4\delta ^3}+𝒪(\frac{1}{\delta ^5})\right],$$
(64)
and the corresponding position of the maximum $`p_M(\delta )`$, the value of the momentum such that $`\mathrm{}[\omega (p,\delta )]\mathrm{}[\omega (p_M,\delta )]`$ for all $`p`$, goes like
$$p_M(\delta )\sqrt{\frac{1}{\delta }}\left[1\frac{1}{\delta ^2}+\frac{47}{16\delta ^4}+𝒪(\frac{1}{\delta ^6})\right].$$
(65)
Of course, as the largest growth rate of the instability seems to occur for $`\eta =0`$, this was our choice in Eqs. (64) and (65).
In sum, we have found that for $`\delta \sqrt{2}`$ the steady state of the atom-molecule system is stable, and for any $`\delta >\sqrt{2}`$ it is unstable. Although we have reported only on a specific stationary solution (54), possibly one out of three, the modes we have not considered explicitly are all unstable. But $`\delta =\sqrt{2}`$ is also the watershed, in that below $`\delta =\sqrt{2}`$ the steady state is all molecules ($`\mathrm{\Phi }_0=0`$), and above $`\delta =\sqrt{2}`$ some atoms are involved. All told, the steady state with everything in molecules is stable for $`\delta \sqrt{2}`$, but any steady state involving atoms is always unstable. The most unstable situation occurs around $`\delta 0`$, with about an equal mixture of atoms and molecules. For $`\delta 1`$ atoms take over, and the time scale for the instability grows longer.
### IV.4 Fate of unstable system
Linearized stability analysis has shown that a joint atom-molecule condensate is unstable in the presence of photoassociating light, i.e., there are small-excitation modes that grow exponentially. But when a small deviation from steady state grows exponentially, eventually it is not a small deviation anymore and the linearized analysis breaks down. To investigate the fate of the system once the instability has set in, we integrate the full GPE numerically as in Ref. JAV99a . Our aim is a qualitative demonstration, so we proceed in $`1+1`$ dimensions, one spatial coordinate $`x`$ and time $`t`$. However, there is nothing in our method that would not immediately work in higher spatial dimensions. The restriction is merely a matter of computer time.
Specifically, in our algorithm we discretize a stretch $`L`$ of the line into equidistant points $`x_i`$, and seek to represent the fields at these discrete points only. For convenience, we use periodic boundary conditions, so that the value of all functions repeats over $`L`$. The GPE is integrated over a time step $`h`$ in two moves. First we ignore the position derivatives in the GPE altogether. This implies that the stripped-down version of the GPE is local; for each $`x_i`$, $`\mathrm{\Phi }(x_i)`$ and $`\mathrm{\Psi }(x_i)`$ only couple to $`\mathrm{\Phi }(x_i)`$ and $`\mathrm{\Psi }(x_i)`$. We integrate the local version of the GPE over the time $`h`$ separately for each $`x_i`$ as two coupled differential equations using a second-order Runge-Kutta step. To prevent a numerical drift of the norm, we complete the initial step of the algorithm by normalizing the fields analogously to Eq. (35). Second, we ignore anything but the position derivatives in the GPE. The resulting partial version of the GPE is nonlocal, but in exchange it is linear and does not mix the fields $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$. To integrate over the same (sic!) time step as in the initial part of the algorithm, we first take the Fourier transform of the result from the first part using the Fast Fourier Transformation (FFT). In Fourier space position derivatives become local, so, to propagate the fields over the step $`h`$, we simply multiply their Fourier transforms by what is now the local, exact (within FFT), linear time evolution operator. Finally, in preparation for the next step, we transform back to real space.
This split-step algorithm is an obvious variation of the time honored split-operator methods for linear partial differential equations SPLITSTEP , and has been described before at least by the group of Firth SKR98 . Nonetheless, it comes with a fair dose of heuristics. It is therefore gratifying that we have been able to verify a good rate of convergence by using successively smaller time steps to integrate over a fixed interval of time.
We present an example of our simulations in Fig. 5, showing the absolute square of the atomic field $`|\mathrm{\Phi }|^2`$ as a function of position $`x`$ and time $`t`$. We use 128 points $`x_i`$. Because of the periodic boundary conditions, the left and right edges of $`x`$ wrap around and are actually the same. The plot is for the parameters $`\delta =0`$, $`\eta =0`$, the range of position $`x`$ is 24.6282 (in units of $`r_0`$) corresponding to three wavelengths of the most unstable excitation mode for these parameters, and time $`t`$ runs from 0 to 127 (in units of $`t_0`$). The unit of $`|\mathrm{\Phi }|^2`$, atom density, is such that for a homogeneous gas with everything in atoms, the density would be $`|\mathrm{\Phi }|^21`$. The system starts at time $`t=0`$ in the steady state appropriate for $`\delta =0`$, $`\eta =0`$, except that we add a small amount of Gaussian noise to each of the points specifying the initial state. Figure 5 is otherwise the same as Fig. 2 in Ref. JAV99a , except that a different seed was used for the random number generator that added the noise. As befits an instability, this innocuous change has lead to a totally different quantitative result at long times.
From the numerical simulations we see that the nature of the instability is such that the atoms and the molecules together combine into dense clumps. These clumps move around, and tend to join when they collide. As far as we can tell, within the present model only one big, dense clump will remain in the end.
One might wonder what is the mechanism behind the clumping. We present a heuristic guess. We are not talking about a thermodynamic system, so minimization of energy is a dubious principle to begin with; and besides, we are dealing with quasienergies in a rotating frame, not real energies. With these warnings out of the way, let us surmise that the system nevertheless attempts to minimize its energy. The atomic and molecular condensates make something akin to the two-level system in quantum optics. When light is added, one gets a dressed two-level system. The energy of the lower one of the two dressed states decreases with increasing Rabi frequency, which is the product of the electric field strength of the laser and the dipole moment. But the analog of Rabi frequency for the two-level system of atomic and molecular condensates, $`\mathrm{\Omega }`$, is also proportional to the square root of density, so the present system may also decrease its energy by increasing the density. Maybe this is what the instability is about.
## V Experimental considerations
The models we have discussed until now have been rather rudimentary. We next take up two types of complications that may come up in experiments. First, in Sec. V.1 enters an angle that is in principle included in our many-body Hamiltonian, although we have not yet considered it expressly: photodissociation of condensate molecules to states outside of the atomic condensate. It turns out to limit the rate at which one can achieve coherent atom-molecule conversion in adiabatic passage. Second, in Sec. V.2 we discuss a number of aspects that have so far been missing from our models: spontaneous emission from the photoassociated state, atom-atom interactions, trapping of atoms and molecules, and various level shifts. Spontaneous emission can be ameliorated by resorting to a two-color photoassociation scheme. Provided that photoassociation is speedy enough, which we believe is technically possible, the rest of these complications may be minor nuisances rather than dominant features of an experiment. What it takes to make photoassociation speedy enough is the subject of Sec. V.3, where we discuss the characteristic Rabi frequency for photoassociation $`\mathrm{\Omega }`$ for various alkalis. Finally, in Sec. V.4, we analyze the published experiment of Ref. WYN00 from our viewpoint of coherent photoassociation.
### V.1 Rogue photodissociation
As we noted already, in the case of photoassociation of an infinite homogeneous condensate, momentum conservation uniquely determines the state of the ensuing molecule. The converse, however, does not hold. When a molecule photodissociates into two atoms, momentum conservation does not force the atoms to return to the atomic condensate. Bose enhancement favors recombination of atoms with the condensate; the characteristic frequency is $`\mathrm{\Omega }`$ for both photoassociation and photodissociation between atomic and molecular condensates. Nonetheless, atoms winding up elsewhere are lost for coherent photoassociation. We have coined the term “rogue photodissociation” for photodissociation processes that send atoms outside the atomic condensate.
One might think that energy conservation in the cycle of photoassociation and photodissociation is the additional constraint that guarantees that the atoms return to the condensate. But this need not be a compelling argument. Any time dependence in the system interferes with energy conservation. For instance, suppose that photodissociation proceeds to the noncondensate states at the same rate $`\mathrm{\Gamma }_0`$ (per atom) as it would in the case of a nondegenerate gas of molecules, so that after a time $`\mathrm{\Gamma }_0^1`$ coherent photoassociation ceases. The photodissociation rate $`\mathrm{\Gamma }_0`$ in itself furnishes a time scale such that energy has to be conserved only to within $`\mathrm{}\mathrm{\Gamma }_0`$. The time evolution involved in nonlinear Rabi flopping would also interfere with energy conservation.
We present here a rudimentary model for rogue photodissociation for the special case when the detuning is swept in order to convert an atomic condensate to a molecular condensate. The key assumption is that we may employ the standard Markov approximation in the analysis of photodissociation. This entails that rogue photodissociation has no memory, but is characterized at each instant of time by a rate of exponential decay. Such an assumption seems dubious in particular when the laser is tuned to the close vicinity of the photodissociation threshold RZA82 . However, we know of no near-threshold case of this kind in which the breakdown of the Markov approximation has proven relevant in an experiment.
Evidently, only the condensate mode, one of very many atomic modes, is strongly affected by Bose enhancement. We take rogue photodissociation to proceed at the rate that would be appropriate for a nondegenerate gas at the given detuning. Second, we model the dependence of the photodissociation rate on detuning using the Wigner threshold law, so that we write
$$\mathrm{\Gamma }(\delta )=\theta (\delta )\sqrt{\frac{\delta }{\mathrm{\Omega }}}\mathrm{\Gamma }_0.$$
(66)
For convenience we have chosen the photoassociation frequency scale $`\mathrm{\Omega }`$ as the reference detuning for photodissociation rate; $`\mathrm{\Gamma }_0`$ is the photodissociation rate for the detuning $`\delta =\mathrm{\Omega }`$. As the third quantitative element of the model, we take the probability that a given atom is in the molecular condensate to be twice the square of the molecular field amplitude as solved from the classical field theory, and normalized as in (50). Explicitly, in dimensional units and for $`\delta 0`$, this probability is found from Eq. (53) as
$$P_M(\delta )=\frac{1}{18}\left[\left(\frac{\delta }{\mathrm{\Omega }}\right)+\sqrt{6+\left(\frac{\delta }{\mathrm{\Omega }}\right)^2}\right]^2$$
(67)
Suppose now that the detuning is swept as $`\delta =\xi \mathrm{\Omega }^2t`$, as in our rapid adiabatic passage example. Ignoring the depletion of the condensates due to the very same rogue photodissociation, we find the total probability for rogue photodissociation
$`P`$ $`=`$ $`{\displaystyle _{\delta (t)0}}𝑑tP_M[\delta (t)]\mathrm{\Gamma }_0\sqrt{{\displaystyle \frac{\delta (t)}{\mathrm{\Omega }}}}`$ (68)
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_0}{\xi \mathrm{\Omega }}}\left[{\displaystyle \frac{2^{3/4}}{3^{1/4}}}{\displaystyle _0^{\mathrm{}}}𝑑\tau \sqrt{\tau }\left(\sqrt{1+\tau ^2}\tau \right)^2\right]`$
$`=`$ $`\alpha {\displaystyle \frac{\mathrm{\Gamma }_0}{\xi \mathrm{\Omega }}},`$
where the numerical constant has the value $`\alpha =4.03197`$. Except for the numerical factor, this is the same expression we already used in a qualitative estimate in Ref. JAV99 .
Since the photodissociation rate grows linearly with light intensity, $`I`$, and the photoassociation characteristic frequency $`\mathrm{\Omega }`$ is proportional to the field strength of the laser, $`\sqrt{I}`$, in the end rogue photodissociation wins out as the light intensity is increased. Qualitatively, when $`P=1`$, rogue photodissociation has overtaken coherent conversion of atoms to molecules. To study this borderline case, we first set $`P=1`$ in Eq. (68) and solve $`\mathrm{\Gamma }_0`$ as a function of $`\mathrm{\Omega }`$. We then insert the result into Eq. (27), thus eliminating $`\mathrm{\Gamma }_0`$. Moreover, the velocity $`v_0`$ in Eq. (27) is then the relative velocity of the dissociated atoms corresponding to the detuning of the laser that gave the photoassociation rate $`\mathrm{\Gamma }_0`$, in this case $`\delta =\mathrm{\Omega }`$. Therefore we have $`v_0^2/2\mu =\mathrm{}\mathrm{\Omega }`$, and $`v_0`$ may be eliminated as well. We finally solve for the borderline value $`\mathrm{\Omega }`$ as a function of the problem parameters. After simple manipulations the ensuing characteristic frequency scale for photoassociation may be written
$$\widehat{\mathrm{\Omega }}=\left(\frac{8\sqrt{2}\pi \xi }{\alpha }\right)^{2/3}(\rho \mathrm{\lambda ̄}^3)^{2/3}ϵ_R.$$
(69)
Here we have introduced $`\mathrm{\lambda ̄}`$, wavelength of the light divided by $`2\pi `$, and the familiar recoil frequency for laser cooling
$$ϵ_R=\frac{\mathrm{}}{2m\mathrm{\lambda ̄}^2}.$$
(70)
The main finding is that rogue photodissociation restricts the light intensity that one may profitably use for adiabatic atom-molecule conversion. This means that there is also a maximum usable photoassociation frequency, or a minimum possible time scale for adiabatic atom-molecule conversion. The way we have written our estimate (69), the frequency scale is provided by the photon recoil frequency, and the corresponding time scale is in the ballpark of $`ϵ_R^1`$. The density dependence of photoassociation is encapsulated in the parameter $`\rho \mathrm{\lambda ̄}^3`$, the usual dimensionless parameter that governs the coupling of light with matter in a dense medium. For present-day condensates $`\rho \mathrm{\lambda ̄}^31`$ is a reasonable rule of thumb. Finally, we have a numerical constant that depends on the rate of sweeping of the detuning, but which may also be set equal to one in a rough estimate. Altogether, when in need of a qualitative number for the photoassociation frequency $`\mathrm{\Omega }`$, we resort to $`\mathrm{\Omega }ϵ_R`$.
Although our estimate of the minimum time scale for coherent photoassociation was developed for rapid adiabatic passage, we believe that (with $`\xi 1`$) it also applies to Rabi flopping. This is because the dimensional parameters of the problem are the same in both cases. In fact, as it comes to dimensional quantities, the minimum time scale for coherent photoassociation is equivalently written FOOT
$$\tau \frac{m}{\mathrm{}\rho ^{2/3}}.$$
(71)
This is the essentially unique quantity with the dimension of time that can be put together using the quantities characterizing a homogeneous, noninteracting, quantum mechanical, zero-temperature BEC; density, atom mass, and $`\mathrm{}`$.
### V.2 Physics missing from model
#### V.2.1 Spontaneous emission
We have discussed a one-color model for photoassociation. The physical drawback is that, where there is a strong dipole matrix element for photoassociation/dissociation driven by external light, there is also a strong dipole matrix element for spontaneous emission. For instance, if a photon is absorbed in photoassociation, then there is also a reverse spontaneous decay of the molecule. The molecule may end up in bound vibrational states, either in the same electronic manifold where the atoms started from, or in some other electronic manifold. Alternatively, the photoassociated molecule may decay back to a dissociation continuum in a process known as radiative escape. Either way, usually the probability is small that the system returns to the same two-atom state in which is started. After each process of spontaneous emission, two atoms are typically lost for any profitable use.
The analogous problem of an unstable excited state is standard fare in quantum optics, and so is the solution: add another laser-driven transition from the spontaneously decaying state to a stable state. In the same way, two-color Raman photoassociation, a free-bound transition followed by a bound-bound transition of the molecule, may take place between (nearly) non-decaying atomic and molecular states BOH96 ; MAC99 ; BOH99 . Here we study Raman photoassociation as a means of achieving an effective two-mode scheme; genuine three-mode phenomena such as STIRAP are discussed elsewhere MAC00
Our present notation for this scheme is sketched in Fig. 6. Suppose the first step of photoassociation takes place with absorption of a photon, then one sets up another laser to force, say, induced emission from the primary photoassociated state to a (more) stable bound molecular state. We call the Rabi frequency in the second step $`\chi `$, the detuning of the laser from resonance in the second transition $`\mathrm{\Delta }`$, and the spontaneous decay rate of the primary photoassociation state $`\mathrm{\Gamma }_\mathrm{s}`$. In this context $`\delta `$ stands for the two-photon detuning, the total energy mismatch for light-induced transition from the initial atoms to the final stable molecular state, including appropriate photon recoil corrections.
Let us model our three-level $`\mathrm{\Lambda }`$ scheme using a variation of the two-mode model (15) as
$$\frac{H}{\mathrm{}}=\delta g^{}g+(\mathrm{\Delta }+\delta )b^{}b\chi (b^{}g+g^{}b)\frac{1}{2}\kappa (b^{}aa+ba^{}a^{}),$$
(72)
where $`g`$, for “ground”, is the annihilation operator for the stable molecular state. Once more, we have added a multiple of the conserved particle number, this time in such a way that the atoms have zero energy. The term $`\chi `$ describes transitions between the bound molecular states. Spontaneous losses from the intermediate state are ignored for the time being.
The Heisenberg equations of motion from Hamiltonian (72) are
$`\dot{g}`$ $`=`$ $`i\delta g+i\chi b,`$ (73)
$`\dot{b}`$ $`=`$ $`i(\mathrm{\Delta }+\delta )b+i\chi g+\frac{1}{2}i\kappa aa,`$ (74)
$`\dot{a}`$ $`=`$ $`i\kappa a^{}b.`$ (75)
We next eliminate the intermediate state adiabatically with the assumption that $`\mathrm{\Delta }`$ is the largest evolution frequency in the system. We thus formally set $`\dot{b}=0`$, and obtain
$$b\frac{\chi g+{\scriptscriptstyle \frac{1}{2}}\kappa aa}{\mathrm{\Delta }}.$$
(76)
Inserting this into Eqs. (73) and (75), we have
$`\dot{g}`$ $`=`$ $`i\left(\delta {\displaystyle \frac{\chi ^2}{\mathrm{\Delta }}}\right)g+\frac{1}{2}i{\displaystyle \frac{\kappa \chi }{\mathrm{\Delta }}}aa,`$ (77)
$`\dot{a}`$ $`=`$ $`i{\displaystyle \frac{\kappa \chi }{\mathrm{\Delta }}}a^{}g+\frac{1}{2}i{\displaystyle \frac{\kappa ^2}{\mathrm{\Delta }}}a^{}aa.`$ (78)
These are Heisenberg equations of motion for the effective Hamiltonian
$$\frac{H_{\mathrm{eff}}}{\mathrm{}}=\left(\delta \frac{\chi ^2}{\mathrm{\Delta }}\right)g^{}g\frac{1}{4}\frac{\kappa ^2}{\mathrm{\Delta }}a^{}a^{}aa\frac{1}{2}\frac{\kappa \chi }{\mathrm{\Delta }}(a^{}a^{}g+g^{}aa).$$
(79)
The effective Hamiltonian describes a two-level system with the two-photon detuning $`\delta `$ and two-photon Rabi frequency $`\kappa \chi /\mathrm{\Delta }`$ in lieu of the usual detuning and Rabi frequency. There are two additional twists to the story. First, the two-photon resonance experiences a light shift $`\chi ^2/\mathrm{\Delta }`$, an old acquaintance from quantum optics. Second, we have an effective atom-atom interaction proportional to $`\kappa ^2/\mathrm{\Delta }`$. An analogous interaction in a nondegenerate gas was discussed earlier in Ref. SHL96 . Other than these tweaks, everything we have said about the two-mode model applies as before.
The adiabatic elimination (76) has a dark side hidden by a notational trick, namely, that it does not preserve boson commutators. Had we written the right-hand side in Eq. (75) as $`ba^{}`$ instead of $`a^{}b`$ and substituted (76), the atom-atom interaction term in the effective Hamiltonian would have displayed the operator ordering $`aaa^{}a^{}`$ instead of $`a^{}a^{}aa`$. The difference is immaterial for a large atom number, the limit we are studying anyway, but a more careful investigation of the adiabatic elimination would be in order if the atom number were not large.
As excessive care with operator products is not warranted, we write from the adiabatic assumption the number of atoms in the intermediate state $`b`$ qualitatively as
$$b^{}b\left(\frac{\chi }{\mathrm{\Delta }}\right)^2g^{}g+\left(\frac{\kappa }{2\mathrm{\Delta }}\right)^2(a^{}a)^2.$$
(80)
Let us scale the operators by $`\sqrt{N}`$, i.e., write $`b=\sqrt{N}\beta `$, and so forth. Then, within a factor of two, the quantum expectation value $`\beta ^{}\beta `$ is the probability that an atom is in the intermediate state, and so on. Equation (80) becomes
$$\beta ^{}\beta \left(\frac{\chi }{\mathrm{\Delta }}\right)^2\gamma ^{}\gamma +\left(\frac{\mathrm{\Omega }}{2\mathrm{\Delta }}\right)^2(\alpha ^{}\alpha )^2.$$
(81)
The fraction of atoms lost per unit time to spontaneous emission from the intermediate state equals $`2\mathrm{\Gamma }_\mathrm{s}\beta ^{}\beta `$.The intermediate detuning $`\mathrm{\Delta }`$ suppresses losses by a factor $`1/\mathrm{\Delta }^2`$, whereas the effective Rabi frequency scales as $`1/\mathrm{\Delta }`$. In principle, and at the present level of the physical model, it is possible to get rid of the harmful spontaneous emission to any desired degree by increasing the intermediate detuning.
#### V.2.2 Interactions between atoms and molecules
Atoms interact among themselves, molecules interact with molecules, and atoms even interact with molecules. For a dilute gas, the atom-atom interaction is often described by the effective two-body potential
$$U(𝐫_1,𝐫_2)=\frac{4\pi \mathrm{}^2a}{m}\delta (𝐫_1𝐫_2),$$
(82)
where $`a`$ is the $`s`$-wave scattering length for the atoms as before. One may write analogous models for atom-molecule and molecule-molecule interactions. If atom-atom interactions were suspected to be a factor, one could add to the field theory the usual two-body atom-atom interactions as
$$_{AA}=\frac{2\pi \mathrm{}^2a}{m}\varphi ^{}(𝐫)\varphi ^{}(𝐫)\varphi (𝐫)\varphi (𝐫),$$
(83)
and so on. Inelastic collisions, such as quenching of the molecules by collisions, may also prove important. At a phenomenological level, they could be described by using a complex scattering length for atom-molecule collisions.
Nonetheless, we have considered neither elastic nor inelastic collisions explicitly in this paper. The motivation is mainly pragmatic. Photoassociation is the novelty of this work anyway. Second, at this point in time virtually nothing is known about the scattering lengths for cases other than atom-atom interactions. Third, as pointed out in Sec. V.1, we anticipate a characteristic frequency scale for photoassociation, $`\mathrm{\Omega }`$, to be of the order of photon recoil frequency $`ϵ_R`$ of laser cooling, say, ten kilohertz. This is larger than a typical frequency scale associated with collisions, $`4\pi \mathrm{}a\rho /m`$, in many of the present alkali experiments. Photoassociation should dominate the action at least over short time scales.
More formally, we see from Eqs. (69) and (71) that the maximum usable photoassociation frequency scales with atom density as $`\rho ^{2/3}`$, while the rate of binary collisions scales as $`\rho `$. In principle and at this level of modeling, it is always possible to make photoassociation win out by decreasing the density.
Of course, not all of our discussions are for short times only. Notably, in the case of the instability of a joint atom-molecule condensate, it may well happen that collisional interactions eventually play a role in the clumping. We plan to return to collisional effects in a future publication, inasmuch as something worthwhile emerges from this front.
#### V.2.3 Trapping of atoms
In the current alkali vapor experiments one does not see infinite homogeneous condensates, but the condensate is ordinarily confined to a magnetic trap with a (practically quadratic) potential $`V(𝐫)`$. This may be taken into account in the field theory by adding a term in the Hamiltonian density,
$$_{\mathrm{VA}}(𝐫)=\varphi ^{}(𝐫)V(𝐫)\varphi (𝐫).$$
(84)
A similar additional term would describe trapping of molecules.
In fact, the beauty of the field theoretical formulation is that it does not depend on any given one-particle basis to describe the motion of the atoms and molecules. The photoassociation term was originally discussed using plane-wave states, which makes the derivation easy, but at the level of field theory there is no manifest vestige of plane waves anymore. Even if we add trapping potentials, there is no need to tamper with the photoassociation term. This should be contrasted with the more delicate situation that emerges if one tries to consider photoassociation directly using the eigenstates of the trap, or indeed some states that would take into account both trapping and atom-atom interactions.
Once more, if the photoassociation frequency scale is of the order $`ϵ_R`$, it is still vastly larger than the typical frequency scales associated with magnetic trapping of atoms, and the corresponding length scale for photoassociation, $`\mathrm{\lambda ̄}`$, is far smaller than the size of a typical condensate. Over short times the condensates behave locally as if they were homogeneous. One just applies the theory of an infinite condensates at each local density, and averages the results over the trap.
On the other hand, there are cases in our formulation where the trapping will matter. For $`\mathrm{\Omega }ϵ_R`$ the modulational instability may well set in as in a homogeneous condensate, but the motion of the atoms and molecules due to the trapping forces will certainly have a long-term effect on the atom-molecule clumps. We do not discuss this issue here, but plan to return to trapping in a future publication.
#### V.2.4 Level shifts
We have already mentioned a few mechanisms that can change the position of the photoassociation resonance. Atom-atom and molecule-molecule interactions alter the energy per atom or per molecule, and thereby modify the resonance condition for photoassociation. Since the atom-molecule ratio conversely depends on the detuning, the makings of bistability and hysteresis are in principle there. We have also discussed the light shift and the many-body shift in a two-color, three-mode configuration.
One more shift we have brought up before JAV98 ; MAC99 , but not yet in this paper, arises because the dissociation continuum is not flat. Given the initial state, the dipole matrix element (or more precisely, the square of the dipole matrix element per unit energy) depends on the final continuum state. The result is that the photodissociation rate picks up an imaginary part. That, of course, amounts to a shift of the photodissociating state with respect to the continuum. The shift is proportional to light intensity, and in a qualitative estimate is comparable to the photodissociation rate; see, e.g., Refs. jav .
Moreover, there are additional shifts due to the presence of each and every discrete state dipole-coupled to the initial state. The sum of all light shifts is actually finite only because the dipole coupling eventually tends to zero when one goes high enough in the energy of the coupled states. The implication is that, a priori, all dipole coupled states, even those off resonance by perhaps several photon energies, have to be considered explicitly. If one finds a significant contribution to the light shift from one far-off resonance state, chances are that one has to consider all of them.
We will not attempt to address continuum and non-resonant light shifts explicitly. Nonetheless, on the basis of the atomic case discussed in Refs. jav , we believe that, at least above the photodissociation threshold, the light shift should be reasonably independent of where exactly the laser is tuned. This would mean that, in the rapid adiabatic passage type atom-molecule conversion, the light shift merely gives a constant bias to the detuning, and is virtually inconsequential.
### V.3 Numerical examples
#### V.3.1 From rate to Rabi frequency
One-color photoassociation has been analyzed by various groups PA-groups , in particular for alkalis at finite temperature. A typical outcome is the photoassociation rate $`_v^{}`$ (in s<sup>-1</sup>) or the photoabsorption rate coefficient $`\alpha _v^{}^{PA}`$ (in cm<sup>5</sup>) for a given bound level $`v^{}`$ of the excited electronic state. The latter quantity is independent of two experimental parameters, namely atom density $`\rho `$ and laser intensity $`I`$. The rate of photoassociation $`_v^{}`$ is obtained via
$$_v^{}=\rho \phi \alpha _v^{}^{PA},$$
(85)
where the photon flux (photons/s cm<sup>2</sup>) is given by $`\phi =I/\mathrm{}\omega _L`$, $`\omega _L`$ being the photon frequency.
But we also know from our earlier work JAV98 ; MAC99 that the photoassociation rate in a thermal sample is given in terms of our detuning $`\delta `$, temperature $`T`$, and photoassociation rate $`\mathrm{\Gamma }`$ as
$$_v^{}=e^{\frac{\mathrm{}\delta }{k_BT}}\rho \lambda _D^3\mathrm{\Gamma },$$
(86)
where
$$\lambda _D=\left(\frac{2\pi \mathrm{}^2}{\mu k_BT}\right)^{1/2}$$
(87)
is the usual thermal de Broglie wavelength, albeit calculated using the reduced mass of the colliding atoms. In what follows, we write the detuning parameter for photoassociation as $`\delta =\omega _{\mathrm{}}\omega _L\mathrm{\Delta }_v^{}`$, where $`\mathrm{}\omega _{\mathrm{}}`$ is the asymptotic energy difference between the electronic curves and $`\mathrm{\Delta }_v^{}`$ is the red-detuning of level $`v^{}`$ from its asymptote. The binding energy of the molecular state is thus equal to $`\mathrm{}\mathrm{\Delta }_v^{}`$. Combining Eqs. (85) and (86) with Eq. (27), we have an expression for the characteristic frequency of coherent photoassociation to the level $`v^{}`$,
$$\mathrm{\Omega }_v^{}=e^{\frac{\mathrm{}\delta }{2k_BT}}\sqrt{\frac{2\pi \mathrm{}^2\alpha _v^{}^{PA}\phi \rho }{v\mu ^2\lambda _D^3}}.$$
(88)
In all of our discussion below we choose the detuning in such a way that $`\mathrm{}\delta =\frac{1}{2}k_BT`$, which gives the corresponding resonance velocity $`v=\sqrt{k_BT/\mu }`$. Strictly speaking, Eq. (27) requires the limit $`v0`$, or equivalently, $`\delta 0`$. However, we always assume, without explicitly checking this assumption, that the temperature is already low enough to bring the system into the region of validity of the Wigner threshold law. The quotient $`\mathrm{\Gamma }/v`$ in Eq. (27) has then supposedly reached the $`v0`$ limit.
We have already introduced the usual density parameter for light-matter coupling $`\mathrm{\lambda ̄}^3\rho `$ to characterize atom density, and the recoil frequency $`ϵ_R`$ as the frequency scale. In the same vein, we write the Rabi frequency for photoassociation as
$$\frac{\mathrm{\Omega }_v^{}}{ϵ_R}=\sqrt{\left(\frac{I}{I_v^{}}\right)}\sqrt{(\mathrm{\lambda ̄}^3\rho )},$$
(89)
where $`I_v^{}`$ is the characteristic light intensity for coherent photoassociation. In explicit numbers, we find from Eq. (88)
$$I_v^{}=\frac{1.46245\times 10^{26}}{\left[\frac{m}{\mathrm{u}}\right]^2\left[\frac{T}{\mathrm{K}}\right]\left[\frac{\lambda }{\mathrm{nm}}\right]^2\left[\frac{\alpha _v^{}^{PA}}{\mathrm{cm}^5}\right]}\frac{\mathrm{W}}{\mathrm{cm}^2},$$
(90)
which also displays the units used to express atomic mass $`m`$, temperature $`T`$, wavelength of light $`\lambda `$, and photoabsorption coefficient $`\alpha _v^{}^{PA}`$.
#### V.3.2 Calculated photoassociation rates
One can estimate the photoabsorption rate coefficient for a pair of atoms. At low temperatures, only $`s`$-wave scattering contributes to the process, and one finds COT95 ; cote-pra
$$\alpha _v^{}^{PA}(\omega ,T)\frac{4\pi ^2\omega }{3c}\lambda _D^3e^{\frac{\mathrm{}\delta }{k_BT}}|D_v^{}(\mathrm{}\delta )|^2,$$
(91)
In the low-temperature limit, the dipole matrix element $`|D_v^{}(\mathrm{}\delta )|^2`$ can be approximated by
$`|D_v^{}(E)|^2`$ $``$ $`|v^{},J=1|D|\mathrm{}\delta ,J=0,`$ (92)
$``$ $`\left({\displaystyle \frac{2\mu k}{\pi \mathrm{}^2}}\right)|D_0|^2(aR_v^{})^2S_v^{}^2`$
with $`E=\mathrm{}^2k^2/2\mu =\mathrm{}\delta `$, where $`D_0=|𝐝|`$ is the asymptotic dipole moment, $`a`$ is the scattering length, $`R_v^{}`$ is the classical outer turning point of the excited level $`v^{}`$, and $`S_v^{}`$ is a dimensionless parameter representing the fraction of the bound wave function contained in the last node COT95 ; cote-pra ; cote-jms .
For example, given <sup>7</sup>Li atoms in a triplet state at 1 mK, a detailed calculation cote-jms showed that the excited level $`v^{}=58`$ has the best Franck-Condon factor with the continuum and the highest-lying bound level $`v^{\prime \prime }=10`$ of the lower triplet electronic state, with $`\alpha _{v^{}=58}^{PA}=2.0\times 10^{32}`$ cm<sup>5</sup> note-rate . The corresponding characteristic intensity is $`I_{v^{}=58}=33\mathrm{W}\mathrm{cm}^2`$.
Beyond this, in Fig. 7 we give the calculated rate coefficients for high-lying vibrational states for both stable isotopes of Li, for both the singlet and the triplet excited states. Correspondingly, Table 1 presents a few numerical examples of the rate coefficients and characteristic intensities.
Much larger photoassociation rates and correspondingly smaller characteristic intensities can be obtained if the excited levels $`v^{}`$ are closer to the dissociation limit (i.e., have smaller binding energies). In fact, the $`v^{}=58`$ level is deeply bound, with a binding energy of $`\mathrm{}\mathrm{\Delta }_{v^{}=58}=5.5\times 10^4`$ hartree (or 122.22 cm<sup>-1</sup> note-hulet ). For example, $`v^{}=86`$ with a binding energy of $`0.87\mathrm{cm}^1`$ has a rate coefficient 1050 times larger, namely $`\alpha _{v^{}=86}^{PA}=2.1\times 10^{29}`$ cm<sup>5</sup>, which gives the characteristic intensity $`I_{v^{}=86}=0.032\mathrm{W}\mathrm{cm}^2`$. The high-lying excited levels $`v^{}`$ have larger rate coefficients because the overlap of the excited bound wave function and the continuum ground wave function is larger. For high $`v^{}`$, the overlap scales like $`\mathrm{\Delta }_v^{}^{1/2}`$.
Similar results for sodium and cesium are available. For small binding energies, the rate coefficient $`\alpha _v^{}^{PA}`$ at 1 mK is of the order of $`10^{28}`$ cm<sup>5</sup> for Na, and $`10^{27}`$ cm<sup>5</sup> for Cs more-rates . These translate to characteristic intensities as small as $`\mu \mathrm{W}\mathrm{cm}^2`$. Notice that for lithium, the rate coefficients for singlet transitions are smaller than for triplet transitions, reflecting the sign of the scattering length: this is especially significant for <sup>6</sup>Li, where the triplet scattering length is negative and enormous COT95a .
Other expressions based on semi-classical treatment have been developed pillet ; more-rates . For example, in pillet , the rate for $`s`$-wave contribution is
$$_v^{}\rho \lambda _D^3e^{\frac{\mathrm{}\delta }{k_BT}}\sqrt{\frac{\pi ^2}{3\delta \mathrm{\Delta }_v^{}}}\mathrm{sin}^2[k(R_v^{}A(k))]\kappa ^2;$$
(93)
where $`\kappa =𝐝𝐄/2\mathrm{}`$ is the Rabi frequency of the laser of intensity $`I`$, $`A(k)`$ represents a shift from the free solution, and $`\mathrm{\Delta }_v^{}`$ (again) is the detuning of the level $`v^{}`$ from the dissociation limit. This expression is valid for $`s`$-wave scattering and at low detunings. In the limit $`k0`$, $`A(k_R)a`$, and the maximum of $`_v^{}`$ is found to be at $`\mathrm{}\delta k_BT/2`$, so that
$$_v^{}\rho \lambda _D^3e^{1/2}\kappa ^2\sqrt{\frac{2\pi ^2\mu ^2}{3\mathrm{}^3\mathrm{\Delta }_v^{}}}\sqrt{k_BT}(R_v^{}a)^2.$$
(94)
This equation is similar to Eqs.(85)-(92). Notice that since $`\lambda _D1/\sqrt{T}`$, the rate scales as $`1/T`$ and $`\mathrm{\Delta }_v^{}^{1/2}`$ for small detuning. For lithium at 0.140 mK (corresponding to the Doppler temperature), assuming a laser intensity of 1000 W/cm<sup>2</sup> (so that $`\kappa =0.8\times 10^{10}`$ s<sup>-1</sup>: see pillet for details), and a detuning $`\mathrm{\Delta }_v^{}1`$ cm$`{}_{}{}^{1}30`$ GHz (or $`4.5\times 10^6`$ hartree), we have $`R_v^{}=135a_0`$ and $`\lambda _D=1489a_0=7.875\times 10^6`$ cm. For a density $`\rho =10^{11}`$ cm<sup>-3</sup>, this gives $`\rho \lambda _D^3=4.88\times 10^5`$ and, neglecting the scattering length $`a`$, one gets $`_v^{}9.6\times 10^4`$ s<sup>-1</sup>. If we scale the rate to $`T=1`$ mK, we get $`1.3\times 10^4`$ s<sup>-1</sup>, less than twice the value 7140 $`s^1`$ obtained from Eq. (85) using the calculated value of $`\alpha _{v^{}=86}^{PA}`$ given above. Notice that many assumptions are made in these estimates: nonetheless, expression (94) gives good order of magnitude for the photoassociation rate for small detuning and low temperatures.
Using the Doppler temperature as a typical temperature, and averaging the rate using a linewidth corresponding to 5 MHz pillet , one finds typical rates for small detunings by scaling the numbers of Table 2 with the appropriate temperature, density, laser intensity, and detunings, according to
$$_v^{}\frac{I}{I_0}\frac{\rho }{\rho _0}\frac{T_0}{T}\sqrt{\frac{\mathrm{\Delta }_0}{\mathrm{\Delta }_v^{}}}\overline{}.$$
(95)
Here, $`I_0=1000`$ W/cm<sup>2</sup>, $`\rho _0=10^{11}`$ cm<sup>-3</sup>, $`\mathrm{\Delta }_0=1`$ cm<sup>-1</sup>, and $`T_0`$ and $`\overline{}`$ are listed in Table 2. Correspondingly, $`\alpha ^{PA}`$ scales as
$$\alpha _v^{}^{PA}\frac{T_0}{T}\sqrt{\frac{\mathrm{\Delta }_0}{\mathrm{\Delta }_v^{}}}\overline{\alpha }^{PA}.$$
(96)
The characteristic intensities should then scale with the square root of the detuning $`\mathrm{\Delta }_v^{}`$, so that we have
$$I_v^{}=\sqrt{\frac{\mathrm{\Delta }_v^{}}{\mathrm{\Delta }_0}}I_0.$$
(97)
The values of $`\lambda `$, $`\phi `$ and $`\overline{\alpha }^{PA}`$ are listed in Table 2. To compare the rates for the various alkali metals, it is convenient to express them for the same parameters. Assuming $`I_0=1000`$ W/cm<sup>2</sup>, $`\rho =10^{11}`$ cm<sup>-3</sup>, $`T=100`$ $`\mu `$K, and $`\mathrm{\Delta }_0=1`$ cm<sup>-1</sup>, we obtain the values listed in Table 3. The photoassociation rates vary between $`1.2\times 10^5`$ s<sup>-1</sup> (or $`\alpha ^{PA}2.9\times 10^{28}`$ cm<sup>5</sup>) for Cs and $`6.3\times 10^5`$ s<sup>-1</sup> (or $`\alpha ^{PA}1.8\times 10^{27}`$ cm<sup>5</sup>) for Li. From the expressions for $`_v^{}`$, one expects the rates to scale like $`1/2\mu `$. If we were to multiply the rates by $`M`$, the mass number, we notice that, except for Li, the scaled rates indeed are similar. The variations left are due to slightly different Rabi frequencies pillet .
Finally, one has to be careful when using the values of Table 3, in which the effect of scattering lengths are not taken into account. These effects can be significant, like in the case of <sup>6</sup>Li, <sup>85</sup>Rb or <sup>133</sup>Cs, where large negative scattering lengths induce larger rates. Also, at larger detunings, one probes deeper region of the excited electronic state, and shorter distances of the lower state continuum wave function, where the exact nodal structure will play an important role (see COT95 ; cote-pra ; cote-jms and Fig. 7).
As already noted, the values for the saturation intensity in Table 3 are to be construed as estimates only. Also, as they are, they apply to a very weakly bound molecular state ($`\mathrm{\Delta }_v^{}=1\mathrm{cm}^1`$). Such high-lying states are likely to be easily perturbed by atom-molecule collisions, so that in practice one might resort to more bound molecular states. Nonetheless, these saturation intensities are remarkably low, occasionally much lower than our initial generic estimate of $`10\mathrm{Wcm}^2`$ JAV99 ; JAV99a . From the experimental viewpoint this is encouraging, as the requirements on laser intensity are greatly moderated. The flip side is that the usable intensity appears to be limited by rogue photoassociation; for densities such that $`\mathrm{\lambda ̄}^3\rho 1`$, the maximum photoassociation Rabi frequency is of the order of the recoil frequency $`ϵ_R`$ . For heavier alkalis the saturation intensities may indeed be low, but small recoil frequency and slow coherent photoassociation are the prices to pay.
### V.4 Discussion of an experiment
In the experiment of Wynar et al., Ref. WYN00 , the authors studied photoassociation of a <sup>87</sup>Rb condensate in a two-color Raman configuration. From our standpoint, the reported data fall into two categories. First, the width and shift of the two-photon resonance line was studied as a function of the intensities of the two lasers. We are going to use these measurements to estimate the bound-bound Rabi frequency. Second, the authors measured the width and the shift of the same resonance line as a function of the density of the condensate. At low intensities, elastic atom-atom and atom-molecule collisions turn out to have a major effect on the resonance parameters. However, these effects were modeled quantitatively in Ref. WYN00 , and one may then get at the quantities pertaining to photoassociation. We use this type of data to determine the characteristic frequency of coherent photoassociation.
#### V.4.1 Adapting the theory
We first cast our theoretical approach in a form that allows for direct juxtapositions with the analysis of experimental data in Ref. WYN00 . We ignore elastic collisions because a comparison with Ref. WYN00 turns out to be possible even without considering them explicitly, and inelastic collisions because their presence was not conclusively established in Ref. WYN00 .
We begin by adding the spontaneous decay of the intermediate state at the rate $`\mathrm{\Gamma }_s`$ as a nonhermitian term in the Hamiltonian for the $`\mathrm{\Lambda }`$ system, Eq. (72). This gives
$$\frac{H}{\mathrm{}}=\delta g^{}g+(\mathrm{\Delta }+\delta )b^{}b+\mathrm{}i\frac{\mathrm{\Gamma }_s}{2}b^{}b.$$
(98)
It is known in quantum optics that introducing an imaginary part to the energy of the decaying state gives correct results as long as the unstable state decays irreversibly only to states that are outside of the state space included in the model.
We then repeat much of the analysis of Sec. V.2.1 for the Hamiltonian (98), with a number of tricks and approximations. First, we take the limit of large intermediate detuning, $`|\mathrm{\Delta }|\mathrm{\Gamma }_s,|\delta |`$. Second, we assume that the bound-bound Rabi frequency $`\chi `$ is much larger than the free-bound Rabi frequency $`\mathrm{\Omega }`$. This implies that in the estimates of the line widths and decay rates, such as those following from Eq. (81), only the $`\chi ^2`$ terms need be kept. Third, we ignore light-induced atom-atom interactions. Fourth, we scale the the operators as before, $`a=\sqrt{N}\alpha `$, etc., which brings out the photoassociation frequency $`\mathrm{\Omega }`$ explicitly. Fifth, after deriving the equations of motion for the atomic operator $`\alpha `$ and final-state molecule operator $`\gamma `$ by adiabatic elimination of the intermediate-state operate $`\beta `$, in the spirit of semiclassical approximation we replace the operators $`\alpha `$, $`\beta `$, and $`\gamma `$ with $`c`$ numbers. At this stage we have, in analogy with Eqs. (77) and (78), the equations
$`\dot{\gamma }`$ $`=`$ $`i\left(\delta {\displaystyle \frac{\chi ^2}{\mathrm{\Delta }}}\right)\gamma {\displaystyle \frac{\mathrm{\Gamma }_s\chi ^2}{2\mathrm{\Delta }^2}}\gamma +\frac{1}{2}i{\displaystyle \frac{\mathrm{\Omega }\chi }{\mathrm{\Delta }}}\alpha ^2,`$ (99)
$`\dot{\alpha }`$ $`=`$ $`i{\displaystyle \frac{\mathrm{\Omega }\chi }{\mathrm{\Delta }}}\alpha ^{}\gamma .`$ (100)
To develop the rate approximation, we first note that the probability that an atom belongs to the stable state $`g`$, $`P_g=2|\gamma |^2`$, satisfies the equation of motion
$$\dot{P}_g=\frac{\mathrm{\Gamma }_s\chi ^2}{\mathrm{\Delta }^2}P_g+i\frac{\mathrm{\Omega }\chi }{\mathrm{\Delta }}(CC^{}).$$
(101)
The “coherence” $`C=\alpha ^2\gamma ^{}`$ has the equation of motion
$$\dot{C}=i\frac{\mathrm{\Omega }\chi }{\mathrm{\Delta }}[P_gP_a\frac{1}{2}P_a^2]+\left[i\left(\delta \frac{\chi ^2}{\mathrm{\Delta }}\right)\frac{\mathrm{\Gamma }_s\chi ^2}{2\mathrm{\Delta }^2}\right]C,$$
(102)
where, in turn $`P_a=|\alpha |^2`$ is the probability that an atom remains in the system unassociated.
We could go on and derive a corresponding equation of motion for $`P_a`$. However, here we assume that $`P_a1`$, and that correspondingly $`P_g1`$. We solve for the coherence $`C`$ adiabatically by setting $`\dot{C}=0`$, and insert the result into Eq. (101) to obtain a rate equation for the population of the stable molecular state $`g`$,
$$\dot{P}_g=\frac{\mathrm{\Gamma }_s\chi ^2}{\mathrm{\Delta }^2}P_g+\frac{\left[\frac{\mathrm{\Omega }\chi }{\mathrm{\Delta }}\right]^2\left[\frac{\mathrm{\Gamma }_s\chi ^2}{2\mathrm{\Delta }^2}\right]}{\left[\delta \frac{\chi ^2}{\mathrm{\Delta }}\right]^2+\left[\frac{\mathrm{\Gamma }_s\chi ^2}{2\mathrm{\Delta }^2}\right]^2}P_a^2.$$
(103)
This displays a shift of the ground state by
$$\mathrm{\Delta }\epsilon _L=\frac{\chi ^2}{\mathrm{\Delta }},$$
(104)
a linewidth
$$\frac{\gamma _L}{2}=\frac{\mathrm{\Gamma }_s\chi ^2}{2\mathrm{\Delta }^2},$$
(105)
and a corresponding direct decay of bound molecules at the rate $`\gamma _L`$.
With the assumption $`P_g1`$ we have made the rate equation (103) essentially one way, i.e., ignored transitions from molecules back to the atoms. The semiclassical approximation in general requires that there are many molecules in a quantum state, so that the molecules make a condensate. However, we believe that, by ignoring transitions from molecules back to atoms, we have divorced the rate equation (103) from any coherence requirement for the molecules. As the experiments do not claim a molecular condensate, this is essential for the analysis that follows.
#### V.4.2 Comparison with experiments
Wynar et al., Ref. WYN00 , discuss the light shift of the resonance line. In their case the frequencies of the two lasers are quite close so that both of them cause line shifts and broadenings, but this is included in their analysis. Translating back to our case with only one laser frequency on the bound-bound transition, the comparison between our formulation and Ref. WYN00 reads
$$\frac{\chi ^2}{\mathrm{\Delta }}=\frac{\beta I_2}{\mathrm{\Delta }_1},$$
(106)
where $`I_2`$ is the intensity (in W cm<sup>-2</sup>) of the bound-bound light, $`\mathrm{\Delta }_1\mathrm{\Delta }=2\pi \times 150\mathrm{MHz}`$ is the intermediate detuning, and $`\beta =3.5\times 10^9\mathrm{m}^2\mathrm{W}^1\mathrm{s}^2`$ is a parameter that we deduce from their Fig. 3a. This gives the approximate formula
$$\chi 0.94\sqrt{\frac{I_2}{\mathrm{Wcm}^2}}\times 2\pi \mathrm{MHz}.$$
(107)
Wynar et al., Ref. WYN00 , also present an analysis of line shapes at low intensity. The rate coefficient for two-photon photoassociation on resonance, $`K_0=9\times 10^{14}\mathrm{cm}^3\mathrm{s}^1`$, is one of the fitted parameters. The resonance photoassociation rate from our Eq. (103) and the corresponding expression in Ref. WYN00 , ignoring inelastic processes and with $`n\rho `$, boil down to
$$\frac{\mathrm{\Omega }^2}{\mathrm{\Gamma }_s}=nK_0.$$
(108)
Given the spontaneous decay rate of this molecular level, $`\mathrm{\Gamma }_s=12\times 2\pi \mathrm{MHz}`$ WYN00 , and noting that their cited value of $`K_0`$ was for the intensity of the photoassociation laser $`I_1=0.51\mathrm{W}\mathrm{cm}^2`$, we find the scaling formula
$$\mathrm{\Omega }5.8\sqrt{\frac{\rho }{10^{14}\mathrm{cm}^3}\frac{I_1}{\mathrm{Wcm}^2}}\times 2\pi \mathrm{kHz},$$
(109)
which gives the characteristic intensity $`I_v^{}=0.08\mathrm{Wcm}^2`$.
In the experiment, the binding energy of the primary photoassociated state was 23 cm<sup>-1</sup>. In contrast, for $`\mathrm{\Delta }_v^{}=1\mathrm{cm}^1`$, Table 3 gives the characteristic intensity $`0.07\mathrm{mWcm}^2`$. The theoretical and experimental characteristic intensities differ by three orders of magnitude, but they are also for different binding energies. The difference in characteristics intensities is no cause for concern, as we do not know at all how the detuning 23 cm<sup>-1</sup> for Rb is related to the minima of the photoassociation rates like those shown in Fig. 7 for Li.
#### V.4.3 Reaching coherence
Suppose that the density of the <sup>87</sup>Rb gas equals $`\rho =\mathrm{\lambda ̄}^3`$, a factor of two larger than quoted in Ref. WYN00 , the laser intensities are $`I_1=I_2=10\mathrm{W}\mathrm{cm}^2`$, typical values in photoassociation experiments, and the intermediate detuning is as in Ref. WYN00 . Then the two-photon Rabi frequency becomes $`\mathrm{\Omega }\chi /\mathrm{\Delta }=1\times 2\pi \mathrm{kHz}`$, while the effective linewidth is $`\chi ^2\mathrm{\Gamma }_s/2\mathrm{\Delta }^2=2\times 2\pi \mathrm{kHz}`$. For coherent phenomena we need a Rabi frequency at least comparable to the linewidth, a condition that is not quite satisfied. But if one were simply to take the intensities $`I_1=100\mathrm{W}\mathrm{cm}^2`$ for the photoassociating light and $`I_2=1\mathrm{W}\mathrm{cm}^2`$ for the bound-bound light, Rabi frequency remains unchanged and linewidth goes down by a factor of 100. One is then in principle deep in the regime of coherent photoassociation. Also, the effective Rabi frequency would be smaller than the recoil frequency, so that rogue photoassociation need not yet be a fatal problem.
On the other hand, eyeballing from Fig. 2 of Ref WYN00 , the linewidth due to elastic atom-atom collisions would be of the order of $`5\times 2\pi \mathrm{kHz}`$, so elastic collisions would interfere with coherent photoassociation. Rabi oscillations would be seriously impeded, rapid adiabatic passage maybe less. One could compensate by lowering the density by two orders of magnitude and at the same time increasing the intensity of the photoassociating laser by another two orders of magnitude to keep the effective Rabi frequency constant, but with lowering of the density one would wind up in trouble with rogue photoassociation.
Experimentally, one can think of varying the intensity of the two lasers, the intermediate state, the intermediate detuning, and the density of the sample. The conditions one must watch out for are that effective Rabi frequency be larger than the effective damping rate and line broadening due to elastic (and inelastic) collisions, yet not so large as to cause excessive rogue photoassociation. If the parameters could be varied freely, within our physical model a solution is always possible. Unfortunately, one must live with some practical limitations on, say, laser intensities. The present experiment of Ref. WYN00 is tantalizingly close to coherent photoassociation, but to really get there may require a careful optimization of the experimental parameters and a solid understanding of what is technically feasible in a given experiment. We will not attempt to enter such discussions.
## VI Concluding remarks
We have developed theory for coherent photoassociation of a Bose-Einstein condensate of atoms, which typically leads to a Bose-Einstein condensate of molecules. Phenomena analogous to coherent optical transients in few-level systems, such as Rabi flopping and rapid adiabatic passage between atomic and molecular condensates are expected. Coherent phenomena depend on the Bose enhancement of the dipole moment matrix element, and cannot occur in a nondegenerate gas, at least not in the thermodynamic limit. Bose-Einstein condensation and photoassociation are both currently front-line research and the marriage thereof might still be a technological challenge, but coherent optical transients in photoassociation should eventually be feasible in experiments.
###### Acknowledgements.
This work is supported in part by the NSF, Grant No. PHY-9801888, and by NASA, Grant No. NAG8-1428. Additional support was provided by the Connecticut Space Grant College Consortium. One of us \[JJ\] thanks Sándor Varró for interesting discussions; in particular, for pointing out Ref. VOLREF .
## Appendix A Dipole matrix element
In our bare-bones example we consider the situation in which we either have one molecule at rest and no atoms, or no molecule and two atoms in states $`𝐤`$ and $`𝐤`$. Then the relative momentum becomes $`\mathrm{}𝐤`$. This peculiarity, and also the reason for the factor $`\frac{1}{2}`$ in Eq. (14), is because the relative momentum for two entities with momenta $`𝐩_1`$ and $`𝐩_2`$ should be defined as $`\frac{1}{2}(𝐩_1𝐩_2)`$ to make it the conjugate of the usual relative position $`𝐫_1𝐫_2`$. We thus have the possible state vectors
$$|\psi =\left(\beta b^{}+\underset{𝐤}{}\alpha _𝐤a_𝐤^{}a_𝐤^{}\right)|0,$$
(110)
where $`\beta `$ and $`\alpha _k`$ are complex amplitudes to be determined, and $`|0`$ stands for the vacuum of atoms and molecules. The states for $`\alpha _𝐤`$ and $`\alpha _𝐤`$ are the same, so the sum over $`𝐤`$ runs only over half of the possible $`𝐤`$ values; say, those with $`k_x>0`$.
Following Ref. JAV99 , we write the time dependent Schrödinger equation from the Hamiltonian (12) in terms of the coefficients $`\beta `$ and $`\alpha _k`$, assuming a plane wave of light. In order to avoid an inessential complication, we ignore photon recoil and simply set $`𝐪=0`$. The result is
$`\dot{\alpha }_𝐤`$ $`=`$ $`i\left({\displaystyle \frac{\mathrm{}𝐤^2}{2\mu }}\delta _0\right)\alpha _𝐤+i{\displaystyle \frac{𝐝^{}𝐄^{}}{2\mathrm{}}}\beta ,`$ (111)
$`\dot{\beta }`$ $`=`$ $`i{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{𝐝𝐄}{2\mathrm{}}}\alpha _𝐤.`$ (112)
The QC approach discussed in Refs. JAV98 and MAC99 displays precisely the same structure, evidently identifying $`𝐝𝐄/2\mathrm{}`$ as the QC Rabi frequency $`\kappa `$ of those papers.
There is a catch, however. Because of the Bose-Einstein statistics, states with the relative momenta for two atoms $`𝐤`$ and $`𝐤`$ are the same, so the density of states for the relative motion is half of what it was for Maxwell-Boltzmann atoms. To obtain the same photodissociation rate, we therefore have to make the matrix element a factor of $`\sqrt{2}`$ larger than in the case of distinguishable atoms. This is the $`\sqrt{2}`$ in Eq. (14). The rate of photodissociation is proportional to the square of the matrix element and picks up a factor of two, which compensates for the missing half of phase space density.
The second-quantized Hamiltonian is chosen so that it gives the right photodissociation rate. What we did not fully realize in Ref. JAV99 is that the same factor of $`\sqrt{2}`$ also leads to twice the photoassociation rate that one would obtain for Maxwell-Boltzmann atoms, given the Franck-Condon factor from the standard calculations or as deduced from the photodissociation rate. It appears that in the literature the rate for $`s`$-wave photoassociation is usually calculated low by a factor of two. Besides, by a simple extension of the statistics arguments, $`p`$-wave photoassociation should vanish for bosons. Are these problems?
We think not. In photoassociation experiments no particular care is usually taken to polarize the atoms. They are not all in the same internal state, and therefore not all photoassociation events are between indistinguishable atoms. This should reduce the factor-of-two boson enhancement of the photoassociation rate for even partial waves, and call forth odd partial waves. The remaining discrepancies with the conventional calculations may well be within experimental and theoretical uncertainties.
It would be of some interest to investigate photoassociation of a polarized low-temperature gas experimentally. We believe that the statistics effects are real, and should be observable. After all, few colleagues seem to have a problem with the assertion that there are no $`s`$-wave collisions for polarized fermions. Photoassociation experiments in a polarized gas should make an interesting test of our phenomenological Hamiltonian.
## Appendix B Kernel for photoassociation in field theory
The process we followed in the main text in deriving the integral kernel $`𝐝(𝐫)`$ \[as in, e.g., Eq. (23)\] was a phenomenological mix between formally pure field theory and collisional physics.
In fact, to obtain the dipole matrix elements $`𝐝_{\mathrm{𝐤𝐤}^{}}`$ in Eq. (12), we carry out an integral of the form
$$𝐝_{\mathrm{𝐤𝐤}^{}}𝐝(𝐤𝐤^{})=𝐝d^3r\overline{\psi }^{}(𝐫)\overline{\varphi }_{𝐤𝐤^{}}(𝐫),$$
(113)
where $`\overline{\psi }`$ is the bound-state molecular wave function and $`\overline{\varphi }_{𝐤𝐤^{}}`$ describes the state of two atoms with the relative momentum $`\mathrm{}(𝐤𝐤^{})`$. Now, in the transformation to field theory (18) we used plane waves as the wave functions for states with a given wave vector. In the matrix element (113) we should correspondingly use plane waves to describe the relative motion, as in
$$\overline{\varphi }_{𝐤𝐤^{}}(𝐫)=\frac{1}{\sqrt{V}}e^{i(𝐤𝐤^{})𝐫}.$$
(114)
However, in our derivation we did not not use plane waves, but distorted plane waves that take into account atom-atom interactions, and only asymptotically (at large distances) corresponds to the relative momentum $`𝐤𝐤^{}`$.
This makes a difference. For instance, if the matrix element (113) is computed using pure plane waves, the final photoassociation kernel reads
$$𝐝(𝐫)=\sqrt{2}𝐝\overline{\psi }^{}(r),$$
(115)
instead of something akin to (LABEL:DR1). Depending on the value of the scattering length, there could be a substantial difference in the overall numerical factor also in the contact-interaction form for atom-molecule coupling, Eq. (24). The question is, which method is “correct”; pure or distorted plane waves?
First of all, atoms do interact, which in general will have an effect on photoassociation. One who wishes to analyze photoassociation quantitatively will have to take atom-atom interactions into account. A field theorist pursuing pure plane waves would have to consider atom-atom interactions explicitly, and work out the consequences in photoassociation ab initio. This is probably a major chore. Our hope is that, with our distorted-plane wave trickery, we have phenomenologically captured most of the effect of atom-atom interactions on photoassociation.
On the other hand, if we use the distorted plane waves to calculate the photoassociation matrix element and in addition introduce a contact interaction model for atoms as in (83), the suspicion arises that we are double-counting atom-atom interactions. Further approximations, such as the contact-interaction form for photoassociation or classical field theory, could either exacerbate or ameliorate the double-counting.
Disentangling atom-atom interactions and photoassociation might make an interesting theoretical exercise, but we do not attempt it in this paper. Instead, we proceed with a model that allows for both distorted-plane wave matrix elements for photoassociation, and an explicit atom-atom interaction. In this way, we at least get the limiting cases of photoassociation in a dilute gas and atom-atom interactions in the absence of photoassociation basically right with little effor
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# Investigating the structure of the accretion disk in WZ Sge from multi-wave-band, time-resolved spectroscopic observations: Paper II
## 1 Introduction
Dwarf-novae, among the cataclysmic variables (CVs), are semi-detached binary star systems where a white-dwarf (primary star) accretes gas from a cool companion (secondary star). The gas from the secondary star accretes through the inner Lagrangian point (L1), into the primary Roche lobe and slowly spirals in toward the white-dwarf. The in-flowing gas forms an accretion disk which orbits around the white-dwarf before actually being accreted onto the primary star. During the time the gas is within the accretion disk, current theory tells us that half of the gravitational energy of the gas is radiated outward while that which remains heats the accretion disk. Indeed most of the visible and IR light output by a typical dwarf-novae is from the accretion disk (?, ?). Emission lines seen in the spectra of a dwarf-nova arise from the vicinity of the accretion disk and can be double peaked due to the Doppler motions of the disk within the binary \[Horne & Marsh 1986\]. The velocity of the material within the disk, together with the emitted flux distribution over the disk, determines the line profile shape. Consequently, spectral line analysis is a fundamental method employed to understand the physical state of the accretion disk and its spatial flux distribution. The most common emission lines used in such analyses are those due to hydrogen and helium.
In this paper, we continue the analysis of the accretion disk in WZ Sge which was started in ? (hereafter PAPER I). PAPER I presented Doppler tomography of optical and IR disk emission lines as well as ratioed Doppler maps. General conclusions reached there indicated different optical depth at the hot-spot and the accretion disk, and an accretion disk structure dissimilar to those predicted by models. In this paper (PAPER II), we perform spectral analysis by direct measurements of the emission line profiles. We organize the present paper in three main sections. Section 2 describes the orbital phase dependence of H$`\alpha `$, H$`\beta `$, and IR emission lines. Section 3 presents newly determined radial velocity curves using optical and IR emission lines. We show that the amplitude and phase zero of the radial velocity curves are systematically biased. We show that this is due to different gas opacities in the hot-spot and the accretion disk, the hot-spot having a larger opacity than the accretion disk. We further develop the idea of increasing hot-spot contamination in systems where the difference between the optical depth in the hot-spot and the accretion disk is large, these systems being short orbital period, low mass transfer rate dwarf-novae. In section 4 we present our determinations for the accretion disk radii and compare our values with previously reported determinations.
Complete details concerning the data set used here (see table 1), are in PAPER I. All data are phased according to the photometric eclipse ephemeris stated in table 7, PAPER I. Optical spectra of WZ Sge were grouped into 50 phase bins per orbit in order to improve the signal-to-noise ratio of the data.
## 2 Emission lines
We analyze the emission lines evolution by inspecting their profile, measuring the fluxes and computing the Balmer decrement (last analysis on optical spectra only). Figure 4 summarizes the results from each analysis.
### 2.1 The line profiles
Figures 1 and 2 present a sample of optical and IR spectra respectively to show the complex emission line profiles and their dramatic evolution throughout an orbital period. Both optical and IR emission line profiles evolve similarly and owe most of their changes to the variable strength of the hot-spot around the orbit. Phases 0.3, 0.5, 0.84 and 0.63 are particularly suitable to point out the hot-spot evolution over the orbit.
Inspection shows that phase 0.3 is when the hot-spot, in both H$`\alpha `$ and H$`\beta `$ emission lines, is at the center of the accretion disk emission, i.e. has zero radial velocity with respect to an observer co-moving with the white-dwarf. From the system geometry determined in PAPER I we expect the hot-spot in the middle of the accretion disk emission line at phase $``$0.3 only if the gas at the stream-disk impact region has just the stream velocity (see figure 4). This condition implies that the accretion disk density is low enough such that the gas at the stream-disk impact region keeps its stream velocity with little effect from the Keplerian motion of the underlying accretion disk material. In the IR spectra the phase of the hot-spot zero radial velocity occurs slightly later, i.e. at phase 0.33. Phase binning of the optical spectra with the same time resolution of the IR spectra has showed that such a delay is significant and implies different velocity of the hot-spot gas in different wave-bands. We explain such a phase delay with the hot-spot emission from different lines spreading along the stream trajectory (figure 4).
The hot-spot appears again at the center of the emission line at phase 0.84, both in the optical and IR spectra. The expected phase for the optical spectra is 0.80 in the assumption of symmetric and isotropic hot-spot emission. Further evidence for asymmetric and anisotropic hot-spot emission is in the hot-spot profile which appears double peaked at phase 0.84. This feature is real as it is common to both H$`\alpha `$ and H$`\beta `$ binned spectra and is visible for a phase interval of $``$0.5 (see also the trailed spectrograms in figure 1 of PAPER I). A detailed study of the double structure is beyond the scope of this paper. The lower spectral resolution in the J and K bands does not resolve such a feature if present (the hot-spot peak separation is $``$9Å in the Balmer lines).
Phase $``$0.53 is when we expect the maximum hot-spot blue-shift and enhanced blue peak emission in the accretion disk lines. Observations do not confirm this expectation and show red and blue shifted peaks of roughly equal strength between phase 0.5 and 0.56, both in the optical and IR spectra. We explain this observation with little or no contribution from the hot-spot to the overall emission line flux. This is also evidence of the anisotropic hot-spot emission.
Asymmetries in the accretion disk are evident in the H$`\alpha `$ emission line which displays a stronger red-shifted peak at phase 0.63. Doppler maps in PAPER I show asymmetric accretion disks also in the H$`\beta `$, Br$`\gamma `$, He I, and Br$`\delta `$ lines. The accretion disk emission line profile evolves from a shallow U-shaped profile to a V-shaped profile, between phase 0.5 and 0.63. This is particularly evident in the optical spectra; among IR observations, only the J band spectra show similar evolution. ?) predict U-shaped emission line profiles for optically thin accretion disks, and V-shaped profiles for the optically thick case. This might suggest that the changing profile of the emission lines is evidence that some areas of the disk have optically thin line emission and some have optically thick.
### 2.2 Line flux
The emission line flux was measured by manually marking the position of line wings with the cursor and summing up the flux from each pixel in between; the continuum flux corresponding to a linear interpolation of the same two cursor position was automatically subtracted during the line flux measurements. Each of the emission lines measured (with the exception of the Pa$`\beta `$ line) show quite strong modulations apparently correlated with the line profile evolution (figure 3). The modulations display maxima at phase 0.84 and 0.3, and a minimum at phase 0.5. Thus, the hot-spot shows stronger emission when viewed from the outside (phase 0.84) and from the inside (0.3), while it does not contribute significantly to the emission line flux when viewed from upstream (phase 0.5). Assuming that the flux at phase 0.5 is due only to the accretion disk (see figure 1 and figure 2, and previous subsection) we can compute the relative contribution by the hot-spot to the total emission line flux. At phase 0.84 the hot-spot emits 50$`\%`$ of the total line flux in H$`\beta `$ and $`30\%`$ in H$`\alpha `$ and the IR lines. When, viewed form the inside the hot-spot relative contribution decreases to 38$`\%`$ in H$`\beta `$ and to $`20\%`$ in H$`\alpha `$ and Br$`\gamma `$. The He I 2.06$`\mu `$m line does not display an evident hump at phase 0.3, thus the relative hot-spot contribution was not computed there. Assuming a spherical hot-spot with radius $`S=3.1\times 10^8`$cm \[Smak 1993\], cylindrical accretion disk with thickness $`2H3.8\times 10^8`$cm <sup>1</sup><sup>1</sup>1The value was derived computing the scale height H as given in \[Williams & Shipman 1988\] assuming a uniform temperature T=6000K in the accretion disk and a white dwarf mass of M$`{}_{WD}{}^{}=0.82`$M (PAPER I)., and inclination $`i=75^o`$ \[Smak 1993\], the decrease of the hot-spot flux from phase 0.84 to phase 0.30 due to occultation by an optically thick disk is 70%. We observe a decrease of only 30% in H$`\beta `$ and $``$50% in H$`\alpha `$ and Br$`\gamma `$. Thus either the accretion disk gas is at least partially transparent to the hot-spot emission line flux, or the assumed geometry inappropriate.
The continuum light curve is expected to be double-humped as already observed by ?) in the IR, and by ?) and ?) in the optical. Right panel of figure 3 shows that only our K-band continuum measurements match with previous observations. The optical continuum shows large uncertainties on each data point and poor modulation throughout the orbit as the method used to measure line and continuum flux does not take into account the underlying white-dwarf absorption. The IR spectra, which do not show any evidence of white-dwarf features have the expected double-humped structure and smaller uncertainties in the continuum light curves. Fluxes from the Pa$`\beta `$ emission line and its underlying continuum have not been reported in figure 3 and table 2 because the data show variations throughout the run of observation. The run covered about four cycles and variations between cycles do not appear to be correlated. We cannot say whether this is due to intrinsic variability of the emitting source or to external causes.
### 2.3 The Balmer decrement
We measured the Balmer decrement for both the accretion disk and the hot-spot region to check for different opacities. The Balmer decrement is defined by the ratio of the line intensity in frequency units $`D_\nu (H\alpha /H\beta )I_\nu (H\alpha )/I_\nu (H\beta )`$. The Balmer decrement, as measured from spectra, relies on line fluxes, because the line intensity is broadened by effects due to both the intrinsic properties of the emitting gas and the instrumentation used during the observation. The comparison of measured flux-ratios to tabulated intensity-ratios assumes that the emission lines of interest have identical profiles, each broadened in a similar way. We observe Balmer lines which are far from being identical in shape, however, table 6 of PAPER I show that time averaged disk profiles have identical velocity broadening. De-blending of the components is not possible without a complex model of hot-spot shape and accretion disk structure, thus we measured the peak intensity for each component, on continuum-subtracted spectra. We measured the disk emission at each phase using only the red or blue line peak, whichever was unaffected by the hot-spot emission at the given phase. The hot-spot line intensity has always some underlying disk emission blending, except near phase $`0.3`$ and $`0.84`$, when the hot-spot emission is at the center of the emission line and the relative contribution from the accretion disk is negligible. The top panel of figure 5 shows that the Balmer decrement within the accretion disk is approximately constant throughout the orbit with an average value of $`D_\nu (H\alpha /H\beta )=3.82`$, which is larger than the predicted value of $`D_\nu (H\alpha /H\beta )1`$ for an optically thick accretion disk (?, ?). $`I_\nu (H\alpha )`$ and $`I_\nu (H\beta )`$ each show phase dependent modulations, implying a non-uniform accretion disk emission, but their ratio remains constant. However, despite the uncertainties, the Balmer decrement in the accretion disk shows a clear trend around phase 0.3 toward higher opacities. The bottom panel of figure 5 shows the hot-spot Balmer decrement vs orbital period. The Balmer decrement here is strongly phase dependent and shows two deep minima at phase intervals 0.2-0.3 and 0.8-0.84, i.e. when observing perpendicular to the stream trajectory at the hot-spot position and measuring the hot-spot line intensities without any underlying accretion disk emission. The hot-spot Balmer decrement is $`D_\nu (H\alpha /H\beta )`$2. The value of $`D_\nu (H\alpha /H\beta )`$ outside phases 0.2-0.3 and 0.8-0.84 is $``$ 3 and approaches the average accretion disk Balmer decrement when the hot-spot emission is weak. In the hot-spot, $`I_\nu (H\beta )`$ shows a stronger modulation around the orbit than $`I_\nu (H\alpha )`$ which is largely biased by the underlying accretion disk emission.
The Balmer decrement provides information about the gas temperature and density where the emission line forms; it gives a unique gas density for each assumed temperature. However, care must be taken because different temperatures may give the same value of the Balmer decrement depending on the assumption of optically thick or optically thin emission lines. In the interpretation of our results we mainly use models by ?) who modeled emission lines in optically thin accretion disks. He computed H$`\beta `$ strengths and Balmer decrements for a grid of temperatures, orbital inclinations, and mid-plane nucleon densities. In figure 6 we plot the results of ?) together with the Balmer decrement we measure at the hot-spot (lower horizontal line) and the accretion disk (upper horizontal line) in WZ Sge.
The value of $`D_\nu (H\alpha /H\beta )`$=2 for the hot-spot region corresponds to a black-body optically thick emission with a temperature of T$`<`$5000K. This decrement apparently matches a number of observed Balmer decrements in CVs \[Williams 1983\]. However, because we observe the He I emission line at 2.06$`\mu `$m, parts of the hot-spot region must be at least 15000K (see section 3). An optically thin gas with a temperature of 15000K and a Balmer decrement of 2 corresponds to a gas density of $`\mathrm{log}N_o=11.5`$ in Williams’ model (see figure 6), which is lower than the derived gas density within the accretion disk (see below). The gas density in the stream is expected to be 3 to 4 times larger than the gas density within the accretion disk \[Lubow & Shu 1976\]. We thus conclude that the hot-spot is at least partially optically thick and characterized by a steep temperature gradient. The Balmer emission lines form down-stream from the initial impact region from a relatively low temperature (T$``$5000K) optically thick gas, and the He I 2.06$`\mu `$m emission line forms in the outer stream-impact region where the temperature is higher (see section 3).
The Balmer decrement $`D_\nu (H\alpha /H\beta )`$ = 3.82 of the gas within the accretion disk would correspond to a black-body of temperature T$``$3400K. However, we expect the temperature to be higher as in an optically thick $`\alpha `$-disk the temperature is never expected to fall below 6000K over a wide range of mass accretion rate \[Williams 1980\]. Moreover, Ratioed Doppler tomograms and radial disk profile of WZ Sge (see PAPER I) suggest a partially optically thin accretion disk. We conclude that the accretion disk is optically thin in the Balmer emission lines. According to ?) a Balmer decrement of $`D_\nu (H\alpha /H\beta )`$ = 3.82 corresponds to an optically thin gas of temperature 6000K and density $`\mathrm{log}N_o`$ 13.5 (see figure 6). A roughly similar value of $`\mathrm{log}N_e13`$ for gas at a temperature of T=5000K may be derived from ?) models. ?), using the escape probability approach, determined Balmer decrements and hydrogen lines strength for an infinite slab of gas at various opacities, temperatures, and densities. We can use these estimates to determine the optically thin accretion disk mass assuming: i) a uniform gas density of $`\mathrm{log}N_o=13.2`$, ii) a cylindrical accretion disk with radii as determined in PAPER I (see also table 6), and height 2H as already defined in subsection 2.2. The result is an accretion disk mass of M$`{}_{d}{}^{}5.7\times 10^{15}M_{}`$, which is about 5 orders of magnitude smaller than the total mass accreted during the inter-outburst period ($`10^9M_{}`$, ?). It is also $``$4 order of magnitude smaller than the critical mass (i.e. the maximum accretion disk mass before an outburst) expected in low mass transfer rate systems <sup>2</sup><sup>2</sup>2We derived the value M$`{}_{crit}{}^{}4\times 10^{11}M_{}`$ either using average values for the accretion disk radii and surface density as predicted by ?) ($`R_{in}4.5\times 10^9`$cm, $`R_{out}1.1\times 10^{10}`$cm, $`\mathrm{log}\mathrm{\Sigma }2.4`$, for $`\dot{M}10^{11}M_{}yr^1`$), or formula (18) in ?) with $`\alpha 0.003`$ as predicted for TOADs by ?).. However, typical values for the mass transfer rates and critical masses have been computed using the assumption of $`\alpha `$-disk models which are optically thick and cannot therefore be compared directly with optically thin gas models. ?) predicts an accretion disk that is only partially optically thin at a mass transfer rate of $`\dot{M}10^{11}M_{}yr^1`$, the optically thin region starting from a radius of $`5\times 10^9`$cm, outward. This value seems in quite good agreement with the radius of $`r2.6\times 10^9`$cm at which we observe the transition of the Balmer line ratio from optically thick to optically thin case (see section 4 and figure 7 in PAPER I).
## 3 Radial Velocity Curves
Optical and IR spectra were used to measure radial velocity curves for H$`\beta `$, H$`\alpha `$, Pa$`\beta `$, Br$`\gamma `$ and He I ($`\lambda 2.06\mu `$m) emission lines. Radial velocity curve measurements were not produced for the two emission lines He I $`\lambda 1.08\mu `$m and Pa$`\gamma `$ as they were strongly blended together.
Figure 7 presents our radial velocity curves and their best fit for the five emission lines listed above. Each data point corresponds to the radial velocity of the emission line central wavelength. Central wavelengths were computed by using a technique similar to the Pogson’s method to determine the mid-time of an eclipsing binary with an asymmetric eclipse \[Glasby 1970\]. The radial velocity measurements were fit by a grid-search method minimizing $`\chi ^2`$. The sinusoidal fitting function used was:
$$v_i=\gamma +K_1\times \mathrm{sin}[2\pi (\varphi _i+0.50\varphi _o)]$$
(1)
where $`v_i`$ is the measured velocity and $`\varphi _i`$ the observed photometric phase. The fitting parameters $`\gamma `$, K1, and $`\varphi _o`$ are respectively the systemic velocity, the Keplerian velocity of the white-dwarf around the system’s center of mass, and the red-to-blue crossing in the radial velocity curve, corresponding to secondary inferior conjunction. We added the phase shift of 0.50 in Eq.1 in order to have $`\varphi _o`$ exactly matching the secondary star inferior conjunction, and not the primary star inferior conjunction.
Our best fit parameters from each emission line represent an inconsistent set of values (see table LABEL:tab:rvfit). The systemic velocity, $`\gamma `$, runs from negative to positive values, spanning a range of $``$ 300km sec<sup>-1</sup> wide. The white-dwarf Keplerian velocity, K1, runs from $``$46 to 120 km sec<sup>-1</sup>; while the red-to-blue crossing, $`\varphi _o`$, swings about phase 0.19, covering a range $`\mathrm{\Delta }\varphi `$0.1 in phase. Interestingly, a similarly inconsistent set of system parameters was determined for the optical radial velocity curves of the SU UMa type dwarf-nova VY Aqr \[Augusteijn 1994\].
When fitting and interpreting the sinusoidal radial velocity curves, it is assumed that: i) the disk orbits with Keplerian velocity around the white-dwarf, and ii) the accretion disk follows the white-dwarf motion around the binary’s center of mass. These approximations fail when hot-spot emission and disk asymmetries are not negligible. ?), shows that the hot-spot can induce a measured K1 larger than the true value and a spuriously delayed moment for inferior conjunction. ?) explain a larger K1 amplitude with non Keplerian motion.
We believe that our measured radial velocity curves result from combination of two different motions: 1) the motion of the accretion disk as described at points i) and ii) of the previous paragraph; and 2) the motion of the gas within the hot-spot region which moves along the ballistic trajectory at the stream velocity. Combination of the two velocity components in points 1) and 2) will change the amplitude and phase of a sinusoidal radial velocity curve, the largest departures occurring in systems where the hot-spot emission represents the largest fraction of the total emission line flux. An anisotropically emitting hot-spot may affect the $`\gamma `$ velocity measurements. Optically thin accretion disks are likely to be low density gas regions (see section 3.1), which neither affects the stream velocity of the in-falling gas, nor absorbs the flux emitted at the impact region. Thus, emission line profiles from optically thin accretion disks are expected to be heavily biased by the hot-spot emission.
The phase offset is found to be correlated with the energy required to produce each emission line. Figure 8 shows the observed phase offset versus the excitation potential of each H emission line as given by the standard formula. Figure 8 clearly shows a linear relationship between the observed phase offset and the emission line energy and implies a smeared out hot-spot region with decreasing excitation energy and temperatures as the gas moves down-stream. Results similar to those in figure 8 may be obtained using the data in ?) on VY Aqr. We fit the phase offset computed by Augusteijn for three hydrogen emission lines and found a linear relationship with almost identical slope and y-intercept as in figure 8.
The He I 2.06$`\mu `$m emission line was not considered in figure 8. The He I line 2.06$`\mu `$m is a transition strongly coupled with the resonant line at 584Å. Both lines decay from the same 2<sup>1</sup>P level but the 584Å transition has a probability of occurrence which is $`10^3`$ times larger \[Najarro et al. 1994\]. Strong He I emission at 2.06$`\mu `$m is rarely seen in astronomical objects and observed only in the extended atmosphere, wind and/or disk-like formations of WN and Be super-giant stars, where they form because of the high temperature, velocity and density of the out-flowing gas (15000-30000K, $``$700 Km sec<sup>-1</sup>, and $`10^5`$ M yr<sup>-1</sup>, respectively, ?). The presence of the He I 2.06$`\mu `$m emission requires the He I 584Å line to form in an optically thick, high temperature 15000-30000K region, which, in short period dwarf-novae, can only occur in the shock heated stream-disk impact region. Formation of He I 2.06$`\mu `$m in a hot corona is excluded by the Doppler tomogram in PAPER I figure 3.
We interpret our results above to provide a picture for the quiescent accretion disk in WZ Sge (figure 4), where each hot-spot emission line is likely to be aligned along the stream trajectory into the accretion disk, with the higher energy emissions arising at the higher temperature outer edge of the accretion disk impact region, and the lower energy lines primarily forming further down-stream.
### 3.1 Phase offset as a measure of the hot-spot bias
White-dwarf and/or hot-spot eclipses are present in the broad band photometric light curves of many dwarf-novae. In the case of the white-dwarf eclipse, the observed eclipse minimum defines the time of the secondary star inferior conjunction. In systems displaying only a bright spot eclipse, the phase offset between binary phase zero and the hot-spot eclipse is small ($``$0.05 in phase, ?), so that we still have a good idea of the time of the true binary phase zero. Similarly, when measuring radial velocity curves, we expect the spectroscopic ephemeris to match the photometric one in systems with perfectly symmetric accretion disk, and a phase offset larger than zero in systems having asymmetries in their disk. In particular, a larger phase offset is expected in systems where the hot-spot represents the main asymmetry in the accretion disk and a considerable fraction of the total emission line flux. This is the case of WZ Sge given that our observations average out possible short term variation and Doppler tomograms in PAPER I show accretion disk asymmetries always considerably fainter than the hot-spot emission; moreover in section 2 we showed the accretion disk to be optically thin. To test the hypothesis of phase offset as a measure of the hot-spot bias and the optical properties of the accretion disk gas we searched the literature for measured dwarf-novae phase offsets and plot these values against the binary orbital period. Theory predicts low mass transfer rates in short orbital period systems \[Howell, Rappaport & Politano 1997\], implying low density possibly optically thin accretion disks.
Phase offsets have been observed in several different CVs (?, ?, and ?), but until now there has been no explicit or systematic search for its cause. ?), observed phase offsets in both nova-like and SU UMa stars, claiming that the hot-spot cannot be the cause of such an offset. ?) concluded that the best candidate to explain the phase offset observed in five dwarf-novae is only the hot-spot. Finally, ?) collected most of the observed phase offsets in the literature and plotted them vs. the binary orbital period but did not provide a possible explanation.
The eleven dwarf-novae listed in table 4 have been carefully selected in order to provide a homogeneous sample. This was achieved by selecting: i) high inclination dwarf-novae that show photometric eclipses, ii) only Balmer line radial velocity curves, and iii) publications reporting both photometric and spectroscopic analysis. Restrictions for point i) are to avoid biases other than the hot-spot and to reduce unknown inclination effects. Point ii) reduces the scatter in the measured phase offset by focusing on emission lines arising from regions with similar physical conditions, and finally, point iii) guarantees consistently measured phase offset and reduces any bias due to possible changes in the quiescent accretion state. Figure 9 plots our selected sample which provides evidence for the expected correlation: the phase offset is larger in short orbital period systems.
## 4 Accretion disk size
It is possible to estimate the inner and outer accretion disk radii in a dwarf-nova by measuring the emission line wing and peak separations respectively, and assuming a Keplerian velocity. For an annulus of gas in Keplerian motion at a distance $`r`$ from a white-dwarf of mass M<sub>WD</sub>, orbiting at an angle $`i`$ to the line of sight, the Keplerian velocity is:
$$V_{kep}=\sqrt{\frac{GM_{WD}}{r}}\mathrm{sin}i$$
(2)
In a double peaked emission line from an accretion disk, the Keplerian velocity of the outer disk radius corresponds to half the line peak separation (hereafter HPS); while, the Keplerian velocity at the inner accretion disk radius corresponds to half the line wing separation (hereafter HWS).
It is important both to have reasonable system parameters and to measure HWS and HPS on unbiased spectra, i.e. spectra in which the wings and peaks of the accretion disk emission line are not affected by hot-spot emission (figure 1, phase 0.3 and 0.84). Spectra in which the hot-spot emission clearly contaminates wings and/or peaks of the accretion disk emission line profile may provide a measure of the hot-spot bias. Both HWS and HPS were measured reading the cursor position with the IRAF task splot. HWS values from each spectra are the average over 5 measurements.
### 4.1 Measurements on uncontaminated spectra
We selected a sample of uncontaminated phase binned H$`\alpha `$ and H$`\beta `$ spectra. The assumption is that the hot-spot emission does not affect the wings of the accretion disk emission lines. This selection criteria produced three binned H$`\alpha `$ spectra (at phases 0.28, 0.30 and 0.32) and six binned H$`\beta `$ spectra (at phases 0.28, 0.30, 0.32, 0.34, 0.82 and 0.84). We applied the same selection criteria to IR spectra and selected two uncontaminated spectra in both the J and K band (at phases 0.33 and $``$0.84 in each case). In case the hot-spot feature was not evident (figure 2, J-band, phase 0.3) we select those spectra which are the closest in phase to the uncontaminated optical spectra listed above. HWS and HPS were first measured on each bin separately, and then averaged. HWS and HPS measured values are in table 5. We see that values from optical and IR emission lines are in good agreement each other, the exceptions being the He I $`2.06\mu `$m and Pa$`\beta `$ lines which apparently have smaller HPS value and larger HWS values than in the H$`\alpha `$, H$`\beta `$, and Br$`\gamma `$ measurements. We applied a t-student test (to a confidence level of $`\alpha =0.05`$), to check for statistical differences between pairs of values. The test showed that the HWS values have to be considered statistically identical in all the emission lines, i.e. the same inner disk radii; while the H$`\beta `$ HPS is statistically larger than the measured HPS from H$`\alpha `$, Br$`\gamma `$, He I $`2.06\mu `$m, and Pa$`\beta `$, thus implying a smaller outer disk radius ($`10\%`$).
Velocities from uncontaminated spectra were used in equation 2 to compute the accretion disk radii for a variety of primary masses and inclinations (see table 6). It is evident that small white-dwarf masses (M$`{}_{WD}{}^{}<0.5`$M), imply a small accretion disk outer radius ($`<50\%`$ of the primary Roche lobe radius) and a relatively large ratio R<sub>WD</sub>/R<sub>in</sub> ($`<0.5`$); while larger masses (M$`{}_{WD}{}^{}1`$M) give a large outer disk radius (up to 90% of the primary Roche lobe radius) and R<sub>WD</sub>/R$`{}_{in}{}^{}0.07`$, implying an annulus shaped accretion disk. The input value of the white-dwarf mass does not affect the ratio, R=R<sub>in</sub>/R<sub>out</sub>, which is equal to 0.2 in each case.
### 4.2 Measurements on the averaged spectra
Optical and IR average spectra show relatively symmetric accretion disk emission lines. This is expected because the hot-spot motion, i.e. the S-wave, smoothes out the hot-spot features at each phase, and evenly distributes its emission throughout the orbit. Average spectra also show broader emission lines than single spectra. This happens for two reasons: i) the common orbital motion of the accretion disk and the white-dwarf around the binary center of mass will broaden the disk emission line by a term K1, i.e. the white-dwarf Keplerian velocity; and ii) the wings of the hot-spot emission blend with the blue and red-shifted accretion disk emission lines. The S-wave velocity is usually smaller than the outer disk edge Keplerian velocity (?; hereafter GKS), thus, we expect average spectra with broader accretion disk emission peaks and smaller peak separation than observed in single spectra. HPS and HWS from average spectra are given in table 5. HPS values from optical and IR emission lines agree fairly well and confirm –with the exception of He I $`2.06\mu `$m and Pa$`\beta `$– the expectation. HWS values are in good agreement too, emission lines with larger HWS value are probably affected by high velocity components from the hot-spot. HWS and HPS measured in the averaged spectra can give information about the S-wave bias when compared with wing and peak separations from uncontaminated spectra. They also test the assumptions used in the accretion disk radii computations.
We compared HWS values from averaged and uncontaminated spectra after correction for the white-dwarf orbital motion. HWS from uncontaminated spectra, and HWS from averaged spectra, K1 subtracted, should be equal in the absence of hot-spot high velocity components. Table 7 shows that HWS values from uncontaminated spectra and K1 subtracted averaged spectra agree very well each other in the case of H$`\beta `$ and Br$`\gamma `$ emission lines; while H$`\alpha `$ and He I $`2.06\mu `$m have HWS values from K1 subtracted average spectra which are $``$ 450 and 250 Km sec<sup>-1</sup> larger than expected in the absence of high velocity components in the hot-spot. Such components bias both radial velocity curve measures and the accretion disk radii computation. The fact that we observe the largest hot-spot high velocity component in the H$`\alpha `$ line is consistent with figure 4 and the less energetic emission forming downstream in the accretion disk. High velocity hot-spot components in the Pa$`\beta `$ emission line should be taken cautiously, because of larger uncertainties in the measurement.
It is worth noting that each average spectra was K1 subtracted by the correspondent white-dwarf orbital motion as determined in section 3. The K1 values from radial velocity curve fitting are biased by the hot-spot and do not reflect the true orbital motion of the white-dwarf. However, because the hot-spot bias is believed to increase the true K1 value (see section 3, ?, and ?), then the applied correction to our averaged spectra should be over-estimated.
Table 5 also reports HWS and HPS as we measured by hand from the average H$`\alpha `$ spectra in figure 1 of GKS. Comparison of the GKS determinations with those from our averaged spectra gives a chance to check for changes in the quiescent accretion state of WZ Sge. The data in GKS were obtained almost one year after the last outburst of WZ Sge in Dec 1978 and ours were obtained about 18 years after the outburst. We measured larger HWS and HPS values in the GKS spectra than in our averaged spectra. In particular, in the GKS spectra the HWS is $``$200 Km sec<sup>-1</sup> larger than in our averaged spectra; while, the HPS matches the values we measured on our uncontaminated spectra. Therefore, we may conclude that: i) the H$`\alpha `$ hot-spot emission high velocity component was possibly larger one year after the super-outburst than 18 years later; ii) the hot-spot emission was probably weaker at the time of GKS observations such that it didn’t affect the accretion disk peak separation; iii) the peak separation in GKS spectra may be considered free from hot-spot bias and implies constant outer accretion disk radius in the H$`\alpha `$ emission.
Particularly interesting is the conclusion of point iii) which claims a constant accretion disk outer radius. Indeed, the same HPS value has been observed in all quiescent spectra of WZ Sge obtained in the last 40 years (see table 8). One possible explanation is in the assumption of the outer accretion disk radius at the 3:1 tidal resonance radius. A constant accretion disk outer radius has some important implications relating to current outburst theory. Some disk outburst theories \[Meyer-Hofmeister, Meyer & Liu 1998, Wynn et al. 2000\] predict that, during quiescence, the outer radius slowly increases up to the 3:1 resonance radius. Once reached at the 3:1 resonance, quiescent super-humps should be present.
### 4.3 Smak’s method of determining R<sub>in</sub>/R<sub>out</sub>
The accretion disk radii ratio, R=R<sub>in</sub>/R<sub>out</sub>, is often the only value provided by authors because of the large uncertainties affecting the determination of the orbital inclination and the binary star masses. A commonly applied method to compute R is the one developed by ?).
Smak’s method assumes axially symmetric accretion disk in Keplerian motion and a power law flux distribution ($`fr^\alpha `$) for the disk emission. The R and $`\alpha `$ values depend on the parameters U, A<sub>84</sub>, and A<sub>41</sub> which are defined as follow: U=$`res/\mathrm{\Delta }\lambda `$, with $`res`$=instrumental resolution, A$`{}_{84}{}^{}\mathrm{log}W_{0.8}\mathrm{log}W_{0.4}`$, and A$`{}_{41}{}^{}\mathrm{log}W_{0.4}\mathrm{log}W_{0.1}`$, where W<sub>0.8</sub>, W<sub>0.4</sub>, and W<sub>0.1</sub> are the emission line width at the fractions 0.8, 0.4, and 0.1 of the peak height above the continuum, respectively.
?) applied Smak’s method to WZ Sge and determined A$`{}_{84}{}^{}=0.07`$, A$`{}_{41}{}^{}=0.1`$, and R=0.3 (see table 9). A ratio R=0.3 is 50% larger than the value 0.2 we determined in section 4.1. We investigated such a difference to understand whether it reflects changes in the quiescent accretion disk size or it implies inconsistencies between the Smak’s method and the method described in section 4.1. We applied Smak’s method to our H$`\alpha `$ and H$`\beta `$ uncontaminated spectra and found A<sub>84</sub> and A<sub>41</sub> values matching those by Mennickent and Arenas (see table 9). However, we expect a ratio $`0.2R<0.3`$. Then, we may conclude that: i) the two methods are consistent and, in particular, the assumptions of axisymmetric disk and power law flux distribution required by the Smak’s method are not fundamental to the determination of R and Smak’s method may be applied also to asymmetric accretion disks such as WZ Sge; ii) the two accretion disk radii ratio, R=0.3 by Mennickent and Arenas and R=0.2 by us, may correspond to a real change in WZ Sge’s quiescent accretion state and indicate an inward motion of the inner accretion disk radius toward the white-dwarf surface (in section 4.2 we showed that the outer disk radius remains constant during quiescence).
## 5 Summary and Conclusions
We observed both hot-spot and accretion disk line emission to vary in shape and strength throughout an orbital period, and conclude that both the hot-spot and the accretion disk are asymmetric and anisotropically emitting. We determined the gas at the impact region to have just the stream velocity, and the accretion disk is a low density gas with little drag effects on the in-falling material from the secondary star. There are different Balmer decrements at the impact region and in the rest of the accretion disk, and we conclude that the gas has a different opacity, $`\tau `$, within the two regions. In particular, the hot-spot is optically thick in the lines, while the accretion disk appears to be optically thin in the lines.
Our measured radial velocity curves, from five emission lines in the optical and in the infrared are found to present an inconsistent set of system parameters. We explain this result as a bias due to the hot-spot emission. In particular, we showed that the hot-spot emission delays the apparent time of the secondary inferior conjunction with respect to photometric phase zero. We also show that the phase offset depends both on the excitation potential energy of the considered emission line and on the optical depth of the gas in the accretion disk. Our extended analysis of a selected sample of dwarf-novae shows increasing phase offset with shorter orbital period, consistent with the idea of optically thin accretion disks in low mass transfer rate systems.
Summarizing our results from the emission line profile analysis and radial velocity curve computations presented here, we can formulate a probable structure for the WZ Sge accretion disk. The accretion disk is generally a low density, optically thin gas. The accretion disk is neither symmetric, nor has it a uniform gas density or temperature structure (see also PAPER I). The hot-spot is optically thick in the emission lines, not symmetric in shape, and not an isotropic emitter. The hot-spot emission arises from a multicomponent extended region in which each emission component is not visible at all phases and the hot-spot emission varies in strength throughout the orbital period. These variable hot-spot emissions lead to an extended hot-spot region which does not emit strongly when viewed from the down-stream direction (phase 0.5 and following), and shows a temperature gradient along the stream trajectory. Figure 4 summarizes our conclusion.
The qualitative picture of WZ Sge provided by our observations and data analysis shows a peculiar quiescent accretion disk dissimilar to others observed at present. WZ Sge is the first short orbital period dwarf-nova showing evidence of optically thin emission both in the continuum (?, and ?), and the lines. Multicolor photometry of OY Car, Z Cha, and HT Cas (?, ?, ?) shows evidence of quiescent accretion disk optically thin in the continuum but thick in the emission lines. This work establishes a first step toward an understanding of the spatial flux and material distribution within the WZ Sge accretion disk and other TOAD (Tremendous Outburst Amplitude Dwarf-novae, ?) candidates.
The accretion disk radii determined in sec.4.1 show that determination of the inner and outer radius strongly depend on the assumed white-dwarf mass and the binary orbital inclination. Unfortunately, white-dwarf mass and radius determinations are still affected by large uncertainties. The uncertainties on the WZ Sge white-dwarf parameters prevent us from uniquely determing the actual size of the accretion disk, the fraction of the primary Roche lobe it fills, or its exact shape. Whether the accretion disk extends down to the white-dwarf surface or is a ring-like accretion disk as suggested by some previous studies (?, ?, ?), cannot be uniquely determined. We determined the accretion disk radii ratio in WZ Sge, R=R<sub>in</sub>/R<sub>out</sub>, to be 0.2, and observed that the outer accretion disk radius does not vary during the quiescence period. One explanation for the constant accretion disk outer radius is the assumption that the accretion disk extends to the 3:1 tidal resonance radius. One possible test of this idea would be to search for and find quiescent super-humps.
## Acknowledgments
The optical data used in this study were obtained by Henke Spruit and Rene Rutten with the William Herschel Telescope of the ING at La Palma. The William Herschel Telescope is operated on the island of La Palma by the Isaac Newton Group in the Spanish Observatorio del Roque de los Muchachos of the Instituto the Astrophysica de Canarias. The United Kingdom Infrared Telescope is operated by the Joint Astronomy Centre on behalf of the U.K. Particle Physics and Astronomy Research Council. SBH acknowledges partial support of this research from the NSF grant AST 98-19770 and from the University of Wyoming office of research. We thank Dr. Peter Tamblyn for his help and information concerning He I emission.
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# 1 Introduction
## 1 Introduction
Since the (2+1) dimensional O(3) nonlinear sigma model (NLSM) was first discussed by Polyakov and Belavin , there have been lots of attempts to improve this soliton model associated with the homotopy group $`\pi _2(S^2)=Z`$. The SU(N) invariant NLSM or $`CP^{N1}`$ model was later introduced in terms of $`N`$ complex fields $`Z_\alpha `$ ($`\alpha =1,\mathrm{},N`$) satisfying a constraint $`Z_\alpha ^{}Z_\alpha 1=0`$. In addition, one can impose a local U(1) invariance
$$Z_\alpha (x)e^{i\alpha (x)}Z_\alpha (x)$$
(1.1)
for arbitrary space-time dependent $`\alpha (x)`$ . The $`CP^{N1}`$ model for $`N=2`$ was shown to be equivalent to the O(3) NLSM where Polyakov and Belavin found instantons . Moreover, it is well known that the topological charge in the O(3) NLSM is equivalent to that in the SU(2) invariant NLSM .
On the other hand, the Dirac method is a well known formalism to quantize physical systems with constraints. In this method, the Poisson brackets in a second-class constraint system are converted into Dirac brackets to attain self-consistency. The Dirac brackets, however, are generically field-dependent, nonlocal and contain problems related to ordering of field operators. These features are unfavorable for finding canonically conjugate pairs. However, if a first-class constraint system can be constructed, one can avoid introducing the Dirac brackets and can instead use Poisson brackets to arrive at the corresponding quantum commutators.
To overcome the above problems, Batalin, Fradkin, and Tyutin (BFT) developed a method which converts the second-class constraints into first-class ones by introducing auxiliary fields. Recently, this BFT formalism has been applied to several models of current interest . In particular, the relation between the Dirac and BFT schemes, which had been obscure and unsettled, was clarified in the SU(2) Skyrmion model. Very recently, in Ref. the BFT Hamiltonian method was also applied to the SU(2) Skyrmion model to directly obtain the compact form of a first-class Hamiltonian, via the construction of the BFT physical fields. Meanwhile, this BFT approach was applied to $`CP^1`$ model , but these studies fell short of obtaining the desired compact form of a first-class Hamiltonian and as a result further developments have been deterred.
The motivation of this paper is to systematically apply the standard Dirac quantization method, the BFT scheme , the Batalin, Fradkin and Vilkovisky (BFV) method and the Becci-Rouet-Stora-Tyutin (BRST) method to the $`CP^1`$ model . In section 2 we convert the second-class constraints into first-class ones according to the BFT method to construct first-class BFT physical fields and directly derive the compact expression of a first-class Hamiltonian in terms of these fields. The existing approach, used in Ref. to derive a first-class Hamiltonian in the SU(2) Skyrmion model, involves an infinite iteration procedure. Our approach avoids this. We then investigate some properties of the Poisson brackets of these BFT physical fields to obtain the Dirac brackets in the limit of vanishing auxiliary fields. We construct in section 3 a BRST-invariant gauge fixed Lagrangian in the BFV scheme through the standard path-integral procedure. Exploiting collective coordinates, in section 4 we perform a semiclassical quantization and in section 5 we explicitly show the equivalence between the $`CP^1`$ model and O(3) nonlinear sigma model (NLSM) at the canonical level, by using the Hopf bundle . In section 6 we show that the energy spectrum of rigid rotator in the $`CP^1`$ model obtained by the standard Dirac method with the suggestion of generalized momenta is consistent with that of the BFT scheme.
## 2 First-class constraints and Hamiltonian
In this section we apply the BFT scheme to the $`CP^1`$ model, which is a second-class constraint system. We start with the $`CP^1`$ model Lagrangian of the form
$$L=\mathrm{d}^2x\left[_\mu Z_\alpha ^{}^\mu Z_\alpha (Z_\alpha ^{}_\mu Z_\alpha )(Z_\beta ^\mu Z_\beta ^{})\right]$$
(2.1)
where $`Z_\alpha =(Z_1,Z_2)`$ is a multiplet of complex scalar fields with a constraint
$$\mathrm{\Omega }_1=Z_\alpha ^{}Z_\alpha 10.$$
(2.2)
One notes here that, as discussed before, this model is invariant under a local U(1) gauge symmetry transformation (1.1). By performing the Legendre transformation, one can obtain the canonical Hamiltonian
$$H_c=\mathrm{d}^2x\left[\mathrm{\Pi }_\alpha ^{}\mathrm{\Pi }_\alpha +_iZ_\alpha ^{}_iZ_\alpha (Z_\alpha ^{}_iZ_\alpha )(Z_\beta _iZ_\beta ^{})\right]$$
(2.3)
where $`\mathrm{\Pi }_\alpha `$ and $`\mathrm{\Pi }_\alpha ^{}`$ are the canonical momenta conjugate to the complex scalar fields $`Z_\alpha `$ and $`Z_\alpha ^{}`$, respectively, given by
$`\mathrm{\Pi }_\alpha `$ $`=`$ $`\dot{Z}_\alpha ^{}Z_\alpha ^{}Z_\beta \dot{Z}_\beta ^{}`$
$`\mathrm{\Pi }_\alpha ^{}`$ $`=`$ $`\dot{Z}_\alpha Z_\alpha Z_\beta ^{}\dot{Z}_\beta .`$ (2.4)
The time evolution of the constraint $`\mathrm{\Omega }_1`$ yields an additional secondary constraint
$$\mathrm{\Omega }_2=Z_\alpha ^{}\mathrm{\Pi }_\alpha ^{}+Z_\alpha \mathrm{\Pi }_\alpha 0$$
(2.5)
and $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ form a second-class constraint algebra
$$\mathrm{\Delta }_{kk^{}}(x,y)=\{\mathrm{\Omega }_k(x),\mathrm{\Omega }_k^{}(y)\}=2ϵ^{kk^{}}Z_\alpha ^{}Z_\alpha \delta (xy)$$
(2.6)
with $`ϵ^{12}=ϵ^{21}=1`$.
Following the BFT formalism which systematically converts the second class constraints into first class ones, we introduce two auxiliary fields $`\mathrm{\Phi }^i`$ according to the number of second class constraints $`\mathrm{\Omega }_i`$ with the Poisson brackets
$$\{\mathrm{\Phi }^i(x),\mathrm{\Phi }^j(y)\}=ϵ^{ij}\delta (xy).$$
(2.7)
The first class constraints $`\stackrel{~}{\mathrm{\Omega }}_i`$ fulfilling the simplest closed algebra
$$\{\stackrel{~}{\mathrm{\Omega }}_i(x),\stackrel{~}{\mathrm{\Omega }}_j(y)\}=0$$
(2.8)
are then constructed as follows
$$\stackrel{~}{\mathrm{\Omega }}_i(x)=\mathrm{\Omega }_i(x)+\mathrm{d}^2yX_{ij}(x,y)\mathrm{\Phi }^j(y)$$
(2.9)
where the matrix $`X_{ij}`$ satisfies the relation
$$\mathrm{\Delta }_{ij}(x,y)+\mathrm{d}^2zX_{ik}(x,z)ϵ^{kl}X_{jl}(z,y)=0.$$
(2.10)
The solution of Eq. (2.10) is for instance given as
$$X_{ij}(x,y)=\left(\begin{array}{cc}2& 0\\ 0& Z_\alpha ^{}Z_\alpha \end{array}\right)\delta (xy)$$
(2.11)
to yield the first class constraints with the redefinition of the two auxiliary fields $`\mathrm{\Phi }^i=(\theta ,\pi _\theta )`$
$`\stackrel{~}{\mathrm{\Omega }}_1`$ $`=`$ $`\mathrm{\Omega }_1+2\theta =Z_\alpha ^{}Z_\alpha 1+2\theta ,`$
$`\stackrel{~}{\mathrm{\Omega }}_2`$ $`=`$ $`\mathrm{\Omega }_2Z_\alpha ^{}Z_\alpha \pi _\theta =Z_\alpha ^{}\mathrm{\Pi }_\alpha ^{}+Z_\alpha \mathrm{\Pi }_\alpha Z_\alpha ^{}Z_\alpha \pi _\theta .`$ (2.12)
Here one notes that the physical fields $`Z_\alpha `$ are geometrically constrained to reside on the $`S^3`$ hypersphere with the modified norm $`Z_\alpha ^{}Z_\alpha =12\theta (x)`$.
Now, we consider the uniqueness of the first class constraints. In fact, according to the Dirac terminology , the first class constraints $`\stackrel{~}{\mathrm{\Omega }}_i`$ are defined to satisfy the following Lie algebra<sup>1</sup><sup>1</sup>1In the case of the nonvanishing $`C_{ij}^k`$, Eq. (2.10) is modified as $`\mathrm{\Delta }_{ij}(x,y)+\mathrm{d}^2zX_{ik}(x,z)ϵ^{kl}X_{jl}(z,y)=C_{ij}^k(x,y)\stackrel{~}{\mathrm{\Omega }}_k(y)`$.
$$\{\stackrel{~}{\mathrm{\Omega }}_i,\stackrel{~}{\mathrm{\Omega }}_j\}=C_{ij}^k\stackrel{~}{\mathrm{\Omega }}_k.$$
(2.13)
Since the first class constraints $`\stackrel{~}{\mathrm{\Omega }}_i`$ are ”strongly zero” $`\stackrel{~}{\mathrm{\Omega }}_i=0`$ to yield $`\{\stackrel{~}{\mathrm{\Omega }}_i(x),\stackrel{~}{\mathrm{\Omega }}_j(y)\}|\mathrm{phy}=0`$ from Eq. (2.13), one does not have any difficulties in construction of the quantum commutators and in quantization of the given physical system. In that sense, one has degrees of freedom in taking a set of the first class constraints, without any criterion. For instance our set of the first class constraints (2.12) is a specific choice to satisfy the minimal Lie algebra (2.8) with $`C_{ij}^k=0`$. Moreover even in this minimal case, we have an equivalent family of the first class constraints governed by SO(2) rotation group, under which the matrix $`X_{ij}`$ transforms as
$$XX^{}=RXR^T$$
(2.14)
where $`R`$ is an orthogonal $`2\times 2`$ matrix satisfying the condition $`RR^T=1`$. Since the matrix $`ϵ`$ is invariant under the SO(2) rotation, namely, $`RϵR^T=ϵ`$, one can easily check that the above rotated $`X^{}`$ also satisfy the relation (2.10) to yield the equivalent family of first class constraints $`\stackrel{~}{\mathrm{\Omega }}_i^{}`$ given by inserting the $`X^{}`$ into Eq. (2.9). For the more important case of the uniqueness of the first class Hamiltonian, we will discuss later in details.
Next, we construct the first class BFT physical fields $`\stackrel{~}{}=(\stackrel{~}{Z}_\alpha ,\stackrel{~}{\mathrm{\Pi }}_\alpha )`$ corresponding to the original fields $`=(Z_\alpha ,\mathrm{\Pi }_\alpha )`$. The $`\stackrel{~}{}`$’s, which reside in the extended phase space, are obtained as a power series in the auxiliary fields $`(\theta ,\pi _\theta )`$ by demanding that they are strongly involutive: $`\{\stackrel{~}{\mathrm{\Omega }}_i,\stackrel{~}{}\}=0`$. After some algebra, we obtain the first class physical fields as
$`\stackrel{~}{Z}_\alpha `$ $`=`$ $`Z_\alpha \left({\displaystyle \frac{Z_\beta ^{}Z_\beta +2\theta }{Z_\beta ^{}Z_\beta }}\right)^{1/2},`$
$`\stackrel{~}{\mathrm{\Pi }}_\alpha `$ $`=`$ $`\left(\mathrm{\Pi }_\alpha {\displaystyle \frac{1}{2}}Z_\alpha ^{}\pi _\theta \right)\left({\displaystyle \frac{Z_\beta ^{}Z_\beta }{Z_\beta ^{}Z_\beta +2\theta }}\right)^{1/2}`$ (2.15)
As discussed in Ref. , any functional $`𝒦(\stackrel{~}{})`$ of the first class fields $`\stackrel{~}{}`$ is also first class, namely, $`\stackrel{~}{𝒦}(;\mathrm{\Phi })=𝒦(\stackrel{~}{})`$. Using the property, we construct a first-class Hamiltonian in terms of the above BFT physical variables. The result is
$$\stackrel{~}{H}=\mathrm{d}^2x\left[\stackrel{~}{\mathrm{\Pi }}_\alpha ^{}\stackrel{~}{\mathrm{\Pi }}_\alpha +(_i\stackrel{~}{Z}_\alpha ^{})(_i\stackrel{~}{Z}_\alpha )(\stackrel{~}{Z}_\alpha ^{}_i\stackrel{~}{Z}_\alpha )(\stackrel{~}{Z}_\beta _i\stackrel{~}{Z}_\beta ^{})\right].$$
(2.16)
We then directly rewrite this Hamiltonian in terms of the original as well as auxiliary fields<sup>2</sup><sup>2</sup>2In deriving the first class Hamiltonian $`\stackrel{~}{H}`$ of Eq. (LABEL:hct), we have used the conformal map condition, $`Z_\alpha ^{}_iZ_\alpha +Z_\alpha _iZ_\alpha ^{}=0`$. to obtain
$`\stackrel{~}{H}`$ $`=`$ $`{\displaystyle }\mathrm{d}^2x[(\mathrm{\Pi }_\alpha ^{}{\displaystyle \frac{1}{2}}Z_\alpha \pi _\theta )(\mathrm{\Pi }_\alpha {\displaystyle \frac{1}{2}}Z_\alpha ^{}\pi _\theta ){\displaystyle \frac{Z_\beta ^{}Z_\beta }{Z_\beta ^{}Z_\beta +2\theta }}`$
$`+(_iZ_\alpha ^{})(_iZ_\alpha ){\displaystyle \frac{Z_\beta ^{}Z_\beta +2\theta }{Z_\beta ^{}Z_\beta }}(Z_\alpha ^{}_iZ_\alpha )(Z_\beta _iZ_\beta ^{})\left({\displaystyle \frac{Z_\gamma ^{}Z_\gamma +2\theta }{Z_\gamma ^{}Z_\gamma }}\right)^2].`$
We observe that the forms of the first two terms in this Hamiltonian are exactly the same as those of the O(3) NLSM .
Here $`\stackrel{~}{H}`$ is strongly involutive with the first class constraints $`\{\stackrel{~}{\mathrm{\Omega }}_i,\stackrel{~}{H}\}=0`$. A problem with $`\stackrel{~}{H}`$ in Eq. (LABEL:hct) is that it does not naturally generate the first-class Gauss law constraint from the time evolution of the constraint $`\stackrel{~}{\mathrm{\Omega }}_1`$. Therefore, by introducing an additional term proportional to the first class constraints $`\stackrel{~}{\mathrm{\Omega }}_2`$ into $`\stackrel{~}{H}`$, we obtain an equivalent first class Hamiltonian
$$\stackrel{~}{H}^{}=\stackrel{~}{H}+\frac{1}{2}\mathrm{d}^2x\pi _\theta \stackrel{~}{\mathrm{\Omega }}_2$$
(2.18)
which naturally generates the Gauss law constraint
$$\{\stackrel{~}{\mathrm{\Omega }}_1,\stackrel{~}{H}^{}\}=\stackrel{~}{\mathrm{\Omega }}_2,\{\stackrel{~}{\mathrm{\Omega }}_2,\stackrel{~}{H}^{}\}=0.$$
(2.19)
One notes here that $`\stackrel{~}{H}`$ and $`\stackrel{~}{H}^{}`$ act in the same way on physical states, which are annihilated by the first-class constraints. Similarly, the equations of motion for observables remain unaffected by the additional term in $`\stackrel{~}{H}^{}`$. Furthermore, in the limit $`(\theta ,\pi _\theta )0`$, our first class system is exactly reduced to the original second class one.
Now it is appropriate to comment on the uniqueness in construction of the first class Hamiltonians. Similar to the above discussions on the uniqueness of the first class constraints, one can have the degrees of freedom in construction of the first class Hamiltonian $`\stackrel{~}{H}`$ or $`\stackrel{~}{H}^{}`$ where $`\stackrel{~}{H}^{}`$ is equivalent to $`\stackrel{~}{H}`$ up to the additional term $`\stackrel{~}{\mathrm{\Omega }}_2`$ which does not affect the vacuum structure as discussed above. However, imposing the condition that in the limit $`(\theta ,\pi _\theta )0`$ the first class system is exactly reduced to the original second class one, one can exploit the degrees of the freedom to uniquely fix the specific form of first class Hamiltonian $`\stackrel{~}{H}^{}`$ in Eq. (2.18) which fulfills the Gauss law constraint.
Next, we consider the Poisson brackets of the fields in the extended phase space $`\stackrel{~}{}`$ and identify the Dirac brackets by taking the vanishing limit of auxiliary fields. After some algebraic manipulation starting from Eq. (2.15), one can obtain the commutators
$`\{\stackrel{~}{Z}_\alpha (x),\stackrel{~}{Z}_\beta (y)\}`$ $`=`$ $`\{\stackrel{~}{Z}_\alpha ^{}(x),\stackrel{~}{Z}_\beta (y)\}=0,`$
$`\{\stackrel{~}{Z}_\alpha (x),\stackrel{~}{\mathrm{\Pi }}_\beta (y)\}`$ $`=`$ $`(\delta _{\alpha \beta }{\displaystyle \frac{\stackrel{~}{Z}_\alpha \stackrel{~}{Z}_\beta ^{}}{2\stackrel{~}{Z}_\gamma ^{}\stackrel{~}{Z}_\gamma }})\delta (xy),`$
$`\{\stackrel{~}{Z}_\alpha (x),\stackrel{~}{\mathrm{\Pi }}_\beta ^{}(y)\}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{Z}_\alpha \stackrel{~}{Z}_\beta }{2\stackrel{~}{Z}_\gamma ^{}\stackrel{~}{Z}_\gamma }}\delta (xy),`$
$`\{\stackrel{~}{\mathrm{\Pi }}_\alpha (x),\stackrel{~}{\mathrm{\Pi }}_\beta (y)\}`$ $`=`$ $`{\displaystyle \frac{1}{2\stackrel{~}{Z}_\gamma ^{}\stackrel{~}{Z}_\gamma }}(\stackrel{~}{\mathrm{\Pi }}_\alpha \stackrel{~}{Z}_\beta ^{}\stackrel{~}{Z}_\alpha ^{}\stackrel{~}{\mathrm{\Pi }}_\beta )\delta (xy),`$
$`\{\stackrel{~}{\mathrm{\Pi }}_\alpha (x),\stackrel{~}{\mathrm{\Pi }}_\beta ^{}(y)\}`$ $`=`$ $`{\displaystyle \frac{1}{2\stackrel{~}{Z}_\gamma ^{}\stackrel{~}{Z}_\gamma }}(\stackrel{~}{\mathrm{\Pi }}_\alpha \stackrel{~}{Z}_\beta \stackrel{~}{Z}_\alpha ^{}\stackrel{~}{\mathrm{\Pi }}_\beta ^{})\delta (xy).`$ (2.20)
In the vanishing auxiliary field limit, the above Poisson brackets in the extended phase space exactly reproduce the corresponding Dirac brackets in the previous works
$`\{\stackrel{~}{Z}_\alpha (x),\stackrel{~}{Z}_\beta (y)\}_{(\theta ,\pi _\theta )=0}`$ $`=`$ $`\{Z_\alpha (x),Z_\beta (y)\}_D,`$
$`\{\stackrel{~}{Z}_\alpha ^{}(x),\stackrel{~}{Z}_\beta (y)\}_{(\theta ,\pi _\theta )=0}`$ $`=`$ $`\{Z_\alpha ^{}(x),Z_\beta (y)\}_D,`$
$`\{\stackrel{~}{Z}_\alpha (x),\stackrel{~}{\mathrm{\Pi }}_\beta (y)\}_{(\theta ,\pi _\theta )=0}`$ $`=`$ $`\{Z_\alpha (x),\mathrm{\Pi }_\beta (y)\}_D,`$
$`\{\stackrel{~}{Z}_\alpha (x),\stackrel{~}{\mathrm{\Pi }}_\beta ^{}(y)\}_{(\theta ,\pi _\theta )=0}`$ $`=`$ $`\{Z_\alpha (x),\mathrm{\Pi }_\beta ^{}(y)\}_D,`$
$`\{\stackrel{~}{\mathrm{\Pi }}_\alpha (x),\stackrel{~}{\mathrm{\Pi }}_\beta (y)\}_{(\theta ,\pi _\theta )=0}`$ $`=`$ $`\{\mathrm{\Pi }_\alpha (x),\mathrm{\Pi }_\beta (y)\}_D,`$
$`\{\stackrel{~}{\mathrm{\Pi }}_\alpha (x),\stackrel{~}{\mathrm{\Pi }}_\beta ^{}(y)\}_{(\theta ,\pi _\theta )=0}`$ $`=`$ $`\{\mathrm{\Pi }_\alpha (x),\mathrm{\Pi }_\beta ^{}(y)\}_D`$ (2.21)
where
$$\{A(x),B(y)\}_D=\{A(x),B(y)\}d^2zd^2z^{}\{A(x),\mathrm{\Omega }_k(z)\}\mathrm{\Delta }^{kk^{}}\{\mathrm{\Omega }_k^{}(z^{}),B(y)\}$$
(2.22)
with $`\mathrm{\Delta }^{kk^{}}`$ being the inverse of $`\mathrm{\Delta }_{kk^{}}`$ in Eq. (2.6). It is also noteworthy that the Poisson brackets of $`\stackrel{~}{}`$’s in Eq. (2.20) have exactly the same form as those of the Dirac brackets of the field $``$. In other words, the functional $`\stackrel{~}{𝒦}`$ in $`\stackrel{~}{𝒦}(;\mathrm{\Phi })=𝒦(\stackrel{~}{})`$ corresponds to the Dirac brackets $`\{A,B\}|_D`$ and hence $`\stackrel{~}{𝒦}`$ corresponding to $`\{\stackrel{~}{A},\stackrel{~}{B}\}`$ becomes
$$\{\stackrel{~}{A},\stackrel{~}{B}\}=\{A,B\}_D|_{A\stackrel{~}{A},B\stackrel{~}{B}}.$$
(2.23)
This kind of situation happens again when one considers the first-class constraints (2.12). More precisely, these first-class constraints in the extended phase space can be rewritten as
$`\stackrel{~}{\mathrm{\Omega }}_1`$ $`=`$ $`\stackrel{~}{Z}_\alpha ^{}\stackrel{~}{Z}_\alpha 1,`$
$`\stackrel{~}{\mathrm{\Omega }}_2`$ $`=`$ $`\stackrel{~}{Z}_\alpha ^{}\stackrel{~}{\mathrm{\Pi }}_\alpha ^{}+\stackrel{~}{Z}_\alpha \stackrel{~}{\mathrm{\Pi }}_\alpha ,`$ (2.24)
which are form-invariant with respect to the second-class constraints (2.2) and (2.5).
## 3 BRST symmetries
In this section we introduce two canonical sets of ghosts and anti-ghosts together with auxiliary fields in the framework of the BFV formalism , which is applicable to theories with the first-class constraints:
$$(𝒞^i,\overline{𝒫}_i),(𝒫^i,\overline{𝒞}_i),(N^i,B_i),(i=1,2)$$
which satisfy the super-Poisson algebra
$$\{𝒞^i(x),\overline{𝒫}_j(y)\}=\{𝒫^i(x),\overline{𝒞}_j(y)\}=\{N^i(x),B_j(y)\}=\delta _j^i\delta (xy).$$
Here the super-Poisson bracket is defined as
$$\{A,B\}=\frac{\delta A}{\delta q}|_r\frac{\delta B}{\delta p}|_l(1)^{\eta _A\eta _B}\frac{\delta B}{\delta q}|_r\frac{\delta A}{\delta p}|_l,$$
where $`\eta _A`$ denotes the number of fermions, called the ghost number, in $`A`$ and the subscript $`r`$ and $`l`$ denote right and left derivatives, respectively.
In the $`CP^1`$ model, the nilpotent BRST charge $`Q`$ and the BRST invariant minimal Hamiltonian $`H_m`$ are given by
$`Q`$ $`=`$ $`{\displaystyle \mathrm{d}^2x(𝒞^i\stackrel{~}{\mathrm{\Omega }}_i+𝒫^iB_i)},`$
$`H_m`$ $`=`$ $`\stackrel{~}{H}^{}{\displaystyle \mathrm{d}^2x𝒞^1\overline{𝒫}_2},`$ (3.1)
which satisfy the relations
$$\{Q,H_m\}=0,Q^2=\{Q,Q\}=0.$$
(3.2)
Our next task is to fix the gauge, which is crucial to identify the BFT auxiliary field $`\theta `$ with the Stueckelberg field. The desired identification follows if one chooses the fermionic gauge fixing function $`\mathrm{\Psi }`$ as
$$\mathrm{\Psi }=\mathrm{d}^2x(\overline{𝒞}_i\chi ^i+\overline{𝒫}_iN^i),$$
(3.3)
with the unitary gauge
$$\chi ^1=\mathrm{\Omega }_1,\chi ^2=\mathrm{\Omega }_2.$$
(3.4)
Here note that the $`\mathrm{\Psi }`$ satisfies the following identity
$$\{\{\mathrm{\Psi },Q\},Q\}=0.$$
(3.5)
The effective quantum Lagrangian is then described as
$$L_{eff}=\mathrm{d}^2x(\mathrm{\Pi }_\alpha ^{}\dot{Z}_\alpha ^{}+\mathrm{\Pi }_\alpha \dot{Z}_\alpha +\pi _\theta \dot{\theta }+B_2\dot{N}^2+\overline{𝒫}_i\dot{𝒞}^i+\overline{𝒞}_2\dot{𝒫}^2)H_{tot}$$
(3.6)
where $`H_{tot}=H_m\{Q,\mathrm{\Psi }\}`$ and the terms $`\mathrm{d}^2x(B_1\dot{N}^1+\overline{𝒞}_1\dot{𝒫}^1)=\{Q,\mathrm{d}^2x\overline{𝒞}_1\dot{N}^1\}`$ have been suppressed by replacing $`\chi ^1`$ with $`\chi ^1+\dot{N}^1`$.
Now we perform path integration over the fields $`B_1`$, $`N^1`$, $`\overline{𝒞}_1`$, $`𝒫^1`$, $`\overline{𝒫}_1`$ and $`𝒞^1`$, by using the equations of motion. This leads to the effective Lagrangian of the form
$`L_{eff}`$ $`=`$ $`{\displaystyle }\mathrm{d}^2x[\mathrm{\Pi }_\alpha ^{}\dot{Z}_\alpha ^{}+\mathrm{\Pi }_\alpha \dot{Z}_\alpha +\pi _\theta \dot{\theta }+B\dot{N}+\overline{𝒫}\dot{𝒞}+\overline{𝒞}\dot{𝒫}`$ (3.7)
$`(\mathrm{\Pi }_\alpha ^{}{\displaystyle \frac{1}{2}}Z_\alpha \pi _\theta )(\mathrm{\Pi }_\alpha {\displaystyle \frac{1}{2}}Z_\alpha ^{}\pi _\theta ){\displaystyle \frac{Z_\gamma ^{}Z_\gamma }{Z_\gamma ^{}Z_\gamma +2\theta }}`$
$`(_iZ_\alpha ^{})(_iZ_\alpha ){\displaystyle \frac{Z_\gamma ^{}Z_\gamma +2\theta }{Z_\gamma ^{}Z_\gamma }}+(Z_\alpha ^{}_iZ_\alpha )(Z_\beta _iZ_\beta ^{})\left({\displaystyle \frac{Z_\gamma ^{}Z_\gamma +2\theta }{Z_\gamma ^{}Z_\gamma }}\right)^2`$
$`{\displaystyle \frac{1}{2}}\pi _\theta \stackrel{~}{\mathrm{\Omega }}_2+2Z_\alpha ^{}Z_\alpha \pi _\theta \overline{𝒞}𝒞+\stackrel{~}{\mathrm{\Omega }}_2N+B\mathrm{\Omega }_2+\overline{𝒫}𝒫]`$
with the redefinitions: $`NN^2`$, $`BB_2`$, $`\overline{𝒞}\overline{𝒞}_2`$, $`𝒞𝒞^2`$, $`\overline{𝒫}\overline{𝒫}_2`$, $`𝒫𝒫_2`$.
After performing the routine variation procedure and identifying $`N=B+\frac{\dot{\theta }}{(12\theta )}`$ we arrive at the effective Lagrangian of the form
$`L_{eff}`$ $`=`$ $`{\displaystyle }\mathrm{d}^2x[{\displaystyle \frac{1}{(12\theta )}}(_\mu Z_\alpha ^{})(^\mu Z_\alpha )(12\theta )^2(B+2\overline{𝒞}𝒞)^2`$
$`{\displaystyle \frac{1}{(12\theta )^2}}(Z_\alpha ^{}_\mu Z_\alpha )(Z_\beta ^\mu Z_\beta ^{}){\displaystyle \frac{1}{12\theta }}_\mu \theta ^\mu B+_\mu \overline{𝒞}^\mu 𝒞]`$
which is invariant under the BRST-transformation
$`\delta _BZ_\alpha `$ $`=`$ $`\lambda Z_\alpha 𝒞,\delta _B\theta =\lambda (12\theta )𝒞,`$
$`\delta _B\overline{𝒞}`$ $`=`$ $`\lambda B,\delta _B𝒞=\delta _BB=0.`$ (3.9)
## 4 Collective coordinate quantization
In this section, we perform a semi-classical quantization of the unit topological charge $`Q=1`$ sector of the $`CP^1`$ model by exploiting the collective coordinates to consider physical aspects of the theory.
As a first approximation to the quantum ground state we could quantize zero modes responsible for classical degeneracy by introducing collective coordinates as follows
$`Z_1`$ $`=`$ $`e^{i(\alpha +\varphi )/2}\mathrm{cos}{\displaystyle \frac{F(r)}{2}},`$
$`Z_2`$ $`=`$ $`e^{i(\alpha +\varphi )/2}\mathrm{sin}{\displaystyle \frac{F(r)}{2}},`$
where $`(r,\varphi )`$ are the polar coordinates and $`\alpha (t)`$ is the collective coordinates. Here, in order to ensure the case of $`Q=1`$, we have used the fact the profile function $`F(r)`$ satisfies the boundary conditions: $`lim_r\mathrm{}F(r)=\pi `$ and $`F(0)=0`$.
Using the above soliton configuration, we obtain the unconstrained Lagrangian of the form
$$L=E+\frac{1}{2}\dot{\alpha }^2,$$
(4.2)
where the soliton static mass and the moment of inertia are given by
$`E`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}{\displaystyle _0^{\mathrm{}}}drr\left[\left({\displaystyle \frac{\mathrm{d}F}{\mathrm{d}r}}\right)^2+{\displaystyle \frac{\mathrm{sin}^2F}{r^2}}\right],`$
$``$ $`=`$ $`\pi {\displaystyle _0^{\mathrm{}}}drr\mathrm{sin}^2F.`$ (4.3)
Introducing the canonical momentum conjugate to the collective coordinate $`\alpha `$
$$p_\alpha =\dot{\alpha },$$
(4.4)
we then have the canonical Hamiltonian
$$H=E+\frac{1}{2}p_\alpha ^2.$$
(4.5)
At this stage, one can associate the Hamiltonian (4.5) with the previous one (2.3), which was given by the canonical momenta $`\pi ^a`$. Given the soliton configuration (LABEL:conf) one can obtain the relation between $`\pi ^a`$ and $`p_\alpha `$ as follows
$$\mathrm{\Pi }_\alpha ^{}\mathrm{\Pi }_\alpha =\frac{\mathrm{sin}^2F}{4^2}p_\alpha ^2$$
(4.6)
to yield the integral
$$\mathrm{d}^2x\mathrm{\Pi }_\alpha ^{}\mathrm{\Pi }_\alpha =\frac{1}{2}p_\alpha ^2.$$
(4.7)
Since the spatial derivative term in (2.3) yields nothing but the soliton energy $`E`$, one can easily see, together with the relation (4.7), that the canonical Hamiltonian (2.3) is equivalent to the other one (4.5), as expected.
Now, let us define the angular momentum operator $`J`$ as follows
$$J=\mathrm{d}^2xϵ_{ij}x^iT^{oj},$$
(4.8)
where the symmetric energy-momentum tensor is given by
$`T^{\mu \nu }`$ $`=`$ $`^\mu Z_\alpha ^{}^\nu Z_\alpha +^\mu Z_\alpha ^\nu Z_\alpha ^{}(Z_\alpha ^\mu Z_\alpha ^{})(Z_\beta ^{}^\nu Z_\beta )`$
$`(Z_\alpha ^{}^\mu Z_\alpha )(Z_\beta ^\nu Z_\beta ^{})g^{\mu \nu }(_\sigma Z_\alpha ^{})(^\sigma Z_\alpha )+g^{\mu \nu }(Z_\alpha ^{}_\sigma Z_\alpha )(Z_\beta ^\sigma Z_\beta ^{}).`$
Then, substituting the configuration (LABEL:conf) into Eq. (LABEL:tt), we obtain the angular momentum operator of the form
$$J=\dot{\alpha }=p_\alpha =i\frac{}{\alpha }$$
(4.10)
to yield the Hamiltonian of the form
$$H=E+\frac{1}{2}J^2.$$
(4.11)
Here one notes that the above Hamiltonian can be interpreted as mass spectrum of a rigid rotator in the $`CP^1`$ model.
Next, let us consider the zero modes in the extended phase space by introducing the soliton configuration
$`Z_1`$ $`=`$ $`(12\theta )^{1/2}e^{i(\alpha +\varphi )/2}\mathrm{cos}{\displaystyle \frac{F(r)}{2}},`$
$`Z_2`$ $`=`$ $`(12\theta )^{1/2}e^{i(\alpha +\varphi )/2}\mathrm{sin}{\displaystyle \frac{F(r)}{2}},`$ (4.12)
which satisfy the first class constraint $`Z_\alpha ^{}Z_\alpha =12\theta `$ of Eq. (2.12). <sup>3</sup><sup>3</sup>3Here one can easily see that the first class physical fields $`\stackrel{~}{Z}_\alpha `$ of Eq. (2.15) satisfy the corresponding first class constraint $`\stackrel{~}{Z}_\alpha ^{}\stackrel{~}{Z}_\alpha =1`$ of Eq. (2.24). In this configuration from Eqs. (2.1) and (5.11) we then obtain
$$L_{eff}=E+\frac{1}{2}\dot{\alpha }^2,$$
(4.13)
which is remarkably the Lagrangian (4.2) given in the original phase space. Consequently the quantization of zero modes in the extended phase space reproduces the same energy spectrum (4.11). This phenomenon originates from the fact that the collective coordinates $`\alpha `$ in the Lagrangian (4.2) are not affected by the constraints (2.2) and (2.12) for the complex scalar fields $`Z_\alpha `$. Here one notes that in the SU(2) Skyrmion model the collective coordinates themselves are constrained to yield the modified energy spectrum in contrast to the case of the $`CP^1`$ model.
## 5 Connection to O(3) nonlinear sigma model
In this section, we will demonstrate the equivalence of the $`CP^1`$ model and O(3) NLSM at the canonical level. In the O(3) NLSM, the dynamical physical fields $`n^a`$ are mappings from the spacetime manifold, which is assumed to be the direct product of a compact two-dimensional Riemann surface $`𝖬^2`$ and the time dimension $`R^1`$, to the two-sphere $`S^2`$, namely $`n^a:𝖬^2R^1S^2`$. On the other hand, the dynamical physical fields of the $`CP^1`$ model are $`Z_\alpha `$ which map the spacetime manifold $`𝖬^2R^1`$ into $`S^3`$, namely $`Z_\alpha :𝖬^2R^1S^3`$. Here one notes that $`S^3`$ is homeomorphic to SU(2) group manifold.
Since the $`CP^1`$ model is invariant under a local U(1) gauge symmetry, which consists of a redefinition of the phase of $`Z_\alpha `$ as in Eq. (1.1), the physical configuration space of the $`CP^1`$ model are the gauge orbits which form the coset $`S^3/S^1=S^2=CP^1`$. In order to associate the physical fields of the $`CP^1`$ model with those of the O(3) NLSM, we exploit the projection from $`S^3`$ to $`S^2`$, namely the Hopf bundle <sup>4</sup><sup>4</sup>4Here one notes that in order to eliminate all the unphysical degrees of freedom one can also supply a gauge fixing condition such as the Coulomb gauge: $`Z_\alpha ^{}_i_iZ_\alpha Z_\alpha _i_iZ_\alpha ^{}=0`$.
$$n^a=Z_\alpha ^{}\sigma ^aZ_\alpha $$
(5.1)
with the Pauli matrices $`\sigma ^a`$, so that we can see that the $`CP^1`$ model Lagrangian (2.1) is equivalent to the O(3) NLSM
$$L=\mathrm{d}^2x\left[\frac{1}{4}(_\mu n^a)(^\mu n^a)\right]$$
(5.2)
where $`n^a`$ ($`a`$=1,2,3) is a multiplet of three real scalar field with a constraint
$$\mathrm{\Omega }_1=n^an^a10.$$
(5.3)
Here note that the topological charge $`Q=Z`$ sector of the O(3) NLSM is guaranteed by the homotopy group $`\pi _2(S^2)=Z`$.
Moreover the collective coordinates (LABEL:conf) of the $`CP^1`$ model can be consistently obtained via the Hopf bundle (5.1) from those of the well known O(3) NLSM
$`n^1`$ $`=`$ $`\mathrm{cos}(\alpha (t)+\varphi )\mathrm{sin}F(r),`$
$`n^2`$ $`=`$ $`\mathrm{sin}(\alpha (t)+\varphi )\mathrm{sin}F(r),`$
$`n^3`$ $`=`$ $`\mathrm{cos}F(r),`$ (5.4)
to yield the rigid rotator energy spectrum (4.11), which is exactly the same as that of the O(3) NLSM with the same soliton static mass $`E`$ and moment of inertia $``$ defined in Eq. (4.3).
Now one can introduce the other bundle for the canonical momenta
$$\pi ^a=\frac{1}{2}(\mathrm{\Pi }_\alpha \sigma ^aZ_\alpha +Z_\alpha ^{}\sigma ^a\mathrm{\Pi }_\alpha ^{})$$
(5.5)
to reproduce the following secondary constraint from the corresponding $`CP^1`$ model one (2.5)
$$\mathrm{\Omega }_2=n^a\pi ^a0.$$
(5.6)
Exploiting the above bundles (5.1) and (5.5) one can easily show that the canonical Hamiltonian (2.3) of the $`CP^1`$ model is reduced to that of the O(3) NLSM
$$H_c=\mathrm{d}^2x\left(\pi ^a\pi ^a+\frac{1}{4}_in^a_in^a\right).$$
(5.7)
Here one notes that in contrast to the Banerjee case , where the reduced Hamiltonian has an additional term proportional to a first class constraint, we have obtained the exactly same Hamiltonian as shown in (5.7).
Similarly, introducing the bundles for the first class physical fields $`\stackrel{~}{n}^a`$ and $`\stackrel{~}{\pi }^a`$
$`\stackrel{~}{n}^a`$ $`=`$ $`\stackrel{~}{Z}_\alpha ^{}\sigma ^a\stackrel{~}{Z}_\alpha `$
$`\stackrel{~}{\pi }^a`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{~}{\mathrm{\Pi }}_\alpha \sigma ^a\stackrel{~}{Z}_\alpha +\stackrel{~}{Z}_\alpha ^{}\sigma ^a\stackrel{~}{\mathrm{\Pi }}_\alpha ^{})`$ (5.8)
one can also find the equivalence between the $`CP^1`$ model and the O(3) NLSM in the extended phase space at the classical level.
## 6 Connection to consistent Dirac quantization
Now we consider consistent connection of the quantization of $`CP^1`$ model in the improved Dirac scheme to that of the standard Dirac one, where one can obtain the quantum commutators via Eqs. (2.20) and (2.21)
$`[Z_\alpha (x),Z_\beta (y)]`$ $`=`$ $`[Z_\alpha ^{}(x),Z_\beta (y)]=0,`$
$`[Z_\alpha (x),\mathrm{\Pi }_\beta (y)]`$ $`=`$ $`i\left(\delta _{\alpha \beta }{\displaystyle \frac{Z_\alpha Z_\beta ^{}}{2Z_\gamma ^{}Z_\gamma }}\right)\delta (xy),`$
$`[Z_\alpha (x),\mathrm{\Pi }_\beta ^{}(y)]`$ $`=`$ $`{\displaystyle \frac{i}{2Z_\gamma ^{}Z_\gamma }}Z_\alpha Z_\beta \delta (xy),`$
$`[\mathrm{\Pi }_\alpha (x),\mathrm{\Pi }_\beta (y)]`$ $`=`$ $`{\displaystyle \frac{i}{2Z_\gamma ^{}Z_\gamma }}\left(\mathrm{\Pi }_\alpha Z_\beta ^{}Z_\alpha ^{}\mathrm{\Pi }_\beta \right)\delta (xy),`$
$`[Pi_\alpha (x),\mathrm{\Pi }_\beta ^{}(y)]`$ $`=`$ $`{\displaystyle \frac{i}{2Z_\gamma ^{}Z_\gamma }}\left(\mathrm{\Pi }_\alpha Z_\beta Z_\alpha ^{}\mathrm{\Pi }_\beta ^{}\right)\delta (xy),`$
where the quantum operator for the canonical momenta are given as
$`\mathrm{\Pi }_\alpha `$ $`=`$ $`i(\delta _{\alpha \beta }{\displaystyle \frac{Z_\alpha ^{}Z_\beta }{2Z_\beta ^{}Z_\beta }})_\beta `$
$`\mathrm{\Pi }_\alpha ^{}`$ $`=`$ $`i(\delta _{\alpha \beta }{\displaystyle \frac{Z_\alpha Z_\beta ^{}}{2Z_\beta ^{}Z_\beta }})_\beta ^{}`$ (6.2)
with the short hands $`_\alpha =\frac{}{Z_\alpha }`$ and $`_\alpha ^{}=\frac{}{Z_\alpha ^{}}`$.
Now we observe that without loss of generality the generalized momenta $`\mathrm{\Pi }_\alpha `$ fulfilling the structure of the commutators (LABEL:commst2) is of the form
$`\mathrm{\Pi }_\alpha ^c`$ $`=`$ $`i(\delta _{\alpha \beta }{\displaystyle \frac{Z_\alpha ^{}Z_\beta }{2Z_\beta ^{}Z_\beta }})_\beta {\displaystyle \frac{icZ_\alpha ^{}}{2Z_\beta ^{}Z_\beta }}`$
$`\mathrm{\Pi }_\alpha ^c`$ $`=`$ $`i(\delta _{\alpha \beta }{\displaystyle \frac{Z_\alpha Z_\beta ^{}}{2Z_\beta ^{}Z_\beta }})_\beta ^{}{\displaystyle \frac{icZ_\alpha }{2Z_\beta ^{}Z_\beta }}`$ (6.3)
with an arbitrary parameter $`c`$ to be fixed later.
On the other hand, the energy spectrum of the rigid rotators in the $`CP^1`$ model can be obtained in the Weyl ordering scheme where the Hamiltonian (2.3) is modified into the symmetric form
$$H_N=E+\mathrm{d}^2x\mathrm{\Pi }_\alpha ^N\mathrm{\Pi }_\alpha ^N$$
(6.4)
where
$`\mathrm{\Pi }_\alpha ^N`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left[(\delta _{\alpha \beta }{\displaystyle \frac{Z_\alpha ^{}Z_\beta }{2Z_\beta ^{}Z_\beta }})_\beta +_\beta (\delta _{\alpha \beta }{\displaystyle \frac{Z_\alpha ^{}Z_\beta }{2Z_\beta ^{}Z_\beta }})+{\displaystyle \frac{cZ_\alpha ^{}}{Z_\beta ^{}Z_\beta }}\right]`$
$`\mathrm{\Pi }_\alpha ^N`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left[(\delta _{\alpha \beta }{\displaystyle \frac{Z_\alpha Z_\beta ^{}}{2Z_\beta ^{}Z_\beta }})_\beta ^{}+_\beta ^{}(\delta _{\alpha \beta }{\displaystyle \frac{Z_\alpha Z_\beta ^{}}{2Z_\beta ^{}Z_\beta }})+{\displaystyle \frac{cZ_\alpha }{Z_\beta ^{}Z_\beta }}\right].`$ (6.5)
After some algebra, one can obtain the Weyl ordered $`\mathrm{\Pi }_\alpha ^N\mathrm{\Pi }_\alpha ^N`$ as follows
$$\mathrm{\Pi }_\alpha ^N\mathrm{\Pi }_\alpha ^N=_\alpha ^{}_\alpha +\frac{3Z_\alpha ^{}Z_\beta }{4Z_\gamma ^{}Z_\gamma }_\alpha ^{}_\beta +\frac{3}{4Z_\gamma ^{}Z_\gamma }Z_\alpha _\alpha \frac{1}{16Z_\gamma ^{}Z_\gamma }(2c1)(2c+3)$$
to yield the modified quantum energy spectrum of the rigid rotators
$$H_N=E+\frac{1}{2}J^2\mathrm{d}^2x\frac{(2c+3)(2c1)}{16}.$$
(6.6)
Now, in order for the Dirac bracket scheme to be consistent with the BFT one, the adjustable parameter $`c`$ in Eq. (6.6) should be fixed with the values
$$c=\frac{1}{2},\frac{3}{2}.$$
(6.7)
Here one notes that these values for the parameter $`c`$ relate the Dirac bracket scheme with the BFT one to yield the desired quantization in the $`CP^1`$ model so that one can achieve the unification of these two formalisms.
## 7 Conclusions and Discussions
In summary, we have constructed first-class BFT physical fields and, in terms of them we have obtained a first-class Hamiltonian, consistent with the Hamiltonian with the original fields and auxiliary fields. The Poisson brackets of the BFT physical fields are also constructed and these Poisson brackets are shown to reproduce the corresponding Dirac brackets in the limit of vanishing auxiliary fields. Subsequently, we have obtained, in the Batalin, Fradkin and Vilkovisky (BFV) scheme, a BRST-invariant gauge fixed Lagrangian including the (anti)ghost fields, and BRST transformation rules under which the effective Lagrangian is invariant. On the other hand, by performing the semiclassical quantization with the collective coordinates, we have obtained the spectrum of rigid rotator, which is shown to be consistent with that obtained by the standard Dirac method with the introduction of generalized momenta. Next, using the Hopf bundle , we have shown that the $`CP^1`$ model is exactly equivalent to the O(3) NLSM. Through further investigation it will be interesting to include the Chern-Simons or Hopf term in the $`CP^1`$ model since there are still subtle ambiguities in the literatures .
Now, it is appropriate to comment on the dynamical aspects of the $`CP^{N1}`$ model which possesses a local U(1) gauge invariance. As discussed in Eq. (5.1), one has traded in two degrees of freedom (three components of $`n^a`$ minus one constraint (5.3)) for the three degrees of freedom (four components of $`Z`$ minus one constraint (2.2)). In fact, the $`Z_\alpha `$ field possesses only two degrees of freedom since an overall phase transformation (1.1) does not change $`n^a`$ and hence the Lagrangian (2.1). In order to check that the Lagrangian (2.1) is invariant under the U(1) local gauge transformation, one can rewrite the Lagrangian as
$$L=\mathrm{d}^2x\left[\frac{2}{g^2}(D_\mu Z_\alpha )^{}(D^\mu Z_\alpha )+\lambda (Z_\alpha ^{}Z_\alpha 1)\right]$$
(7.8)
with the covariant derivative $`D_\mu =_\mu iA_\mu `$ and the auxiliary gauge field $`A_\mu =iZ_\alpha ^{}_\mu Z_\alpha `$. Here we have explicitly included the primary constraint $`\mathrm{\Omega }_1`$ in Eq. (5.3) together with the Lagrangian multiplier field $`\lambda `$, and the coupling constant $`g^2`$ . Under the U(1) local gauge transformation (1.1), one can have
$$A_\mu A_\mu +_\mu \alpha $$
(7.9)
under which the Lagrangian (7.8) is invariant. At the classical level, the gauge field $`A_\mu `$ associated with the U(1) symmetry is a redundant one which can be eliminated by using the equations of motion .
On the other hand, it was shown that in a stationary phase approximation the expectation value $`\lambda `$ does not vanish for $`g^2`$ large enough . In this case, one can replace the field $`\lambda (x)`$ in the Lagrangian (7.8) with a constant $`\lambda =\lambda `$ to yield the effective Lagrangian of the disordered phase where $`Z_\alpha `$ is an effective field no longer subjected to a constraint on its magnitude . Thus, at the quantum level, the $`Z_\alpha `$ field acquires a mass and the spin-spin correlation function becomes short ranged. Here we have effectively traded in the constraint for a mass. On the other hand, the gauge field $`A_\mu `$ may also acquire the kinetic term to become dynamical at the quantum level . Moreover, if the coupling constant $`g^2`$ becomes larger than some critical values, the symmetric phase appears with massless vector boson pole. Since the O(3) NLSM has no local U(1) symmetry and may become singular due to a divergence in a composite vector boson channel, the above equivalence between the $`CP^1`$ model and O(3) NLSM breaks down at the quantum level where one needs to take into account properly the dynamical gauge boson effects.
In the BFT scheme, at the quantum level, a similar situation happens to yield the quantum effects and the corresponding breakdown of the equivalence between the $`CP^1`$ model and the O(3) NLSM. Moreover, through further investigation, it will be interesting to study a new term $`\lambda \theta `$ in $`\lambda \stackrel{~}{\mathrm{\Omega }}_1`$ associated with the first class constraint $`\stackrel{~}{\mathrm{\Omega }}_1`$ in Eq. (2.12), which may play a role in quantum level phenomenology.
STH would like to thank the University of South Carolina for the warm hospitality during his visit. STH and YJP would like to thank B.H. Lee for helpful discussions and acknowledge financial support in part from the Korean Ministry of Education, BK21 Project No. D-1099 and Grant No. 2000-2-11100-002-5 from the Basic Research Program of the Korea Science and Engineering Foundation. The work of KK and FM are supported in part by NSF (USA), Grant No. 9900756 and No. INT-9730847.
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# 1 Introduction
## 1 Introduction
There is considerable interest in understanding hadronic decays involving $`\eta `$ and $`\eta ^{}`$ in the final state. The phenomenological study of hadronic processes involving flavour singlet pseudoscalar mesons makes assumptions about their composition. Here we address the issue of the nature of the $`\eta `$ and $`\eta ^{}`$ from QCD directly, making use of lattice techniques.
Lattice QCD directly provides a bridge between the underlying quark description and the non-perturbative hadrons observed in experiment. The amplitudes to create a given meson from the vacuum with a particular operator made from quark fields are measurable, an example being the determination of $`f_\pi `$. It also allows a quantitative study of the disconnected quark contributions that arise in the flavour singlet sector. The lattice approach provides other information such as that obtained by varying the number of quark flavours and their masses.
In the case of pseudoscalar mesons, the chiral perturbation theory approach also provides links between a quark description and the hadronic states. For the pion, this has the well known consequence that the decay constant $`f_\pi `$ describes quantitatively both the $`\mu \nu `$ and $`\gamma \gamma `$ decays. For the flavour singlet states ($`\eta `$ and $`\eta ^{}`$), the situation is more complicated . The axial anomaly now involves a gluonic component and the definition of decay constants is not straightforward. From the chiral perturbation theory description, one expects the mixing of $`\eta `$ and $`\eta ^{}`$ to be most simply described in a quark model basis. In the flavour singlet sector, for pseudoscalar mesons, we then have contributions to the mass squared matrix with quark model content $`(u\overline{u}+d\overline{d})/\sqrt{2}`$ and $`s\overline{s}`$ (which we label as $`nn`$ and $`ss`$ respectively):
$$\left(\begin{array}{cc}m_{nn}^2+2x_{nn}& \sqrt{2}x_{ns}\\ \sqrt{2}x_{ns}& m_{ss}^2+x_{ss}\end{array}\right)$$
(1)
Here $`m`$ corresponds to the mass of the flavour non-singlet eigenstate and is the contribution to the mass coming from connected fermion diagrams while $`x`$ corresponds to the contribution from disconnected fermion diagrams. In the limit of no mixing (all $`x=0`$, the OZI suppressed case), then we have the quenched QCD result that the $`\eta `$ is degenerate with the $`\pi `$ meson and the $`\eta ^{}`$ would correspond to the $`s\overline{s}`$ pseudoscalar meson. This is not the case, of course, and the mixing contributions $`x`$ are important.
Using as input $`m_{nn}`$, $`m_{ss}`$, $`m_\eta `$ and $`m_\eta ^{}`$, the three mixing parameters $`x`$ cannot be fully determined. It is usual to express the resulting one parameter freedom in terms of a mixing angle, here defined by
$$\eta =\eta _{nn}\mathrm{cos}\varphi \eta _{ss}\mathrm{sin}\varphi \eta ^{}=\eta _{nn}\mathrm{sin}\varphi +\eta _{ss}\mathrm{cos}\varphi $$
(2)
We show the resulting values of the mixing parameters $`x`$ in the figure (the input value for $`m_{ss}`$ will be discussed later).
The $`\eta `$ and $`\eta ^{}`$ mesons are often described in an SU(3) motivated quark basis, namely $`\eta _8=(u\overline{u}+d\overline{d}2s\overline{s})/\sqrt{6}`$, $`\eta _1=(u\overline{u}+d\overline{d}+s\overline{s})/\sqrt{3}`$. The mixing angle $`\theta `$ in this basis would be given by $`\varphi 54.7^0`$ in a lowest order chiral perturbation theory. In order to have $`f_Kf_\pi `$, one needs higher order terms in the chiral perturbation theory treatment and then the mixing scheme becomes more complicated in this basis with more than one angle needed.
In the SU(3) symmetric limit, $`m_{nn}=m_{ss}=m`$ and $`x_{nn}=x_{ns}=x_{ss}=x`$, so that only one mixing parameter is relevant and the mixing matrix simplifies considerably to a diagonal form with elements $`m^2`$ (octet) and $`m^2+3x`$ (singlet). Previous lattice studies have used degenerate quarks, so have explored this case and have found that the mixing parameter $`x`$ is of a magnitude which can explain qualitatively the observed splitting between the $`\eta `$ and $`\eta ^{}`$ mesons.
Here we undertake a non-perturbative study in QCD from first principles which will be able to establish the values of the mixing parameters $`x`$, including the pattern of SU(3) breaking. This more comprehensive study would take into account the different masses of the light ($`u`$ and $`d`$) quarks and the heavier $`s`$ quark. Within the lattice approach, it is not at present feasible to evaluate using quarks as light as the nearly massless $`u`$ and $`d`$ quarks and also it is more tractable to use an even number of degenerate quarks in the vacuum. As we shall show, despite these restrictions, a thorough study of the mixing between $`\eta `$ and $`\eta ^{}`$ is possible.
Our lattice study uses dynamical configurations with $`N_f=2`$ flavours of sea quarks of type 1 and we consider the properties of pseudoscalar mesons made of either quark 1 or quark 2, where quark 2 corresponds to a heavier quark. Thus quark 2 is treated as partially quenched. Here we have in mind exploring a situation which will be relevant to treating strange quark propagation (quark 2) in a vacuum containing only lighter quarks (quark 1). We focus here on the results of lattice evaluations, for background to the methods used see ref.. We address three topics where lattice input permits us to construct a firm foundation for the $`\eta `$, $`\eta ^{}`$ mixing:
* From comparing pseudoscalar meson masses with valence quarks of two different masses (namely meson masses $`m_{11}`$, $`m_{12}`$ and $`m_{22}`$), we can estimate the mass $`m_{ss}`$ of the unmixed $`\overline{s}s`$ meson, given the observed $`m_{ns}`$ and $`m_{nn}`$ masses (ie $`K`$ and $`\pi `$ respectively).
* From measuring the mixing parameters $`x_{11}`$, $`x_{12}`$ and $`x_{22}`$ between initial and final flavour singlet states consisting of either quark 1 or 2 with different masses as above, we can establish the pattern of SU(3) breaking in the mixing.
* For $`N_f=2`$ degenerate flavours of quark, we determine the pseudoscalar decay constants for the flavour singlet ($`P_0`$) and non-singlet ($`P_1`$) meson. This input allows us to discuss the relation between the observed $`\gamma \gamma `$ decay modes of $`\pi ^0`$, $`\eta `$ and $`\eta ^{}`$ and the underlying quark content.
## 2 Lattice results
### 2.1 The $`s\overline{s}`$ pseudoscalar mass
Chiral symmetry considerations lead to the expectation that the pseudoscalar meson composed of quarks of mass $`M_q`$ has mass squared $`m^2`$ which behaves linearly with $`M_q`$ at small quark mass. However, at large quark mass ($`c`$ and $`b`$ quarks for instance), one expects the meson mass to vary approximately linearly with the quark mass. Here we are not concerned with the region of very small quark mass where chiral logs are important , so we summarise this behaviour by
$$m^2=bM_q+cM_q^2+O(M_q^3)$$
(3)
For a pseudoscalar meson made of two different quarks of mass $`M_n`$ and $`M_s`$, we shall assume its mass only depends on $`(M_n+M_s)/2`$ and not on $`(M_sM_n)/2`$ as found in lattice studies and in lowest order chiral perturbation theory. If eq. 3 were valid with just the linear term in the quark mass(ie $`c=0`$), then one directly obtains the required mass of the pseudoscalar meson composed of $`s`$ quarks, $`m_{ss}^2=2m_{sn}^2m_{nn}^2`$, that is $`2K^2\pi ^2`$, leading to $`m_{ss}=0.687`$ GeV.
This can be explored on a lattice by measuring the pseudoscalar meson mass for valence quarks in combinations 11, 22 and 12. Then, for small $`c/b`$, we have
$$\frac{c}{4b^2}=\frac{\frac{1}{2}(m_{11}^2+m_{22}^2)m_{12}^2}{(m_{22}^2m_{11}^2)^2}$$
(4)
This has been studied in the quenched approximation giving evidence for a positive coefficient $`c`$ in eq. 3. In the quenched approximation, however, the chiral behaviour at small quark mass is anomalous since the theory is not unitary. A better way to study this issue on the lattice is then to use dynamical configurations with sea quarks of type 1 and to consider the propagation of mesons made of either quark 1 or quark 2, where quark 2 corresponds to a heavier quark.
We present results from UKQCD configurations with $`N_f=2`$ flavours of sea quark with SW-clover coefficient $`C_{SW}=1.76`$, lattice size $`12^324`$, and with sea quarks having $`\kappa =0.1398`$, corresponding to sea quarks of mass around the strange quark mass ($`m_P/m_V=0.67`$). Then we take the heavier valence quark (with $`\kappa =0.1380`$) as corresponding to approximately twice the strange mass ($`m_P/m_V=0.81`$).
The fits with two states to a $`4\times 4`$ matrix of mesonic correlators for $`t`$ range 3 to 10 give results for the spectrum shown in Table 1.
Taking account of the correlation among the errors, we obtain the dimensionless ratio
$$m_{11}^2c/(4b^2)=0.011(3)$$
(5)
which indicates a statistically significant curvature from the $`c`$ term. Setting the scale using $`a^1=1.47`$ GeV then the value of $`c/b^2`$ in physical units can be obtained from $`m_{11}=698`$ MeV.
Applying this value of $`c`$ to the determination of the $`m_{ss}`$ mass from the $`\pi `$ and K masses, gives a relative shift upwards due to the curvature term ($`c`$) of 1.1(3)%, corresponding to a value of $`m_{ss}=0.687+0.008`$ GeV.
Note that this value also helps us to identify the meson mass ratio corresponding to strange quarks, namely $`m_P/m_V=m_{ss}/m_\varphi =0.682`$.
This lattice study thus answers the question of the likely deviation in the pseudoscalar mass formula from the result given by the lowest order chiral expression.
### 2.2 Flavour-singlet mixing
The mass splitting between flavour non-singlet and singlet mesons can be measured using lattice evaluation of disconnected quark propagators. This is not an easy task: the contamination from excited states is difficult to remove and the statistical errors turn out to be relatively large. Initial studies have been in the quenched approximation . Here, although there is no flavour splitting of the masses, the mass splitting matrix element $`x`$ can be evaluated. It is, however, preferable to be able to study the mass splitting directly and hence here we focus on results from full QCD simulations .
The study of the mass spectrum of flavour singlet ($`P_0`$) and non-singlet pseudoscalar meson ($`P_1`$) using $`N_f=2`$ flavours of sea quark 1 leads to singlet mass $`m_0=(m_{11}^2+2x_{11})^{1/2}`$ and non-singlet mass $`m_1=m_{11}`$ respectively which allows $`x_{11}`$ to be extracted. We shall also be interested in the dependence of $`x`$ on quark masses and on non-diagonal mixings. These can be studied with a little less rigour as we discuss later.
We use the UKQCD lattices introduced in the previous section. The disconnected diagrams were evaluated using a variance reduction method which uses all the data available with no dilution from the stochastic method used. We use as many different operators for the pseudoscalar meson as possible to have the largest basis in which to extract the ground state - local and non-local (fuzzed) in space with both $`\gamma _5`$ and $`\gamma _5\gamma _4`$ spin structure. Unfortunately, even with this basis of four operators, we are unable to determine the singlet mass precisely. For example for both valence and sea quarks having $`\kappa =0.1398`$, as shown in Table 1, we obtain $`am_0=0.56(4)`$ from a one state fit to $`t=3`$ to 7 with a $`4\times 4`$ matrix of meson correlators. A two state fit to a wider $`t`$ range (2-7) gives a similar mass value. The corresponding non-singlet pseudoscalar mass is also given in Table 1, so the determination of $`x`$ from $`m_1^2+2x_{11}=m_0^2`$ has relatively large errors ($`x_{11}=0.10(4)`$ GeV<sup>2</sup> using $`a^1=1.47`$ GeV). As an alternative, we also fit the ratio of the singlet to non-singlet correlators directly to a ground state mass difference. Using the $`t`$ range 2-7 and local and fuzzed pseudoscalar operators, we obtain $`am_0am_1=0.12(7)`$. This method gives a slightly larger mass value (indicating $`x_{11}=0.13(8)`$ GeV<sup>2</sup>) but has even larger errors. These results indicate that much larger ensembles of gauge configurations will be needed to make more precise this approach of determining $`x`$ from masses.
If one studies correlations of meson operators made from valence quarks of type 2 in a sea of quarks of type 1, one will find the ground state pseudoscalar meson to be that composed of quarks of type 1 (assuming type 2 quarks are heavier than type 1). Because of this, we need to explore in more detail to study the SU(3) breaking of the mixing parameters.
To get a first look at this issue, we consider a quenched lattice and measure the ratio of the disconnected to connected diagrams for pseudoscalar meson propagation. We present results for $`\beta =5.7,C_{SW}=1.57,12^3\times 24`$ with 100 configurations with $`\kappa =0.14077`$ and 0.13843. The non-singlet spectrum at these parameters was studied previously giving $`m_V/m_P`$ values of 0.65 and 0.78 which correspond approximately to strange quarks and quarks twice as heavy as strange. The scale was set as $`a^1=1.2`$ GeV. The disconnected meson correlator was determined using a stochastic method with variance reduction .
Assuming dominance by ground state meson contributions, the ratio of disconnected to connected diagrams at time separation $`t`$ is
$$\frac{D_{ij}}{C_{ij}}=\frac{N_fx_{ij}(t+1)}{2(m_{ii}m_{jj})^{1/2}}$$
(6)
with flavour non-singlet pseudoscalar mass $`m_{ii}`$ for quarks of type $`i`$. The factor of $`t+1`$ comes from the number of lattice sites at which the disconnected diagram can be split. To clarify the pattern of SU(3) breaking, we also study the non-diagonal case where we also measure the disconnected to connected ratio (here $`C_{ij}`$ is taken as $`(C_{ii}C_{jj})^{1/2}`$). In extracting $`x_{12}`$ we can make an additional correction for the contribution from the propagation of mesons with different masses $`m_{11}`$ and $`m_{22}`$ although, in practice, this correction is very small.
Using local meson operators, we obtain for $`x`$ the values in Table 2. The largest $`t`$ value has least contributions from excited state contamination and the consistency of the results versus $`t`$ suggests that such contamination is small. As was found previously , $`x`$ increases as the quark mass is decreased. Moreover, we can check to see if there is a factorisation of $`x`$ as expected in some chiral perturbation theory descriptions , namely $`x_{12}^2=x_{11}x_{22}`$, and we find that $`x_{12}`$ lies somewhat below the value given by this assumption.
We now revert to discussing the more realistic (partially quenched) case: with heavier quarks of type 2 in a sea of two flavours of quarks of type 1. The method described above for the quenched case can be applied here too. In principle this method is now only valid for small $`N_fxt`$ for the propagation of quark 1. From this analysis of the measured $`D/C`$ values, we get the $`x`$ values shown in Table 3. The values of $`x_{11}`$ are similar to those obtained above (with larger errors) directly from the rigorous method of using the mass differences. This suggests that the strong assumptions made in determining $`x`$ directly from $`D/C`$ are actually reasonable in practice. This, and the consistency of values from different $`t`$ and different mesonic operators, gives us confidence to use the $`x`$ values from quarks of type 2 (which are partially quenched anyway) as a guide to the quark mass dependence of $`x`$. The $`x`$ values again show an increase with decreasing quark mass and also approximate factorisation.
Setting the quark mass to strange (since $`m_P/m_V=0.682`$ in nature for $`s`$ quarks) in both quenched and $`N_f=2`$ evaluations leads to a consistent lattice estimate of $`x_{ss}`$ in the range 0.09 to 0.13 GeV<sup>2</sup>. This value is also consistent with that reported from a study of $`N_f=2`$ by the CP-PACS collaboration with $`m_P/m_V=0.69`$ and $`a^1=1.29`$ GeV giving values of $`x_{ss}=0.10`$ GeV<sup>2</sup> and 0.14 GeV<sup>2</sup> (depending on using $`t_{\mathrm{min}}=2,3`$ in fits, respectively). These lattice values are obtained at quite coarse lattice spacings and there may be some additional systematic error arising from the extrapolation to the continuum limit. We have, however, chosen to use a clover improved fermion action to minimise this extrapolation error.
We are unable to determine the mixing strengths $`x`$ for lighter quarks than strange. So we assume that the value of $`x`$ continues to increase as the quark mass is decreased below strange in a similar way to the decrease we see from twice strange (type 2) to strange (type 1).
Consider now the consequence of this determination of the mixing. We use input masses $`m_{nn}=0.137`$ GeV, $`m_{ss}=0.695`$ GeV (as discussed above) and aim to have $`x`$ values in line with our results above, namely $`x_{ss}0.12`$ GeV<sup>2</sup>, $`x_{ns}^2x_{nn}x_{ss}`$ and we also expect, though with big errors from the extrapolation, $`x_{nn}/x_{ss}2`$. The figure shows the $`x`$ values needed to reproduce the known $`\eta `$ and $`\eta ^{}`$ masses for each mixing angle $`\varphi `$. The lattice determination of $`x_{ss}`$ is shown by the dotted horizontal band. Keeping close to this band while satisfying the other lattice constraints is possible for the mixing illustrated by the vertical line. This has $`x_{nn}=0.292,x_{ns}=0.218,x_{ss}=0.13`$ GeV<sup>2</sup> which gives a description of the observed $`\eta `$ and $`\eta ^{}`$ masses while being consistent with our QCD inspired evidence about the mixing strengths. This assignment corresponds to a mixing angle $`\varphi `$ in the $`\eta _{nn},\eta _{ss}`$ basis of 44.5<sup>0</sup>. Note that this is almost maximal which implies that the quark content (apart from the relative sign) of the $`\eta `$ and $`\eta ^{}`$ meson is the same. The corresponding mixing angle in the $`\eta _8`$, $`\eta _1`$ basis (modulo comments above) is a value of $`\theta `$ of $`10.2^0`$.
### 2.3 Flavour-singlet decay constants
The decays of $`\pi ^0,\eta `$ and $`\eta ^{}`$ to $`\gamma \gamma `$ are expected to proceed via the quark triangle diagram. The quark model gives a decay proportional to $`Q_i^2`$ for the contribution from a quark of charge $`Q_i`$. Thus for the $`\pi ^0`$ meson and the flavour-singlet $`nn`$ and $`ss`$ mesons, the quark charge contributions to the decay amplitudes would be in the ratio $`1:5/3:\sqrt{2}/3`$. The experimental reduced decay amplitudes for $`\pi ^0`$, $`\eta `$, and $`\eta ^{}`$ are in the ratio $`1.0:1.00(10):1.27(7)`$. This information can be used to analyse the quark content of the pseudoscalar mesons subject to a quantitative understanding of the decay mechanisms.
The conventional approach assumes that the decay constants for the decays of the three mesons are the same and then the relative decay amplitudes give information on the quark content. This suggests a mixing angle of $`\theta 20^0`$ is preferred .
We now address the issue of determining these decay constants directly from QCD using lattice methods. Our study uses 2 flavours of degenerate quark and we define the decay constants by
$$0|A^\mu |P_1(q)=f_1q^\mu 0|A^\mu |P_0(q)=f_0q^\mu $$
(7)
For the isospin 1 state $`P_1`$ ($`\pi `$ \- like), this is on a firm footing because of the anomaly hence $`f_1`$ will be scale invariant. For the flavour singlet pseudoscalar meson $`P_0`$, the decay constant defined as above will not be scale invariant because of gluonic contributions to the anomaly . In this exploratory study we determine the decay constants with lattice regularisation and we shall compare the singlet and non-singlet values.
These decay constants can be thought of as giving the quark wave function at the origin of the pseudoscalar meson. Since the mass splitting between singlet and non-singlet is not reproduced directly in quenched QCD, it is essential to use lattice studies that do include sea quark effects in this study of decay matrix elements.
Results were obtained using fits to full (connected and disconnected) meson propagation with 4 different types of meson creation and destruction operator. These are local and fuzzed operators with either $`\gamma _5`$ or $`\gamma _4\gamma _5`$ couplings, so giving $`4\times 4`$ matrix of pseudoscalar correlators. We used the $`N_F=2`$ UKQCD configurations referred to above. For the disconnected correlators, the variance reduction technique is essential to get a reasonable signal to noise ratio, particularly for the operators involving the $`\gamma _4\gamma _5`$ factor.
The lattice result for $`f`$ with various valence quark masses with fixed sea quark mass (quark 1) as above is shown in Table 1. For the non-singlet results using $`a^1=1.47`$ GeV and the tadpole-improved perturbative value of $`Z`$ of 0.81 (and of $`c_A`$ which is involved in mixing of the lattice pseudoscalar and axial currents but has a very small effect in practice) we get $`f_{11}=198(8)`$ MeV. Since this corresponds to strange quarks, it is in reasonable agreement with experiment assuming a steady increase from $`f_{nn}=131`$ MeV and $`f_{ns}=160`$ MeV to $`f_{ss}`$. We do see evidence for this increase in $`f`$ with quark mass directly on the lattice going from quarks of type 1 (strange) to type 2 (twice strange) as shown in Table 1.
The flavour singlet results are shown in Table 1. They are determined by fits to the appropriate (connected plus disconnected) meson correlators which are a $`4\times 4`$ matrix at each $`t`$ value. Despite this extensive data set, the determinations of $`f`$ have relatively large statistical errors and the systematic error from changing the type of fit is also comparable. For our case with $`N_f=2`$ degenerate quarks, the comparison of the flavour singlet and non-singlet shows that the singlet decay constants appear to be somewhat larger, though the errors are too big to substantiate this.
Combining the mass dependence we find in the flavour non-singlet sector with the near equality of singlet and non-singlet decay constants, we can deduce properties of the physical case with three light quarks. Thus, in terms of the conventional treatment , we would expect $`f_\eta /f_\pi >1`$ and $`f_\eta ^{}/f_\pi >1`$. One way to minimise the effects of mixing is to consider $`X=(a_\eta ^2+a_\eta ^{}^2)/a_\pi ^2`$ where $`a`$ refers to the reduced decay amplitude. Using the conventional formulae for the decay amplitudes would then give a value of $`X=3r^2`$ (where $`r`$ is a suitably weighted average of $`f_\eta /f_\pi `$ and $`f_\eta ^{}/f_\pi `$ which are both greater than 1). Thus the conventional treatment gives $`X>3`$ which is significantly larger than the experimental value of 2.64(24). Thus it appears unlikely that the conventional treatment (with the decay to $`\gamma \gamma `$ being given by the analogue of the formula for pions) is correct for any mixing angle.
We conclude that there is no support for the conventional assumption that the singlet decays are given by a similar expression to the non-singlet. As has been pointed out by many authors , this is plausible for at least two reasons: (i) the $`\eta `$ and $`\eta ^{}`$ mesons are heavier and therefore less likely to dominate the axial current or, equivalently, higher order corrections to chiral perturbation theory will be more important (ii) the flavour-singlet axial anomaly has a gluonic component which will give additional contributions to any hadronic process.
## 3 Conclusion
From our careful non-perturbative study of mass formulae for flavour non-singlet pseudoscalar mesons made of different quarks, we deduce that the $`ss`$ state lies at 695 MeV. We then determine the pattern of mixing for the flavour singlet sector, obtaining $`x_{ss}0.12`$ GeV<sup>2</sup>, $`x_{nn}/x_{ss}2`$ and $`x_{ns}^2x_{nn}x_{ss}`$. These conditions are indeed consistent and point to a mixing close to maximal ($`\varphi =45\pm 2^0`$) in the $`nn`$, $`ss`$ basis (this corresponds to a conventional ($`\eta _8`$, $`\eta _1`$) mixing $`\theta `$ of $`10\pm 2^0`$). We are able to explore the decay constants for singlet pseudoscalar mesons for the first time. Our results show similar decay constants for singlet and non-singlet states of the same mass but with quite large errors.
We have not addressed here the issue of the origin of these mixing parameters $`x`$. Lattice studies have the capability to relate them to topological charge density fluctuations or to other vacuum properties.
Our lattice studies have been hampered by two constraints. One is that the disconnected quark diagrams needed for a study of singlet mesons are intrinsically noisy. Much larger data sets (tens of thousands of gauge configurations) will be needed to increase precision. Another constraint is that we are unable to work with sea quarks substantially lighter than strange. We have also not attempted a continuum limit extrapolation of our lattice results. Although we are using a lattice formalism that should improve this extrapolation, it would be safer to test it directly. The lattice non-perturbative results do, however, show clearly the structure of the mixing in the singlet pseudoscalar mesons.
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# Time-dependent Gutzwiller approximation for the Hubbard model
\[
## Abstract
We develop a time-dependent Gutzwiller approximation (GA) for the Hubbard model analogous to the time-dependent Hartree-Fock (HF) method. The formalism incorporates ground state correlations of the random phase approximation (RPA) type beyond the GA. Static quantities like ground state energy and double occupancy are in excellent agreement with exact results in one dimension up to moderate coupling and in two dimensions for all couplings. We find a substantial improvement over traditional GA and HF+RPA treatments. Dynamical correlation functions can be easily computed and are also substantially better than HF+RPA ones and obey well behaved sum rules.
\]
The Gutzwiller (GW) trial wave function is probably the most popular variational approach to the Hubbard model which incorporates correlation effects beyond the Hartree-Fock (HF) approximation. Since the Hubbard model describes the competition between hopping and correlation induced localization of the charge carriers, the idea is to apply a projector to a Slater determinant (SD) which reduces the number of doubly occupied sites. The optimum double occupancy probability is determined variationally.
Similar to HF best results are obtained if one allows for unrestricted charge and spin distributions which are determined also variationally. For example for the half-filled Hubbard model a SD with long range antiferromagnetic order is favoured.
A formal diagrammatic solution of the GW variational problem has been given by Metzner and Vollhardt, however, for the most part of practical purposes one approximates the corresponding expectation values using the so-called Gutzwiller approximation (GA).
The GA can be derived using a variety of methods. In particular it is recovered at the mean-field level (saddle-point) of the four-slave boson functional integral method introduced by Kotliar and Ruckenstein (KR). The latter offers the possibility of going beyond the Gutzwiller result as for example the inclusion of transversal spin degrees of freedom . In addition it provides a scheme to include fluctuations beyond the mean-field (MF) solution. Expansions around the slave-boson saddle point have been performed for homogeneous systems in Ref. in order to calculate correlation functions in the charge and longitudinal spin channels. However, the expansion of the KR hopping factor $`z^{SB}`$ is a highly nontrivial task both with respect to the proper normal ordering of the bosons and also with respect to the correct continuum limit of the functional integral .
The complexity of the expansions around the slave boson saddle point have severely hampered practical computations of dynamical quantities within this formalism. One of the few successful attempts is the computation of the optical conductivity in the paramagnetic state in Ref. . Remarkably although the starting SD describes a paramagnetic system, spectral weight on the Hubbard bands appears as an effect of fluctuations. As far as we know this approach has not been pursued in broken symmetry states due to technical difficulties, including the fact that the KR choice for the $`z^{SB}`$ hopping factor does not lead to controlled sum rules.
In this work we introduce a simple scheme to compute fluctuation corrections around the GA to dynamical and static correlation functions and the ground state energy. The method can be viewed as a time-dependent GA in the same way as the random phase approximation (RPA) method on top of a HF solution (HF+RPA) can be viewed as time-dependent HF approximation in the limit of small amplitude oscillations . For this reason we label the method as GA+RPA. It is also a generalization of the method of Ref. in order to describe the low temperature Fermi liquid regime. The method incorporates ground state correlations beyond the ones of the Gutzwiller type just as HF+RPA takes into account ground state correlations not present in the HF wave function.
The GA+RPA ground state energy of the one-band Hubbard model is in excellent agreement with exact results up to moderate coupling in one dimension (1d) and for all couplings in a 2d system (Fig. 1). The optical conductivity of a Hubbard chain is in much better agreement with numerical results than HF+RPA (Fig. 2). In addition sum rules are well behaved in the HF+RPA sense.
We consider the one-band Hubbard model
$$H=\underset{ij,\sigma }{}t_{ij}c_{i,\sigma }^{}c_{j,\sigma }+U\underset{i}{}n_{i,}n_{i,}$$
(1)
where $`c_{i,\sigma }`$ destroys an electron with spin $`\sigma `$ at site $`i`$, and $`n_{i,\sigma }=c_{i,\sigma }^{}c_{i,\sigma }`$. $`U`$ is the on-site Hubbard repulsion and $`t_{ij}`$ denotes the transfer parameter between sites $`i`$ and $`j`$. In the numerical computations below we take only nearest neighbour matrix elements $`t_{ij}=t`$ to be non-zero.
Our starting point is an energy functional $`E[\rho ,D]`$ of the GA type. Here $`\rho `$ is the density matrix of an associated Slater determinant $`|SD>`$, i.e. $`\rho _{i\sigma ,j\sigma ^{}}=SD|c_{i\sigma }^{}c_{j\sigma ^{}}|SD`$ and $`D`$ is a vector of the GA double occupancy parameters $`D_i`$ at site $`i`$. In order to consider arbitrary fluctuations the charge and spin distribution of $`\rho `$ and the distribution of $`D`$ should be completely unrestricted. For simplicity we consider only solutions where the associated SD is an eigenstate of the z-component of the total spin operator ($`\rho _{i\sigma ,j\sigma ^{}}\delta _{\sigma ,\sigma ^{}}\rho _{ij\sigma }`$).
$`E[\rho ,D]`$ can be obtained by exploiting the equivalence between the KR saddle point solution and the GA. It is given by:
$$E[\rho ,D]=\underset{ij\sigma }{}t_{ij}z_{i\sigma }z_{j\sigma }\rho _{ij\sigma }+U\underset{i}{}D_i$$
(2)
and
$$z_{i\sigma }=\frac{\sqrt{(1\rho _{ii}+D_i)(\rho _{ii\sigma }D_i)}+\sqrt{D_i(\rho _{ii,\sigma }D_i)}}{\sqrt{\rho _{ii\sigma }(1\rho _{ii\sigma })}}.$$
(3)
with $`\rho _{ii}=_\sigma \rho _{ii\sigma }`$. The stationary solution $`\rho ^{(0)},D^{(0)}`$ is determined by minimizing the energy functional with respect to $`\rho `$ and $`D`$.
The variation with respect to the density matrix has to be constrained to the subspace of Slater determinants by imposing the projector condition $`\rho ^2=\rho `$. Within this subspace we now consider small time-dependent amplitude fluctuations of the density matrix $`\rho (t)`$. We add a weak time-dependent field to Eq. (2) of the form: $`F(t)=_{i\sigma ,j\sigma ^{}}(f_{i\sigma ,j\sigma ^{}}e^{i\omega t}c_{i\sigma }^{}c_{j\sigma ^{}}+h.c.)`$. This produces small amplitudes oscillations $`\delta \rho (t)`$ around the stationary density i.e. $`\delta \rho (t)\rho (t)\rho ^{(0)}`$.
We assume that at each instant of time the double occupancy parameter is at the minimum of the energy functional compatible with the corresponding $`\rho (t)`$; i.e. the double occupancy parameters $`\{D\}`$ adjust antiadiabatically to the time evolution of the density matrix. This is reasonable since the double occupancy involves processes which are generally high in energy and hence fast. We anticipate that for the cases explored, this approximation works well up to energies as large as the Hubbard band in the optical conductivity (Fig. 2). As the density varies the double occupancy shifts from $`D^{(0)}`$ to satisfy the antiadiabaticity constraint. We define $`\delta D(t)=D(t)D^{(0)}`$ and $`\delta \rho (t)`$ and $`\delta D(t)`$ are linear in $`f`$.
The formal complication of the present approach as compared to the standard RPA has its origin in the proper adjustment of $`D`$ to the time evolution of $`\rho (t)`$, i.e. the determination of $`\delta D(t)`$. This step is achieved by expanding the energy functional Eq. (2) up to second order in $`\delta \rho `$ and $`\delta D`$ around the saddle point:
$`E[\rho ,D]`$ $`=`$ $`E_0+\overline{h}_0^{}\delta \overline{\rho }+{\displaystyle \frac{1}{2}}\delta \overline{\rho }^{}L_0\delta \overline{\rho }`$ (4)
$`+`$ $`\delta DS_0\delta \overline{\rho }+{\displaystyle \frac{1}{2}}\delta D^tK_0\delta D`$ (5)
where the bar indicates that we are treating a matrix as a column vector and the not indicates evaluation in the stationary state. Here we have defined an effective one-particle Gutzwiller Hamiltonian:
$$h_{ji\sigma }=\frac{E}{\rho _{ij\sigma }}$$
(6)
and the matrices
$`L_{ij\sigma ,kl\sigma ^{}}`$ $`=`$ $`{\displaystyle \frac{^2E}{\rho _{ij\sigma }^{}\rho _{kl\sigma ^{}}}}`$ (7)
$`S_{k,ij\sigma }`$ $`=`$ $`{\displaystyle \frac{^2E}{D_k\rho _{kl\sigma }}}`$ (8)
$`K_{k,l}`$ $`=`$ $`{\displaystyle \frac{^2E}{D_kD_l}}.`$ (9)
Using the condition of antiadiabaticity
$$\frac{E}{\delta D}=0$$
(10)
in Eq. (4) we obtain a linear relation between $`\delta \rho `$ and $`\delta D`$. Eliminating $`\delta D`$ from Eq. (4) finally yields an expansion of the energy as a functional of $`\delta \rho `$ alone $`\stackrel{~}{E}[\rho ]E[\rho ,D(\rho )]`$,
$`\stackrel{~}{E}[\rho ]=E_0+\overline{h}_0^{}\delta \overline{\rho }+{\displaystyle \frac{1}{2}}\delta \overline{\rho }^{}(L_0S_0^{}K_0^1S_0)\delta \overline{\rho }`$ (11)
This can be regarded as the expansion of an effective interacting energy functional in which the interaction potential between particles is density dependent. This kind of functional often appears in the context of nuclear physics and a well developped machinery exist to compute the RPA fluctuations induced by the interaction. We will only briefly outline here the corresponding formalism (for details see Ref. ). The advantage of this method with respect to other methods (eg. diagrammatic) is that the present derivation is solely based on the knowledge of an energy functional of a SD density matrix which is precisely what the GA provides.
At the saddle-point $`h`$ and $`\rho `$ can be diagonalized simultaneously. As a result one obtains $`(h_0)_{kl}=\delta _{kl}ϵ_k`$ and the density matrix has eigenvalue 1 below the Fermi level and eigenvalue 0 above the Fermi level. We will notate states below the Fermi level as hole ($`h`$) states and the states above the Fermi level as particle states ($`p`$).
Up to linear order the density matrix obeys the equation of motion:
$$i\mathrm{}\delta \dot{\rho }=[\stackrel{~}{h}_0,\delta \rho ]+[\frac{\stackrel{~}{h}}{\rho }.\delta \rho ,\rho ^{(0)}]+[f^{GA},\rho ^{(0)}]$$
(12)
where $`\stackrel{~}{h}`$ is defined as in Eq. (6) but with $`\stackrel{~}{E}`$ instead of $`E`$ (Note that $`\stackrel{~}{h}_0=h_0`$). $`\frac{\stackrel{~}{h}}{\rho }.\delta \rho `$ is a short hand notation for
$$\underset{ph}{}\left(\frac{\stackrel{~}{h}}{\rho _{ph}}|_{\rho =\rho ^{(0)}}\delta \rho _{ph}+\frac{\stackrel{~}{h}}{\rho _{hp}}|_{\rho =\rho ^{(0)}}\delta \rho _{hp}\right).$$
(13)
$`f^{GA}`$ is the GA version of $`f`$, i.e. it includes the $`z_0`$ factors for intersite matrix elements.
One can show that particle-particle and hole-hole matrix elements of Eq. (12) are zero and the particle-hole ($`ph`$) matrix elements of $`\delta \rho `$ satisfy the well known RPA eigenvalue equation. The RPA dynamical matrix can be obtained from Eqs. (11),(12). Upon diagonalizing the RPA matrix by a Bogoliubov transformation one obtains the eigenvectors $`V^{(\lambda )}=(X_{ph}^{(\lambda )},Y_{ph}^{(\lambda )})`$ and eigenvalues $`E^{(\lambda )}`$ where the latter correspond to the excitation energies of the system. An explicit expression for the response functions and a discussion of sum rules can be found in Ref. which apply straightforwardly to our case.
The present formalism is well suited for the calculation of charge excitations in inhomogeneous doped systems which will be presented elsewhere. In the following we restrict ourselves to the half-filled Hubbard model in the antiferromagnetic Néel state. The double occupancy at the RPA level is given by: $`D_{RPA}=𝑑\omega _\lambda 0|n_i|\lambda \lambda |n_i|0\delta (\omega E^{(\lambda )})`$ where the integrand is the Lehmann representation of an appropriately defined density-density correlation function. The matrix elements $`0|n_{}|\lambda `$ for $`\lambda >0`$ can be computed in terms of the eigenvectors $`V^{(\lambda )}`$.
In the inset of Fig. 1, we show the GA+RPA double occupancy compared with exact results and other approximations in a 1d system. For small $`U`$ long range magnetic order is not enough to reduce substantially the HF double occupancy from the noninteracting value. RPA on top of HF corrects for this but because the starting point is quite far from the exact result the correction is not so accurate and one gets that HF+RPA overstrikes the exact double occupancy. On the contrary for the GA only a small correction is needed and RPA performs remarkably well. Note that $`U(D_{RPA}D_{HF})`$ is a measure for the residual interaction in HF+RPA. In the GA+RPA approach such a simple relation is lost but clearly a smaller correction of the MF double occupancy suggests a smaller residual interaction.
From the interaction energy $`UD_{RPA}`$ we compute the correction to the ground state energy using the coupling constant integration trick. We find very good agreement with the exact results as shown in Fig. 1. This holds in 1d up to intermediate values of $`U/t`$ and in a 4x4 2d cluster for all $`U/t`$. The improvement with dimensionality is expected as in any MF + RPA computation.
In order to examine the quality of dynamical correlation functions we have studied the optical conductivity in the GA+RPA approach. As in the HF+RPA method the f-sum rule is exactly satisfied with the following prescription. The optical conductivity on one side of the equality should be computed at the GA+RPA (HF+RPA) level and the expectation value on the other side (essentially the kinetic energy in our case) computed at GA (HF) level.
In this regard the f-sum rule provides also an encouraging argument that the GA+RPA dynamical correlation functions are much more accurate than those obtained via the corresponding HF+RPA method. We have compared the exact kinetic energy of a 4x4 lattice with unrestricted GA and HF results for various hole concentrations and have found that over a wide range of doping and on-site correlation $`U`$ there is almost perfect agreement between the GA method and exact results. On the other hand the HF kinetic energy has an error that for example for $`U=4t`$ is at least 40 times larger. This is not surprising since GA takes into account the correlation induced reduction of kinetic energy in a much better way than HF.
Fig. 2 displays $`\sigma (\omega )`$ for a 32-site Hubbard ring and half-filling in case of $`U/t=4`$. The onset of excitations across the Mott-Hubbard gap is signalled by the appearance of a large peak in $`\sigma (\omega )`$. We find excellent agreement between Monte-Carlo (MC) and GA+RPA whereas the excitation energy in HF is clearly overestimated. Note that the MC data display an additional hump at approximately twice the energy of the first peak. Both particle-particle scattering processes (not included at RPA level) and the failure of the antiadiabaticity assumption for the double occupancy at high energies can explain the discrepancy. We see however that for energies of the order of the Hubbard band the method performs very well.
In conclusion we have presented a time-dependent GA for the calculation of dynamical and static quantities in the Hubbard model. The approach is conceptually very simple and leads to much better agreement with exact results than previous approximations.
As in any computation of fluctuations we are dealing with the residual interaction between particles beyond the mean-field level. Roughly speaking since the GA contains ground state correlations not included in a HF wave function the residual interaction is a smaller perturbation to the MF state and hence it is natural that RPA works much better in this case.
It is interesting to remark that the computation of the ground state energy presented here is reminiscent of the evaluation of the one in the uniform electron gas based on Hubbard type dielectric functions. Also there the computation goes through a density-density correlation function and the coupling constant integration trick. Furthermore the Hubbard local field correction to the dielectric function takes into account the correlation hole in the uniform electron gas whereas the GW projector method takes into account similar correlations in the Hubbard model. The connection between these two approaches deserves further investigation as it may lead to a unified approach to strongly correlated systems.
We greatfully acknowledge many useful discussions with C. Di Castro, C. Castellani, M. Grilli and R. Raimondi. We acknowledge partial financial support from INFM. G.S. acknowledges financial support from the Deutsche Forschungsgemeinschaft as well as hospitality and support from the Dipartimento di Fisica of Università di Roma “La Sapienza”where part of this work was carried out.
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# Universality of the local spacing distribution in certain ensembles of Hermitian Wigner matrices
## 1. Introduction and main results
Consider a probability measure $`_N`$ on the space of all $`N\times N`$ hermitian matrices. We will be interested in the statistical properties of the spectrum as $`N`$ becomes large, in particular in features that are insensitive to the details of the particular sequence of probability measures we are considering. It is believed, on the basis of numerical simulations, that for many types of hermitian random matrix ensembles, i.e. choices of $`_N`$, the local statistical properties of the eigenvalues are the same as for the Gaussian Unitary Ensemble (GUE), where $`d_N(M)=Z_N^1\mathrm{exp}(\frac{N}{2}\text{Tr }M^2)dM`$. Here $`dM`$ is Lebesgue measure on the space $`_N^{N^2}`$ of all $`N\times N`$ hermitian matrices. The asymptotic eigenvalue density as $`N\mathrm{}`$ (density of states) is given by the Wigner semicircle law $`\rho (t)=\frac{1}{2\pi }\sqrt{(4t^2)_+}`$. Let $`\rho _N(x_1,\mathrm{},x_N)`$ be the induced probability density on the eigenvalues. The semicircle law is the limit of the one-dimensional marginal density as $`N\mathrm{}`$. The m - point correlation function
(1.1)
$$R_m^{(N)}(x_1,\mathrm{},x_m)=\frac{N!}{(Nm)!}_{^{Nm}}\rho _N(x)𝑑x_{m+1}\mathrm{}𝑑x_N,$$
is given by, ch. 5, ,
(1.2)
$$R_m^{(N)}(x_1,\mathrm{},x_m)=det(K_N(x_i,x_j))_{i,j=1}^m,$$
where the kernel $`K_N(x,y)`$ is given by
(1.3)
$$K_N(x,y)=\frac{\kappa _{N1}}{\kappa _N}\frac{p_N(x)p_{N1}(y)p_{N1}(x)p_N(y)}{xy}e^{N(x^2+y^2)/4}.$$
Here $`p_N(x)=\kappa _Nx^N+\mathrm{}`$ are the normalized orthogonal polynomials with respect to the weight function $`\mathrm{exp}(Nx^2/2)`$ on $``$ (rescaled Hermite polynomials). From these formulas, and Plancherel-Rotach asymptotics for the Hermite polynomials it follows that
(1.4)
$$\underset{N\mathrm{}}{lim}\frac{1}{(N\rho (u))^m}R_m^{(N)}(u+\frac{t_1}{N\rho (u)},\mathrm{},u+\frac{t_m}{N\rho (u)})=det\left(\frac{\mathrm{sin}\pi (t_it_j)}{\pi (t_it_j)}\right)_{i,j=1}^m$$
if $`\rho (u)>0`$. It has been proved, , , , that this is also true in other invariant ensembles of the form $`d_N(M)=Z_N^1\mathrm{exp}(N\text{Tr }V(M))dM`$. The orthogonal polynomials in (1.3) are then replaced by polynomials orthogonal with respect to $`\mathrm{exp}(NV(x))`$ on $``$. That the ensemble is invariant means that the probability measure is invariant under the conjugation $`MU^1MU`$, with a unitary matrix $`U`$. Sufficient control of the limit (1.4) for all $`m1`$, makes it possible to determine the asymptotic spacing distribution, i.e. distances between nearest neighbour eigenvalues, see . More precisely, let $`\{t_N\}`$ be a sequence such that $`t_N\mathrm{}`$ but $`t_N/N0`$ as $`N\mathrm{}`$ and define, , , $`S_N(s,x)`$, $`s0`$, $`x^N`$, to be the symmetric function, which for $`x_1<\mathrm{}<x_N`$ is defined by
(1.5)
$$S_N(s,x)=\frac{1}{2t_N}\mathrm{\#}\{1jN1;x_{j+1}x_j\frac{s}{N\rho (u)},|x_ju|\frac{t_N}{N\rho (u)}\}.$$
Given an hermitian matrix $`M`$ let $`x_1(M)<\mathrm{}x_N(M)`$ be its eigenvalues; we write $`x(M)=(x_1(M),\mathrm{},x_N(M))`$. Then it is proved in that
(1.6)
$$\underset{N\mathrm{}}{lim}𝔼_N[S_N(s,x(M))]=_0^sp(\sigma )𝑑\sigma ,$$
for a large class of invariant ensembles. Here $`p(\sigma )`$ is the density of the $`\beta =2`$ local spacing distribution, the Gaudin distribution, given by the probability density
(1.7)
$$p(s)=\frac{d^2}{ds^2}det(IK)_{L^2(0,s)},$$
where $`K`$ is the operator on $`L^2(0,s)`$ with kernel $`K(t,s)=\mathrm{sin}\pi (ts)/\pi (ts)`$, the sine kernel, see .
The aim of the present paper is to extend (1.4) and (1.6) to other, non-invariant ensembles. It is conjectured, see p.9, that (1.4) and (1.6) should hold also for so called Wigner matrices where the elements are independent but not necessarily Gaussian variables. In this case the probability measure is not invariant under conjugation by unitary matrices. For other results on Wigner matrices see for example , , , , and . In particular, in the universality of the fluctuations of the largest eigenvalue is established. To be more precise, consider the complex random variables $`w_{jk}`$, $`1jk`$ with independent laws $`P_{jk}=P_{jk}^RP_{jk}^I`$, where $`P_{jj}^I=\delta _0`$. Let $`𝒲^p`$, a class of Wigner ensembles, denote the class of all $`\{P_{jk}\}_{1jk}`$ which satisfy
(1.8)
$$zdP_{jk}(z)=0,|z|^2dP_{jk}(z)=\sigma ^2$$
for all $`1jk`$, and furthermore
(1.9)
$$\underset{j,k}{sup}|z|^p𝑑p_{jk}(z)<\mathrm{}.$$
If $`w_{kj}=\overline{w}_{jk}`$, $`W=(w_{jk})_{j,k=1}^N`$ is an $`N\times N`$ hermitian Wigner matrix .
Fix $`a>0`$ and let $`\varphi _a(t)=(\pi a^2)^{1/2}\mathrm{exp}(t^2/a^2)`$ be a Gaussian density function. Define $`Q_{jk}^{R,I}=\varphi _aP_{jk}^{R,I}`$, $`1j<k`$, $`Q_{jj}^R=\varphi _{a\sqrt{2}}P_{jj}^R`$, $`j1`$ and $`Q_{jj}^I=\delta _0`$. Then $`Q`$ is also a Wigner ensemble and we let $`𝒲_a^p`$ denote the subclass of $`𝒲^p`$ obtained in this way. Note that although $`𝒲_a^p`$ does not contain all Wigner ensembles it does contain cases where the distribution of the matrix elements have very different shapes, so in this sense it is rather broad, and proving universality in $`𝒲_a^p`$ clearly shows that the universality is not restricted to the invariant ensembles. Another way to describe this ensemble of random matrices is as follows. Let $`V`$ be a GUE-matrix with the probability measure $`Z_N^1\mathrm{exp}(\frac{1}{2}\text{Tr }V^2)dV`$ and let $`W`$ be an $`N\times N`$ Wigner matrix with distribution $`P𝒲^p`$, i.e. the law of $`w_{jk}`$ is $`P_{jk}`$. We will assume that the variance $`\sigma ^2=1/4`$, which can always be achieved by rescaling. Then $`W+aV`$ has the distribution $`Q`$, and we write
$$M=\frac{1}{\sqrt{N}}(W+aV).$$
We can think of this in terms of Dyson’s Brownian motion model, , $`W+aV`$ is obtained from $`W`$ by letting the matrix elements execute a Brownian motion for a time $`a^2`$, see sect. 2. If $`P𝒲^p`$ and $`W`$ is is an $`N\times N`$ hermitian matrix we let $`P^{(N)}`$ denote the distribution of $`H=W/\sqrt{N}=(h_{jk})`$, i.e.
$$dP^{(N)}(H)=\underset{1jkN}{}dP_{jk}(\sqrt{N}h_{jk}).$$
The matrix $`M`$ has the distribution $`Q^{(N)}`$, which is given by
(1.10)
$$dQ^{(N)}(M)=2^{N/2}\left(\frac{N}{\pi a^2}\right)^{N^2/2}\left(__Ne^{\frac{N}{2a^2}\text{Tr }(MH)^2}𝑑P^{(N)}(H)\right)dM,$$
and this is the measure we will study. The asymptotic distribution of the eigenvalues $`x_1,\mathrm{},x_N`$ of $`M`$ is the semicircle law
(1.11)
$$\rho (u)=\frac{2}{\pi (1+4a^2)}\sqrt{(1+4a^2u^2)_+}.$$
The following proposition will be proved in sect. 2 using an argument from , .
###### Proposition 1.1.
The symmetrized eigenvalue measure on $`^N`$ induced by $`Q^{(N)}`$ has a density
(1.12)
$$\rho _N(x)=__N\rho _N(x;y(H))𝑑P^{(N)}(H)$$
where
(1.13)
$$\rho _N(x;y)=\left(\frac{N}{2\pi a^2}\right)^{N/2}\frac{\mathrm{\Delta }_N(x)}{\mathrm{\Delta }_N(y)}det(e^{\frac{N}{2a^2}(x_jy_k)^2})_{j,k=1}^N$$
and $`\mathrm{\Delta }_N(x)=_{1i<jN}(x_ix_j)`$ is the Vandermonde determinant.
The main result of the present paper is that for Wigner ensembles from $`𝒲_a^p`$ we can prove (1.4) and (1.6), and thus extend the universality to a rather broad class of Wigner matrices.
###### Theorem 1.2.
Fix $`a>0`$ and assume that $`|u|\sqrt{1/2+2a^2}`$. Let $`R_M^{(N)}(x_1,\mathrm{},x_m)`$ be the correlation functions, defined by (1.1), of the eigenvalue measure $`\rho _N`$, (1.12), for $`Q^{(N)}`$, (1.10). Let $`fL_c^{\mathrm{}}(^m)`$, the set of all $`L^{\mathrm{}}`$ functions on $`^m`$ with compact support, and set for $`x^N`$
$$(Sf)(x)=\underset{i_1,\mathrm{},i_m}{\overset{}{}}f(x_{i_1},\mathrm{},x_{i_m}),$$
where the sum is over all distinct indices from $`\{1,\mathrm{},N\}`$. If $`Q𝒲_a^p`$ with $`p>2(m+2)`$, then
(1.14) $`\underset{N\mathrm{}}{lim}{\displaystyle __N}(Sf)(N\rho (u)(x_1(M)u),\mathrm{}N\rho (u)(x_N(M)u))𝑑Q^{(N)}(M)`$
$`\underset{N\mathrm{}}{lim}{\displaystyle _^m}f(t_1,\mathrm{},t_m){\displaystyle \frac{1}{(N\rho (u))^m}}R_m^{(N)}(u+{\displaystyle \frac{t_1}{N\rho (u)}},\mathrm{},u+{\displaystyle \frac{t_m}{N\rho (u)}})d^mt`$
$`={\displaystyle _^m}f(t_1,\mathrm{},t_m)det({\displaystyle \frac{\mathrm{sin}\pi (t_it_j)}{\pi (t_it_j)}})_{i,j=1}^md^mt.`$
The condition on $`u`$ is made just to simplify the saddle-point argument in sect. 3; the result should hold for any $`u`$ with $`\rho (u)>0`$.
We can also prove that the spacing distribution is the same as for GUE.
###### Theorem 1.3.
Fix any $`a>0`$ and assume that $`Q𝒲_a^{6+ϵ}`$, $`ϵ>0`$. Let $`S_N(s,x)`$ be defined by (1.5). Then, for any $`s0`$,
(1.15)
$$\underset{N\mathrm{}}{lim}__NS_N(s,x(M))𝑑Q^{(N)}(M)=_0^sp(\sigma )𝑑\sigma ,$$
where $`p(s)`$ is given by (1.7).
The theorems will be proved in sect. 4 after the preparatory work in sect. 2 and 3.
## 2. The correlation functions
We will start by proving Proposition 1.1 using the Harish-Chandra/Itzykson-Zuber formula following , . After that we will give a formula for the correlation functions of $`\rho _N(x;y)`$, which is very close to the formula in , but our derivation will be different. A central role will be played by non-intersecting one-dimensional Brownian motions and we will use the formulas of Karlin and McGregor. Also we will discuss the relation to Dyson’s Brownian motion model. This connection can be found in and we will only give an outline.
Proof. Let $`F(x)`$ be a continuous symmetric function on $`^N`$. By Fubini’s theorem
(2.1)
$$__NF(x(M))𝑑Q^{(N)}(M)=c_N^{(1)}__N\left(__NF(x(M))e^{\frac{N}{2a^2}\text{Tr }(MH)^2}𝑑M\right)𝑑P^{(N)}(H)$$
with $`c_N^{(1)}=2^{N/2}(N/\pi a^2)^{N^2/2}`$. In the right hand side of (2.1) we make the substitution $`M=U^1RU`$, with $`UU(N)`$ and $`R_N`$, and then integrate over $`U(N)`$. If we use Fubini’s theorem again, we obtain
$$c_N^{(1)}__N\left(__NF(x(R))\left(_{U(N)}e^{\frac{N}{2a^2}\text{Tr }(U^1RUH)^2}𝑑U\right)𝑑R\right)𝑑P^{(N)}(H).$$
Here we have also used the fact that $`dM=dR`$. The integral over $`U(N)`$ can now be evaluated using the Harish-Chandra/Itzykson-Zuber formula, ,, see also A.5. We obtain the integral
$$c_N^{(1)}c_N^{(2)}__N\left(__NF(x)\frac{1}{\mathrm{\Delta }_N(x)\mathrm{\Delta }_N(y)}det(e^{N(x_jy_k)^2/2a^2})_{j,k=1}^NdR\right)𝑑P^{(N)}(H),$$
where $`y_1,\mathrm{},y_N`$ are the eigenvalues of $`H`$ and $`c_N^{(2)}=(a^2/N)^{N(N1)/2}_{j=1}^Nj!`$. The integrand in the middle integral depends only on the eigenvalues $`x`$ of $`R`$ and hence we can integrate out the other degrees of freedom in the standard way, ch. 3, and obtain, after using Fubini’s theorem,
$`{\displaystyle __N}F(x(M))𝑑Q^{(N)}(M)`$
(2.2) $`=c_N^{(1)}c_N^{(2)}c_N^{(3)}{\displaystyle __N}\left({\displaystyle _^N}F(x){\displaystyle \frac{\mathrm{\Delta }_N(x)}{\mathrm{\Delta }_N(y)}}det(e^{N(x_jy_k)^2/2a^2})_{j,k=1}^Nd^Nx\right)𝑑P^{(N)}(H)`$
with $`c_N^{(3)}=\pi ^{N(N1)/2}_{j=1}^N(j!)^1`$. We see that $`c_N^{(1)}c_N^{(2)}c_N^{(3)}=(N/2\pi a^2)^{N/2}`$ and since (2) holds for arbitrary bounded, continuous, symmetric $`F(x)`$ we have proved that the symmetrized eigenvalue measure is given by 1.12. This proves proposition 1.1. $`\mathrm{}`$
Let $`p_t(x,y)`$ be the transition probability of a Markov process $`X(t)`$ on $``$ with continuous paths. Consider $`N`$ independent copies of the process $`(X_1(t),\mathrm{},X_N(t))`$ and assume that this is a strong Markov process in $`^N`$. Suppose that the particles start at positions $`y_1<\mathrm{}<y_N`$ at time 0. The probability density that they are at positions $`x_1<\mathrm{}<x_N`$ at time $`S`$ given that their paths have not intersected anytime during the time interval $`[0,S]`$ is, by a theorem of Karlin and McGregor, ,
$$det(P_S(y_j,x_k))_{j,k=1}^N.$$
Hence, the conditional probability density that the particles are at positions $`y_1<\mathrm{}<y_N`$ at time 0, at positions $`x_1<\mathrm{}<x_N`$ at time $`S`$, at positions $`z_1<\mathrm{}<z_N`$ at time $`S+T`$, given that their paths have not intersected in the time interval $`[0,S+T]`$ is
(2.3)
$$q_{S,T}(x;y;z)\frac{1}{𝒵_N}det(P_S(y_j,x_k))_{j,k=1}^Ndet(P_T(x_j,z_k))_{j,k=1}^N,$$
where
$$𝒵_N=_{x_1<\mathrm{}<x_N}det(P_S(y_j,x_k))_{j,k=1}^Ndet(P_T(x_j,z_k))_{j,k=1}^Nd^Nx;$$
we assume that $`𝒵_N>0`$. Note that the expression (2.3) is a symmetric function of $`x_1,\mathrm{},x_N`$, so we can regard it as a probability measure on $`^N`$. Our next lemma shows that we can obtain $`\rho _N(x;y)`$ defined by (1.13) as a limit of the measure in (2.3).
###### Lemma 2.1.
Let $`z_j=j1`$, $`1jN`$ and let $`p_t(x,y)=(2\pi t)^{1/2}\mathrm{exp}((xy)^2/2t)`$ be the transition probability for Brownian motion. Then, for any $`x^N`$ and $`y_1<\mathrm{},y_N`$,
(2.4)
$$\underset{T\mathrm{}}{lim}q_{S,T}(x;y;z)=\frac{1}{(2\pi S)^{N/2}}\frac{\mathrm{\Delta }_N(x)}{\mathrm{\Delta }_N(y)}det(e^{(x_jy_k)^2/2S})_{j,k=1}^Nq_S(x;y).$$
Note that $`\rho _N(x;y)=q_{a^2/N}(x;y)`$.
Proof. Write
$`det(P_S(y_j,x_k))_{j,k=1}^Ndet(P_T(x_j,z_k))_{j,k=1}^N`$
(2.5) $`={\displaystyle \frac{1}{(2\pi )^N(TS)^{N/2}}}det(e^{(x_jy_k)^2/2S})_{j,k=1}^N{\displaystyle \underset{j=1}{\overset{N}{}}}e^{\frac{x_j^2+z_j^2}{2T}}det(e^{x_jz_k/2T})_{j,k=1}^N.`$
Note that $`𝒵_N`$ is the conditional probability density of going from $`y_1<\mathrm{},y_N`$ to $`z_1<\mathrm{},z_N`$ without collosions, i.e.
$`𝒵_N`$ $`=det(p_{S+T}(y_j,z_k))_{j,k=1}^N`$
(2.6) $`={\displaystyle \frac{1}{(2\pi )^{N/2}(S+T)^{N/2}}}{\displaystyle \underset{j=1}{\overset{N}{}}}e^{\frac{y_j^2+z_j^2}{2(S+T)}}det(e^{y_jz_k/2(S+T)})_{j,k=1}^N.`$
Now, since $`z_j=j1`$, we have two Vandermonde determinants in (2) and (2). If we evaluate these, take the quotient between (2) and (2) and then take the limit $`T\mathrm{}`$, we obtain the right hand side of (2.4). $`\mathrm{}`$
Proposition 1.1 and lemma 2.1 establish a link between the eigenvalue distribution of $`M=(W+aV)/\sqrt{N}`$ and the non-intersecting Brownian paths. If we set $`S=a^2/N`$, then the right hand side of (2.4) and (1.13) are identical; $`y_1<\mathrm{}<y_N`$ are the eigenvalues of $`H=W/\sqrt{N}`$. This relation can also be seen in another way, which we will now outline. Let $`X(t)=(x_{jk}(t))_{j,k=1}^N`$ be an $`N\times N`$ Hermitian matrix, where $`\text{Re }x_{jk}(t)`$, $`\text{Im }x_{jk}(t)`$, $`jk`$ are independent Brownian motions with variance $`(1+\delta _{jk})/2`$. Assume that $`X(0)=H`$ is distributed according to $`P^{(N)}`$. Then the distribution of $`X(a^2/N)`$ is the same as that of $`M=(W+aV)/\sqrt{N}`$. Following Dyson, , see also , it is possible to derive a stochastic differential equation for the eigenvalues $`\lambda _1(t),\mathrm{},\lambda _N(t)`$ of $`X(t)`$,
(2.7)
$$d\lambda _i=dB_i+\underset{ki}{}\frac{1}{\lambda _i\lambda _k}dt,$$
where $`B_i`$ are independent standard Brownian motions on $``$, and with the intial conditions $`\lambda _i(0)=y_i`$, $`1iN`$. We can also consider the problem of non-intersecting Brownian motions in a different way than that of Karlin and McGregor. Namely, let $`K=\{x^N;x_1<\mathrm{}<x_N\}`$ and consider Brownian motion in $`^N`$ starting at $`yK`$ and conditioned to remain in $`K`$ forever. As proved in , see also , , if $`\lambda _i`$ are the components of the $`N`$-dimensional conditioned Brownian motion they satisfy the stochastic differential equation (2.7) with the same initial conditions. This gives another way to obtain (1.13) without using the Harish-Chandra, Itzykson/Zuber formula. Actually, we can turn the argument around and give a proof of this formula.
We turn now to the computation of the correlation functions of the right hand side of (2.4), but we start more generally with (2.3). This can be analyzed using the techniques of , compare the analysis of the Schur measure, , in , and see also . For completeness, let us outline the result we need from . Let $`(\mathrm{\Omega },\mu )`$ be a measure space. Assume that $`\varphi _j,\psi _jL^2(\mathrm{\Omega },\mu )`$, $`1jN`$, and $`fL^{\mathrm{}}(\mathrm{\Omega },\mu )`$. Set
$$Z_N[f]=\frac{1}{N!}_{\mathrm{\Omega }_N}det(\varphi _j(x_k))_{j,k=1}^Ndet(\psi _j(x_k))_{j,k=1}^N\underset{j=1}{\overset{N}{}}f(x_j)d\mu (x_j)$$
and
$$A=\left(_\mathrm{\Omega }\varphi _j(x)\psi _j(x)𝑑\mu (x)\right)_{j,k=1}^N.$$
###### Proposition 2.2.
() Assume that $`Z_N[1]0`$. Then $`A`$ is invertible and we can define
(2.8)
$$K_N(t,s)=\underset{j,k=1}{\overset{N}{}}\psi _k(t)(A^1)_{kj}\varphi _j(s).$$
Then, for any $`gL^{\mathrm{}}(\mathrm{\Omega },\mu )`$,
(2.9)
$$\frac{Z_N[1+g]}{Z_N[1]}=det(I+K_Ng)_{L^2(\mathrm{\Omega })}.$$
If we define a density on $`\mathrm{\Omega }^N`$ by
(2.10)
$$u_N(x)=\frac{1}{N!Z_N[1]}det(\varphi _j(x_k))_{j,k=1}^Ndet(\psi _j(x_k))_{j,k=1}^N,$$
then it has the correlation functions
(2.11)
$$\frac{N!}{(NM)!}_{\mathrm{\Omega }^{Nm}}u_N(x)𝑑x_{m+1}\mathrm{}𝑑x_N=det(K_N(x_i,x_j))_{i,j=1}^N.$$
Proof. We will indicate the main steps in the proof of (2.9) of which (2.11) is a consequence, see . Set
$$B=\left(_\mathrm{\Omega }\varphi _j(x)\psi _j(x)g(x)𝑑\mu (x)\right)_{j,k=1}^N.$$
Then, by the formula
$$Z_N[f]=det\left(_\mathrm{\Omega }\varphi _j(x)\psi _j(x)f(x)𝑑\mu (x)\right)_{j,k=1}^N,$$
which goes all the way back to , and which is not difficult to prove by expanding the determinants, we see that $`detA=Z_N[1]0`$, so $`A`$ is invertible and
(2.12)
$$\frac{Z_N[1+g]}{Z_N[1]}=\frac{det(A+B)}{detA}=det(I+A^1B).$$
Now,
$$(A^1B)_{jk}=_\mathrm{\Omega }\psi _k(x)\left(\underset{\mathrm{}=1}{\overset{N}{}}(A^1)_j\mathrm{}\varphi _{\mathrm{}}(x)g(x)\right)𝑑\mu (x),$$
and we define $`T:^NL^2(\mathrm{\Omega },d\mu )`$ and $`S:L^2(\mathrm{\Omega },d\mu )^N`$ by the kernels $`T(x,k)=\psi _k(x)`$ and $`S(j,x)=_{\mathrm{}=1}^N(A^1)_j\mathrm{}\varphi _{\mathrm{}}(x_g(x)`$. Then, by (2.12) and a determinant identity,
$`{\displaystyle \frac{Z_N[1+g]}{Z_N[1]}}`$ $`=det(I+ST)_^N=det(I+TS)_{L^2(\mathrm{\Omega },d\mu )}`$
$`=det(I+K_Ng)_{L^2(\mathrm{\Omega },d\mu )},`$
with $`K_N`$ given by (2.8). Note that $`K_Ng`$, which means first multiplication by $`g`$ and then application of the operator on $`L^2(\mathrm{\Omega },d\mu )`$ with kernel $`K_N`$, is a finite rank operator. $`\mathrm{}`$
Observe now that if we take $`\mathrm{\Omega }=`$, $`d\mu (x)=dx`$, $`\varphi _j(x)=P_T(x,z_j)`$ and $`\psi _j(x)=p_S(y_j,x)`$, then (2.3) is a probability density of the form (2.10) and we can apply the proposition. Note that
$$(A)_{jk}=_{}p_T(x,z_j)p_S(y_k,x)𝑑x=p_{S+T}(y_k,z_j).$$
The kernel which gives the correlation functions is
$$K_N^{S,T}(u,v)=\underset{k=1}{\overset{N}{}}p_S(y_k,v)\left(\underset{j=1}{\overset{N}{}}(A^1)_{jk}p_T(u,z_j)\right).$$
Let $`A_k(v)`$ be the matrix we obtain from $`A`$ by replacing column $`k`$ by $`(p_T(v,z_1)\mathrm{}p_T(v,z_N))^T`$. Then, by Kramers’ rule,
(2.13)
$$K_N^{S,T}(u,v)=\underset{k=1}{\overset{N}{}}p_S(y_k,v)\frac{detA_k(v)}{detA}.$$
This formula and proposition 2.2 is the basis for the next proposition. The result is closely related to the result derived in by different methods.
###### Proposition 2.3.
The correlation functions for $`q_S(x;y)`$ defined by (2.4) are given by
(2.14) $`R_m^N(x_1,\mathrm{},x_m;y)`$ $`{\displaystyle \frac{N!}{(Nm)!}}{\displaystyle _{^{Nm}}}q_S(x;y)𝑑x_{m+1}\mathrm{}𝑑x_N`$
$`=det(K_N^S(x_i,x_j;y))_{i,j=1}^m,`$
where
(2.15) $`K_N^S(u,v;y)={\displaystyle \frac{e^{(v^2u^2)/2S}}{(vu)S(2\pi i)^2}}{\displaystyle _\gamma }𝑑z{\displaystyle _\mathrm{\Gamma }}𝑑w(1e^{(vu)z/S})`$
$`\times {\displaystyle \frac{1}{z}}\left(w+zvS{\displaystyle \underset{j}{}}{\displaystyle \frac{y_j}{(wy_j)(zy_j)}}\right)e^{(w^22vwz^2+2uz)/2S}.`$
Here $`\gamma `$ is the union of the curves $`tt+i\omega `$, $`t`$ and $`tti\omega `$, $`t`$ with a fixed $`\omega >0`$, and $`\mathrm{\Gamma }:tit`$.
Proof. We have to show that with $`p_t(u,v)=(2\pi t)^{1/2}\mathrm{exp}((uv)^2/2t)`$ and $`z_j=j1`$ the limit of the right hand side of (2.13) as $`T\mathrm{}`$ can be written as (2.15). The result then follows from lemma 2.1, proposition 2.2 and the dominated convergence theorem. We see that
(2.16)
$$detA=\frac{1}{(2\pi (S+T))^{N/2}}\underset{j=1}{\overset{N}{}}e^{\frac{z_j^2+y_j^2}{2(S+T)}}\underset{1i<jN}{}(e^{\frac{y_j}{S+T}}e^{\frac{y_i}{S+T}})$$
by the formula for a Vandermonde determinant. Let $`\mathrm{\Gamma }_M^{}`$ be the curve $`tt+iM`$, $`t`$, $`M`$ fixed. Then
$$p_T(z_j,v)=\frac{1}{\sqrt{2\pi T}}e^{\frac{z_j^2}{2(S+T)}\frac{v^2}{2T}}\frac{1}{\sqrt{2\pi }}_{\mathrm{\Gamma }_M^{}}e^{\frac{\tau ^2}{2}+z_j(\frac{v}{T}+i\tau \sqrt{\frac{S}{2T(S+T)}})}𝑑\tau $$
Hence,
$`detA_k(v)`$ $`={\displaystyle \frac{1}{2\pi \sqrt{T}}}{\displaystyle \frac{1}{(2\pi (S+T))^{(N1)/2}}}\left({\displaystyle \underset{j=1}{\overset{N}{}}}e^{\frac{z_j^2}{2(S+T)}}\right)\left({\displaystyle \underset{jk}{}}e^{\frac{y_j^2}{2(S+T)}}\right)`$
$`\times e^{\frac{v^2}{2T}}{\displaystyle _{\mathrm{\Gamma }_M^{}}}e^{\frac{\tau ^2}{2}}det\stackrel{~}{A}_k(v)d\tau ,`$
where $`\stackrel{~}{A}_k(v)`$ is the matrix we get from $`(\mathrm{exp}(\frac{z_jy_k}{S+T}))_{j,k=1}^N`$ by replacing column $`k`$ by $`(\mathrm{exp}(z_j(\frac{v}{T}+i\tau \sqrt{\frac{S}{2T(S+T)}})))_{j=1}^N`$. Since $`z_j=j1`$ we have a Vandermonde determinant and we obtain
$`detA_k(v)`$ $`=\sqrt{{\displaystyle \frac{S+T}{T}}}{\displaystyle \frac{1}{(2\pi (S+T))^{N/2}}}\left({\displaystyle \underset{j=1}{\overset{N}{}}}e^{\frac{z_j^2}{2(S+T)}}\right)\left({\displaystyle \underset{jk}{}}e^{\frac{y_j^2}{2(S+T)}}\right)`$
(2.17) $`\times e^{\frac{v^2}{2T}}{\displaystyle \frac{1}{\sqrt{2ı}}}{\displaystyle _{\mathrm{\Gamma }_M^{}}}e^{\frac{\tau ^2}{2}}{\displaystyle \underset{1i<jN}{}}(e^{\frac{y_j}{S+T}}e^{\frac{y_i}{S+T}})d\tau ,`$
where $`y_k`$ should be replaced by $`(S+T)(\frac{v}{T}+i\tau \sqrt{\frac{S}{2T(S+T)}})`$. Take the quotient of (2.16) and (2) and let $`T\mathrm{}`$. This gives
$$\underset{T\mathrm{}}{lim}\frac{detA_k(v)}{detA}=\frac{1}{\sqrt{2\pi }}_{\mathrm{\Gamma }_M^{}}e^{\frac{\tau ^2}{2}}\underset{jk}{}\left(\frac{v+i\sqrt{S}\tau y_j}{y_ky_j}\right)d\tau .$$
Choose $`M`$ so that $`v\sqrt{S}M=L`$, where $`L`$ is given, and make the change of variables $`w=v+i\sqrt{S}\tau `$. Then
$$\underset{T\mathrm{}}{lim}\frac{detA_k(v)}{detA}=\frac{1}{i\sqrt{2\pi S}}_{\mathrm{\Gamma }_L}e^{\frac{(wv)^2}{2S}}\underset{jk}{}\left(\frac{wy_j}{y_ky_j}\right)dw,$$
where $`\mathrm{\Gamma }_L:tL+it`$, $`t`$. Thus, using (2.13),
$$K_N^S(u,v;y)=\frac{1}{2\pi iS}\underset{k=1}{\overset{N}{}}e^{(y_ku)^2/2S}_{\mathrm{\Gamma }_L}e^{\frac{(wv)^2}{2S}}\underset{jk}{}\left(\frac{wy_j}{y_ky_j}\right)dw.$$
Let $`\gamma `$ be a curve surrounding $`y_1,\mathrm{},y_N`$ and choose $`L`$ so large that $`\gamma `$ and $`\mathrm{\Gamma }`$ do not intersect. The residue theorem gives
$$\frac{1}{2\pi i}_\gamma \frac{e^{(zu)^2/2S}}{wz}\underset{j=1}{\overset{N}{}}\frac{wy_j}{zy_j}dz=\underset{k=1}{\overset{N}{}}e^{(y_ku)^2/2S}\underset{jk}{}\left(\frac{wy_j}{y_ky_j}\right)$$
for all $`w\mathrm{\Gamma }_L`$. Thus,
(2.18)
$$K_N^S(u,v;y)=\frac{e^{\frac{v^2u^2}{2S}}}{(2\pi i)^2S}_\gamma 𝑑z_{\mathrm{\Gamma }_L}𝑑we^{\frac{1}{2S}(w^22vwz^2+2uz)}\frac{1}{wz}\underset{j=1}{\overset{N}{}}\frac{wy_j}{zy_j}.$$
In (2.18) we make the change of variables $`zbz`$, $`wbw`$ with $`b`$ close to 1. This will modify the contours but we can use Cauchy’s theorem to deform back to $`\gamma `$ and $`\mathrm{\Gamma }_L`$. Now, take the derivative with respect to $`b`$ and then put $`b=1`$. This gives the equation
$`0`$ $`=K_N^S(u,v;y)+{\displaystyle \frac{e^{\frac{v^2u^2}{2S}}}{(2\pi i)^2S^2}}{\displaystyle _\gamma }𝑑z{\displaystyle _{\mathrm{\Gamma }_L}}𝑑w{\displaystyle \frac{1}{wz}}e^{\frac{1}{2S}(w^22vwz^2+2uz)}`$
$`\times \left[w^2z^2+uzvw+S{\displaystyle \underset{j=1}{\overset{N}{}}}\left({\displaystyle \frac{w}{wy_j}}{\displaystyle \frac{z}{zy_j}}\right)\right]{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{wy_j}{zy_j}}.`$
This can be written
$`{\displaystyle \frac{}{u}}((uv)K_N^S(u,v;y))={\displaystyle \frac{e^{\frac{v^2u^2}{2S}}}{(2\pi i)^2S^2}}{\displaystyle _\gamma }𝑑z{\displaystyle _{\mathrm{\Gamma }_L}}𝑑w`$
$`\left[w+zvS{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{y_j}{(wy_j)(zy_j)}}\right]e^{(w^22vwz^2)/2S}e^{uz/S}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{wy_j}{zy_j}},`$
and integration of this formula gives (2.15). In this last formula we can choose $`L`$ arbitrarily and take $`\gamma `$ to be the curve in the proposition by using Cauchy’s formula, This completes the proof. $`\mathrm{}`$
We now take $`S=a^2/N`$ and set
(2.19)
$$𝒦_N(u,v;y)=e^{\frac{N(u^2v^2)}{2a^2}+\omega (uv)}K_N^{a^2/N}(u,v;y),$$
where $`\omega `$ is a constant that will be specified later. Note that we can replace $`K_N^{a^2/N}`$ with $`𝒦_N`$ in (2.14) without changing the correlation functions, so we can just as well work with $`𝒦_N`$. Set
$`f_N(z)={\displaystyle \frac{1}{2a^2}}(z^22uz)+{\displaystyle \frac{1}{N}}{\displaystyle \underset{j=1}{\overset{N}{}}}\mathrm{log}(zy_j)`$
$`g_N(z,w)={\displaystyle \frac{1}{a^2z}}\left(w+zu{\displaystyle \frac{a^2}{N}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{y_j}{(wy_j)(zy_j)}}\right)`$
$`h(z,w)={\displaystyle \frac{e^{\omega (uv)}}{N\rho (u)(vu)}}e^{\frac{N}{a^2}(uv)w}(e^{N(uv)w/a^2}e^{N(uv)(wz)/a^2}),`$
so that
(2.20)
$$𝒦_N(u,v;y)=N\rho (u)_\gamma \frac{dz}{2\pi i}_\mathrm{\Gamma }\frac{dw}{2\pi i}h(z,w)g_N(z,w)e^{N(f_N(w)f_N(z))}.$$
These are the formulas we will use in the asymptotic analysis. A straightforward computation shows that
(2.21)
$$g_N(z,w)=\frac{1}{z}f_N^{}(z)+\frac{f_N^{}(z)f_N^{}(w)}{zw}.$$
## 3. Asymptotics
The eigenvalues $`y_1,\mathrm{},y_N`$ of the Wigner matrix $`H`$ converge to the semicircle law
(3.1)
$$\sigma (t)=\frac{2}{\pi }\sqrt{1t^2},|t|1.$$
In order to be able to perform the saddle point analysis of (2.20) we need uniform control of the convergence of $`f_N(z)`$ to its limit
(3.2)
$$f(z)=\frac{1}{2a^2}(z^22uz)+_1^1\mathrm{log}(zt)\sigma (t)𝑑t.$$
In order to show this we must start with some probability estimates. Write $`\mathrm{\Omega }_{R,\eta }=\{z;|\text{Re }z|R,\eta |\text{Im }z|R\}`$.
###### Lemma 3.1.
Let $`FL^{\mathrm{}}(^N)`$ be symmetric and let $`\eta >0`$ and $`R>0`$ be given. Assume that $`P𝒲^p`$, $`p>4`$ and $`0<\xi <\mathrm{min}(\frac{1}{2}\frac{2}{p},\frac{1}{16})`$. Then, there is a probability measure $`\stackrel{~}{P}^{(N)}`$ on $`_N`$ such that
(3.3)
$$\left|__NF(x(H))𝑑P^{(N)}(H)__NF(x(H))𝑑\stackrel{~}{P}^{(N)}(H)\right|N^{2p(\frac{1}{2}\xi )}F_{\mathrm{}},$$
and
(3.4)
$$\underset{z\mathrm{\Omega }_{R,\eta }}{sup}\left|\frac{1}{N}\text{Tr }\mathrm{log}(zH)_1^1\mathrm{log}(zt)\sigma (t)𝑑t\right|CN^\xi $$
a.s. with respect to $`\stackrel{~}{P}^{(N)}`$.
Proof. Given $`P𝒲^p`$ we introduce a cut-off $`L>0`$ and define a new probability measure $`P_L𝒲^p`$ by
$$dP_{L,jk}^{R,I}(t)=\frac{1}{d_{L,jk}}\chi _{[L,L]}(t)dP_{jk}^{R,I}(t),1jk$$
where $`d_{L,jk}`$ is a normalization constant. Note that $`P_{L,jk}`$ is supported in $`K=[L,L]^2`$. Set $`d_L^{(N)}=_{1jkN}d_{L,jk}`$. Then,
$`\left|{\displaystyle __N}F(x(H))𝑑P^{(N)}(H){\displaystyle __N}F(x(H))𝑑P_L^{(N)}(H)\right|`$ $`F_{\mathrm{}}(1d_L^{(N)})(1+{\displaystyle \frac{1}{d_L^{(N)}}})`$
(3.5) $`{\displaystyle \frac{CN^2}{L^p}}F_{\mathrm{}}`$
for some constant $`C`$. The last estimate follows from
(3.6)
$$1d_L^{(N)}=P[\text{some }|W_{jk}|L]N^2\underset{1jk}{sup}\frac{E[|W_{jk}|^p]}{L^p}\frac{CN^2}{L^p}$$
by (1.9). Set $`D_N=\mathrm{\Omega }_{R,\eta }\frac{1}{N}^2`$ and note that $`\mathrm{\#}D_NCN^2`$ for some constant $`C`$ that only depend on $`R,\eta `$. For a given function $`f`$ set
$$A_N(f;\delta )=\{H_N;|\frac{1}{N}\text{Tr }(f(H))_1^1f(t)𝑑\sigma (t)|\delta \},$$
where $`\sigma (t)`$ is the semicircle law (3). Set
(3.7)
$$A_N(\delta )=\underset{zD_N}{}A_N(f_z,\delta ),$$
where $`f_z(t)=\mathrm{log}(zt)`$ (principal branch). To estimate the probability of $`A_N(\delta )`$ under $`P^{(N)}`$ we will use a result of Guionnet and Zeitouni, . Let
$$|f|_{}=\underset{t,s}{sup}\frac{|f(x)f(y)|}{|xy|},$$
and $`f_{}=f_{\mathrm{}}+|f|_{}`$. Then, by , corollary 1.6a), and the discussion before this corollary, given $`ϵ>0`$, there are positive constants $`C_0(ϵ)`$, $`C_1`$ and $`C_2`$ such that if we write
(3.8)
$$\delta _1(N)=C_1L^2|f|_{}N^1+C_2(ϵ)f_{}N^{1/4+ϵ},$$
then
(3.9)
$$P_L^{(N)}\left[|\frac{1}{N}\text{Tr }f(H)_1^1f(t)\sigma (t)𝑑t|\delta \right]4\mathrm{exp}\left[\frac{C_2N^2}{L^4|f|_{}}(\delta \delta _1(N))^2\right]$$
for any $`\delta >\delta _1(N)`$. Since under $`P_L^{(N)}`$ all $`|H_{jk}|\sqrt{2}(L/\sqrt{N})`$, the spectral radius is $`2L`$. Thus, the left hand side of (3.9) is unchanged if we replace $`f=f_z`$ with $`f=f_z^L(t)`$, where $`f_z^L(t)=\mathrm{log}(zt)`$ if $`|t|2L`$, $`f_z^L(t)=\mathrm{log}(z2L)`$ if $`t>2L`$ and $`f_z^L(t)=\mathrm{log}(z+2L)`$ if $`t<2L`$. Now, $`f_z^L(t)`$ is Lipschitz and there is a constant $`C_3`$, independent of $`L`$, such that $`|f_z^L(t)|_{}C_3`$ and $`f_z^L(t)_{}C_3(1+\mathrm{log}L)`$ for all $`z\mathrm{\Omega }_{z,\eta }`$. Take $`L=L_N=N^{1/2\xi }`$ and $`ϵ=1/6`$ in (3.8). Then $`\delta _1(N)CN^{2\xi }`$ and if we choose $`\delta =N^\xi `$ in (3.9) we obtain
(3.10)
$$P_L^{(N)}\left[|\frac{1}{N}\text{Tr }f_z(H)_1^1f_z(t)\sigma (t)𝑑t|N^\xi \right]c_1\mathrm{exp}(c_2N^{2\xi })$$
for some positive constants $`c_1,c_2`$. If we use (3.10) we see that the probability of the complement of the event in (3.7) can be estimated as
(3.11)
$$P_{L_N}^{(N)}[A_N(N^\xi )^c]CN^2e^{c_2N^{2\xi }}.$$
Set
$$d\stackrel{~}{P}^{(N)}(H)=(P_{L_N}^{(N)}[A_N(N^\xi )])^1\chi _{A_N(N^\xi )}(H)dP_{L_N}^{(N)}.$$
Note that $`N^2/L_N^p=N^{2p(1/2\xi )}`$, so combining (3), (3.7) and (3.11) we obtain the estimate (3.3). From the definition of $`A_N(\delta )`$ we see that (3.4) holds for $`zD_N`$, but then a straightforward approximation argument extends it to all $`z\mathrm{\Omega }_{R,\eta }`$. This completes the proof of lemma 3.1. $`\mathrm{}`$
We now come to the central asymptotic result.
###### Lemma 3.2.
Let $`\mathrm{\Omega }_{R,\eta }`$ be as above, let $`\xi (0,1/2]`$ and let $`K`$ be a compact subset of $``$. Also let $`u_N`$ be a sequence such that $`u_Nu`$ as $`N\mathrm{}`$. Furthermore, let $`Y_{R,\eta }`$ be the set of all $`y^N`$ such that
(3.12)
$$\underset{z\mathrm{\Omega }_{R,\eta }}{sup}\left|\frac{1}{N}\underset{j=1}{\overset{N}{}}\mathrm{log}(zy_j)_1^1\mathrm{log}(zt)\sigma (t)𝑑t\right|CN^\xi $$
for some constant $`C`$ and all $`N1`$, where $`\sigma (t)`$ is given by (3.1). Then, we can find $`R_0>0`$, $`\eta _0>0`$ and a constant $`C`$ such that for all $`yY_{R_0,\eta _0}`$, $`\tau K`$, $`|u|\sqrt{1/2+2a^2}`$ and $`N1`$,
(3.13)
$$\left|\frac{1}{N\rho (u)}𝒦_N(u_N,u_N+\frac{\tau }{N\rho (u)};y)\frac{\mathrm{sin}\pi \tau }{\pi \tau }\right|C(|uu_N|+N^\xi ),$$
where $`\rho (u)`$ is given by (1.11).
Proof. It follows from the formula (2.20) that
(3.14)
$$\frac{1}{N\rho (u)}𝒦_N(u_N,u_N+\frac{\tau }{N\rho (u)};y)=N_\gamma \frac{dz}{2\pi i}_\mathrm{\Gamma }\frac{dw}{2\pi i}h(z,w)g_N(z,w)e^{N(f_N(w)f_N(z))},$$
where $`g_N(z,w)`$ is given by (2.21),
$$f_N(z)=\frac{1}{2a^2}(z^22u_Nz)+\frac{1}{N}\underset{j=1}{\overset{N}{}}\mathrm{log}(zy_j)$$
and
$$h(z,w)=\frac{e^{\omega _0\tau }}{\tau }\left(e^{\tau w/a^2\rho (u)}e^{\tau (wz)/a^2\rho (u)}\right)$$
We have taken $`\omega =\omega _0/N\rho (u)`$, where $`\omega _0`$ is given by (3.23) below. The integral in (3.14) will be analyzed using a saddle point argument. It follows from (3.12) and Cauchy’s integral formula that there is a constant $`C`$ such that for all $`N1`$, $`\tau K`$, $`yY_{R/2,2\eta }`$ and $`|u|\sqrt{1/2+2a^2}`$,
(3.15) $`|f_N^{}(z)f^{}(z)|C(N^\xi +|uu_N|)`$
$`|f_N^{\prime \prime }(z)f^{\prime \prime }(z)|CN^\xi .`$
A computation shows that,
$$f^{}(z)=\frac{1}{a^2}(zu)+2(z\sqrt{z^21}).$$
Set $`S(w)=(w+1/w)/2`$ with inverse $`S^1(z)=z+\sqrt{z^21}`$, where $`\sqrt{z^21}=\sqrt{z1}\sqrt{z+1}`$ (principal argument). The function $`S`$ maps $`\{|w|>1\}`$ to $`[1,1]`$ and $`|w|=1`$ is mapped to $`[1,1]`$. Note that
$$f^{}(S(w))=\frac{w}{2a^2}+(2+\frac{1}{2a^2})\frac{1}{w}\frac{u}{a^2}.$$
Write $`u=\sqrt{1+4a^2}\mathrm{cos}\theta _c`$, where $`\theta _c[0,\pi ]`$. Our assumption on $`u`$ means that $`|\mathrm{cos}\theta _c|1/2`$. Note that $`f^{}(S(w))=0`$ has the solutions $`w_c^\pm =\sqrt{1+4a^2}\mathrm{exp}(\pm i\theta _c)`$. Hence the critical points for $`f`$ are $`z_c^\pm =S(w_c^\pm )`$.
We will now define some contours that we will use. Pick $`\delta >0`$ (small), see below. Set, for some $`ϵ>0`$ (small), $`\gamma _1^+(t)=S(\sqrt{1+4a^2}e^{i\delta }t)`$, $`\mathrm{}<t0`$, $`\gamma _2^+(t)=S(\sqrt{1+4a^2}e^{it})`$, $`\delta t\theta _cϵ`$, $`\gamma _3^+(t)=S(\sqrt{1+4a^2}e^{it})`$, $`\theta _cϵt\theta _c+ϵ`$, $`\gamma _4^+(t)=S(\sqrt{1+4a^2}e^{it})`$, $`\theta _c+ϵt\pi \delta `$ and $`\gamma _5^+(t)=S(\sqrt{1+4a^2}e^{i(\pi \delta )}t)`$, $`0t<\mathrm{}`$. Also, set $`\gamma _j^{}(t)=\overline{\gamma _j^+(t)}`$, $`1j5`$. Then, we can take $`\gamma =_{j=1}^5(\gamma _j^+\gamma _j^{})=\gamma ^+\gamma ^{}`$ in (3.14). Let $`t_0(1/\sqrt{1+4a^2},1)`$ be such that $`\text{Im }S(t_0w_c^+)=\eta `$, and write $`\alpha =\text{Re }S(t_0w_c^+)`$. Set, for some $`ϵ>0`$ (small), $`\mathrm{\Gamma }_1^+(t)=\alpha +it`$, $`0t\eta `$, $`\mathrm{\Gamma }_2^+(t)=S(tw_c^+)`$, $`t_0t1ϵ`$, $`\mathrm{\Gamma }_3^+(t)=S(tw_c^+)`$, $`1ϵt1+ϵ`$ and $`\mathrm{\Gamma }_4^+(t)=S(tw_c^+)`$, $`1+ϵt`$. Also, set $`\mathrm{\Gamma }_j^{}(t)=\overline{\mathrm{\Gamma }_j^+(t)}`$, $`1j4`$. We can then take $`\mathrm{\Gamma }=_{j=1}^4(\mathrm{\Gamma }_j^+\mathrm{\Gamma }_j^{})=\mathrm{\Gamma }^+\mathrm{\Gamma }^{}`$ in (3.14). Set
(3.16)
$$L_N^{bd}(\tau ;y)=N_{\gamma _3^b}\frac{dz}{2\pi i}_{\mathrm{\Gamma }_3^d}\frac{dw}{2\pi i}h(z,w)g_N(z,w)e^{N(f_N(w)f_N(z))},$$
where $`b,d\{+,\}`$ and write $`L_N=L_N^{++}L_N^+L_N^++L_N^{}`$.
###### Claim 3.3.
We can choose $`R_0>0`$, $`\eta _0>0`$ and $`ϵ,\delta >0`$, so that $`\gamma _3^++\gamma _3^{}+\mathrm{\Gamma }_3^++\mathrm{\Gamma }_3^{}`$ lies in a neighbourhood of $`z_c^\pm `$ which is included in $`\mathrm{\Omega }_{R_0/2,2\eta _0}`$ and for all $`N1`$, $`\tau K`$, $`yY_{r/2,2\eta }`$ and $`|u|\sqrt{1/2+2a^2}`$,
(3.17)
$$\left|\frac{1}{N\rho (u)}𝒦_N(u_N,u_N+\frac{\tau }{N\rho (u)};y)L_N(\tau ;y)\right|Ce^{cN}$$
with $`c>0`$
The claim will be proved below. We will now use the claim to finish the proof of lemma 3.2. It follows from (3.15) that there are critical points $`z_N^\pm =S(w_N^\pm )`$ for $`f_N(z)`$ such that
(3.18)
$$|z_N^\pm z_c^\pm |C(N^\xi +|uu_N|).$$
We can deform $`\gamma _3^\pm `$ ($`\mathrm{\Gamma }_3^\pm `$) into contours $`\gamma _N^\pm `$ ($`\mathrm{\Gamma }_N^\pm `$) such that the endpoints are unchanged, $`\gamma _N^\pm (0)=\mathrm{\Gamma }_N^\pm (0)=z_N^\pm `$ and $`\gamma _N^\pm `$ ($`\mathrm{\Gamma }_N^\pm `$) have $`C^1`$-distance $`C(N^\xi +|uu_N|)`$ to $`\gamma _3^\pm `$ ($`\mathrm{\Gamma }_3^\pm `$). We can also asume that these contours are chosen so that $`\gamma _N^\pm (t)=S(w_N^\pm e^{\pm it})`$ and $`\mathrm{\Gamma }_N^\pm (t)=S(w_N^\pm (1+t))`$ for $`|t|ϵ`$.
We can now proceed in the standard way with a local saddle point argument in (3.16) and prove that there is a constant $`C`$ such that
(3.19) $`\left|L_N^{bd}(\tau ;y)h(z_N^b,z_N^d)g_N(z_N^b,z_N^d){\displaystyle \frac{2\pi }{(2\pi i)^2}}{\displaystyle \frac{(\gamma _N^b)^{}(0)(\mathrm{\Gamma }_N^d)^{}(0)e^{N(f_N(z_N^b)f_N(z_N^d))}}{\sqrt{f_N^{\prime \prime }(z_N^b)(\gamma _N^b)^{}(0)^2}\sqrt{f_N^{\prime \prime }(z_N^d)(\mathrm{\Gamma }_N^d)^{}(0)^2}}}\right|`$
$`{\displaystyle \frac{C}{\sqrt{N}}}`$
for all $`N1`$, $`\tau K`$, $`yY_{R_0,\eta _0}`$ and $`|u|\sqrt{1/2+2a^2}`$. Note that $`z_N^+=\overline{z_N^{}}`$ and $`f_N(z_N^+)f_N(z_N^{})`$ is purely imaginary. Now, $`(\gamma _N^b)^{}(0)=biS^{}(w_N^b)`$, $`(\mathrm{\Gamma }_N^b)^{}(0)=w_N^bS^{}(w_N^b)`$ and a computation shows that
$$f_N^{\prime \prime }(z_N^b)(\gamma _N^b)^{}(0)^2=f_N^{\prime \prime }(z_N^b)(\mathrm{\Gamma }_N^b)^{}(0)^2=f_N^{\prime \prime }(z_N^b)S^{}(w_N^b)^2(w_N^a)^2,$$
which has a positive real part by (3.15) and the fact that $`f^{\prime \prime }(z_c^b)S^{}(w_c^b)^2(w_c^a)^2`$ has a positive real part. From (2.21) we see that $`g_N(z_N^b,z_N^d)=0`$ if $`bd`$ and $`g_N(z_N^b,z_N^b)=f_N^{\prime \prime }(z_N^b)`$. It follows that
$$\frac{g_N(z_N^b,z_N^b)(\gamma _N^b)^{}(0)(\mathrm{\Gamma }_N^b)^{}(0)}{\sqrt{f_N^{\prime \prime }(z_N^b)(\gamma _N^b)^{}(0)^2}\sqrt{f_N^{\prime \prime }(z_N^b)(\mathrm{\Gamma }_N^b)^{}(0)^2}}=bi.$$
Also, from (3.18) it follows that $`|h(z_N^b,z_N^b)h(z_c^b,z_c^b)|C(N^\xi +|uu_N|)`$, and thus (3.19) yields
(3.20)
$$|L_N^{bd}(\tau ;y)|\frac{C}{\sqrt{N}}$$
if $`bd`$ and
(3.21)
$$\left|L_N^{bb}(\tau ;y)+\frac{bh(z_c^b,z_c^b)}{2\pi i}\right|C(N^\xi +|uu_N|).$$
Combining (3.16), (3.20) and (3.21) we obtain
(3.22)
$$\left|L_N(\tau ;y)+\frac{h(z_c^+,z_c^+)h(z_c^{},z_c^{})}{2\pi i}\right|C(N^\xi +|uu_N|).$$
Now,
$$h(z_c^\pm ,z_c^\pm )=\frac{e^{\omega _0\tau }}{\tau }\left(e^{\tau z_c^\pm /a^2\rho (u)}1\right)$$
and a computation shows that
(3.23)
$$\frac{z_c^\pm }{a^2\rho (u)}=\pi \frac{1+2a^2}{2a^2}\mathrm{cot}\theta _c\pm \pi i\omega _0\pm \pi i.$$
Thus (3.22) becomes
$$\left|L_N(\tau ;y)\frac{\mathrm{sin}\pi \tau }{\pi \tau }\right|C(N^\xi +|uu_N|).$$
If we combine this estimate with (3.17) we see that the lemma is proved. $`\mathrm{}`$
It remains to prove claim 3.3.
Proof. Let $`\gamma _{}^\pm =_{j3}\gamma _j^\pm `$ and $`\mathrm{\Gamma }_{}^\pm =_{j3}\mathrm{\Gamma }_j^\pm `$. We have to estimate
$$I_1^{bd}=N_{\gamma _{}^b}|dz|_{\mathrm{\Gamma }^d}|dw||h(z,w)||g_N(z,w)|e^{N\text{Re }(f_N(w)f_N(z_c^d))N\text{Re }(f_N(z)f_N(z_c^b))},$$
and
$$I_2^{bd}=N_{\gamma ^b}|dz|_{\mathrm{\Gamma }_{}^d}|dw||h(z,w)||g_N(z,w)|e^{N\text{Re }(f_N(w)f_N(z_c^d))N\text{Re }(f_N(z)f_N(z_c^b))},$$
where $`b,d\{+,\}`$. Note that $`f_N(z_c^+)f_N(z_c)`$ is purely imaginary. We will concentrate on $`I_1^{++}`$ since the other cases are similar.
Using the inequality
$$\left|\frac{wy_j}{zy_j}\right|=\left|1+\frac{wz}{zy_j}\right|1+C(|w|+|z|)$$
it is not difficult to see that there are constants $`C_1`$ and $`C_2`$ such that
(3.24)
$$|h(z,w)||g_N(z,w)|e^{N\text{Re }(f_N(w)f_N(z))}C_1E^{C_2N(|z|+|w|)+N(\text{Re }(w^22uw)\text{Re }(z^22uz))/2a^2}$$
for all $`y^N`$, $`\tau K`$ and $`|u|\sqrt{1/2+2a^2}`$. Note that $`|\text{Im }z|c>0`$ for all $`z\gamma `$. (The constant $`c`$ depends on the $`\delta `$ in the definition of $`\gamma `$, but as we will see below $`\delta `$ depends only on the parameter $`a`$ in the problem.) From the estimate (3.24) it follows that by picking $`R=R_0`$ sufficiently large, the contribution to $`I_1^{++}`$ from $`z`$ and/or $`w`$ outside $`\mathrm{\Omega }_{R_0,0}`$ is $`e^N`$. Thus we can assume that $`z,w\mathrm{\Omega }_{R_0,0}`$. Next, we will derive the other estimates we will need to prove the claim.
Assume that $`z\mathrm{\Omega }_{R_0,\eta }`$ and $`w\mathrm{\Gamma }_1^+`$. Then,
$`|g_N(z,w)e^{Nf_N(w)}|`$
$`C\left(1+{\displaystyle \frac{1}{N}}{\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \frac{1}{|wy_k|}}\right){\displaystyle \underset{j=1}{\overset{N}{}}}|wy_j|e^{N\text{Re }(w^22uw)/2a^2}`$
$`C\left(1+{\displaystyle \frac{1}{N}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{1}{|\alpha +i\eta y_j|}}\right){\displaystyle \underset{j=1}{\overset{N}{}}}|\alpha +i\eta y_j|e^{N\text{Re }(w^22uw)/2a^2}`$
$`Ce^{N[\text{Re }f_N(\alpha +i\eta )+\text{Re }(w^22uw)((\alpha +i\eta )^22u(\alpha +i\eta ))]/2a^2}.`$
If we use (3.12) and the definition of $`f_N`$ we obtain
$`|g_N(z,w)e^{N(f_N(w)f_N(z_c^+))}|`$
(3.25) $`Ce^{cN(N^\xi +|uu_N|)+N\eta ^2/2a^2+N\text{Re }(f(\alpha +i\eta )f(z_c^+))/2a^2}`$
for $`z\mathrm{\Omega }_{R_0,\eta }`$ and $`w\mathrm{\Gamma }_1^+`$.
We will now compute how $`\text{Re }f(z)`$ changes along $`\gamma `$. Assume that $`\theta _c0`$, the other case is analogous. Consider $`\gamma (\theta )=S(\sqrt{1+4a^2}e^{i\theta })`$, $`\delta \theta \pi \delta `$. A computation, using the fact that $`f^{}(\gamma (\theta _c))=0`$ gives $`\text{Re }\frac{d}{d\theta }f(\gamma (\theta ))=\frac{1+2a^2}{2a^2}\mathrm{sin}\theta (\mathrm{cos}\theta _c\mathrm{cos}\theta )`$. From this we see that there is a constant $`c_0>0`$ such that
(3.26)
$$\text{Re }(f(\sqrt{1+4a^2}e^{i\theta })f(z_c^+))c_0(\theta \theta _c)^2.$$
Next, consider $`\gamma _1(t)=S(\sqrt{1+4a^2}e^{i\delta }t)`$, $`t0`$. If we write $`\omega _\delta =\sqrt{1+4a^2}e^{i\delta }`$, then
$$\frac{d}{dt}f(\gamma _1(t))=\frac{1}{4a^2}[\omega _\delta t2u+\frac{1+4a^2}{\omega _\delta t}][1\frac{1}{(\omega _\delta t)^2}].$$
Set $`\omega _\delta t=s(t)e^{i\theta (t)}`$. A computation shows that
(3.27) $`\text{Re }{\displaystyle \frac{d}{dt}}f(\gamma _1(t))={\displaystyle \frac{1}{4a^2\sqrt{1+4a^2}}}\{[(s(t)+{\displaystyle \frac{1}{s(t)}})\mathrm{cos}\theta (t)2\mathrm{cos}\theta _c]`$
$`\times [1+4a^2{\displaystyle \frac{1}{s(t)^2}}\mathrm{cos}2\theta (t)]{\displaystyle \frac{1}{s(t)^2}}\mathrm{sin}2\theta (t)(s(t){\displaystyle \frac{1}{s(t)}})\mathrm{sin}\theta (t)\}.`$
Note that $`\mathrm{sin}\theta (t)=s(t)^1\sqrt{1+4a^2}\mathrm{sin}\delta `$. It follows that the right hand side of (3.27) equals
(3.28) $`{\displaystyle \frac{1}{4a^2\sqrt{1+4a^2}}}\{[(s(t)+{\displaystyle \frac{1}{s(t)}})(1+4a^2{\displaystyle \frac{1}{s(t)^2}}+2{\displaystyle \frac{(1+4a^2)\mathrm{sin}^2\delta }{s(t)^4}})`$
$`2{\displaystyle \frac{(1+4a^2)\mathrm{sin}^2\delta }{s(t)^4}}(s(t){\displaystyle \frac{1}{s(t)}})]\mathrm{cos}\theta (t)2(1+4a^2{\displaystyle \frac{1}{s(t)^2}}+2{\displaystyle \frac{(1+4a^2)\mathrm{sin}^2\delta }{s(t)^4}})\mathrm{cos}\theta _c\}`$
and this is
$$\frac{1}{4a^2\sqrt{1+4a^2}}(1+4a^2\frac{1}{s(t)^2})\left[(s(t)+\frac{1}{s(t)})\mathrm{cos}\theta (t)2(1+\frac{1+4a^2}{2a^2}\mathrm{sin}^2\delta )\mathrm{cos}\theta _c\right],$$
since $`s(t)1`$. Choose $`\delta \theta _c/4`$ so that
$$(1+\frac{1+4a^2}{2a^2}\mathrm{sin}^2\delta )\mathrm{cos}\theta _c\mathrm{cos}\frac{\theta _c}{2}.$$
Since $`s(t)+1/s(t)2`$ and $`\theta (t)\delta `$ we see that there is a constant $`c_0>0`$ such that
(3.29)
$$\text{Re }\frac{d}{dt}f(\gamma _1(t))c_0.$$
For $`\gamma _5(t)=S(\sqrt{1+4a^2}e^{i(\pi \delta )}t)`$, $`t0`$, we still have the formula (3.28) with $`\gamma _1(t)`$ replaced by $`\gamma _5(t)`$ and, since $`\pi \delta \theta (t)\pi `$, we see that the right hand side is
(3.30)
$$\frac{1}{\sqrt{1+4a^2}}[(s(t)+\frac{1}{s(t)})\mathrm{cos}(\pi \theta (t))+2\mathrm{cos}\theta _c]$$
and consequently there is a constant $`c_0>0`$ such that
(3.31)
$$\text{Re }\frac{d}{dt}f(\gamma _5(t))c_0.$$
Consider now how $`\text{Re }f(w)`$ changes along $`\mathrm{\Gamma }^+`$. Set $`\mathrm{\Gamma }(t)=S(tw_c^+)`$, $`tt_0`$. A computation gives
$$\text{Re }\frac{d}{dt}f(S(tw_c))=\frac{1t}{2a^2t^2}[1+t(1+4a^2)(t^2(1+4a^2)+\frac{1}{t})\mathrm{cos}2\theta _c].$$
Now, since $`|u|\sqrt{1/2+2a^2}`$, it follows that $`\mathrm{cos}2\theta _c0`$ and thus
$`\text{Re }{\displaystyle \frac{d}{dt}}f(S(tw_c)){\displaystyle \frac{1t}{2a^2t^2}}(1+t(1+4a^2))\text{if }t_0t1`$
(3.32) $`\text{Re }{\displaystyle \frac{d}{dt}}f(S(tw_c)){\displaystyle \frac{1t}{2a^2t^2}}(1+t(1+4a^2))\text{if }t1.`$
The first of these estimates can be used to show that if we pick $`\eta =\eta _0`$ sufficiently small, then
$$\eta ^2+\text{Re }(f(\alpha +i\eta )f(z_c^+))c_0$$
for some positive $`c_0`$. If we use this in (3) we obtain
(3.33)
$$|g_N(z,w)e^{N(f_N(w)f_N(z_c^+))}|Ce^{c_0^{}N}$$
for some positive $`c_0^{}`$. We can now use (3.26), (3.29), (3.31), (3) and (3.33) to estimate $`I_1^{++}`$ and see that it is $`Ce^{cN}`$ for some positive $`c`$. $`\mathrm{}`$
## 4. Proof of the theorems
We start with the proof of theorem 1.2.
Proof. By proposition 1.1 and Fubini’s theorem the integral in the left hand side of (1.14) can be written
(4.1)
$$__N\left(_^N\rho _N(x,y(H))(Sf)(N\rho (u)(x_1u),\mathrm{},N\rho (u)(x_Nu))d^Nx\right)𝑑P^{(N)}(H)$$
Note that $`S(f)_{\mathrm{}}N^mf_{\mathrm{}}`$. Since $`\rho _N(x,)`$ is a probability density on $`^N`$ we can use lemma 3.1 to replace the expression in (4.1) by
(4.2)
$$__N\left(_^N\rho _N(x,y(H))(Sf)(N\rho (u)(x_1u),\mathrm{},N\rho (u)(x_Nu))d^Nx\right)𝑑\stackrel{~}{P}^{(N)}(H)$$
with an error $`CN^mf_{\mathrm{}}N^{2p(1/2\xi )}=o(1)`$, since $`p>2(m+2)`$, provided we choose $`\xi `$ small enough. Now, since $`\rho _N(x,)`$ is symmetric it follows from (1.13), (2.4), (2.14) and (2.19) that the expression in (4.2) can be written
(4.3) $`{\displaystyle __N}{\displaystyle _^m}`$ $`f(t_1,\mathrm{},t_m)`$
$`\times det({\displaystyle \frac{1}{N\rho (u)}}𝒦(u+{\displaystyle \frac{t_i}{N\rho (u)}},u+{\displaystyle \frac{t_j}{N\rho (u)}};y(H)))_{i,j=1}^md^mtd\stackrel{~}{P}^{(N)}(H).`$
Since $`f`$ has compact support and we know that (3.4) holds a.s. $`[\stackrel{~}{P}^{(N)}]`$ it follows from lemma 3.2, with $`u_N=u+t_i/N\rho (u)`$, $`\tau =t_jt_i`$, that
$$\left|𝒦(u+\frac{t_i}{N\rho (u)},u+\frac{t_j}{N\rho (u)};y(H))\frac{\mathrm{sin}\pi (t_it_j)}{\pi (t_it_j)}\right|CN^\xi ,$$
for a.a. $`[\stackrel{~}{P}^{(N)}]`$ and all $`(t_1,\mathrm{},t_m)`$ in the support of $`f`$. Thus we can take the limit as $`N\mathrm{}`$ in (4.3) and obtain the right hand side of (1.14). This completes the proof. $`\mathrm{}`$
Before proving theorem 1.3 we need some preliminary results on the level spacing distribution. Let $`\rho _N(x)`$ be a symmetric probability density on $`^N`$ with correlation functions defined by (1.1). Assume that $`R_1^{(N)}/N\rho (t)`$ (weakly) as $`N\mathrm{}`$, so that $`\rho (t)`$ is the asymptotic density. Let $`u`$ be a given point such that $`\rho (u)>0`$, and let $`t_N`$ be a sequence such that $`t_N\mathrm{}`$ but $`t_N/N0`$ as $`N\mathrm{}`$. Set, for $`|r|1/2`$,
$$_m^{(N)}(\sigma _1,\mathrm{},\sigma _m;r)=\frac{1}{(N\rho (u))^m}R_m^{(N)}(u+\frac{2t_Nr+\sigma _1}{N\rho (u)},u+\frac{2t_Nr+\sigma _m}{N\rho (u)})$$
and let $`_m(\sigma _1,\mathrm{},\sigma _m)`$ be the limiting correlation functions, which we assume are continuous, symmetric and translation invariant. Assume that, for each $`s0`$,
(4.4)
$$D_N(s)=\underset{m=N+1}{\overset{\mathrm{}}{}}\frac{s^m}{m!}\underset{|\sigma _j|s}{sup}|_m(\sigma _1,\mathrm{},\sigma _m)|<\mathrm{}.$$
Set
$$H(s)=\underset{m=0}{\overset{\mathrm{}}{}}\frac{(1)^m}{m!}_{[0,s]^m}_m(\sigma _1,\mathrm{},\sigma _m)d^m\sigma $$
(the probability of no particle in $`[0,s]`$), which is well defined by (4.4). Also, set
(4.5)
$$ϵ_m^{(N)}=\underset{|\sigma _j|s,|r|1/2}{sup}|_m^{(N)}(\sigma _1,\mathrm{},\sigma _m;r)_m(\sigma _1,\mathrm{},\sigma _m)|.$$
###### Proposition 4.1.
Let $`S_N(s,x)`$ be defined by (1.5). Then
(4.6)
$$\left|_^NS_N(s,x)\rho _N(x)d^Nx_0^sH^{\prime \prime }(u)𝑑u\right|D_N(s)+\underset{m=2}{\overset{N}{}}\frac{s^{m1}}{(m1)!}ϵ_m^{(N)}.$$
Proof. We first show that
(4.7)
$$_0^sH^{\prime \prime }(u)𝑑u=\underset{m=2}{\overset{N}{}}\frac{s^{m1}}{(m1)!}_{[0,s]^{m1}}_m(0,\tau _2,\mathrm{},\tau _m)𝑑\tau _2\mathrm{}𝑑\tau _m,$$
see . Since $`_m`$ is translation invariant and symmetric by assumption, we have
(4.8) $`H^{}(u)`$ $`=\underset{ϵ0}{lim}{\displaystyle \frac{1}{ϵ}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}{\displaystyle _{[ϵ,u]^m[0,u]^m}}_m(x_1,\mathrm{},x_m)d^mx`$
$`=\underset{ϵ0}{lim}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}{\displaystyle \frac{1}{ϵ}}\left(m{\displaystyle _{[ϵ,0]\times [0,u]^{m1}}}_m(x_1,\mathrm{},x_m)d^mx\right)`$
$`={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{(m1)!}}{\displaystyle _{[0,u]^{m1}}}_m(0,x_2,\mathrm{},x_m)d^{m1}x,`$
where we have also used (4.4) and the continuity of $`_m`$. Continuing in the same way we see that $`H(u)`$ is actually a $`C^{\mathrm{}}`$ function, in particular $`H^{\prime \prime }(u)`$ is well defined and continuous. From (4.8) we get
$$H^{}(s)=_m(0)+\underset{m=2}{\overset{\mathrm{}}{}}\frac{(1)^m}{(m1)!}_{[0,s]^{m1}}_m(0,x_2,\mathrm{},x_m)d^{m1}x.$$
Hence $`H^{}(0)=_m(0)`$ and we see that the right hand side of (4.7) equals $`H^{}(u)H^{}(0)`$, which is what we wanted to prove.
It is proved in , using a result from , that
$`{\displaystyle _^N}S_N(s,x)\rho _N(x)d^Nx`$
$`={\displaystyle \underset{m=2}{\overset{N}{}}}{\displaystyle \frac{(1)^m}{(m1)!}}{\displaystyle _{1/2}^{1/2}}𝑑r{\displaystyle _{[0,\mathrm{min}(s,(12r)t_N)]^{m1}}}_m^{(N)}(0,\sigma _2,\mathrm{},\sigma _m;r)d^{m1}\sigma .`$
Hence, the estimate (4.6) follows from (4.4), (4.5) and (4.7). $`\mathrm{}`$
We turn now to the proof of theorem 1.3.
Proof. Just as in the proof of theorem 1.2 above we see that since $`P𝒲^{6+ϵ}`$ and $`S_N_{\mathrm{}}N/2t_N`$,
(4.9) $`\left|{\displaystyle __N}S_N(s,x(M))𝑑Q^{(N)}(M){\displaystyle __N}\left({\displaystyle _^N}S_N(s,x)\rho _N(x;y(H))d^Nx\right)𝑑\stackrel{~}{P}^{(N)}(H)\right|`$
$`C{\displaystyle \frac{N}{t_N}}N^{2(6+ϵ)(1/2\xi )}{\displaystyle \frac{C}{t_N}},`$
if we take $`\xi `$ sufficiently small, and also that (3.4) holds. From proposition 1.1, (2.4) and proposition 2.3 we know the correlation functions of $`\rho _N(x;y)`$, and if we take $`u_N=u+(2t_Nr+\sigma _i)(N\rho (u))^1`$ in lemma 3.2 we see that
(4.10) $`\left|{\displaystyle \frac{1}{N\rho (u)}}𝒦(u+{\displaystyle \frac{2t_Nr+\sigma _i}{N\rho (u)}},u+{\displaystyle \frac{2t_Nr+\sigma _j}{N\rho (u)}};y(H)){\displaystyle \frac{\mathrm{sin}\pi (\sigma _i\sigma _j)}{\pi (\sigma _i\sigma _j)}}\right|`$ $`C({\displaystyle \frac{t_N}{N}}+N^\xi )`$
$`\omega _N`$
for a.a. $`H`$ $`[\stackrel{~}{P}^{(N)}]`$. Thus, the limiting correlation functions are
$$_m(\sigma _1,\mathrm{},\sigma _m)det\left(\frac{\mathrm{sin}\pi (\sigma _i\sigma _j)}{\pi (\sigma _i\sigma _j)}\right)_{i,j=1}^m.$$
Since the matrix in the determinant is positive definite it follows from the Hadamard inequality that
$$D_N(s)\underset{m=N+1}{\overset{\mathrm{}}{}}\frac{s^m}{m!}.$$
Also, since
$$_m^{(N)}(\sigma _1,\mathrm{},\sigma _m;y)=det\left(\frac{1}{N\rho (u)}𝒦(u+\frac{2t_Nr+\sigma _i}{N\rho (u)},u+\frac{2t_Nr+\sigma _j}{N\rho (u)};y)\right)_{i,j=1}^m$$
it follows from (4.10), the multilinearity of the determinant and Hadamard’s inequality that
$$|_m^{(N)}(\sigma ;y)_m(\sigma )|m(1+\omega _N)^{m1}\omega _Nm^{m/2},$$
and hence $`ϵ_m^{(N)}m(1+\omega _N)^{m1}\omega _Nm^{m/2}`$. Now, by proposition 4.1, Stirling’s formula and the fact that $`\omega _N0`$,
(4.11) $`\left|{\displaystyle _^N}S_N(s,x)\rho _N(x;y(H))d^Nx{\displaystyle _0^s}H^{\prime \prime }(u)𝑑u\right|`$
$`{\displaystyle \underset{m=N+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{s^m}{m!}}+\omega _N{\displaystyle \underset{m=2}{\overset{N}{}}}{\displaystyle \frac{s^m}{(m1)!}}(1+\omega _N)^{m1}m^{(m+2)/2}=o(1)`$
as $`N\mathrm{}`$, for a.a. $`H`$ $`[\stackrel{~}{P}^{(N)}]`$. If we combine (4.9) and (4.11) we see that the theorem is proved. $`\mathrm{}`$
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# Nonlinear wave equation, nonlinear Riemann problem, and the twistor transform of Veronese webs
## 0. Introduction
We denote the $`d`$-dimensional coordinate vector space over the base field by $`𝕍^d`$, and the corresponding projective space by $`^{d1}`$. For $`(v_1,\mathrm{},v_d)𝕍^d\left\{0\right\}`$ we denote by $`\left(v_1:\mathrm{}:v_d\right)`$ the corresponding element of $`^{d1}`$. By $`𝔹_r^d𝕍^d`$ we denote the open ball of radius $`r`$ centered at the origin. Then $`\left(𝔹_r^1\right)^d`$ is a cube in the real case and a polydisk in the complex case.
As a convention, put $`|\mathrm{}|=\mathrm{}`$, so that $`\left\{|z|>1\right\}`$ includes $`z=\mathrm{}`$.
The word “smooth” can have 3 different meanings: in the case of the base field $``$ it can mean either $`C^{\mathrm{}}`$-smooth or real-analytic, in the case of the base field $``$ it means complex-analitic. When only some of these cases work, we use more specific terms.
In this paper we study a special family of nonlinear wave equations. Elements of this family are parameterized by numbers $`A,B,C`$ which satisfy
(0.1)
$$A0\text{}B0\text{}C0\text{}A+B+C=0.$$
Given such numbers, the equation is
$$Aw_xw_{yz}+Bw_yw_{xz}+Cw_zw_{xy}=0;$$
here $`w(x,y,z)`$ is a function of three variables, If we need to specify $`A,B,C`$, we may call this equation the $`(A,B,C)`$-equation. Whenever we mention an $`(A,B,C)`$-equation we assume that $`A,B,C`$ satisfy (0.1).
In this paper we study only those solutions of $`(A,B,C)`$-equations which are in general position, according the the following
###### Definition 0.1.
Say that a function $`w(x,y,z)`$ is non-degenerate if $`w_x0`$, $`w_y0`$, $`w_z0`$ whenever $`w(x,y,z)`$ is defined.
###### Definition 0.2.
Say that two functions $`w(x,y,z)`$ and $`w^{}(x,y,z)`$ are gauge transforms of each other, if $`w=\tau w^{}`$ for an appropriate invertible scalar function $`\tau `$ of one variable.
The first target of this paper is the following statement:
###### Theorem 0.3.
Suppose that triples $`(A,B,C)`$ and $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$ satisfy conditions (0.1). Consider equations
(0.2) $`Aw_xw_{yz}+Bw_yw_{xz}+Cw_zw_{xy}`$ $`=0,`$
(0.3) $`\stackrel{~}{A}v_xv_{yz}+\stackrel{~}{B}v_yv_{xz}+\stackrel{~}{C}v_zv_{xy}`$ $`=0,`$
the system of equations
(0.4) $`A\stackrel{~}{B}w_xv_y`$ $`=\stackrel{~}{A}Bw_yv_x`$
$`A\stackrel{~}{C}w_xv_z`$ $`=\stackrel{~}{A}Cw_zv_x,`$
and the equation (with $`\alpha =A/\stackrel{~}{A}`$, $`\beta =B/\stackrel{~}{B}`$, $`\gamma =C/\stackrel{~}{C}`$)
(0.5)
$$(v_x,v_y,v_z)(\alpha w_x,\beta w_y,\gamma w_z),$$
here for two vector-functions we write $`ab`$ if $`a(x,y,z)=\psi (x,y,z)b(x,y,z)`$ for an appropriate nowhere-0 scalar function $`\psi `$. Then locally near $`(x,y,z)=(0,0,0)`$
1. For non-degenerate functions $`w`$, $`v`$ System (0.4) is equivalent to Equation (0.5);
2. Given a solution $`(w,v)`$ of System (0.4) with non-degenerate $`w`$ and $`v`$ and any gauge transforms $`w_1`$ of $`w`$ and $`v_1`$ of $`v`$ the pair $`(w_1,v_1)`$ is a solution of System (0.4);
3. Suppose that $`A\stackrel{~}{B}\stackrel{~}{A}B.`$ Given a solution $`(w,v)`$ of System (0.4) with non-degenerate $`w`$ and $`v`$, the function $`w`$ satisfies Equation (0.2), the function $`v`$ satisfies Equation (0.3);
4. Suppose that $`A\stackrel{~}{B}\stackrel{~}{A}B.`$ Given a non-degenerate solution $`w`$ of Equation (0.2), there is a non-degenerate function $`v`$ such that the pair $`(w,v)`$ satisfies Equations (0.4). As a corollary, $`v`$ satisfies Equation (0.3);
5. Such a function $`v`$ is defined uniquely up to a gauge transform.
Theorem 0.3 is proved in Section 4. While one could prove this theorem purely analytically, we emphasize the geometric meaning of its statements, thus prove it via relationship to $`3`$-dimensional Veronese webs, which are introduced in Sections 1 and 2.
###### Remark 0.4.
One can consider the last two statements of Theorem 0.3 as statements about existence of non-pointwise relationship between Equations (0.2) and (0.3). Given a solution $`w`$ of Equation (0.2), one obtains a (more or less unique) solution $`v`$ of Equation (0.3) by solving Equations (0.4). Note that the latter equations are equations of lower order than (0.3) when considered as equations in $`v`$.
In other words, System (0.4) provides a Bäcklund–Darboux transform of order 1 between two equations (0.2) and (0.3) of order 2. Moreover, this transform is linear in $`v`$.
The second target of this paper is to explicitly solve any $`(A,B,C)`$-equation in complex domain in terms of the nonlinear Riemann problem. This problem is a straightforward nonlinear analogue of the (linear) Riemann conjugation problem:
###### Definition 0.5.
Consider a complex-analytic function $`g(\lambda ,t)`$ defined for $`\epsilon <|\lambda |<1/\epsilon `$ and $`|t|<\delta `$, assume that for any given $`\lambda `$, $`\epsilon <|\lambda |<1/\epsilon `$, the function $`tg(\lambda ,t)`$ is invertible. Suppose that equations
$$\sigma _{}\left(\lambda \right)=g(\lambda ,\sigma _+\left(\lambda \right))\text{ for }\epsilon <|\lambda |<1/\epsilon ,|\sigma _+\left(\lambda \right)|<\delta \text{ for }|\lambda |<1/\epsilon ,$$
uniquely determine complex-analytic functions $`\sigma _+\left(\lambda \right)`$ defined for $`|\lambda |<1/\epsilon `$, and $`\sigma _{}\left(\lambda \right)`$ defined for $`|\lambda |>\epsilon `$. Denote the number $`\sigma _+\left(0\right)`$ by $`_{\epsilon \delta }\left(g\right)`$.
The function $`_{\epsilon \delta }`$ sends a function $`g(\lambda ,t)`$ of two variables to a complex number. Call this function the non-linear Riemann transform. Note that changing $`\epsilon `$ and $`\delta `$ cannot change the value of $`_{\epsilon \delta }\left(g\right)`$ (though this expression can become undefined), thus we are going to drop $`\epsilon `$, $`\delta `$ and denote this function by $``$.
The next step is to define a special family $`g_{x,y,z}(\lambda ,t)`$ of functions of two variables, given one function $`g(\lambda ,t)`$. Use notation $`F_{\text{M},k}\left(\lambda \right)`$ for Lagrange interpolation polynomials on points $`\text{M}=\{\mu _1,\mathrm{},\mu _m\}`$:
$$F_{\text{M},l}\left(\lambda \right)=F_{\text{M}\left\{\mu _l\right\}}\left(\lambda \right)/F_{\text{M}\left\{\mu _l\right\}}\left(\mu _l\right),F_\text{M}\left(\lambda \right)=\underset{\mu \text{M}}{}\left(\lambda \mu _n\right).$$
###### Definition 0.6.
Consider sets of $`k`$ numbers $`\mathrm{\Lambda }=\{\lambda _1,\mathrm{},\lambda _k\}`$ and of $`m`$ numbers $`\text{M}=\{\mu _1,\mathrm{},\mu _m\}`$ satisfying $`|\lambda _l|>0`$, $`|\mu _l|>0`$. Consider a function $`g(\lambda ,t)`$. Let $`\lambda _0=0`$, $`\mathrm{\Lambda }_0=\mathrm{\Lambda }\left\{0\right\}`$, $`F_+=F_{\mathrm{\Lambda }_0,0}`$, $`F_{+,l}=F_{\mathrm{\Lambda }_0,l}`$, $`F_{}\left(\lambda \right)=F_\text{M}\left(\lambda \right)/\lambda ^m`$, $`F_{,l}=F_{\text{M},l}\left(\lambda \right)\mu _l^{m1}/\lambda ^{m1}`$. For collections $`\left\{a_i\right\}`$ and $`\left\{b_i\right\}`$ of $`k`$ and $`m`$ numbers correspondingly denote
$$𝒢_{\mathrm{\Lambda }\text{M},\left\{a_i\right\}\left\{b_i\right\}}(\lambda ,t)=F_{}\left(\lambda \right)^1\left(g(\lambda ,\stackrel{~}{t})\underset{l=1}{\overset{m}{}}b_lF_{,l}\left(\lambda \right)\right),\stackrel{~}{t}=tF_+\left(\lambda \right)+\underset{l=1}{\overset{k}{}}a_lF_{+,l}\left(\lambda \right).$$
Given 3 numbers $`\lambda _1,\lambda _2,\lambda _3`$, let $`g_{x,y,z}(\lambda ,t)\stackrel{\text{def}}{=}𝒢_{\{\lambda _1,\lambda _2\}\left\{\lambda _3\right\},\{x,y\}\left\{z\right\}}(\lambda ,t)`$.
Given a function $`\phi \left(\lambda \right)`$, $`|\lambda |=1`$, define $`\mathrm{ind}\phi `$ as $`\frac{1}{2\pi i}_{|\lambda |=1}\frac{d\phi \left(\lambda \right)}{\phi \left(\lambda \right)}`$.
###### Theorem 0.7.
Consider a complex-analytic function $`g(\lambda ,t)`$ defined for $`\epsilon <|\lambda |<1/\epsilon `$ and $`|t|<\delta `$, such that $`g(\lambda ,0)0`$ and $`\mathrm{ind}\frac{g}{t}(\lambda ,0)=2`$. Fix $`0<r<1`$, $`\lambda _1,\lambda _2,\lambda _3^1`$, $`0<|\lambda _{1,2}|<r`$, $`|\lambda _3|>1/r`$. Then
1. the function
(0.6)
$$w(x,y,z)=\left(g_{x,y,z}\right)$$
is correctly defined for small $`x`$, $`y`$, $`z`$, is complex-analytic and nondegenerate, and satisfies the equation (0.2) with
(0.7)
$$A=\lambda _1\left(\lambda _2\lambda _3\right)\text{}B=\lambda _2\left(\lambda _3\lambda _1\right)\text{}C=\lambda _3\left(\lambda _1\lambda _2\right);$$
2. for any scalar functions $`\psi `$, $`\phi _1`$, $`\phi _2`$, $`\phi _3`$ of one variable which send 0 to 0 the function $`\widehat{w}=\psi \left(w(\phi _1\left(x\right),\phi _2\left(y\right),\phi _3\left(z\right))\right)`$ satisfies the same $`(A,B,C)`$-equation as $`w(x,y,z)`$;
3. for any triple $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$ which satisfies conditions (0.1) one can find $`\lambda _1,\lambda _2,\lambda _3^1`$ and $`T0`$ which satisfy the above inequalities and Equation (0.7) with $`A=T\stackrel{~}{A}`$, $`B=T\stackrel{~}{B}`$, $`C=T\stackrel{~}{C}`$;
4. for any nondegenerate complex-analytic solution $`\widehat{w}(x,y,z)`$ of (0.2) defined near (0,0,0) the function $`g(\lambda ,t)`$ constructed in Theorem 10.1 (for some particular value of the function $`Y\left(x\right)`$) satisfies the conditions above, and $`\widehat{w}=\psi \left(w(x,\phi _2\left(y\right),z)\right)`$; here $`w(x,y,z)`$ is defined by (0.6), $`\phi _2`$ is the inverse function to $`y=Y\left(x\right)`$, and $`\psi \left(t\right)=\widehat{w}(t,Y\left(t\right),0)`$.
This theorem is proved in Section 14.
###### Remark 0.8.
Note that Theorem 10.1 determines the gluing function $`g(\lambda ,t)`$ in terms of the values of $`w`$ and the normal derivative of $`w`$ on a hypersurface. Thus Theorem 0.7 can be considered as a procedure to solve Equation (0.2) basing on the Cauchy initial data.
Such an approach would not gain a lot if the nonlinear Riemann problem were complicated to solve. However, in Section 16 we are going to show that it is as complicated as solving an ODE of high dimension.
Plan. In Sections 1 and 2 we define Veronese webs. In Section 3 we show that constructing a $`3`$-dimensional Veronese web is equivalent to solving an $`(A,B,C)`$-equation. In Section 4 we prove Theorem 0.3, thus construct Bäcklund–Darboux transformations between different $`(A,B,C)`$-equations. In Sections 5 and 6 we show how Veronese webs jump into existence given the statement of Theorem 0.3.
Sections 7 and 8 contain first encounters with the twistor transform of the Veronese web. Although full of technical (and long but simple) statements, these sections enable working with the twistor transform as with a manifold (as opposed to a germ), thus remove many linguistic complications. In Section 9 we introduce convenient coordinate systems on the twistor transform, in Section 10 we describe the gluing functions as solutions of appropriate ODEs.
Section 11 starts dealing with the inverse problem of reconstructing the web by its twistor transform. After recalling what are infinitesimal deformations of submanifolds, we obtain the first solution of the inverse problem, the solution which requires a lot of additional data. Section 12 contains technical results which would allow to drop these additional data in complex-analytic cases: Kodaira–Spencer deformation theory for sections of bundles (Theorem 12.2), and the “inverse” theory (Proposition 12.8) which explicitly constructs a small tubular neighborhood in which the deformation theory works.
Section 13 studies in which cases the “additional data” of the inverse twistor transform can be dropped. We call such webs airy webs, and show that Veronese webs are airy. This section also provides an alternative heuristic for utility of so-called Kronecker webs introduced in : they are airy webs with the parameter space being $`^1`$.
Section 14 completes the full circle by proving Theorem 0.7, thus providing the explicit construction of the inverse twistor transform. Given a non-degenerate solution of the $`(A,B,C)`$-equation, Section 10 had shown how to explicitly calculate gluing functions for the twistor transform via solutions of ODEs. Section 14 shows how to use these gluing data for reconstruction of the initial solution of the $`(A,B,C)`$-equation. Theorem 0.7 provides a way to completely integrate the $`(A,B,C)`$-equation in the non-degenerate case.
The first Appendix (Section 15) connects the results of Section 10 with Turiel classification of Veronese webs of arbitrary dimension . Additionally, we introduce terms using which one can classify arbitrary airy webs of codimension 1. The second Appendix (Section 16) shows that the nonlinear Riemann problem is not harder to solve than Lipschitz ODEs in Hilbert spaces.
## 1. Webs
Recall the definition of a foliation.
###### Definition 1.1.
A prefoliation $``$ of codimension $`r`$ on a manifold $`M`$ is a representation of $`M`$ as a disjoint union of subsets called leaves, each of which is a connected embedded submanifold of codimension $`r`$.
Given an open subset $`UM`$, one can define a restriction $`|_U`$ of $``$ to $`U`$, the leaves of which are connected components of $`LU`$, $`L`$ running through leaves of $``$. Say that $``$ is direct if $`M=N\times F`$ with a connected $`F`$, and leaves are $`\left\{n\right\}\times F`$, $`nN`$. In such a case $`N`$ is called the base of $``$.
The tangent space $`𝒯_m`$ to $``$ at $`mM`$ is the tangent space $`𝒯_mL_m`$ to the leaf $`L_m`$ of $``$ through $`m`$, and the normal space $`𝒩_m`$ at $`mM`$ is $`𝒯_mM/𝒯_m`$. Cotangent space $`𝒯_m^{}`$ and conormal space $`𝒩_m^{}`$ at $`m`$ are defined as dual spaces to the tangent space and the normal space at $`m`$. Clearly, $`𝒩_m^{}`$ can be identified with the orthogonal complement $`\left(𝒯_m\right)^{}`$ to $`𝒯_m𝒯_mM`$ in $`𝒯_m^{}M`$.
###### Definition 1.2.
Say that a prefoliation $``$ is a foliation if every point $`mM`$ has a neighborhood $`U`$ such that $`|_U`$ is diffeomorphic to a direct prefoliation.
Obviously, tangent, cotangent, normal and conormal spaces to a foliation form vector bundles over $`M`$, and $`𝒯𝒯M`$, $`𝒩^{}𝒯^{}M`$ are vector subbundles.
###### Definition 1.3.
A web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ of codimension $`r`$ on a manifold $`M`$ is a family of foliations of codimension $`r`$ on $`M`$, one foliation $`_\lambda `$ per each $`\lambda \mathrm{\Lambda }`$. Say that a web is smooth if $`\mathrm{\Lambda }`$ is a manifold, and the vector subbundle $`𝒩^{}_\lambda 𝒯^{}M`$ depends smoothly on $`\lambda \mathrm{\Lambda }`$ (to be more precise, consider $`𝒩^{}_\lambda `$ as a section of the bundle of Grassmannians $`\mathrm{Gr}_r\left(𝒯^{}M\right)`$).
In what follows we use the shortcut $`_{}`$ for $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ when we are not interested in the set $`\mathrm{\Lambda }`$ of parameters of the web.
###### Definition 1.4.
Say that a web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$ is weakly separating if for any two points $`m_1,m_2M`$ there is $`\lambda \mathrm{\Lambda }`$ such that $`m_1`$ and $`m_2`$ are on different leaves of $`_\lambda `$. Say that a web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$ is weakly separating near $`mM`$ if $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}|_U`$ is weakly separating for an appropriate neighborhood $`Um`$.
Say that a web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ is separating at $`m`$ if for any tangent vector $`v𝒯_mM`$, $`v0`$, there is $`\lambda \mathrm{\Lambda }`$ such that $`v𝒯_m_\lambda `$.
## 2. $`3`$-dimensional Veronese webs
Recall the definition of a Veronese web ().
###### Definition 2.1.
Given a web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$ of codimension $`r`$, and a point $`mM`$, let $`\mathrm{n}_m\left(\lambda \right)𝒯_m^{}M`$, $`\lambda \mathrm{\Lambda }`$, be the normal subspace at $`m`$ to the leaf $`L`$ of $`_\lambda `$ which passes through $`m`$. In the case $`r=1`$ one can consider $`\mathrm{n}_m`$ as a mapping from $`\mathrm{\Lambda }`$ to the projectivization $`𝒯_m^{}M`$ of $`𝒯_m^{}M`$.
###### Definition 2.2.
A Veronese web on a manifold $`M`$ is a smooth separating web $`\left\{_\lambda \right\}_{\lambda ^1}`$ on $`M`$ of codimension 1, such that for any point $`mM`$, $`\mathrm{n}_m`$ is a regular mapping $`^1𝒯_m^{}M`$ of degree $`d=dimM1`$.
###### Remark 2.3.
The condition that $`_{}`$ is separating is equivalent to $`\mathrm{Im}\mathrm{n}_m`$ being not contained in any proper projective subspace of $`𝒯_m^{}M`$. Recall that all regular mappings $`\nu :^1^d`$ of degree $`d`$ which satisfy this property differ only by a projective transformation of $`^d`$. Moreover, the projective transformation $`T:^d^d`$ such that $`T\nu _1=\nu _2`$ is uniquely defined if $`\nu _1`$ and $`\nu _2`$ are two such mappings. A convenient model of such a mapping is given by $`(x:y)(x^d:x^{d1}y:\mathrm{}:xy^{d1}:y^d)`$.
These curves are Veronese curves in the terminology of , or rational normal curves in the terminology of algebraic geometry. The name Veronese web suggests relationship with Veronese curves; in turn, the name Veronese curve was introduce in recognition of the fact that the Veronese surface $`^2^5`$ has the same property: any deformation of it differs by a fraction-linear transformation $`^5^5`$ only.
Restrict our attention to the particular case of $`3`$-dimensional Veronese webs. In this case the only requirement on the family $`\left\{_\lambda \right\}_{\lambda ^1}`$ is that for any $`mM`$ the points $`\mathrm{n}_m\left(\lambda \right)`$, $`\lambda ^1`$ form a smooth (parameterized) quadric in the two-dimensional projective plane $`𝒯_m^{}M`$. Here the parameterization differs from the parameterization given by any stereographic projection by a fraction-linear transformation $`^1^1`$ only. In what follows we consider such parameterizations of quadrics only (any smooth parameterization is such in the complex-geometry case).
###### Lemma 2.4.
A parameterized quadric $`\gamma :^1^2`$ is uniquely determined by $`\gamma \left(\lambda _i\right)`$, $`i=1,2,3,4`$. Here $`\{\lambda _1,\lambda _2,\lambda _3,\lambda _4\}`$ is an arbitrary set of 4 points on $`^1`$.
For any 4 points $`P_i^1`$, $`i=1,2,3,4`$, on $`^2`$ such that no 3 of these points are on the same line one can find a parameterized quadric $`\gamma :^1^2`$ such that $`P_i=\gamma \left(\lambda _i\right)`$, $`i=1,2,3,4`$.
###### Proof.
Recall that given a point $`p^N`$, one can consider a projection $`\pi _p`$ with the center at $`p`$, which sends $`^N\left\{p\right\}`$ onto a projective space $`𝒯_p^N`$ of tangent directions at $`p`$. Here $`\pi _p\left(q\right)`$ is the direction of the line $`\left(pq\right)`$.
Consider compositions $`\pi _{P_i}\gamma :^1𝒯_p^2^1`$, $`i=1,2,3,4`$. Since $`\pi _{P_i}|_{\mathrm{Im}\gamma }`$ is a stereographic projection, these compositions are fraction-linear mappings between projective lines. Thus they are determined by images of any 3 distinct points on $`^1`$. Thus the line $`\left(P_i\gamma \left(\lambda \right)\right)`$ is uniquely determined by $`P_{1,2,3,4}`$. Since $`\gamma \left(\lambda \right)=\left(P_1\gamma \left(\lambda \right)\right)\left(P_2\gamma \left(\lambda \right)\right)`$ if $`\lambda \lambda _{1,2}`$, $`\gamma `$ is uniquely determined by $`P_{1,2,3,4}`$.
To show the existence take any parameterized quadric $`\stackrel{~}{\gamma }:^1^2`$, let $`\stackrel{~}{P}_i=\stackrel{~}{\gamma }\left(\lambda _i\right)`$, $`i=1,2,3,4`$. Then no 3 points out of $`\stackrel{~}{P}_1,\mathrm{},\stackrel{~}{P}_4`$ are on the same line, thus there is a projective mapping $`T:^2^2`$ such that $`\xi \left(\stackrel{~}{P}_i\right)=P_i`$, $`i=1,2,3,4`$. Then $`T\stackrel{~}{\gamma }`$ is the parameterized quadric we need. ∎
###### Corollary 2.5.
Consider a manifold $`M`$, $`dimM=3`$. A Veronese web $`_\lambda `$ on $`M`$ can be reconstructed given 4 foliations $`_{\lambda _i}`$, $`i=1,2,3,4`$, on $`M`$. Here $`\{\lambda _1,\lambda _2,\lambda _3,\lambda _4\}`$ is an arbitrary set of 4 points on $`^1`$.
###### Proof.
Since $`\mathrm{n}_m\left(\lambda _i\right)`$, $`i=1,2,3,4`$, are known for any $`mM`$, by Lemma 2.4 one can find $`\mathrm{n}_m\left(\lambda \right)`$ for any $`\lambda ^1`$ and $`mM`$. This uniquely determines $`_\lambda `$ for any $`\lambda ^1`$. ∎
Fix 4 points $`\{\lambda _1,\lambda _2,\lambda _3,\lambda _4\}^1`$. Given a Veronese web on $`M`$ and a point $`m_0M`$, consider a small neighborhood $`U`$ of $`m_0`$ in $`M`$. One may assume that in $`U`$ the foliations $`_{\lambda _i}`$, $`i=1,2,3,4`$ can be written by equations $`x=\mathrm{const}`$, $`y=\mathrm{const}`$, $`z=\mathrm{const}`$, $`W=\mathrm{const}`$; here $`x,y,z,W`$ are functions on $`U`$. Moreover, $`dx|_{m_0}`$, $`dy|_{m_0}`$ and $`dz|_{m_0}`$ are linearly independent. Indeed, the directions of these 3 vectors are 3 distinct points on a quadric in the projective plane, thus are not on the same line.
Consider $`x,y,z`$ as 3 components of a vector-function $`\phi :U𝕍^3`$, let $`V=\phi \left(U\right)`$. We know that the derivative of this function at $`m_0M`$ is non-degenerate, thus decreasing $`U`$ we may assume that $`\phi `$ gives a diffeomorphism $`UV`$. Then $`w=W\phi ^1`$ is a function on $`V`$, and $`W\left(m\right)=w(x\left(m\right),y\left(m\right),z\left(m\right))`$ if $`mU`$.
###### Lemma 2.6.
The scalar function $`w`$ on $`V𝕍^3`$ and $`\phi \left(m_0\right)V`$ uniquely determine the Veronese web $`_\lambda `$ up to a local diffeomorphism near $`m_0M`$.
###### Proof.
Instead of determining a web up to a local diffeomorphism near $`m_0M`$ it is enough to uniquely determine the diffeomorphic image $`\phi _{}\left(_{}\right)`$ of this web, which is a web on a neighborhood of $`\phi \left(m_0\right)𝕍^3`$. By Corollary 2.5 it is enough to determine $`\phi _{}\left(_{\lambda _i}\right)`$, $`i=1,2,3,4`$. However, leaves of $`\phi _{}\left(_{\lambda _1}\right)`$, $`\phi _{}\left(_{\lambda _2}\right)`$, $`\phi _{}\left(_{\lambda _3}\right)`$ are given by equations $`x=\mathrm{const}`$, $`y=\mathrm{const}`$, $`z=\mathrm{const}`$; here $`(x,y,z)`$ is the standard coordinate system on $`𝕍^3`$. Similarly, leaves of $`\phi _{}\left(_{\lambda _4}\right)`$ are given by the equation $`w(x,y,z)=\mathrm{const}`$. ∎
A change of equations $`x`$, $`y`$, $`z`$ of foliations $`_{\lambda _i}`$, $`i=1,2,3`$, to $`x+C_1`$, $`y+C_2`$, $`z+C_3`$ corresponds to a translation of $`V`$ and $`w`$ by $`(C_1,C_2,C_3)`$, thus one may assume that $`\phi \left(m_0\right)=(0,0,0)`$. Similarly, one may assume that $`w(0,0,0)=0`$.
## 3. Nonlinear wave equation as an integrability condition
Fix a set of 4 points $`\{\lambda _1,\lambda _2,\lambda _3,\lambda _4\}^1`$.
###### Definition 3.1.
Say that a function $`w`$ on an open subset $`M𝕍^3`$ is $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$-admissible if there is a Veronese web $`_\lambda `$ on $`M`$ such that foliations $`_{\lambda _i}`$, $`i=1,2,3,4`$, are given by equations $`x=\mathrm{const}`$, $`y=\mathrm{const}`$, $`z=\mathrm{const}`$, $`w(x,y,z)=\mathrm{const}`$; here $`(x,y,z)`$ is the standard coordinate system on $`𝕍^3`$.
First of all, if $`w`$ is $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$-admissible, Lemma 2.4 implies that for any point $`mM`$ the directions $`dx|_m`$, $`dy|_m`$, $`dz|_m`$ and $`dw|_m`$ are in general position. In other words, $`w_x0`$, $`w_y0`$, $`w_z0`$ everywhere in $`M`$. Thus $`w`$ is non-degenerate (as defined in Section 0. Introduction).
Given non-degeneracy of $`w`$, for any $`\lambda ^1`$ and $`mM`$ the construction of the proof of Corollary 2.5 gives a direction $`\mathrm{n}_m\left(\lambda \right)`$ in the projectivization of $`\left(𝕍^3\right)^{}`$. If $`w`$ is $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$-admissible, then $`m\mathrm{n}_m\left(\lambda \right)`$ coincides with the field of normal directions of the foliation $`_\lambda `$.
Obviously,
###### Lemma 3.2.
Consider a non-degenerate function $`w`$ defined on $`M𝕍^3`$. Suppose that for any $`\lambda ^1`$ the direction field $`\mathrm{n}_m\left(\lambda \right)`$, $`mM`$, given by the construction of the proof of Corollary 2.5 coincides with the field of normal directions of a foliation on $`M`$. Then $`w`$ is $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$-admissible.
Thus to check whether a non-degenerate function $`w`$ is $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$-admissible it is enough to check whether a given direction field coincides with a normal field to a foliation. Such direction fields can be described by the following particular case of the Frobenius integrability condition :
###### Lemma 3.3.
Consider a $`1`$-form $`\omega `$ on a manifold $`M`$ which does not vanish at any point of $`M`$. Call $`\omega `$ Frobenius integrable if there exists a foliation $``$ of codimension 1 on $`M`$ such that $`\omega \left(m\right)`$ is normal to the tangent space at $`m`$ to the leaf $`L_m`$ of $``$ through $`m`$ for any $`mM`$.
Then $`\omega `$ is Frobenius integrable iff $`\omega d\omega =0`$.
###### Proof.
The “only if” part is simple: in an appropriate neighborhood $`U`$ of any given point $`m_0M`$ the foliation $`|_U`$ can be written as $`g=\mathrm{const}`$; here $`g`$ is a function on $`U`$, and $`dg0`$ for any $`mU`$. Thus $`\omega =hdg`$ for an appropriate function $`h`$ on $`U`$, and $`\omega d\omega =hdgdhdg=0`$.
For the “if” part it is enough to show the existence locally on $`M`$, since the foliation is unique if it exists, thus gluing pieces together is not a problem. We may assume that $`M`$ is an open subset of $`𝕍^n`$, and that $`\omega |_{m_0}=dx_n|_{m_0}`$. Say that a tangent vector $`v`$ at $`mM`$ is $`k`$-compatible, $`k=1,\mathrm{},n1`$, if $`\omega |_m,v=0`$ and $`v`$ is of the form $`\frac{}{x_k}+a\frac{}{x_n}`$ with an appropriate number $`a`$. Obviously, in an appropriate neighborhood of any point $`m_0M`$ there is exactly one $`k`$-compatible vector $`v_k\left(m\right)`$ for $`k=1,\mathrm{},n1`$. Define functions $`a_{\left(k\right)}\left(m\right)`$ by $`v_k\left(m\right)=\frac{}{x_k}+a_{\left(k\right)}\left(m\right)\frac{}{x_n}`$. Then the fundamental relationship between commutator and de Rham differential<sup>1</sup><sup>1</sup>1One can easily check this relation in local coordinates.
$$\omega ,[v_k,v_l]=v_k\omega ,v_lv_l\omega ,v_k+d\omega ,v_kv_l$$
implies $`\omega ,[v_k,v_l]=d\omega ,v_kv_l`$. Since $`\omega d\omega =0`$, one can write $`d\omega =\omega \alpha `$; here $`\alpha `$ is a $`1`$-form defined near $`m_0`$. Hence
$$d\omega ,v_kv_l=\omega ,v_k\alpha ,v_l\alpha ,v_k\omega ,v_l=0.$$
Thus $`\omega ,[v_k,v_l]=0`$. On the other hand, $`[v_k,v_l]=\left(v_ka_lv_la_k\right)\frac{}{x_n}`$. Together with $`\omega ,[v_k,v_l]=0`$ this implies $`[v_k,v_l]=0`$. By the principal theorem of the theory of ODE, one can find local coordinates $`(y_1,\mathrm{},y_n)`$ such that $`v_k=\frac{}{y_k}`$, $`k=1,\mathrm{},n1`$. Since $`\omega `$ is orthogonal to $`v_k`$, $`k=1,\mathrm{},n1`$, this implies that $`\omega =h\left(y\right)dy_n`$, thus $`y_n=\mathrm{const}`$ gives a foliation with the required properties. ∎
The next step is to provide an explicit construction of the normal directions $`\mathrm{n}_m\left(\lambda \right)`$ in terms of $`w`$.
###### Lemma 3.4.
Given a Veronese curve $`\gamma \left(\lambda \right)`$ in $`^{n1}`$, one can find polynomials $`p_1\left(\lambda \right),\mathrm{},p_n\left(\lambda \right)`$ of degree $`n1`$ such that $`\gamma \left(\lambda \right)=(p_1\left(\lambda \right):\mathrm{}:p_n\left(\lambda \right))`$ for $`\lambda \mathrm{}`$. Polynomials $`p_k\left(\lambda \right)`$ are defined uniquely up to multiplication by the same constant.
###### Proof.
Any Veronese curve in $`^{n1}`$ is a projective transformation of the closure of the image of the mapping $`\lambda (1:\lambda :\mathrm{}:\lambda ^{n1})`$. A consideration of the corresponding linear transformation of $`𝕍^n`$ provides polynomials $`p_1,\mathrm{},p_n`$.
It is enough to show uniqueness for the curve $`\left(1:\lambda :\mathrm{}:\lambda ^{n1}\right)`$. Obviously, $`p_k\left(\lambda \right)=\lambda ^{k1}p_1\left(\lambda \right)`$. Moreover, since $`\mathrm{deg}p_nn1`$, $`p_1\left(\lambda \right)`$ is a constant. ∎
Thus any Veronese curve in $`^2`$ is a projectivization of a polynomial vector-function $`v\left(\lambda \right)`$ of degree exactly 2. Note that $`v\left(\lambda \right)0`$ for any $`\lambda `$.
This implies that the dependence on $`\lambda `$ of the directions $`\mathrm{n}_m\left(\lambda \right)`$, $`\lambda \mathrm{}`$, can be described by the direction of the $`1`$-form $`\alpha \left(m\right)+\lambda \beta \left(m\right)+\lambda ^2\gamma \left(m\right)`$; here $`\alpha ,\beta ,\gamma `$ are appropriate $`1`$-forms on $`M𝕍^3`$ which are defined up to multiplication by the same function on $`M`$. If $`\lambda 0`$, $`\mathrm{n}_m\left(\lambda \right)`$ is the direction of $`\gamma \left(m\right)+\lambda ^1\beta \left(m\right)+\lambda ^2\alpha \left(\lambda \right)`$, taking the limit $`\lambda \mathrm{}`$ implies that $`\mathrm{n}_m\left(\mathrm{}\right)`$ is the direction of $`\gamma \left(m\right)`$.
###### Lemma 3.5.
Consider vectors $`v_1,v_2,v_3,v_4`$ in $`𝕍^3`$ such that $`v_1,v_2,v_3`$ are linearly independent. Fix a set of 4 points $`\{\lambda _1,\lambda _2,\lambda _3,\lambda _4\}^1\left\{\mathrm{}\right\}`$. There is a unique polynomial vector-function $`v\left(\lambda \right)`$ of degree 2 such that $`v\left(\lambda _4\right)=v_4`$, and $`v\left(\lambda _k\right)`$ is proportional to $`v_k`$, $`k=1,2,3`$.
###### Proof.
Write $`v_4`$ as $`av_1+bv_2+cv_3`$. Since $`v\left(\lambda \right)`$ can be written as $`\alpha \left(\lambda \right)v_1+\beta \left(\lambda \right)v_2+\gamma \left(\lambda \right)v_3`$, we know that $`\alpha \left(\lambda _2\right)=\alpha \left(\lambda _3\right)=0`$, $`\alpha \left(\lambda _4\right)=a`$. This uniquely determines the quadratic polynomial $`\alpha \left(\lambda \right)`$. Proceed similarly for $`\beta \left(\lambda \right)`$ and $`\gamma \left(\lambda \right)`$. ∎
###### Corollary 3.6.
Fix a set of 4 points $`\{\lambda _1,\lambda _2,\lambda _3,\lambda _4\}^1\left\{\mathrm{}\right\}`$. Given a non-degenerate function $`w`$ on $`M𝕍^3`$, the direction $`\mathrm{n}_m\left(\lambda \right)`$ defined by the construction of the proof of Corollary 2.5 coincides with $`\left(p_1\left(\lambda \right)w_x:p_2\left(\lambda \right)w_y:p_3\left(\lambda \right)w_z\right)`$; here
(3.1)
$$p_i\left(\lambda \right)=\left(\lambda _4\lambda _i\right)\left(\lambda \lambda _j\right)\left(\lambda \lambda _k\right),$$
for any permutation $`\left(ijk\right)`$ of (123).
###### Corollary 3.7.
Consider distinct points $`\lambda _i\mathrm{}`$, $`i=1,2,3,4`$. Consider a non-degenerate function $`w`$ on $`M𝕍^3`$. Let
(3.2)
$$\omega _\lambda \stackrel{\text{def}}{=}p_1\left(\lambda \right)w_xdx+p_2\left(\lambda \right)w_ydy+p_3\left(\lambda \right)w_zdz;$$
here $`p_{1,2,3}\left(\lambda \right)`$ are from (3.1). Then the following conditions are equivalent:
1. $`w`$ is $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$-admissible;
2. $`\omega _\lambda d\omega _\lambda =0`$ for any $`\lambda `$;
3. $`\omega _\lambda d\omega _\lambda =0`$ for any 5 distinct values of $`\lambda `$;
4. $`\omega _\lambda d\omega _\lambda =0`$ for any $`\lambda _0\{\lambda _1,\mathrm{},\lambda _4\}`$;
If $`0\{\lambda _1,\mathrm{},\lambda _4\}`$, these conditions are equivalent to
(3.3)
$$\nu _{23}w_xw_{yz}+\nu _{31}w_yw_{xz}+\nu _{12}w_zw_{xy}=0,$$
here $`\nu _{kl}=\lambda _k/\left(\lambda _4\lambda _k\right)\lambda _l/\left(\lambda _4\lambda _l\right)`$.
###### Proof.
Obviously, $`\omega _\lambda |_m0`$ for any $`\lambda \mathrm{}`$ and any $`mM`$. By Lemma 3.3, $`\omega _\lambda d\omega _\lambda =0`$ is equivalent to existence of a foliation to which $`\omega _\lambda `$ is normal. Thus by Lemma 3.2 the first statement implies the second one.
If $`\omega _\lambda d\omega _\lambda =0`$ for any $`\lambda `$, then by Corollary 3.6, the required in Lemma 3.2 foliation exists for $`\lambda \mathrm{}`$. However, $`\stackrel{~}{\omega }_\lambda =\lambda ^2\omega _\lambda `$ is defined for $`\lambda ^1\left\{0\right\}`$, and $`\stackrel{~}{\omega }_\lambda d\stackrel{~}{\omega }_\lambda `$ is a polynomial of degree 4 in $`\lambda ^1`$. Thus $`\stackrel{~}{\omega }_\lambda d\stackrel{~}{\omega }_\lambda =0`$, including $`\lambda =\mathrm{}`$. Moreover, $`\stackrel{~}{\omega }_{\mathrm{}}|_m0`$ for any $`m`$, which implies the existence of $`_\lambda `$ for $`\lambda =\mathrm{}`$ as well. Thus the second statement implies the first one.
Since $`\omega _\lambda `$ is quadratic in $`\lambda `$, $`\omega _\lambda d\omega _\lambda `$ is a polynomial of degree 4 in $`\lambda `$. Thus the second statement is equivalent to the third one. By construction $`\omega _{\lambda _{1,2,3,4}}`$ are proportional to $`dx`$, $`dy`$, $`dz`$, and $`dw`$ correspondingly. This implies that $`\omega _\lambda d\omega _\lambda =0`$ for $`\lambda \{\lambda _1,\lambda _2,\lambda _3,\lambda _4\}`$. Consequently, the fourth statement is equivalent to the third one.
Assume that $`\lambda _0=0\{\lambda _1,\lambda _2,\lambda _3,\lambda _4\}`$. Let $`\mu _k=\lambda _4/\lambda _k1`$, $`k=1,2,3`$. Then
$$\omega _0=\lambda _1\lambda _2\lambda _3\stackrel{~}{\omega },\stackrel{~}{\omega }\stackrel{\text{def}}{=}\mu _1w_xdx+\mu _2w_ydy+\mu _3w_zdz,$$
and $`\stackrel{~}{\omega }d\stackrel{~}{\omega }`$ can be written as
$$\mu _1\mu _2\mu _3\left(\left(\mu _2^1\mu _3^1\right)w_xw_{yz}+\left(\mu _3^1\mu _1^1\right)w_yw_{xz}+\left(\mu _1^1\mu _2^1\right)w_zw_{xy}\right)dxdydz.$$
(It is clear that $`\mu _k0`$ for $`k=1,2,3`$.) Since $`\nu _{kl}=\mu _k^1\mu _l^1`$, the equation $`\omega _0d\omega _0=0`$ is proportional to (3.3), which implies the last statement of the corollary. ∎
Obviously, $`\omega _\lambda d\omega _\lambda =\alpha _{k=1}^4\left(\lambda \lambda _k\right)`$; here $`\alpha `$ is a $`3`$-form on $`M`$ which does not depend on $`\lambda `$. Thus the equations $`\omega _{\lambda _0}d\omega _{\lambda _0}=0`$ for different values $`\lambda _0`$ are proportional, and it does not matter much which value of $`\lambda _0`$ one would use. Consequently, any other choice of $`\lambda _0`$ would lead to an equation which is proportional to (3.3), and one can drop the conditions that $`0\{\lambda _1,\mathrm{},\lambda _4\}`$. Moreover, it is possible to drop the condition $`\mathrm{}\{\lambda _1,\mathrm{},\lambda _4\}`$ as well:
###### Theorem 3.8.
A non-degenerate function $`w`$ on $`M𝕍^3`$ is $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$-admissible iff it satisfies an $`(A,B,C)`$-equation (0.2) with $`A/C=(\lambda _1:\lambda _2:\lambda _3:\lambda _4)`$; here $`(a:b:c:d)=\frac{da}{dc}\frac{bc}{ba}`$ is the cross-ratio of $`a,b,c,d`$.
###### Proof.
Indeed, a direct calculation shows that $`\nu _{12}+\nu _{23}+\nu _{31}=0`$, $`\nu _{12}0`$, $`\nu _{23}0`$, $`\nu _{31}0`$, and $`\nu _{23}/\nu _{12}=(\lambda _1:\lambda _2:\lambda _3:\lambda _4)`$. Thus the statement holds for $`\mathrm{}\{\lambda _1,\mathrm{},\lambda _4\}`$. However, if $`T`$ is a projective transformation, then $`w`$ is $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$-admissible iff it is $`(T\lambda _1,T\lambda _2,T\lambda _3,T\lambda _4)`$-admissible. Since cross-ratio is invariant w.r.t. projective transformations, it is enough to prove the statement for $`(T\lambda _1,T\lambda _2,T\lambda _3,T\lambda _4)`$ with an arbitrary $`T`$. By an appropriate choice of $`T`$ we can ensure that $`\mathrm{}\{\lambda _1,\mathrm{},\lambda _4\}`$ (and additionally $`0\{\lambda _1,\mathrm{},\lambda _4\}`$ if we wish). ∎
###### Remark 3.9.
Since the cross-ratio of 4 distinct points can take any value distinct from $`0,1,\mathrm{}`$, one can momentarily see that for any triple $`(A,B,C)`$ which satisfies (0.1) and for any 3 distinct points $`\lambda _1,\lambda _2,\lambda _3`$ one can find $`\lambda _4`$ such that $`A/C=(\lambda _1:\lambda _2:\lambda _3:\lambda _4)`$. Thus any $`(A,B,C)`$-equation can be interpreted as an integrability condition of a Veronese web: any Veronese web gives rise to a non-degenerate solution of such an equation, and any non-degenerate solution can be represented in this form.
###### Remark 3.10.
One can generalize Corollary 3.7 to the case of Veronese webs of arbitrary dimension. In dimension $`d`$ one still needs one function $`w`$ of $`d`$ variables to completely determine a web up to a local diffeomorphism. The foliation $`_\lambda `$ can be described by a $`1`$-form $`\omega _\lambda `$ which is normal to leaves of $`_\lambda `$, and is given by a formula similar to (3.2).
The $`1`$-form $`\omega _\lambda `$ depends on $`\lambda `$ as a polynomial of degree $`d1`$, and the integrability condition $`\omega _\lambda d\omega _\lambda =0`$ is a polynomial of degree $`2d2`$. Thus a non-degenerate function $`w`$ of $`d`$ variables corresponds to a Veronese web iff $`\omega _\lambda d\omega _\lambda =0`$ for $`2d1`$ different values of $`\lambda `$. By its construction, $`\omega _\lambda d\omega _\lambda =0`$ automatically holds for $`d+1`$ value of $`\lambda `$. Thus a naive generalization (as done in ) of Corollary 3.7 would be that it is enough to require $`\omega _\lambda d\omega _\lambda =0`$ at $`d2`$ “additional” values of $`\lambda `$.
However, contain a much stronger result: if $`\omega _\lambda d\omega _\lambda =0`$ for any “additional” value of $`\lambda `$, then $`\omega _\lambda d\omega _\lambda =0`$ for any $`\lambda `$, thus $`w`$ determines a Veronese web. Unfortunately, this condition is on a $`3`$-form in $`d`$-dimensional space, thus it is still an overdetermined system of partial differential equations on $`w`$, if $`d>3`$. It is very interesting to investigate whether arguments of allow extraction of one equation on $`w`$ which implies $`\omega _\lambda d\omega _\lambda =0`$.
## 4. Bäcklund–Darboux transformations
By Remark 3.9, any solution of an $`(A,B,C)`$-equation gives rise to a Veronese web, which in turn leads to a solution of $`(A^{},B^{},C^{})`$-equation, possibly with different $`(A^{},B^{},C^{})`$.
###### Corollary 4.1.
Let $`w`$ be a non-degenerate solution of $`(A,B,C)`$-equation in a neighborhood of $`(0,0,0)`$, $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$ be numbers such that $`A/C=(\lambda _1:\lambda _2:\lambda _3:\lambda _4)`$. Then for any number $`\lambda `$
1. there is a function $`v(x,y,z)`$ defined in a neighborhood of $`(0,0,0)`$ such that the following identity of vector-functions holds:
(4.1)
$$(v_x,v_y,v_z)=\psi (x,y,z)(\alpha w_x,\beta w_y,\gamma w_z);$$
here $`\alpha ,\beta ,\gamma \alpha =p_1\left(\lambda \right)`$, $`\beta =p_2\left(\lambda \right)`$, $`\gamma =p_3\left(\lambda \right)`$, $`p_k`$ are polynomials given by (3.1), and $`\psi `$ is an appropriate scalar-valued function;
2. if $`\lambda \{\lambda _1,\lambda _2,\lambda _3\}`$, then $`v(x,y,z)`$ can be chosen to be non-degenerate;
3. if $`v`$ is non-degenerate it is $`(\lambda _1,\lambda _2,\lambda _3,\lambda )`$-admissible;
4. if $`\lambda \{\lambda _1,\lambda _2,\lambda _3\}`$, and $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$ satisfy conditions (0.1), and $`\stackrel{~}{A}/\stackrel{~}{C}=(\lambda _1:\lambda _2:\lambda _3:\lambda )`$, then the function $`v(x,y,z)`$ satisfies $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$-equation.
###### Proof.
By Corollary 3.7, $`w`$ is $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$-admissible, thus it corresponds to a web $`_{}`$. Write the leaves of $`_\lambda `$ as $`v(x,y,z)=\mathrm{const}`$, and apply Corollary 3.7 again. ∎
Obviously, the function $`v`$ of the previous corollary is defined uniquely up to a gauge transformation (see Definition 0.2).
Let us find relationships between 9 constants $`(A,B,C)`$, $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$ and $`(\alpha ,\beta ,\gamma )`$ which appear in the statements of this section. Construct $`\stackrel{~}{\nu }_{kl}`$ basing on the $`4`$-tuple $`(\lambda _1,\lambda _2,\lambda _3,\lambda )`$ using the same formula as used to construct $`\nu _{kl}`$ basing on $`(\lambda _1,\lambda _2,\lambda _3,\lambda _4)`$. Let
$$\tau =\left(\lambda _4\lambda _1\right)\left(\lambda _4\lambda _2\right)\left(\lambda _4\lambda _3\right)\frac{\lambda }{\lambda _4}.$$
Then it is easy to check that $`\alpha =\tau \nu _{23}/\stackrel{~}{\nu }_{23}`$, $`\beta =\tau \nu _{13}/\stackrel{~}{\nu }_{13}`$, $`\gamma =\tau \nu _{12}/\stackrel{~}{\nu }_{12}`$. Since simultaneous multiplication of $`\alpha ,\beta ,\gamma `$ by the same non-zero number does not change the meaning of Equation (4.1), we conclude that one can take $`\alpha =A/\stackrel{~}{A}`$, $`\beta =B/\stackrel{~}{B}`$, $`\gamma =C/\stackrel{~}{C}`$.
###### Proof of Theorem 0.3 .
The first statement is obvious, and the second one is the corollary of the first since $`dv_1dv`$ if $`v_1`$ is a gauge transform of $`v`$. The third and the fourth statements are reformulations of parts of Corollary 3.7. The last statement is a direct corollary of the first one and of the following obvious statement:
###### Lemma 4.2.
Given $`(v_x,v_y,v_z)(v_x^{},v_y^{},v_z^{})`$ for two non-degenerate functions $`v`$ and $`v^{}`$ defined in a neighborhood of (0,0,0) in $`𝕍^3`$, one can decrease the neighborhood so that the functions become gauge transforms of each other.
This finishes the proof of Theorem 0.3. ∎
To enhance the statements about Equation (0.5), note the following two lemmas.
###### Lemma 4.3.
Given numbers $`\alpha 0`$, $`\beta 0`$, $`\gamma 0`$ such that $`\alpha \beta `$, $`\alpha \gamma `$, $`\beta \gamma `$, there exist two triples $`(A,B,C)`$ and $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$ which both satisfy conditions (0.1), and $`\alpha =A/\stackrel{~}{A}`$, $`\beta =B/\stackrel{~}{B}`$, $`\gamma =C/\stackrel{~}{C}`$. The numbers $`A,B,C,\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C}`$ are defined uniquely up to multiplication by the same constant.
###### Proof.
Given $`\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C}`$ put $`A=\alpha \stackrel{~}{A}`$, $`B=\beta \stackrel{~}{B}`$, $`C=\gamma \stackrel{~}{C}`$. The conditions (0.1) on $`(A,B,C)`$ can be translated to an additional linear equation $`\alpha \stackrel{~}{A}+\beta \stackrel{~}{B}+\gamma \stackrel{~}{C}=0`$ on $`\stackrel{~}{A}`$, $`\stackrel{~}{B}`$, $`\stackrel{~}{C}`$. This equation is independent of $`\stackrel{~}{A}+\stackrel{~}{B}+\stackrel{~}{C}=0`$, thus there is a unique (up to proportionality) solution $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$ of these two equations. What remains to check is that this solution does not contradict the conditions $`\stackrel{~}{A}0`$, $`\stackrel{~}{B}0`$, $`\stackrel{~}{C}0`$. However, $`\stackrel{~}{A}=0`$ contradicts $`\beta \gamma `$, etc. ∎
###### Lemma 4.4.
Given two triples $`(A,B,C)`$ and $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$ which both satisfy conditions (0.1), put $`\alpha =A/\stackrel{~}{A}`$, $`\beta =B/\stackrel{~}{B}`$, $`\gamma =C/\stackrel{~}{C}`$. Then either $`\alpha =\beta =\gamma `$, or $`\alpha \beta `$, $`\alpha \gamma `$, $`\beta \gamma `$.
This statement is elementary.
## 5. Inverse construction
Of course, Theorem 0.3 can be proven by elementary methods without any reference to Veronese webs. However, Veronese webs are not useful because this theorem can be proven “naturally” by using Veronese webs. In fact Veronese webs appears naturally as reformulations of the statement of this theorem.
Indeed, given a solution of Equation (0.2), consider Systems (0.4) for all possible triples $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$. Since proportional triples $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$ give essentially the same systems, we can enumerate all the triples by the ratio $`\lambda =\stackrel{~}{A}/\stackrel{~}{C}`$, which can be considered as an element of $`^1`$ with the only restrictions being $`\lambda \mathrm{}`$, $`\lambda 0`$, $`\lambda 1`$.
For any such value of $`\lambda `$ one obtains a solution $`v^{\left[\lambda \right]}`$ of Equation (0.3). This solution is defined in a neighborhood $`U_\lambda `$ of (0,0,0), and it is easy to show that this neighborhood may be chosen independently of $`\lambda `$, denote it by $`U`$. The solution $`v^{\left[\lambda \right]}`$ is not unique, but the foliation $`_\lambda `$ of $`U`$ defined by $`v^{\left[\lambda \right]}=\mathrm{const}`$ is uniquely defined. Moreover, $`_\lambda `$ depends smoothly on $`\lambda ^1\{0,1,\mathrm{}\}`$. What remains it to consider what happens near $`\lambda =0`$, near $`\lambda =1`$, and near $`\lambda =\mathrm{}`$.
If $`\lambda 0`$, then $`\stackrel{~}{A}`$ is very small, thus Equation (0.5)
$$(v_x,v_y,v_z)(Aw_x,\stackrel{~}{A}B\stackrel{~}{B}^1w_y,\stackrel{~}{A}C\stackrel{~}{C}^1w_z)$$
becomes close to $`(v_x,v_y,v_z)(Aw_x,0,0)`$, or, in other words, to $`(v_x,v_y,v_z)`$ (1,0,0). The solution to this equation is $`v=v\left(x\right)`$, thus the foliation $`_\lambda `$ has a limit $`x=\mathrm{const}`$ when $`\lambda `$ 0. Similarly, the limit when $`\lambda 1`$ is $`y=\mathrm{const}`$, when $`\lambda \mathrm{}`$ is $`z=\mathrm{const}`$.
Thus an investigation of the statement of Theorem 0.3 directly leads to a family a foliations which depend smoothly on a parameter $`\lambda ^1`$. In the following section we show that the conditions that the normal directions to the foliations span a quadratic cone is also related to the elementary theory of Equation (0.2).
Additionally, the following statement is easy to obtain elementary, but it is an immediate corollary of Theorem 0.3:
###### Corollary 5.1.
If $`w`$ is a non-degenerate solution of Equation (0.2), then any gauge transform of $`w`$ is also a solution of Equation (0.2).
## 6. Linearization
Given a solution $`\overline{\kappa }`$ of a non-linear (system of) equation(s) $`F\left(\kappa \right)=0`$, the linearized equation at $`\overline{\kappa }`$ is the equation $`F\left(\overline{\kappa }+\epsilon \kappa \right)=O\left(\epsilon ^2\right)`$. It is a (system of) linear equation(s) on $`\kappa `$ with the coefficients being partial derivatives of $`F`$ at $`\overline{\kappa }`$.
Obviously, given a solution $`\overline{w}`$ of Equation (0.2), the linearization is
(6.1)
$$A\overline{w}_xw_{yz}+B\overline{w}_yw_{xz}+C\overline{w}_zw_{xy}+A\overline{w}_{yz}w_x+B\overline{w}_{xz}w_y+C\overline{w}_{xy}w_z=0,$$
The left-hand side is a linear differential operator of second order in $`w`$, denote this operator $`l_{\overline{w}}`$ or just $`l`$. The principal symbol of $`l_{\overline{w}}`$ is
(6.2)
$$\mathrm{\Lambda }(x,y,z,\xi ,\eta ,\zeta )=A\overline{w}_x\eta \zeta +B\overline{w}_y\xi \zeta +C\overline{w}_z\xi \eta .$$
This is a non-degenerate quadratic form in $`(\xi ,\eta ,\zeta )`$ iff $`\overline{w}`$ is non-degenerate. Moreover, it vanishes if $`(\xi ,\eta ,\zeta )=(1,0,0)`$, or $`(\xi ,\eta ,\zeta )=(0,0,1)`$, or $`(\xi ,\eta ,\zeta )=(0,0,1)`$. This shows that the linearization is hyperbolic iff $`\overline{w}`$ is non-degenerate. This is why it makes sense to call the equation (0.2) a nonlinear wave equation.
Fix a point $`(x,y,z)`$. Recall that a covector $`(\xi ,\eta ,\zeta )`$ at $`(x,y,z)`$ is characteristic if $`\mathrm{\Lambda }(x,y,z,\xi ,\eta ,\zeta )=0`$. Characteristic covectors of a hyperbolic linear differential equation form a cone in the cotangent space, this cone is called a wave cone. Since our equation is of second order, it is a quadratic cone. Recall that a surface in $`𝕍^3`$ is called characteristic if the normal direction to this surface at any point is characteristic. One can define similar notions for square systems of equations by taking $`det\mathrm{\Lambda }`$ instead of $`\mathrm{\Lambda }`$.
Recall how to construct characteristic surfaces. Consider an expression $`l\left(e^{ik\phi (x,y,z)}\right)`$ when $`k\mathrm{}`$. It can be written as $`\mathrm{\Phi }_\phi (k,x,y,z)e^{ik\phi (x,y,z)}`$; here $`\mathrm{\Phi }_\phi `$ depends polynomially on $`k`$, the degree being 2 or less. Say that $`\phi `$ is an eikonal solution if $`\mathrm{\Phi }_\phi `$ is a polynomial in $`k`$ of degree $`1`$. If $`\mathrm{\Phi }_{\phi ,2}(x,y,z)`$ is the coefficient at $`k^2`$ in $`\mathrm{\Phi }_\phi `$, then the equation
$$\mathrm{\Phi }_{\phi ,2}(x,y,z)=0$$
is a non-linear differential equation of the first order on $`\phi `$. Call this equation the eikonal equation.
Obviously, eikonal solutions coincide with solutions to the eikonal equation. Moreover, it is easy to see that the eikonal equation is equivalent to the surfaces $`\phi =c`$ being characteristic surfaces for any constant $`c`$.
A similar statements holds for square systems of differential equations if one considers $`l\left(ve^{ik\phi (x,y,z)}\right)`$ as a linear function of a vector $`v`$. Then $`\mathrm{\Phi }_\phi `$ becomes a square matrix, and we can consider the degree of $`det\mathrm{\Phi }`$ in $`k`$ instead of the degree of $`\mathrm{\Phi }`$ in $`k`$.
###### Proposition 6.1.
Consider a non-degenerate solution $`w`$ of Equation (0.2). Then $`w`$ is also a solution of the linearized equation (6.1) at $`w`$. Moreover, $`w`$ is also an eikonal solution for this linearized equation.
###### Proof.
To prove the first statement, apply Corollary 5.1. Since $`w`$ is a solution, so is $`w+\epsilon \omega `$ for any $`\epsilon `$. Similarly, since $`w+\epsilon e^{ikw}`$ is a solution for any $`\epsilon `$ and $`k`$, $`w`$ is an eikonal solution as well. ∎
###### Proposition 6.2.
Consider a solution $`(w,v)`$ of System (0.4) with non-degenerate $`w`$ and $`v`$. Let $`l^{\left\{1\right\}}`$, $`l^{\left\{2\right\}}`$ be the linearizations of Equations (0.2), (0.3) at $`w`$, and $`l^{\left\{3\right\}}`$ be the linearization of Equation (0.4) at $`(w,v)`$. Then
1. Characteristic cones of $`l^{\left\{1\right\}}`$, $`l^{\left\{2\right\}}`$, $`l^{\left\{3\right\}}`$ coincide.
2. The function $`v`$ is a solution of the eikonal equation for $`l^{\left\{1\right\}}`$.
###### Proof.
It is easy to check the first claim by a direct calculation. In the second claim we already know that $`v`$ is a solution of the eikonal equation for $`l^{\left\{2\right\}}`$. Since characteristic cones coincide, $`v`$ is also a solution of the eikonal equation for $`l^{\left\{1\right\}}`$. ∎
###### Remark 6.3.
Let us provide a more conceptual heuristic proof of the first claim of the proposition. It is enough to consider characteristic cones for $`l^{\left\{1\right\}}`$ and $`l^{\left\{3\right\}}`$. If $`\phi `$ is a solution of the eikonal equation for $`l^{\left\{3\right\}}`$, then $`l^{\left\{3\right\}}(\stackrel{~}{w},\stackrel{~}{v})=O\left(1\right)`$ when $`k\mathrm{}`$; here
$$\stackrel{~}{w}(x,y,z)=We^{ik\phi (x,y,z)},\stackrel{~}{v}(x,y,z)=Ve^{ik\phi (x,y,z)}$$
and $`W`$ and $`V`$ are appropriate constants. The usual arguments of calculus of asymptotics (see, for example, ) show that by allowing $`W`$ and $`V`$ depend smoothly on $`x,y,z,k^1`$ one can ensure that $`l^{\left\{3\right\}}(\stackrel{~}{w},\stackrel{~}{v})`$ is asymptotically 0 when $`k\mathrm{}`$.
In other words, starting with a solution of the eikonal equation for $`l^{\left\{3\right\}}`$, one can construct an asymptotic solution for $`l^{\left\{3\right\}}`$. Since the relationship between $`l^{\left\{3\right\}}`$ and $`l^{\left\{1\right\}}`$ is a linearization of relation between System (0.4) and Equation (0.2), we conclude that $`W(x,y,z,k^1)e^{ik\phi (x,y,z)}`$ is an asymptotic solution for $`l^{\left\{1\right\}}`$ (as given this argument is heuristic only, one needs to check that the order of taking limits in $`k`$ and in $`\epsilon `$ is correct). Thus $`\phi `$ is also a solution of the eikonal equation for $`l^{\left\{1\right\}}`$. Since the characteristic cone is spanned by differentials of eikonal solutions, the characteristic cone for $`l^{\left\{3\right\}}`$ is a subset of a characteristic cone for $`l^{\left\{1\right\}}`$.
On the other hand, characteristic cones of $`l^{\left\{1\right\}}`$ and $`l^{\left\{3\right\}}`$ are quadratic cones, thus they should coincide.
###### Remark 6.4.
Let us repeat the arguments of Section 5 in the linearized situation. For any $`\lambda ^1`$ we can construct a corresponding triple $`(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})`$ with $`\stackrel{~}{A}/\stackrel{~}{C}=\lambda `$, and a solution $`v^{\left[\lambda \right]}`$ of the corresponding Equation (0.3), thus of $`l^{\left\{2\right\}}`$. Then $`v^{\left[\lambda \right]}`$ is a solution of the eikonal equation for $`l^{\left\{1\right\}}`$. Its level surfaces are characteristic surfaces of $`l^{\left\{1\right\}}`$. For each value of $`\lambda `$ we obtain one characteristic surface passing through a given point.
Moreover, when we vary $`\lambda `$ the coefficients $`A/\stackrel{~}{A}`$, $`B/\stackrel{~}{B}`$, $`C/\stackrel{~}{C}`$ in Equation (0.5) vary as well. They cannot be proportional for different values of $`\lambda `$, thus all the above characteristic surfaces passing through a given point have different directions.
In other words, at a given point we obtain a family of characteristic directions parameterized by $`^1`$. But characteristic directions span a quadratic cone, and the base of this cone is $`^1`$. It easily follows that given a characteristic direction at a given point one can find a value of $`\lambda ^1`$ such that $`dv^{\left[\lambda \right]}`$ at the given point goes in the prescribed direction.
This concludes arguments of Section 5, since using elementary arguments we concluded that results of Theorem 0.3 imply that normal directions to $`_\lambda `$ at a given point should span a quadratic cone.
## 7. $`_{}`$-convex sets and the twistor transform
###### Definition 7.1.
Given a foliation $``$ on $`M`$ and an open subset $`UM`$, we say that $`U`$ is $``$-convex if there is an open subset $`VU`$ such that $`|_V`$ is direct (as defined in Section 1), and for any leaf $`L`$ of $`|_V`$ the set $`LU`$ is connected. Call $`U`$ strictly $``$-convex if additionally the image of $`U`$ under the natural projection $`U𝔅_{|_V}`$ is homeomorphic to a ball.
It is obvious that any point $`mM`$ has a strictly $``$-convex neighborhood. For example, any direct neighborhood (see Section 1) goes.
###### Definition 7.2.
Given a foliation $``$ on $`M`$, denote by $`𝔅_{}`$ the set of leaves of $`|_U`$, and by $`b:M𝔅_{}`$ the natural projection. Given an $``$-convex subset $`U`$, the set $`𝔅_{|_U}`$ has a natural structure of a manifold. Obviously, when one decreases an $``$-convex subset $`U`$, the base $`𝔅_{|_U}`$ decreases as well. In particular, if $`mM`$, then the germ<sup>2</sup><sup>2</sup>2Given a manifold $`M`$ with a closed submanifold $`N`$, an open submanifold $`UM`$ is compatible with $`M`$ if $`UN`$. Extend compatibility relation to an equivalence relation $``$ between manifolds $`M_1N`$. Call equivalence classes $`\stackrel{~}{M}`$ germs near $`N`$. A mapping of germs (or a germ of a mapping) $`(\stackrel{~}{M},N)(\stackrel{~}{M}^{},N^{})`$ is a smooth mapping $`f:MM^{}`$ such that $`f\left(N\right)N^{}`$; here $`M`$, $`M^{}`$ are some representatives of classes $`\stackrel{~}{M}`$, $`\stackrel{~}{M}^{}`$. Such mappings are considered up to the natural equivalence relation induced by restriction to compatible open subsets. of $`𝔅_{|_U}`$ near $`b\left(m\right)`$ does not depend on the $``$-convex neighborhood $`U`$ of $`m`$. Call this germ the local base of the foliation $``$ near $`m`$.
###### Definition 7.3.
Given a web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$, call an open subset $`UM`$ (strictly) $`_{}`$-convex if $`U`$ is (strictly) $`_\lambda `$-convex for all the foliations $`_\lambda `$.
Recall that a section of a mapping $`\pi :MN`$ is a right inverse to $`\pi `$ mappings $`NM`$.
###### Definition 7.4.
Consider a smooth web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$, and an $`_{}`$-convex subset $`UM`$. For any fixed $`\lambda \mathrm{\Lambda }`$ consider the manifold $`𝔅_{_\lambda |_U}`$. Taken together, they form a manifold $`𝔗=𝔗_{_{}}=_{\lambda \lambda }𝔅__\lambda `$ equipped with a projection $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$ (which sends $`𝔅__\lambda \left\{\lambda \right\}`$). Call the pair $`(𝔗,\pi )`$ the twistor transform of $`_{}|_U`$.
Given a point $`mM`$, let $`\mathrm{\Sigma }_m\left(\lambda \right)`$ be a leaf of $`_\lambda `$ which passes through $`m`$. Consider $`\mathrm{\Sigma }_m\left(\lambda \right)`$ as a point of $`𝔗`$. Then $`\mathrm{\Sigma }_m:\mathrm{\Lambda }𝔗`$ is a section of the projection $`\pi `$. If it cannot lead to a confusion, denote the image of this map by the same symbol $`\mathrm{\Sigma }_m`$.
Describe in more details how the bases of $`_\lambda |_U`$ for different $`\lambda `$ fit together inside $`𝔗`$. Call a submanifold $`SM`$ a cross-sections of a foliation $``$ on $`M`$ if $`S`$ is transversal to the leaves of $``$, and each leaf of $``$ intersects $`S`$ at most once. Obviously, cross-sections exist after restriction of $``$ to an appropriate open subset $`U`$, and are identified with open subsets of the base $`𝔅_{|_U}`$.
Moreover, if $`_{}`$ is a smooth web, and $`S`$ is a cross-section to $`_{\lambda _0}`$, then for any point $`mS`$ there is a neighborhood $`US`$, $`Um`$, and a neighborhood $`V\mathrm{\Lambda }`$, $`V\lambda _0`$, such that $`U`$ is a cross-section for $`_\lambda `$, $`\lambda U`$. This gives a local identification of bases of $`_\lambda `$, $`\lambda V`$, thus a structure of a manifold on $`𝔗`$.
###### Remark 7.5.
One can show that for a smooth web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$ with a compact manifold $`\mathrm{\Lambda }`$, any point $`mM`$ has an $`_{}`$-convex neighborhood $`U`$. Different choices of $`U`$ lead to different twistor transforms, but all of them contain $`\mathrm{\Sigma }_m`$. Thus in such a case the germ of $`𝔗_{_{}|_U}`$ near $`\mathrm{\Sigma }_m`$ does not depend on $`U`$.
In fact, this germ is well-defined for any smooth web $`_{}`$. Indeed, the construction with cross-sections allows gluing local bases for $`_\lambda `$ near $`m`$ into a germ of a manifold near $`\mathrm{\Sigma }_m`$.
To simplify the following exposition, we pretend that the twistor transform is well-defined after a restriction of the web to an appropriate small open subset of $`M`$. This is always so if $`\mathrm{\Lambda }`$ is compact. The general case can be always treated honestly by switching to the language of germs.
## 8. Explicit construction of the twistor transform
In the case of codimension 1 the construction of $``$-convex subsets can be easily made explicit. Moreover, such an explicit construction would make statements in the rest of the paper simpler to formulate.
Put $`\mathrm{\Delta }(a,b)=|a|/|b|+|b|/|a|`$, $`\mathrm{\Delta }(a_1,\mathrm{},a_d)=_{1k<ld}\mathrm{\Delta }(a_k,a_l)`$. Consider the following condition on a $`1`$-form $`\alpha `$ on $`U𝕍^d`$, $`0U`$:
(8.1)
$$\underset{\begin{array}{c}k=1\\ l=1\end{array}}{\overset{d}{}}\left|\frac{\alpha _k}{x_l}\right|^2\frac{1}{E\mathrm{\Delta }(a_1,\mathrm{},a_d)^Pr^2}\underset{k=1}{\overset{d}{}}\left|\alpha _k\right|^2$$
here $`E`$, $`r`$ and $`P`$ are numbers, $`\alpha _k(x_1,\mathrm{},x_d)`$, $`k=1,\mathrm{},d`$, are components if $`\alpha `$, and $`a_k=\alpha _k(0,\mathrm{},0)`$, $`1kd`$. This condition makes sense if $`a_1,\mathrm{},a_d0`$, but if $`P=0`$, then it makes sense for any $`\alpha `$.
The following lemma is not surprising:
###### Lemma 8.1.
Fix an integer $`d>0`$. Consider a $`1`$-form $`\alpha `$ defined on $`𝔹_r^d𝕍^d`$ and a foliation $``$ on $`𝔹_r^d`$ of codimension 1. Suppose that $`\alpha |_00`$, and $`\alpha \left(x\right)`$ is normal to $`L_x`$ for any $`x𝔹_r^d`$; here $`L_x`$ is the leaf of $``$ which passes through $`x`$. There are numbers $`D,E>0`$ (which depend on $`d`$ only) such that for any $`0<\rho <r/D`$
1. if $`\alpha `$ satisfies (8.1) with $`P=0`$ in $`𝔹_r^d`$, then $`𝔹_\rho ^d`$ is strictly $``$-convex;
2. if $`\alpha `$ satisfies (8.1) with $`P=2`$ in $`𝔹_r^d`$, and $`a_k0`$, $`1kd`$, then $`\left(𝔹_\rho ^1\right)^d`$ is strictly $``$-convex;
###### Proof.
Transposing coordinates $`x_k`$, one can ensure that $`|a_d||a_k|`$, $`k=1,\mathrm{},d1`$. Changing $`\alpha `$ to $`\alpha /a_d`$ allows us to assume that $`a_d=1`$.
Obviously, one can find $`D`$ and $`E`$ such that the condition above implies that in $`𝔹_r^d`$ one has $`|\alpha _d1|1/2`$ and $`_{k=1}^{d1}|\alpha _k|^22d`$. Consequently, in $`\left(𝔹_{r/D}^{d1}\times 𝕍^1\right)𝔹_r^d`$ one can write any leaf of $``$ which passes through $`(0,\mathrm{},0,c)`$, $`|c|<4\sqrt{d}r/D`$, as $`x_d=\phi _c(x_1,\mathrm{},x_{d1})`$, and $`_{k=1}^{d1}|\phi _c/x_k|^2<3\sqrt{d}`$. Thus one can include $`𝔹_\rho ^d`$ and $`\left(𝔹_\rho ^1\right)^d`$ into a chart-like subset of $`𝔹_r^{d1}\times 𝕍^1`$.
The next step is to show that the leaves intersected with $`𝔹_\rho ^d`$ or $`\left(𝔹_\rho ^1\right)^d`$ are connected. In the case of the ball it is enough to show that
$$N_c(x_1,\mathrm{},x_{d1})=|\phi _c(x_1,\mathrm{},x_{d1})|^2+\underset{k=1}{\overset{d1}{}}|x_k|^2$$
is concave on $`𝔹_{r/D}^{d1}`$ for $`|c|<4\sqrt{d}r/D`$. It is enough to show that the Hessian $`^2|\phi _c|^2/x_kx_l`$ of $`|\phi _c|^2`$ on $`𝔹_{r/D}^{d1}`$ cannot have a large negative eigenvalue under an appropriate choice of constants $`E`$ and $`D`$. This Hessian is a sum of a non-negative part $`2\left(\phi _c/x_k\right)\left(\phi _c/x_l\right)`$ and of $`2\phi _c^2\phi _c/x_kx_l`$.
In turn, it is enough to show that<sup>3</sup><sup>3</sup>3In the complex-analytic case one needs to consider $`\overline{}/x_k\overline{}x_l`$ as well as $`^2/x_kx_l`$. $`|\phi _c|^2_{\begin{array}{c}k=1\\ l=1\end{array}}^{d1}|^2\phi _c/x_kx_l|^2`$ can be made bounded by 1/16. Since $`|\phi _c|`$ can be bounded by $`7\sqrt{d}r/D`$, it is enough if we can bound second derivatives of $`\phi _c`$ as $`O\left(1/r\right)`$.
However, the estimates on $`\alpha _k`$, $`k=1,\mathrm{},n`$, given above allow one to estimate second derivatives of $`\phi _c`$ in terms of derivatives of $`\alpha _k`$. This finishes the proof of $``$-convexity in the case of the ball.
Investigate strict $``$-convexity in the case of the ball. It is clear that one can invert $`\phi _c`$ and write $`c=\psi (x_1,\mathrm{},x_d)`$. It is enough to prove that the $`\psi `$-image of a small ball is convex, which follows from the following simple
###### Lemma 8.2.
There is a number $`E`$ (which depends on $`d`$ only) such that given a function $`\psi `$ on $`𝔹_r^d`$ such that $`d\psi `$ satisfies (8.1) with $`P=0`$ in $`𝔹_r^d`$, then the image $`\psi \left(𝔹_\rho ^d\right)`$ is convex for $`0<\rho <r`$.
Investigate the case of the polydisk. The stronger assumptions we have in the polydisk case allow ensuring $`|\alpha _ka_k|<|a_k|/F`$ for any given $`F>0`$. Now the statement follows from the following
###### Lemma 8.3.
Given $`d`$, there are numbers $`F`$ and $`D`$ which satisfy the following condition. Given a smooth function $`\psi (x_1,\mathrm{},x_d)`$ defined on $`\left(𝔹_r^1\right)^d`$ and any numbers $`a_1,\mathrm{},a_d`$, and $`c`$, if $`\psi `$ satisfies
(8.2)
$$|\psi /x_ka_k|<|a_k|/F\text{}k=1,\mathrm{},d,$$
on $`\left(𝔹_r^1\right)^d`$, then $`\psi \left(\left(𝔹_\rho ^1\right)^d\right)`$ is convex, and $`\psi ^1\left(c\right)\left(𝔹_\rho ^1\right)^d`$ is connected if non-empty for any $`0<\rho <r/D`$.
###### Proof.
The statement is obvious in the real case, so assume complex-analytic situation. Start with the case $`d=1`$. Put $`D=2`$, $`F=4`$. We may assume $`r=1`$, $`a_1=1`$, then $`|\psi ^{\prime \prime }|<1/2`$ on $`𝔹_{1/2}^1`$. Thus the direction of the tangent line $`l_\tau `$ to the curve $`\psi \left(e^{i\tau }/2\right)`$ rotates counterclockwise when $`\tau `$ grows, with the angular velocity being close to 1. This implies convexity of $`\psi \left(𝔹_{r/D}^1\right)`$. The connectivity of $`\psi ^1\left(c\right)`$ is obvious.
In the case $`d>1`$ the convexity follows from similar arguments: the boundary of the image of $`\left(𝔹_\rho \right)^d`$ is the curve $`\mathrm{\Psi }\left(\tau _1\right)=\psi (e^{i\tau _1}\rho ,e^{i\tau _2\left(\tau _1\right)}\rho ,\mathrm{},e^{i\tau _d\left(\tau _1\right)}\rho )`$; here $`\tau _k`$ are appropriate functions, $`d\tau _k/d\tau 1`$, and the direction of the tangent line the curve $`\mathrm{\Psi }\left(\tau \right)`$ behaves as in the case $`d=1`$.
For connectivity proceed by induction in $`d`$. We may assume that $`|a_d||a_k|`$, $`k=1,\mathrm{},d1`$. Increasing $`F`$ and $`D`$, one can ensure that $`\psi ^1\left(c\right)\left(\left(𝔹_\rho ^1\right)^{d1}\times 𝔹_r^1\right)`$ is given by $`x_d=\phi _c(x_1,\mathrm{},x_{d1})`$ if $`c\psi \left(\left(𝔹_\rho ^1\right)^d\right)`$, and $`\phi _c`$ satisfies (8.2) with $`d1`$ taken instead of $`d`$. Thus $`\phi _c^1\left(c_1\right)\left(𝔹_\rho ^1\right)^{d1}`$ is connected if non-empty. On the other hand, $`\psi ^1\left(c\right)\left(𝔹_\rho ^1\right)^d`$ is diffeomorphic to $`\phi _c^1\left(𝔹_\rho ^1\right)\left(𝔹_\rho ^1\right)^{d1}`$. Since $`\phi _c\left(\left(𝔹_\rho ^1\right)^{d1}\right)𝔹_\rho ^1`$ is convex, it is connected, thus $`\psi ^1\left(c\right)\left(𝔹_\rho ^1\right)^d`$ is connected as well. ∎
This finishes the proof of Lemma 8.1. ∎
###### Amplification 8.4.
Consider a $`1`$-form $`\alpha `$ on $`𝔹_r^d`$, and a $`1`$-form $`\stackrel{~}{\alpha }`$ with components $`\stackrel{~}{\alpha }_k=\kappa _k\alpha _k`$, $`k=1,\mathrm{},d`$; here $`\kappa _k`$ are arbitrary numbers, some of which are non-0. Consider a foliation $``$ on $`𝔹_r^d`$ of codimension 1. Suppose that $`\stackrel{~}{\alpha }\left(x\right)`$ is normal to $`L_x`$ for any $`x𝔹_r^d`$; here $`L_x`$ is the leaf of $``$ which passes through $`x`$. Let $`a_k\stackrel{\text{def}}{=}\alpha _k(0,\mathrm{},0)0`$ for $`1kd`$.
There are numbers $`D,E>0`$ (which depend on $`d`$ only) such that for $`0<\rho <r/D`$
1. if $`\alpha `$ satisfies (8.1) with $`P=2`$ in $`𝔹_r^d`$, then $`𝔹_\rho ^d`$ is strictly $``$-convex;
2. if $`\alpha `$ satisfies (8.1) with $`P=4`$ in $`𝔹_r^d`$, then $`\left(𝔹_\rho ^1\right)^d`$ is strictly $``$-convex;
###### Proof.
Proceed similarly to the proof of Lemma 8.1. One may assume that $`\mathrm{max}_k|a_k|=1`$. Let $`A=\mathrm{min}_k|a_k|`$. With the stronger conditions of the amplification one can ensure that $`|\alpha _ka_k|/A`$ is sufficiently small in $`𝔹_r^d`$. Then the condition (8.1) give absolute bounds on derivatives of $`\alpha _k`$, both from above and from below.
Multiplying $`\kappa _k`$ by an appropriate constant, we may assume that $`\mathrm{max}_k|\kappa _k|=1`$. Then given an estimate (8.1) for $`\alpha `$, we can estimate $`_{k=1}^d\left|\stackrel{~}{\alpha }_k\right|^2`$ from below, and $`_{\begin{array}{c}k=1\\ l=1\end{array}}^d\left|\frac{\stackrel{~}{\alpha }_k}{x_l}\right|^2`$ from above in $`𝔹_r^d`$, loosing 2 units in $`P`$. In particular, $`\stackrel{~}{\alpha }`$ satisfies (8.1) with $`P=0`$ or $`P=2`$. ∎
Apply the obtained results to the nonlinear wave equation. Consider the following condition on a function $`w`$ defined on a subset $`V𝕍^3`$, $`0V`$:
(8.3)
$$\underset{\begin{array}{c}k=1\\ l=1\end{array}}{\overset{3}{}}\left|\frac{^2w}{x_kx_l}\right|^2\frac{1}{E\mathrm{\Delta }^Pr^2}\underset{k=1}{\overset{3}{}}\left|\frac{w}{x_k}\right|^2,$$
here $`E`$, $`P`$ and $`r`$ are numbers, and $`\mathrm{\Delta }=\mathrm{\Delta }(w_{x_1}(0,0,0),w_{x_2}(0,0,0),w_{x_3}(0,0,0))`$.
###### Theorem 8.5.
There are numbers $`E,D>0`$ such that given a non-degenerate solution $`w(x,y,z)`$ of Equation (0.2) which satisfies (8.3) in a ball $`𝔹_r^3`$, then there is a neighborhood $`U`$ of (0,0,0) which is strictly $`_{}`$-convex w.r.t. the Veronese web $`_{}`$ which corresponds to $`w`$; here one can take
1. $`U=𝔹_\rho ^3`$ if $`P=2`$, $`0<\rho <r/D`$;
2. $`U=\left(𝔹_\rho ^1\right)^3`$ if $`P=4`$, $`0<\rho <r/D`$.
###### Proof.
Obviously, any ball or polydisk is $``$-convex for 3 exceptional foliations $`\left\{x=\mathrm{const}\right\}`$, $`\left\{y=\mathrm{const}\right\}`$, $`\left\{z=\mathrm{const}\right\}`$ of the web. Other foliations of the web are given by $`\left\{v_{(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})}=\mathrm{const}\right\}`$; here $`v_{(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})}`$ is a non-degenerate solution of (0.5). Application of Amplification 8.4 finishes the proof. ∎
This theorem allows one to explicitly construct the twistor transform of the Veronese web $`_{}`$ associated to $`w`$. Consider the set $`U`$ of the theorem, then the manifold with points enumerating leaves of all the foliations $`_\lambda |_U`$, $`\lambda ^1`$, is the twistor transform of $`_{}`$.
Given an abstract Veronese web $`_{}`$, by Lemma 2.6 one can describe this web by a function $`w(x,y,z)`$ which, by Theorem 3.8, satisfies (0.2) for appropriate $`(A,B,C)`$. Thus one can apply the theorem above to construct the twistor transform of $`_{}`$.
## 9. Sectional coordinates
Recall that a submersion is a smooth mapping of manifolds $`f:MN`$ such that $`df|_m:𝒯_mM𝒯_{f\left(m\right)}N`$ is surjective for any $`mM`$.
###### Lemma 9.1.
Consider a complex manifold $`𝔗`$ with a submersion $`\pi `$ onto a manifold $`\mathrm{\Lambda }`$ and a submanifold $`S𝔗`$ of codimension $`r`$ such that $`\pi |_S`$ is a diffeomorphism. Given a covering $`\left\{V_i\right\}`$ of $`𝔗`$ by Stein submanifolds, there is an open subset $`US`$ and identifications of $`U\pi ^1\left(V_i\right)`$ with $`V_i\times S_i`$, $`S_i^r`$, $`S_i0`$; these identifications intertwine $`\pi `$ with the projections $`V_i\times S_iV_i`$, and send $`S\pi ^1\left(V_i\right)`$ to $`V_i\times \left\{0\right\}`$.
###### Proof.
Suppose that $`r=1`$. Consider any function $`\stackrel{~}{s}_i`$ on a neighborhood of $`S_i\stackrel{\text{def}}{=}S\pi ^1\left(V_i\right)`$ such that the vertical derivative of $`\stackrel{~}{s}_i`$ on $`S_i`$ does not vanish. Put $`s_i\stackrel{\text{def}}{=}\stackrel{~}{s}_i\stackrel{~}{s}_i\mathrm{\Sigma }_m\pi `$. Then $`(\pi ,s_i)`$ gives the required identification of a neighborhood of $`S_i`$ with a subset of $`V_i\times `$.
The existence of such a function $`\stackrel{~}{s}_i`$ follows from the fact that a neighborhood of $`\pi ^1\left(V_i\right)S`$ is Stein if $`V_i`$ is Stein. Indeed, any bundle over a Stein manifold with a fiber isomorphic to a disk $`𝔹_\epsilon ^1`$ is Stein .
In the case $`r>1`$ one needs to consider $`d`$ functions $`\stackrel{~}{s}_{i,k}`$ instead of one, and replaces $`𝔹_\epsilon ^1`$ by $`𝔹_\epsilon ^d`$ (using results of .) ∎
###### Remark 9.2.
These “abstract nonsense” arguments allow the following construction: given a twistor transform $`𝔗`$ of a complex-analytic Veronese web $`Mm`$, cover $`^1`$ by two disks $`V_{1,2}`$, and glue a neighborhood of $`\mathrm{\Sigma }_m`$ from two domains isomorphic to $`V_i\times 𝔹_\epsilon ^1`$ (with $`\pi `$ compatible with projections to $`V_i`$). The gluing function $`g`$ is going to be a mapping $`V_1\times 𝔹_\epsilon ^1(\lambda ,t)(\lambda ,g(\lambda ,t))V_2\times 𝔹_\epsilon ^1`$, with $`g(\lambda ,t)`$ defined on $`\left(V_1V_2\right)\times 𝔹_\epsilon ^1`$. In particular, the function $`g`$ determines the germ of $`𝔗^1`$ near $`\mathrm{\Sigma }_m`$ up to isomorphism. Later, in Theorem 13.12, we will see that this implies that the germ of the Veronese web near $`m`$ is determined by $`g`$ up to isomorphism.
However, if $`𝔗`$ is a twistor transform one can achieve the same result without applying the heavy machinery of complex analysis. One can explicitly construct the required coordinate systems on open subsets of $`𝔗`$.
###### Definition 9.3.
Consider a submanifold $`\gamma `$ of a manifold $`M`$ equipped with a web $`_{}`$ with a twistor transform $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$. Say that an open subset $`U𝔗`$ is compatible with $`\gamma `$, if for any $`m\gamma `$ and any $`\lambda \pi \left(U\right)`$ the leaf of $`_\lambda `$ passing through $`m`$ is in $`U`$.
Obviously, a $`\gamma `$-compatible open subset $`U𝔗`$ is diffeomorphic to $`\pi \left(U\right)\times \gamma `$. In other words, such a subset defines a local trivialization of the bundle $`\pi `$. It is clear that $`\gamma `$ and $`V\stackrel{\text{def}}{=}\pi \left(U\right)`$ determine $`U`$ uniquely.
In the rest of this section we assume that $`_{}`$ is a Veronese web. As Lemma 13.13 will show, for Veronese webs the normal bundles to sections of $`\pi `$ are not trivializable, thus in this case $`\pi \left(U\right)`$ cannot coincide with $`^1`$.
Continue assuming that $`𝔗`$ is not a germ, but a bona fide manifold.
###### Lemma 9.4.
Consider a point $`m`$ on a Veronese web $`_{}`$ on $`M`$ and a curve $`\gamma `$ passing through $`m`$. Let $`V_{m,\gamma }^1`$ consist of points $`\lambda `$ such that $`\gamma `$ is not tangent to $`L_\lambda \left(m\right)`$ at $`m`$; here $`L_\lambda \left(m\right)`$ is the leaf of $`_\lambda `$ which passes through $`m`$. Let an open subset $`V^1`$ be compactly included into $`V_{m,\gamma }`$. Then there is a neighborhood $`\gamma _1`$ of $`m`$ in $`\gamma `$ and a compatible with $`\gamma _1`$ subset $`U𝔗`$ with $`\pi \left(U\right)=V`$.
###### Proof.
If $`\lambda _0V_{m,\gamma }`$, there is a neighborhood $`V`$ of $`\lambda _0`$ and a neighborhood $`W`$ of $`m`$ such that for $`\lambda V`$ the leaves of $`_\lambda `$ are not tangent to $`\gamma `$ at any point of $`W`$, and each leaf intersects $`\gamma W`$ in at most one point. Since $`\overline{V}^1`$ is compact, one can decrease $`W`$ so that this condition is satisfied for any $`\lambda V`$. Taking $`\gamma _1=\gamma W`$, and $`U`$ to consists of leaves of $`_\lambda `$, $`\lambda V`$, which intersect $`\gamma `$ finishes the proof. ∎
###### Lemma 9.5.
The subset $`V_{m,\gamma }^1`$ of Lemma 9.4 is open, depends on $`𝒯_m\gamma `$ only, and $`^1V_{m,\gamma }`$ consists of at most $`dimM1`$ points. Given any subset $`Z^1`$ of at most $`dimM1`$ points and $`mM`$, one can find a curve $`\gamma `$ passing through $`m`$ such that $`V=^1Z`$. Different possible directions $`𝒯_m\gamma `$ correspond $`1`$-to-1 to different ways of assigning multiplicities to points of $`Z`$ with the total being $`dimM1`$.
###### Proof.
The statements of this lemma concern one tangent space $`𝒯_mM`$ only. The tangent spaces $`𝒯_mL_\lambda \left(m\right)𝒯_mM`$ are orthogonal complements to directions $`\mathrm{n}_m\left(\lambda \right)`$ in $`𝒯_m^{}M`$. Thus $`V_{m,\gamma }`$ is determined by $`𝒯_m\gamma `$ and the image of the curve $`\mathrm{n}_m:^1𝒯_m^{}M`$. This is a Veronese curve, and any two such curves are isomorphic. Thus we may replace $`𝒯_m^{}M`$ by an arbitrary vector space $`S`$ with a Veronese curve.
Take $`S`$ to be the symmetric power $`\mathrm{Sym}^{d1}𝕍^2`$, $`dimS=d`$, and let the Veronese curve consists of $`\left(d1\right)`$st powers of elements of $`𝕍^2`$. Then $`S^{}`$ can be identified with homogeneous polynomials of degree $`d1`$ of two variables (two coordinates on $`𝕍^2`$), thus $`𝒯_m\gamma 𝒯_mM=S^{}`$ provides such a polynomial $`p`$ up to a constant.
It is easy to check that $`\lambda V_{m,\gamma }^1=𝕍^2`$ iff $`p`$ does not vanish at the points of $`𝕍^2`$ in the direction of $`\lambda `$. There are at most $`\mathrm{deg}p=d1`$ such directions, and given such directions with appropriate multiplicities, one can find a polynomial $`pS^{}`$ which vanishes at these points. ∎
Now we can implement the program outlined in Remark 9.2:
###### Corollary 9.6.
Given $`mM`$, one can find two curves $`\gamma _1`$, $`\gamma _2`$ passing through $`m`$ and two open subsets $`U_1,U_2𝔗`$ compatible with $`\gamma _1`$, $`\gamma _2`$ correspondingly such that $`U_1U_2`$ is a neighborhood of the section $`\mathrm{\Sigma }_m𝔗`$.
###### Proof.
Indeed, one can find $`\gamma _1`$, $`\gamma _2`$ such that $`^1V_{m,\gamma _1}`$ is contained in a small neighborhood of 0, and $`^1V_{m,\gamma _1}`$ is contained in a small neighborhood of $`\mathrm{}`$. To finish the proof, note that $`\mathrm{\Sigma }_m\pi ^1\left(\pi U\right)U`$ for any subset $`U𝔗`$ which is compatible with a curve $`\gamma `$ passing through $`m`$. ∎
Consider two curves as in Corollary 9.6. Let $`V_1=\pi U_1`$, $`V_2=\pi U_2`$. Then $`U_1V_1\times \gamma _1`$, $`U_2V_2\times \gamma _2`$, thus identifications of $`\gamma _1`$ and $`\gamma _2`$ with $`𝔹_\epsilon ^1`$ lead the gluing function $`g(\lambda ,t)`$ as in the beginning of this section. The other way to look at $`g`$ is to consider it as a family of gluings $`\widehat{g}_\lambda :\gamma _1\gamma _2`$, $`\lambda V_1V_2`$.
Describe these gluings $`\widehat{g}_\lambda `$ in geometric terms. This description does not mention $`𝔗`$ as a manifold, thus one need not assume that $`𝔗`$ exists as a manifold.
###### Corollary 9.7.
Given a point $`m_0`$ on a Veronese web $`M`$, one can find two curves $`\gamma _1`$, $`\gamma _2`$ passing through $`m_0`$, a neighborhood $`WM`$ of $`m_0`$, and two open subsets $`V_1,V_2^1`$ such that
1. For any $`\lambda V_j`$, $`j=1,2`$, and any $`m\gamma _j`$ the leaf of $`_\lambda |_W`$ which passes through $`m`$ intersects $`\gamma _j`$ at exactly one point $`m`$ and is transversal to $`\gamma _j`$;
2. For any $`\lambda V_1V_2`$, and any $`m\gamma _1`$ the leaf of $`_\lambda |_W`$ which passes through $`m`$ intersects $`\gamma _2`$; denote the (unique) point of intersection by $`\widehat{g}_\lambda \left(m\right)`$;
3. $`V_1V_2=^1`$.
The germ near $`\left(V_1V_2\right)\times \left\{m\right\}`$ of the function $`\widehat{g}_{}:\left(V_1V_2\right)\times \gamma _1\gamma _2`$ uniquely determines the germ of the twistor transform $`𝔗`$ of $`M`$ near the section $`\mathrm{\Sigma }_{m_0}`$ and the germ of $`_{}`$ near $`m`$. For any $`0<\epsilon <1`$ one can ensure that $`V_1\{z|z|>\epsilon \}`$, $`V_2\{z|z|<1/\epsilon \}`$.
## 10. Explicit construction of the gluing function
In conditions of Corollary 9.7 identify a neighborhood of $`m_0`$ in $`\gamma _1`$ with $`𝔹_\epsilon ^1`$, and a neighborhood of $`m_0`$ in $`\gamma _2`$ with a subset of $``$. This would make the gluing function $`g(\lambda ,t)`$ into a function $`\left(V_1V_2\right)\times 𝔹_\epsilon ^1`$. A different choice of identifications would lead to $`\stackrel{~}{g}(\lambda ,t)=f\left(g(\lambda ,F\left(t\right))\right)`$ for appropriate invertible functions $`f\left(z\right)`$, $`F\left(z\right)`$.
Describe $`g(\lambda ,t)`$ in terms of the function $`w(x,y,z)`$ which identifies the Veronese web. Later, in Appendix 15, we will see that the gluing function should depend only on the restriction of $`w`$ and first derivatives of $`w`$ to an appropriate surface. Here we prove this only in the case of surfaces of a special form.
###### Theorem 10.1.
Consider a complex-analytic non-degenerate solution $`w(x,y,z)`$ of the nonlinear wave equation (0.2) defined in a neighborhood of (0,0,0). Fix $`0<r<1`$, $`\lambda _1,\lambda _2,\lambda _3^1`$, $`|\lambda _{1,2}|<r`$, $`|\lambda _3|>1/r`$. Let $`Y\left(x\right)`$ be any function such that $`Y^{}\left(x\right)=dY/dx`$ is nowhere 0, and $`Y\left(0\right)=0`$. Consider the following family of ODEs with a parameter $`\mu `$ on a function $`z\left(x\right)`$:
$$\frac{dz}{dx}=\frac{Aw_x(x,Y\left(x\right),z)}{\mu Cw_z(x,Y\left(x\right),z)}\frac{Bw_y(x,Y\left(x\right),z)}{\left(\mu 1\right)Cw_z(x,Y\left(x\right),z)}Y^{}\left(x\right);$$
Let $`g_\mu \left(t\right)`$ be $`z\left(0\right)`$; here $`z\left(x\right)`$ is the solution of this equation with the initial data $`z\left(t\right)=0`$. Then for any $`\epsilon _1>0`$ one can find an appropriate $`\delta >0`$ so that the function $`g_\mu \left(t\right)`$ is correctly defined if $`|\mu |>\epsilon _1`$, $`|\mu 1|>\epsilon _1`$, and $`|t|<\delta `$.
Consider $`\epsilon `$ such that $`r<\epsilon <1`$. Define a surface $`\stackrel{~}{𝔗}`$ by gluing $`𝔹_{1/\epsilon }^1\times 𝔹_\delta ^1`$ and $`\left(^1\overline{𝔹}_\epsilon ^1\right)\times `$ via $`𝔹_{1/\epsilon }^1\times 𝔹_\delta ^1(\lambda ,t)(\lambda ,\stackrel{~}{g}(\lambda ,t))\left(^1\overline{𝔹}_\epsilon ^1\right)\times `$, $`\epsilon <|\lambda |<1/\epsilon `$, $`|t|<\delta `$; here $`\stackrel{~}{g}(\lambda ,t)=g_\mu \left(t\right)`$, $`\mu =(\lambda _1:\lambda _2:\lambda _3:\lambda )`$, and $`\delta `$ corresponds to $`\epsilon _1`$ such that $`|\mu |>\epsilon _1`$ and $`|\mu 1|>\epsilon _1`$ if $`|\lambda |>\epsilon `$. Since $`g_\mu \left(0\right)0`$, $`\stackrel{~}{𝔗}`$ has a section $`\stackrel{~}{\mathrm{\Sigma }}_{(0,0,0)}=\left\{(\lambda ,0)\right\}`$. Coordinates $`\lambda `$ glue into a projection $`𝔗^1`$.
Suppose that
(10.1)
$$\left|\frac{Q\lambda _1\lambda _2}{Q1}\right|>1/\epsilon ,Q=Y^{}\left(0\right)\frac{Bw_y(0,0,0)}{Aw_x(0,0,0)}.$$
Then the germ of $`\stackrel{~}{𝔗}`$ near $`\stackrel{~}{\mathrm{\Sigma }}_{(0,0,0)}`$ is isomorphic to the germ of the twistor transform $`𝔗`$ of the Veronese web associated<sup>4</sup><sup>4</sup>4As in Remark 3.9. to $`w(x,y,z)`$ near $`\mathrm{\Sigma }_{(0,0,0)}`$.
###### Proof.
Consider the $`3`$-dimensional Veronese web $`M`$ associated to $`w(x,y,z)`$ such that the foliations $`\left\{x=\mathrm{const}\right\}`$, $`\left\{y=\mathrm{const}\right\}`$, $`\left\{z=\mathrm{const}\right\}`$ are associated to $`\lambda =\lambda _1`$, $`\lambda =\lambda _2`$, $`\lambda =\lambda _3`$. Take $`m_0=(0,0,0)`$, $`\gamma _2`$ to be the $`z`$-axis. Then the subset $`V_{m_0,\gamma _2}`$ (in notations of Lemma 9.4) is $`^1\{\lambda _1,\lambda _2\}^1\overline{𝔹}_\epsilon ^1`$, since $`\gamma _2`$ is an intersection of a leaf of $`_{\lambda _1}`$ and of a leaf of $`_{\lambda _2}`$. Similarly, for a curve $`\gamma `$ in $`xy`$-plane the subset $`V_{m,\gamma }`$ is $`^1\{\lambda _3,\lambda \left(m\right)\}`$; here $`\lambda \left(m\right)=\lambda _1`$ for the curves $`x=\mathrm{const}`$ in $`xy`$-plane, $`\lambda \left(m\right)=\lambda _2`$ for the curves $`y=\mathrm{const}`$ in $`xy`$-plane. It is clear that for a curve with any other direction $`\lambda \left(m\right)\lambda _1`$ and $`\lambda \left(m\right)\lambda _2`$. In particular, it is so for the curve $`\gamma _1`$ given by $`y=Y\left(x\right)`$. Thus $`V_{m_0,\gamma _1}V_{m_0,\gamma _2}=^1`$. Thus $`\gamma _1`$, $`\gamma _2`$ satisfy conditions of Corollary 9.7, thus one can glue the twistor transform $`𝔗`$ from two open subsets, one being a bundle over $`V_{m_0,\gamma _1}`$, another over $`V_{m_0,\gamma _2}`$.
Moreover, $`V_{m_0,\gamma _1}^1\overline{𝔹}_\epsilon ^1`$, and if $`|\lambda \left(m_0\right)|>1/\epsilon `$, then $`V_{m_0,\gamma _2}𝔹_{1/\epsilon }^1`$. In such a case $`𝔗`$ can be glued from two open subsets, one being a bundle over $`^1\overline{𝔹}_\epsilon ^1`$, another over $`𝔹_{1/\epsilon }^1`$. To describe $`𝔗`$, it is enough to describe the gluing function $`g(\lambda ,t)`$, $`\epsilon <|\lambda |<1/\epsilon `$, for small $`t`$. Taking $`z`$ as the coordinate on $`\gamma _2`$ and $`x`$ as the coordinate on $`\gamma _1`$, one can describe this gluing function in the following way: take a point $`m=(t,Y\left(t\right),0)`$ on $`\gamma _1`$, find the leaf of $`_\lambda `$ which passes through $`m`$, and intersect this leaf with $`\gamma _2`$. Then $`g(\lambda ,t)`$ is the $`z`$-coordinate of the point of intersection.
Consider the surface $`N`$ given by the equation $`y=Y\left(x\right)`$. The foliation $`_\lambda `$ can be described by the equations $`v(x,y,z)=\mathrm{const}`$; here the derivative of $`v`$ is given by Corollary 3.6. The curves cut out by this foliation on $`N`$ have both $`(dY/dx,1,0)`$ and $`(v_x,v_y,v_z)`$ as normal vectors. Thus these curves are tangent to directions
$$(p_3\left(\lambda \right)w_z,p_3\left(\lambda \right)w_zdY/dx,p_1\left(\lambda \right)w_xp_2\left(\lambda \right)w_ydY/dx)$$
(notations as in (3.1)). One can easily check that the ODE of the theorem describes $`xz`$-projections of these curves for $`\mu =(\lambda _1:\lambda _2:\lambda _3:\lambda )`$. Thus $`g(\lambda ,t)=\stackrel{~}{g}(\lambda ,t)`$.
The only thing to prove is $`|\lambda \left(m_0\right)|>1/\epsilon `$. In fact $`\lambda \left(m_0\right)=\frac{Q\lambda _1\lambda _2}{Q1}`$. To check this, it is enough to find the intersection of the leaf of $`_\lambda `$ through (0,0,0) with $`z=0`$. As above, the direction of this curve is given by $`(v_y,v_x,0)=(p_2\left(\lambda \right)w_y,p_1\left(\lambda \right)w_x,0)`$. Again, it is easy to check that this agrees with (10.1). ∎
###### Remark 10.2.
Obviously, the condition (10.1) is satisfied in $`Y^{}\left(0\right)`$ is inside a non-empty disk in $`^1`$. In fact, there is a canonical choice of $`Y\left(x\right)`$ which automatically satisfies (10.1). Indeed, the condition $`\lambda \left(m\right)=\lambda _3`$ gives a direction field on $`xy`$-plane, take an integral curve of this direction field. Explicitly,
(10.2)
$$\frac{dY}{dx}=\frac{Aw_x(x,Y,0)}{Bw_y(x,Y,0)},Y\left(0\right)=0.$$
###### Remark 10.3.
If $`\epsilon `$ with the properties required in the theorem does not exist, by decreasing $`\delta `$ one can ensure that the set of values of $`\lambda `$ for which (10.1) does not hold is in a small disk $`D`$ which does not contain $`\lambda _1`$ and $`\lambda _2`$. If there is a circle on $`^1`$ which separates $`\{\lambda _1,\lambda _2)`$ from $`\lambda _3`$ and $`D`$, then one can use this circle instead of $`\left\{|z|=1\right\}`$ in Theorem 10.1.
If there is no such circle, then $`\frac{BY^{}\left(0\right)w_y(0,0,0)}{Bw_y(0,0,0)Y^{}\left(0\right)Aw_x(0,0,0)}`$ is real and is between 0 and 1. In particular, by a projective transform of $`^1`$ one can make $`\lambda _1,\lambda _2,\lambda _3`$ real, and the disk $`D`$ centered on the real axis between $`\lambda _1`$ and $`\lambda _2`$. If additionally $`w(x,y,z)`$ is real for real $`x,y,z`$, and $`A,B,C`$ are real, then the real $`(A,B,C)`$-equation is hyperbolic near (0,0,0) w.r.t. the surface $`y=Y\left(x\right)`$. Thus this case is of special interest.
In such a case it is hard to describe $`𝔗`$ by representing $`^1`$ as a union of two disks, but one can glue $`𝔗`$ using the same function $`g_\mu \left(t\right)`$ if one covers $`^1`$ by two regions of more complicated form. For example, consider small disks $`D_{1,2,3}`$ centered at $`\lambda _{1,2,3}`$, consider a contour $`L`$ which goes along the line $`\mathrm{Im}\lambda =0`$ with the exceptions of going around $`D_1`$ and $`D_2`$ from above, and around $`D`$ and $`D_3`$ from below. The function $`g(\lambda ,t)`$ is still correctly defined for $`\lambda `$ near $`L`$, thus one can describe $`𝔗`$ by gluing neighborhoods of the regions above $`L`$ and below $`L`$.
Note that for the values of $`\lambda L`$ which are on the real axis the function $`g(\lambda ,t)`$ can be defined in terms of solving a real ODE.
## 11. Equipped twistor transforms and infinitesimal families
For a mapping $`\pi :MN`$ denote by $`\mathrm{\Gamma }(N,\pi )`$ the set of sections of $`\pi `$, i.e., of right inverse mappings to $`\pi `$.
###### Definition 11.1.
Given a web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$ with the twistor transform $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$, consider the family $`\left\{\mathrm{\Sigma }_m\right\}_{mM}`$ of sections of $`\pi `$. The equipped twistor transform of $`_{}`$ is the mapping $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$ together with a family of sections $`\left\{\mathrm{\Sigma }_m\right\}_{mM}`$.
Given such a structure $`(𝔗,\mathrm{\Lambda },\pi ,M,\mathrm{\Sigma }_{})`$, and $`\lambda \mathrm{\Lambda }`$, consider a mapping $`\mathrm{s}_\lambda :M\pi ^1\left(\lambda \right):m\mathrm{\Sigma }_m\left(\lambda \right)`$. If this structure comes from an equipped twistor transform, this mapping is a submersion.
###### Lemma 11.2.
Consider a submersion $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$ together with a family of sections $`\left\{\mathrm{\Sigma }_m\right\}_{mM}`$ parameterized by a manifold $`M`$. If the rank of differential $`d\mathrm{s}_\lambda |_m`$ of the mapping $`\mathrm{s}_\lambda `$ does not depend on $`m`$ and $`\lambda `$, then $`M`$ is equipped with a canonically defined web structure $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$, the leaf $`L_{\lambda ,m_0}`$ of $`_\lambda `$, $`\lambda \mathrm{\Lambda }`$, which passes through $`m_0M`$ consists of points $`mM`$ such that $`\mathrm{\Sigma }_m\left(\lambda \right)=\mathrm{\Sigma }_{m_0}\left(\lambda \right)`$.
If $`(𝔗,\mathrm{\Lambda },\pi ,M,\mathrm{\Sigma }_{})`$ is a twistor transform of a web $`\stackrel{~}{}_{}`$ on $`M`$, then $`\stackrel{~}{}_{}=_{}`$.
###### Proof.
Indeed, mappings with constant rank of the differential are submersions onto their images, thus preimages of points are foliations on $`M`$. The other statements are obvious. ∎
It is clear that in the conditions of the lemma if $`\mathrm{s}_\lambda `$ is not of maximal possible rank (i.e., is not a submersion), then $`𝔗^{}=_\lambda \mathrm{Im}\mathrm{s}_\lambda `$ is a submanifold of $`𝔗`$, and $`(𝔗^{},\mathrm{\Lambda },\pi ^{},M,\mathrm{\Sigma }_{})`$ is the twistor transform of $`_{}`$; here $`\pi ^{}=\pi |_𝔗^{}`$.
Lemma 11.2 shows that one can reconstruct a web on $`M`$ by its equipped twistor transform $`(𝔗,\mathrm{\Lambda },\pi ,M,\mathrm{\Sigma }_{})`$. In fact in many cases to reconstruct the web one needs much less data than $`(𝔗,\mathrm{\Lambda },\pi ,M,\mathrm{\Sigma }_{})`$. Later, in Section 13, we explain when the same information is contained in $`(𝔗,\mathrm{\Lambda },\pi )`$, at least if one considers $`M`$ up to isomorphism. Illustrate this by several weaker statements.
Suppose that the mapping $`\mathrm{\Sigma }_{}:M\mathrm{\Gamma }(\mathrm{\Lambda },\pi ):m\mathrm{\Sigma }_m`$ is injective, in other words, $`_{}`$ is separating. In such cases $`M`$ as a set is identified with $`\mathrm{Im}\mathrm{\Sigma }_{}`$. In fact $`\mathrm{\Gamma }(\mathrm{\Lambda },\pi )`$ has a natural topology, and if $`\mathrm{\Sigma }`$ is a homeomorphism on its image, then the topology on $`M`$ can be also reconstructed basing on $`\mathrm{Im}\mathrm{\Sigma }_{}\mathrm{\Gamma }(\mathrm{\Lambda },\pi )`$. In such a case if we are interested in $`(M,_{})`$ up to homeomorphism, it may be reconstructed given $`(𝔗,\mathrm{\Lambda },\pi ,\mathrm{Im}\mathrm{\Sigma }_{})`$.
One should expect that the same argument will work for diffeomorphisms as far as the differential of $`\mathrm{\Sigma }_{}`$ is injective. However, in general $`\mathrm{\Gamma }(\mathrm{\Lambda },\pi )`$ is not finite-dimensional, thus this question is a little bit more subtle. However, it is relatively easy to describe what is an individual tangent space to $`\mathrm{\Gamma }(\mathrm{\Lambda },\pi )`$. This tangent space is going to be the target of the differential of $`\mathrm{\Sigma }_{}`$.
###### Definition 11.3.
Given a section $`\mathrm{\Sigma }`$ of submersion $`\pi :𝔗\mathrm{\Lambda }`$, the tangent space to $`\mathrm{\Gamma }(\mathrm{\Lambda },\pi )`$ at $`\mathrm{\Sigma }`$ is the vector space $`\mathrm{\Gamma }(S,𝒩S)`$, $`S=\mathrm{Im}\mathrm{\Sigma }`$. Call elements of $`\mathrm{\Gamma }(S,𝒩S)`$ infinitesimal deformations. Given a family $`\left\{\mathrm{\Sigma }_m\right\}_{mM}`$ of sections of $`\pi `$, the infinitesimal family of $`\left\{\mathrm{\Sigma }_m\right\}`$ at $`m_0M`$ is the naturally defined mapping $`d\mathrm{\Sigma }|_{m_0}:𝒯_{m_0}M\mathrm{\Gamma }(\mathrm{\Sigma }_{m_0},𝒩\mathrm{\Sigma }_{m_0})`$. Say that a family $`\left\{\mathrm{\Sigma }_m\right\}`$ is immersive if $`d\mathrm{\Sigma }|_m`$ is a monomorphism for any $`mM`$.
Describe what is $`d\mathrm{\Sigma }|_m`$ and what is the geometric meaning of this definition. To define $`d\mathrm{\Sigma }|_m`$, it is enough to consider the case $`dimM=1`$. A smooth $`1`$-parametric family $`\sigma _t`$, $`tT𝕍^1`$, of sections of $`\pi `$ is a mapping $`\sigma :\mathrm{\Lambda }\times T𝔗`$ such that $`\pi \sigma `$ coincides with the projection $`p_1:\mathrm{\Lambda }\times T\mathrm{\Lambda }`$. Given $`\sigma `$ and $`tT`$, consider the derivatives $`d\sigma |_{(\lambda ,t)}`$ at points of $`\mathrm{\Lambda }\times \left\{t\right\}`$. Clearly, $`d\sigma |_{(\lambda ,t)}`$ maps $`𝒯_\lambda \mathrm{\Lambda }𝒯_t𝕍^1`$ to $`𝒯_{\sigma (\lambda ,t)}𝔗`$. It can be split into a direct sum of a mapping $`d\sigma |_{(\lambda ,t)}^{\left(1\right)}:𝒯_\lambda \mathrm{\Lambda }𝒯_{\sigma (\lambda ,t)}𝔗`$ and $`d\sigma |_{(\lambda ,t)}^{\left(2\right)}:𝒯_t𝕍^1𝒯_{\sigma (\lambda ,t)}𝔗`$.
Note that the condition $`\pi \sigma =p_1`$ determines some components of $`d\sigma |_{(\lambda ,t)}`$. Indeed, consider $`S=\mathrm{Im}\sigma (,t)`$. It is a submanifold of $`𝔗`$. Given $`\lambda \mathrm{\Lambda }`$, the vector space $`𝒯_{\sigma (\lambda ,t)}𝔗`$ can be decomposed into a direct sum of tangent spaces to $`\pi ^1\left(\lambda \right)`$ and to $`S`$. Denote components of $`v𝒯_{\sigma (\lambda ,t)}𝔗`$ in this decomposition by $`v^{\text{vert}}`$ and $`v^{\text{hor}}`$. In particular, the mappings $`d\sigma ^{\left(1\right)}`$, $`d\sigma ^{\left(2\right)}`$ can be further subdivided into $`d\sigma ^{\left(1\right)\text{vert}}`$, $`d\sigma ^{\left(2\right)\text{vert}}`$, $`d\sigma ^{\left(1\right)\text{hor}}`$, $`d\sigma ^{\left(2\right)\text{hor}}`$. It is clear that given two families $`\sigma `$ and $`\stackrel{~}{\sigma }`$, if vertical components of $`d\sigma `$ and $`d\stackrel{~}{\sigma }`$ coincide, then $`d\sigma `$ and $`d\stackrel{~}{\sigma }`$ coincide. Moreover, $`d\sigma ^{\left(1\right)\text{vert}}`$ obviously vanishes. In particular, the only “interesting” part of differential of $`\sigma `$ is $`d\sigma ^{\left(2\right)\text{vert}}`$.
On the other hand, the vertical component of $`v𝒯_{\sigma (\lambda ,t)}𝔗`$ can be also naturally identified with an element of the quotient by the vector subspace of horizontal sections $`𝒯_{\sigma (\lambda ,t)}𝔗/𝒯_{\sigma (\lambda ,t)}S=𝒩_{\sigma (\lambda ,t)}S`$, i.e., with a normal vector to $`S`$ at $`\sigma (\lambda ,t)`$. Since $`d\sigma ^{\left(2\right)\text{vert}}`$ sends $`\delta t𝒯_{t_0}𝕍^1`$ to a normal vector to $`S`$ at $`\sigma (\lambda ,t)`$ for each $`\lambda \mathrm{\Lambda }`$, it associates to $`\delta t`$ a section of the normal bundle $`𝒩S`$.
The following statement is obvious:
###### Lemma 11.4.
The equipped twistor transform of a web is immersive iff the web is separating.
It is clear that for an immersive family $`\mathrm{\Sigma }_m`$, $`mM`$, the mappings $`\mathrm{s}_\lambda `$, $`\lambda \mathrm{\Lambda }`$, separate points on small open subsets of $`M`$ (even infinitesimally). Thus the structure of the manifold on $`M`$ is reconstructed from the mapping of the set $`M`$ to $`\mathrm{\Gamma }(\mathrm{\Lambda },\pi )`$.
###### Corollary 11.5.
Consider a weakly separating and separating web $`_{}`$ on $`M`$. Then $`_{}`$ can be reconstructed up to a diffeomorphism by the twistor transform $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$ of $`_{}`$ together with the subset $`\mathrm{\Gamma }(\mathrm{\Lambda },\pi )`$ consisting of sections which correspond to points of $`M`$.
## 12. Kodaira–Spencer deformation of a section
In the classification of complex-analytic Veronese webs the central role is played by the following corollary<sup>5</sup><sup>5</sup>5Since one-dimensional Cauchy–Riemann equations are not overdetermined, in the case $`dim\mathrm{\Lambda }=1`$ we are most interested in Kodaira–Spencer deformation theory can be replaced by an argument involving an implicit function theorem (in normed spaces). of Kodaira–Spencer deformation theory (for example, see ).
###### Definition 12.1.
Say that a vector bundle $`E`$ over a topological space $`\mathrm{\Lambda }`$ is cohomologically trivial if $`H^k(\mathrm{\Lambda },E)=0`$ for $`k>0`$.
###### Theorem 12.2.
Consider an $`n`$-dimensional complex manifold $`𝔗`$ equipped with a surjective submersion $`\pi :𝔗\mathrm{\Lambda }`$, and with a section $`\mathrm{\Sigma }:\mathrm{\Lambda }𝔗`$ of the projection $`\pi `$. Let $`S=\mathrm{Im}\mathrm{\Sigma }`$, suppose that $`𝒩S`$ is cohomologically trivial, and $`\mathrm{\Lambda }`$ is compact. Then there is a connected complex manifold $`M`$, a mapping $`\sigma :\mathrm{\Lambda }\times M𝔗`$, and a neighborhood $`U`$ of $`S`$ in $`𝔗`$ such that
1. $`\pi \sigma `$ coincides with the projection $`\mathrm{\Lambda }\times M\mathrm{\Lambda }`$;
2. for any section $`s`$ of $`\pi |_U`$ there is unique $`mM`$ such that $`s=\sigma |_{\mathrm{\Lambda }\times \left\{m\right\}}`$; denote by $`m_0M`$ the point which corresponds to $`s=\mathrm{\Sigma }`$;
3. the infinitesimal family<sup>6</sup><sup>6</sup>6See Definition 11.3. $`d\sigma |_{m_0}:𝒯_{m_0}M\mathrm{\Gamma }(S,𝒩S)`$ is a bijection.
###### Remark 12.3.
To translate to the usual formulation of deformation theory, instead of deforming the mapping $`\mathrm{\Sigma }`$, one should deform the submanifold $`S`$. Then the first condition on $`\sigma `$ disappears (is just gives a normalization by identifying the deformed submanifold with $`\mathrm{\Lambda }`$), the second one identifies $`M`$ with the moduli set of those submanifolds in $`U𝔗`$ which project $`1`$-to-1 to $`\mathrm{\Lambda }`$. The fact that the set $`M`$ can be equipped with a structure of a manifold is the most nontrivial part of the statement. If $`𝔗`$ is in fact a total space of a vector bundle $``$ over $`\mathrm{\Lambda }`$, then this statement is trivial, with $`M=\mathrm{\Gamma }(^1,)`$.
Additionally, the existence of the projection on $`\mathrm{\Lambda }`$ (thus of retraction on $`S`$) removes all the bulkiness from the statement on a deformation of an arbitrary submanifold, since one does not need to consider the deformation of the the complex structure on $`S`$.
###### Remark 12.4.
One should interpret the last statement of the theorem as the fact that any infinitesimal deformation is a infinitesimal family of an actual $`1`$-parameter deformation of $`S`$. Compare this with Definition 11.3.
In our discussion we are most interested in the case $`dim𝔗=2`$, $`\mathrm{\Lambda }=^1`$. Then $`𝒩S`$ is a line bundle, thus is isomorphic to $`𝒪\left(d1\right)`$ with $`d0`$, and $`dimM=d`$. In fact we need a particular case $`d=3`$, but for some time we are going to discuss the general case of arbitrary $`d`$, $`𝔗`$ and $`\mathrm{\Lambda }`$.
###### Definition 12.5.
Say that a mapping $`\pi :𝔗\mathrm{\Lambda }`$ of complex manifolds is a disk bundle if $`𝔗`$ is a manifold with $`C^0`$-boundary, $`dim𝔗=dimS+1`$, for any $`\lambda \mathrm{\Lambda }`$ there is a neighborhood $`U\lambda `$ such that $`\pi |_{\pi ^1U}`$ is homeomorphic to the projection $`p_1:U\times DU`$; here $`D`$ is $`\{z|z|1\}.`$
###### Proposition 12.6.
In the conditions of Theorem 12.2 assume that $`dim𝔗=2`$, $`\mathrm{\Lambda }=^1`$, and that $`\pi `$ is a disk bundle. Consider two curves $`\gamma _{1,2}𝔗`$ such that restrictions $`\pi |_{\gamma _1}`$ and $`\pi |_{\gamma _2}`$ are bijections. Suppose that $`d=\mathrm{deg}\left(𝒩\gamma _1\right)+1`$, and $`\gamma _1`$ intersects $`\gamma _2`$ in $`d`$ points. Then $`\gamma _1=\gamma _2`$.
###### Proof.
Suppose $`\gamma _1\gamma _2`$. Let $`X_1,\mathrm{},X_k`$ be the points of intersection of $`\gamma _1`$ and $`\gamma _2`$. Let $`\overline{𝔗}`$ be blow-up of $`𝔗`$ at these points (make repeated blow-ups if needed to remove all the points of intersection). Removing proper preimages of $`\pi ^1\pi \left(X_i\right)`$, $`i=1,\mathrm{},k`$, from $`\overline{𝔗}`$, we obtain a manifold $`\stackrel{~}{𝔗}`$ with a mapping $`\stackrel{~}{\pi }`$ to $`^1`$ such that preimages of points of $`^1`$ are disks, with the exception of the points $`\pi \left(X_i\right)`$, preimages of which are isomorphic to $`^1\{\}`$. Cutting out far-away points of $``$ together with an appropriate neighborhood on $`\stackrel{~}{𝔗}`$, one may ensure that the resulting manifold is a disk bundle over $`^1`$.
Each blow-up decreases the degree of the normal bundle by 1, thus we reduced the statement to the case $`d<0`$, and $`\gamma _1\gamma _2=\mathrm{}`$. Show that this leads to contradiction.
Indeed, topological bundles with the fibers being oriented disks are isomorphic iff their boundaries are isomorphic as bundles with a fiber being oriented circles. In turn, any such bundle is isomorphic to a spherical bundle of a line bundle over $`^1`$, which is determined by its degree up to an isomorphism. We conclude that the topological bundle $`𝔗^1`$ is isomorphic to a neighborhood of $`0`$-section in the total space of $`𝒪\left(n\right)`$, $`n>0`$. However, $`𝒪\left(n\right)`$ has no continuous nowhere-0 sections: indeed, such a section would give a trivialization of the spherical bundle of $`𝒪\left(n\right)`$, thus, due to arguments given above, to an isomorphism of $`𝒪\left(n\right)`$ with $`𝒪\left(0\right)`$. ∎
###### Remark 12.7.
The condition of being a disk bundle is very essential. For example, suppose that $`𝔗`$ is an open subset of $`^1\times ^1`$ with $`\pi `$ being the projection on the first $`^1`$. It is easy to find such an $`𝔗`$ which contains both the “constant” section $`x0`$ of $`\pi `$, and the $`\mathrm{id}`$-section $`xx`$. Moreover, for most points of $`^1`$ the preimage in $`𝔗`$ can be made a disk. Thus a topological argument is required indeed.
The next step is to provide a way to find the subset $`U`$ of Theorem 12.2 if all we new is the family $`\sigma `$.
###### Proposition 12.8.
In the conditions of Theorem 12.2 assume that $`dim𝔗=2`$, $`\mathrm{\Lambda }=^1`$, and that $`\pi `$ is a disk bundle. Let $`\mathrm{deg}\left(𝒩S\right)=d1`$, $`\{\lambda _1,\mathrm{},\lambda _d\}^1`$ be a set of $`d`$ distinct points. Let $`B_k=\pi ^1\lambda _k`$, $`U_k`$ be an open subset of $`B_k`$, $`k=1,\mathrm{},d`$. Let $`\stackrel{~}{\sigma }`$ be a mapping $`^1\times M𝔗`$ such that $`\pi \stackrel{~}{\sigma }`$ is the projection $`^1\times M^1`$. Suppose that for any collection $`X=\left\{X_k\right\}_{k=1}^d𝔗`$ such that $`X_kU_k`$ there is $`m_XM`$ such that $`\mathrm{Im}\left(\stackrel{~}{\sigma }_{m_X}\right)B_k=X_k`$, $`k=1,\mathrm{},d`$.
Let $`U=𝔗\left(_k\left(B_kU_k\right)\right)`$ (in other words, narrow fibers $`B_k`$ over $`X_k`$ to become $`U_k`$). Then for any curve $`\gamma U`$ which projects isomorphically to $`^1`$ there is $`mM`$ such that $`\gamma =\mathrm{Im}\stackrel{~}{\sigma }_m`$.
###### Proof.
Take $`X_k=\gamma B_k`$, and apply Proposition 12.6 to $`\gamma `$ and $`\mathrm{Im}\left(\stackrel{~}{\sigma }_{m_X}\right)`$. ∎
###### Remark 12.9.
This proposition provides a way to check that a given family $`\stackrel{~}{\sigma }`$ and $`U𝔗`$ may work as the family $`\sigma `$ from Theorem 12.2. Note that given $`\stackrel{~}{\sigma }`$ which satisfies the last condition of Theorem 12.2, it is always possible to find the subsets $`U_k`$ with the required properties. Indeed, if $`𝔗`$ is an open subset of the total space of $`𝒪\left(d1\right)`$, then this follows from the fact that a section of $`𝒪\left(d1\right)`$ is uniquely determined by values in $`d`$ different points (compare with Legendre interpolation formula, or Vandermond determinant). In general one needs to apply the implicit function theorem to the mapping $`m\mathrm{\Sigma }_m=\mathrm{Im}\stackrel{~}{\sigma }|_{\mathrm{\Lambda }\times \left\{m\right\}}`$.
Moreover, for the manifolds we are going to consider here (twistor transforms of Veronese webs) we can provide an explicit description of the family $`\sigma `$ and of subsets $`U_k`$.
###### Corollary 12.10.
In the conditions of Theorem 8.5, consider the twistor transform $`𝔗\stackrel{𝜋}{}^1`$ of the $`_{}`$-convex subset $`\left(𝔹_\rho ^1\right)^3`$. Let $`\mathrm{\Sigma }_m`$ be the section of $`𝔗`$ corresponding to $`m\left(𝔹_\rho ^1\right)^3`$. Then the mapping $`\stackrel{~}{\sigma }(m,\lambda )\stackrel{\text{def}}{=}\mathrm{\Sigma }_m\left(\lambda \right)`$, $`m\left(𝔹_\rho ^1\right)^3`$, $`\lambda ^1`$, satisfies the conditions of Proposition 12.8 with $`U_k=B_k`$.
###### Proof.
Let $`\lambda _{1,2,3}`$ be the values of $`\lambda `$ which correspond to exceptional foliations $`\left\{x=\mathrm{const}\right\}`$, $`\left\{y=\mathrm{const}\right\}`$, $`\left\{z=\mathrm{const}\right\}`$ of the web. Then $`B_k`$, $`k=1,2,3`$, are naturally identified with $`𝔹_\rho ^1`$. A choice of $`X_kU_k=B_k`$, $`k=1,2,3`$, corresponds to a choice of 3 leaves of these 3 foliations, or, in other words, to a choice of coordinates $`x`$, $`y`$, $`z`$ such that $`|x|,|y|,|z|<\rho `$. Put $`m=(x,y,z)\left(𝔹_\rho ^1\right)^3`$, then $`\mathrm{\Sigma }_m`$ passes through $`X_k`$, $`k=1,2,3`$. ∎
## 13. Airy webs
Consider a smooth web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$. Suppose that the twistor transform $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$ of $`_{}`$ is well-defined as a manifold.
###### Definition 13.1.
Say that a smooth web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ is strictly airy if for any smooth section $`\mathrm{\Sigma }`$ of $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$ there is a point $`mM`$ such that $`\mathrm{\Sigma }=\mathrm{\Sigma }_m`$. A web is airy if any point has a neighborhood $`U`$ such that $`_{}|_U`$ is strictly airy.
###### Remark 13.2.
This definition requires some modifications if only the germ of $`𝔗=𝔗_{_{}}`$ near $`\mathrm{\Sigma }_{m_0}𝔗`$ is well-defined; here $`m_0M`$. In such a case consider a family $`\sigma _{}:\mathrm{\Lambda }\times T𝔗`$ of sections of $`\pi `$ parameterized by (a germ of) a manifold $`T`$, and such that $`\sigma _{t_0}=\mathrm{\Sigma }_{m_0}`$ for the base point $`t_0T`$. We would require that there is a family $`p_{}:TM`$ of points of $`M`$ such that $`\sigma _t=\mathrm{\Sigma }_{p_t}`$ for $`tT`$ near $`t_0`$.
###### Remark 13.3.
It should be clear that airy webs exist only in complex-analytic situation, otherwise the set of sections is not a finite-dimensional manifold. Moreover, it is reasonable to conjecture that $`\mathrm{\Lambda }`$ cannot be a Stein manifold if $`dim\mathrm{\Lambda }>0`$.
The principal property of airy webs is the following immediate corollary of Corollary 11.5:
###### Theorem 13.4.
Consider a weakly separating and separating strictly airy web. Locally such a web is uniquely determined (up to a local diffeomorphism) by its twistor transform $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$.
###### Proposition 13.5.
In the conditions of Theorem 12.2 suppose that global sections of the vector bundle $`𝒩S`$ span any fiber of $`𝒩S`$. Then a neighborhood $`M_1`$ of $`m_0`$ in $`M`$ is equipped with a web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ of codimension equal to $`\mathrm{codim}\mathrm{\Sigma }`$. This web is separating and airy.
###### Proof.
Deduce the first statement from Lemma 11.2. It is enough to calculate $`\mathrm{rk}d\mathrm{s}_\lambda |_m`$. The condition on global sections is equivalent to $`\mathrm{rk}d\mathrm{s}_\lambda |_{m_0}=\mathrm{codim}\mathrm{\Sigma }`$, thus all we need to show is that this rank does not change if we move to a nearby section $`\mathrm{\Sigma }_m`$ of $`\pi `$.
Consider $`𝒩S`$ as a sheaf of $`𝒪_S`$-modules. For $`PS`$ denote by $`𝒩S\left(P\right)`$ the sheaf of $`𝒪_S`$-modules with local sections being sections of $`𝒩S`$ which vanish at $`P`$. By the Grauert semicontinuity theorem , the Euler characteristic $`\left(1\right)^kdimH^k(S,𝒩S\left(P\right))`$ of $`𝒩S\left(P\right)`$ does not change when $`P`$ changes, and the individual terms $`dimH^k(S,𝒩S\left(P\right))`$ are semicontinuous from above. Similar results hold for $`dimH^k(\mathrm{\Sigma }_m,𝒩\mathrm{\Sigma }_m\left(P\right))`$ considered as functions of $`mM`$ and $`P\mathrm{\Sigma }_m`$.
Consider the exact sequence of sheaves $`0𝒩S\left(P\right)𝒩S\stackrel{v_P}{}\overline{𝒩_PS}`$ 0; here $`\overline{𝒩_PS}`$ is the skyscraper sheaf with the fiber over $`P`$ being $`𝒩_PS`$. Since the mapping $`v`$ of taking the value at $`P`$ is surjective on global sections, the cohomological long exact sequence shows that $`H^k(S,𝒩S\left(P\right))=0`$ for $`k1`$ and $`PS`$. This implies $`H^k(\mathrm{\Sigma }_m,𝒩\mathrm{\Sigma }_m\left(P\right))=0`$ for $`k>1`$ if $`mm_0`$, thus $`dimH^0(\mathrm{\Sigma }_m,𝒩\mathrm{\Sigma }_m\left(P\right))`$ does not depend on $`mm_0`$ and $`P\mathrm{\Sigma }_m`$. Now a consideration of the long exact sequence for $`0𝒩\mathrm{\Sigma }_m\left(P\right)𝒩\mathrm{\Sigma }_m\stackrel{v_P}{}\overline{𝒩_P\mathrm{\Sigma }_m}0`$ shows that $`v_P`$ is surjective for $`mm_0`$ and $`P\mathrm{\Sigma }_m`$. This implies that $`d\mathrm{s}_\lambda |_m`$ is a surjection.
By Lemma 11.2, a neighborhood of $`\mathrm{\Sigma }_{m_0}𝔗`$ is the twistor transform of a web on an open subset $`M_1M`$, $`M_1m_0`$. Since $`d\sigma |_{m_0}`$ (and, by similar arguments, $`d\sigma |_m`$ for any $`mM_1`$) is an injection, this web is separating.
Prove airiness. One can find a neighborhood $`U`$ of $`S`$ in $`𝔗`$ such that $`\mathrm{\Sigma }_mU`$ implies $`mM_1`$. Indeed, let $`U_1=_{mM_1}\mathrm{\Sigma }_m`$. It is a neighborhood of $`S`$, thus one can apply Theorem 12.2 to $`U_1`$ instead of $`𝔗`$. Obviously, the resulting neighborhood $`UU_1`$ of $`S`$ satisfies the requirement above. Let $`M_2=\{mM_1\mathrm{\Sigma }_mU\}`$. Now any section $`\mathrm{\Sigma }`$ of $`U\mathrm{\Lambda }`$ has a form $`\mathrm{\Sigma }=\mathrm{\Sigma }_m`$ for $`mM_1`$. Obviously, this implies also $`mM_2`$. On the other hand, a section of a twistor transform of $`M_2`$ induces a section of $`U\mathrm{\Lambda }`$, thus the restriction of the web on $`M_2`$ is airy. ∎
By Definition 2.1, given a smooth web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$, each point $`mM`$ induces a vector bundle $`\mathrm{n}_m`$ over $`\mathrm{\Lambda }`$, the fiber over $`\lambda `$ being $`\mathrm{n}_m\left(\lambda \right)`$. Obviously,
###### Lemma 13.6.
Consider the section $`\mathrm{\Sigma }_m`$ of the twistor transform $`𝔗\stackrel{𝜋}{}\mathrm{\Lambda }`$ of $`_{}`$. Then $`𝒩\mathrm{\Sigma }_m\pi ^{}\mathrm{n}_m^{}`$.
###### Proposition 13.7.
Consider a complex-analytic separating web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on $`M`$ with compact $`\mathrm{\Lambda }`$. Suppose that $`\mathrm{n}_m`$ is cohomologically trivial for any $`mM`$. Then there is a manifold $`M^{}M`$ with an airy separating web $`\left\{_\lambda ^{}\right\}_{\lambda \mathrm{\Lambda }}`$ on it such that for any $`\lambda \mathrm{\Lambda }`$ and any leaf $`L`$ of $`_\lambda `$ there is a leaf $`L^{}`$ of $`_\lambda ^{}`$ such that $`L=ML^{}`$. The germ of $`M^{}`$ near $`M`$ is canonically defined.
###### Proof.
Due to canonicity of $`M^{}`$ it is enough to prove this statement locally on $`M`$. Thus we may assume that $`_{}`$ is weakly separating and separating. Consider the twistor transform of $`_{}`$. Since $`\mathrm{n}_m`$ may be identified with $`𝒩\mathrm{\Sigma }_m`$, Proposition 13.5 is applicable. (The condition on global sections is automatically satisfied if $`𝔗`$ is a twistor transform.) This provides a construction of $`M^{}`$ and $`_{}^{}`$. ∎
###### Remark 13.8.
This explains the choice of the term airy: it reasonable to imagine that $`M^{}`$ is obtained from $`M`$ by “blowing out” $`M`$. Here leaves of foliations $`_\lambda `$ work as “walls of microscopic air cells” in $`M`$. If $`M^{}`$ is a Veronese web, and we consider just enough foliations $`_{\lambda _k}`$, $`k=1,\mathrm{},K`$, to uniquely determine $`_{}`$ (so $`K=dimM^{}+1`$) then after blow-out each cell becomes a tiny simplex in $`M^{}`$. Before “expansion” each cell is folded into a polytop of smaller dimension.
Proposition 13.7 immediately implies
###### Theorem 13.9.
Consider a complex-analytic separating web $`\left\{_\lambda \right\}_{\lambda \mathrm{\Lambda }}`$ on a connected manifold $`M`$ such that $`\mathrm{n}_m\left(\lambda \right)`$ is a cohomologically trivial vector bundle over $`\mathrm{\Lambda }`$. Then $`_{}`$ is airy iff $`dimM=dim\mathrm{\Gamma }(\mathrm{\Lambda },\mathrm{n}_m)`$ for one (then any) $`mM`$.
###### Remark 13.10.
Note that if $`\mathrm{\Lambda }=^1`$, then $`\mathrm{n}_m\left(\lambda \right)`$ is automatically cohomologically trivial (since by definition this vector bundle is induced from a Grassmannian).
###### Remark 13.11.
In Section 15 we provide a somewhat inverse construction to Proposition 13.7: given $`M^{}`$, we introduce a class of submanifolds $`MM^{}`$ which are equipped with a web having the same twistor transform.
The arguments above give a more detailed proof of one of the principal results of :
###### Theorem 13.12.
Complex-analytic Veronese webs are airy and are (uniquely up to a local diffeomorphism) locally determined by their twistor transform.
This a direct corollary of
###### Lemma 13.13.
Consider a Veronese web $`\left\{_\lambda \right\}_{\lambda ^1}`$ on a $`d`$-dimensional manifold $`M`$. Then $`\mathrm{n}_m𝒪\left(d1\right)`$ for any $`mM`$.
###### Proof.
It is enough to show that $`𝒪\left(d+1\right)`$; here the line bundle $``$ is induced by the Veronese inclusion $`j:^1^{d1}`$ from the tautological line bundle on $`^{d1}`$ (which is isomorphic to $`𝒪\left(1\right)`$). The fiber of $``$ over $`\lambda ^1`$ is the $`1`$-dimensional subspace of $`𝕍^d`$ corresponding to $`j\left(\lambda \right)`$.
A linear function $`l`$ on $`𝕍^d`$ induces a section of $`^{}`$, zeros of this section correspond to points on $`\mathrm{Ker}l\mathrm{Im}j`$. Thus it is enough to construct a hyperplane in $`^{d1}`$ which transversally intersects $`\mathrm{Im}j`$ in $`d1`$ points. Since all the Veronese inclusions are projectively isomorphic, it is enough to consider one given by $`(x:y)(x^{d1}:x^{d2}y:\mathrm{}:xy^{d2}:y^{d1})`$. Let $`\mathrm{\Pi }_{k=1}^{d1}\left(tk\right)=t^{d1}+_{k=0}^{d2}a_kt^k`$. Then the functional $`l`$ with coordinates $`(1,a_{d2},\mathrm{},a_0)`$ satisfies the condition above. ∎
In fact Veronese webs coincide (in complex-analytic situation) with separating airy smooth webs of codimension 1 with $`\mathrm{\Lambda }=^1`$. One can also classify arbitrary separating airy smooth webs with $`\mathrm{\Lambda }=^1`$, the result coincides with Kronecker webs as defined in .
## 14. Non-linear Riemann problem
Consider $`0<\epsilon <1`$, $`\delta >0`$, and a complex-analytic function $`g(\lambda ,t)`$ defined for $`\epsilon <|\lambda |<1/\epsilon `$ and $`|t|<\delta `$. Assume that for any given $`\lambda `$, $`\epsilon <|\lambda |<1/\epsilon `$, the function $`g(\lambda ,t)`$ is invertible. Glue domains $`\left(^1\overline{𝔹}_\epsilon ^1\right)\times `$ and $`𝔹_{1/\epsilon }^1\times 𝔹_\delta ^1`$ together by gluing $`(\lambda ,t_+)𝔹_{1/\epsilon }^1\times 𝔹_\delta ^1`$ to $`(\lambda ,t_{})=(\lambda ,g(\lambda ,t_+))\left(^1\overline{𝔹}_\epsilon ^1\right)\times `$ for $`\epsilon <|\lambda |<1/\epsilon `$ and $`|t|<\delta `$. The result is a $`2`$-dimensional complex manifold $`𝔗`$ with a surjective submersive mapping $`\pi :𝔗^1`$ given by $`(\lambda ,t)\lambda `$.
###### Lemma 14.1.
1. If $`g(\lambda ,0)0`$, then $`\pi `$ has a section given by $`\lambda (\lambda ,0)`$;
2. Sections $`\mathrm{\Sigma }`$ of $`\pi `$ can be identified with pairs of functions $`\sigma _+:𝔹_{1/\epsilon }^1𝔹_\delta ^1`$ and $`\sigma _{}:\left(^1\overline{𝔹}_\epsilon ^1\right)`$ such that $`\sigma _{}\left(\lambda \right)=g(\lambda ,\sigma _+\left(\lambda \right))`$ if $`\epsilon <|\lambda |<1/\epsilon `$;
3. Given a section $`\mathrm{\Sigma }`$ of $`\pi `$ associated to a pair $`\sigma _\pm `$, the degree of the normal bundle of $`\mathrm{\Sigma }`$ in $`𝔗`$ is given by $`\mathrm{ind}\frac{g}{t}(\lambda ,\sigma _+\left(\lambda \right))`$. Here $`\mathrm{ind}\phi \left(\lambda \right)=\frac{1}{2\pi i}_{|\lambda |=1}\frac{d\phi \left(\lambda \right)}{\phi \left(\lambda \right)}`$.
###### Proof.
The first two statements are obvious. On the other hand, the normal bundle of $`\mathrm{\Sigma }`$ is canonically trivialized on $`|\lambda |<1/\epsilon `$ and on $`|\lambda |>\epsilon `$, with the gluing function being $`\frac{g}{t}(\lambda ,\sigma _+\left(\lambda \right))`$. To calculate the degree of a line bundle $``$ it is enough to construct a section $`\tau _+`$ in $`|\lambda |1`$ and a section $`\tau _{}`$ in $`|\lambda |1`$. Suppose that $`\tau _\pm \left(\lambda \right)`$ have no zeros on $`|\lambda |=1`$. Then $`\tau _0=\tau _{}\left(\lambda \right)/\tau _+\left(\lambda \right)`$ is a well-defined function on the unit circle with values in $`^\times `$, and $`\mathrm{deg}=n_++n_{}\mathrm{ind}\tau _0`$; here $`n_\pm `$ are numbers of zeros of $`\tau _\pm `$ (inside and outside of the unit circle correspondingly). In our case $`n_+=n_{}=0`$, and $`\tau _0=\frac{g}{t}(\lambda ,\sigma _+\left(\lambda \right))`$. ∎
###### Theorem 14.2.
Consider a manifold $`K`$ and a function $`g(\lambda ,t,\kappa )`$, $`\kappa K`$, which depends analytically on parameters and such that for any given $`\kappa `$ the function satisfies the condition in the beginning of this section. Suppose that for $`\kappa _0K`$ one has $`g(\lambda ,0,\kappa _0)0`$, and suppose that $`\mathrm{ind}\frac{g}{t}(\lambda ,0,\kappa _0)=1`$. Then there exists $`0<\delta _1<\delta `$ and a neighborhood $`K_1\kappa _0`$, $`K_1K`$, such that for any $`\kappa K_1`$ the conditions
$$\sigma _{,\kappa }\left(\lambda \right)=g(\lambda ,\sigma _{+,\kappa }\left(\lambda \right),\kappa )\text{ for }\epsilon <|\lambda |<1/\epsilon ,|\sigma _{+,\kappa }\left(\lambda \right)|<\delta _1\text{ for }|\lambda |<1/\epsilon ,$$
uniquely determine analytic functions $`\sigma _{+,\kappa }\left(\lambda \right)`$ defined for $`|\lambda |<1/\epsilon `$, and $`\sigma _{,\kappa }\left(\lambda \right)`$ defined for $`|\lambda |>\epsilon `$. Functions $`\sigma _{\pm ,\kappa }\left(\lambda \right)`$ depend analytically on $`\kappa `$.
###### Proof.
Glue domains $`𝔹_{1/\epsilon }^1\times 𝔹_\delta \times K`$ and $`\left(^1\overline{𝔹}_\epsilon ^1\right)\times \times K`$ together by gluing $`(\lambda ,t,\kappa )`$ to $`(\lambda ,g(\lambda ,t,\kappa ),\kappa )`$ for $`\epsilon <|\lambda |<1/\epsilon `$, $`|t|<\delta `$, and $`\kappa K`$. Denote the resulting manifold by $`𝔗`$, denote by $`\pi :𝔗^1`$ the mapping $`(\lambda ,t,\kappa )\lambda `$, by $`\mathrm{\Pi }`$ the natural projection $`𝔗K`$, and by $`\mathrm{\Sigma }`$ the section of $`\pi `$ given by $`\lambda (\lambda ,0,\kappa _0)`$. Consider the normal bundle $`𝒩\mathrm{\Sigma }`$ of $`\mathrm{\Sigma }`$ inside $`𝔗`$. Let $`𝒩^{\left(0\right)}\mathrm{\Sigma }`$ be the normal bundle of $`\mathrm{\Sigma }`$ inside $`\mathrm{\Pi }^1\left(\kappa _0\right)`$. We know that $`\mathrm{deg}𝒩^{\left(0\right)}\mathrm{\Sigma }=1`$. On the other hand, $`𝒩\mathrm{\Sigma }/𝒩^{\left(0\right)}\mathrm{\Sigma }`$ is isomorphic to $`\mathrm{\Pi }^{}𝒯_{\kappa _0}K`$, thus is a trivial vector bundle over $`\mathrm{\Sigma }`$. Thus both $`𝒩^{\left(0\right)}\mathrm{\Sigma }`$ and $`𝒩\mathrm{\Sigma }/𝒩^{\left(0\right)}\mathrm{\Sigma }`$ are cohomologically trivial.
The exact sequence $`\mathrm{}H^k(\mathrm{\Sigma },𝒩^{\left(0\right)}\mathrm{\Sigma })H^k(\mathrm{\Sigma },𝒩\mathrm{\Sigma })H^k(\mathrm{\Sigma },𝒩\mathrm{\Sigma }/𝒩^{\left(0\right)}\mathrm{\Sigma })\mathrm{}`$ shows that $`𝒩\mathrm{\Sigma }`$ is also cohomologically trivial, and $`\mathrm{\Gamma }(\mathrm{\Sigma },𝒩\mathrm{\Sigma })𝒯_{\kappa _0}K`$. In other words, the Kodaira–Spencer theory (Theorem 12.2) is applicable, and there is an mapping $`\mathrm{\Sigma }_{}:^1\times M𝔗`$ and $`m_0M`$, such that $`\mathrm{Im}\mathrm{\Sigma }_{m_0}=\mathrm{\Sigma }`$, and the associated infinitesimal family $`\delta \mathrm{\Sigma }_m:𝒯_{m_0}M\mathrm{\Gamma }(\mathrm{\Sigma },𝒩\mathrm{\Sigma })`$ is a bijection.
On the other hand, $`\mathrm{\Pi }\mathrm{\Sigma }_m:^1K`$ is a deformation of a constant mapping to a point $`\kappa _0K`$, thus is a constant mapping itself. Denote the image-point of this constant mapping by $`\kappa \left(m\right)`$. It is clear that the derivative of $`m\kappa \left(m\right)`$ coincides with the composition $`𝒯_{m_0}M\mathrm{\Gamma }(\mathrm{\Sigma },𝒩\mathrm{\Sigma })H^k(\mathrm{\Sigma },𝒩\mathrm{\Sigma }/𝒩^{\left(0\right)}\mathrm{\Sigma })𝒯_{\kappa _0}K`$, thus $`\kappa \left(m\right)`$ is a local diffeomorphism. Thus we can identify $`M`$ with an open subset of $`K`$. We obtain a family of mappings $`\mathrm{\Sigma }_\kappa :^1𝔗`$, $`\kappa MK`$, such that $`\mathrm{\Pi }\mathrm{\Sigma }_\kappa `$ is the constant mapping to $`\kappa K`$. In other words, $`\mathrm{\Sigma }_\kappa `$ is a section of $`\pi |_{\mathrm{\Pi }^1\kappa }`$, thus induces a pair of functions $`\sigma _{\pm ,\kappa }\left(\lambda \right)`$.
This shows existence of solutions $`\sigma _{\pm ,\kappa }`$, as well as the analytic dependence on parameters. Uniqueness follows from the other parts of Kodaira–Spencer theory (Theorem 12.2). ∎
Using Definition 0.5, one can restate Theorem 14.2 in the following way:
###### Corollary 14.3.
Consider a function $`g(\lambda ,t)`$ defined for $`\epsilon <|\lambda |<1/\epsilon `$ and $`|t|<\delta `$, such that $`g(\lambda ,0)0`$ and $`\mathrm{ind}\frac{g}{t}(\lambda ,0)=1`$. (Obviously, $`\left(g\right)=0`$.) Consider an analytic family $`g_\kappa (\lambda ,t)`$, $`\epsilon <|\lambda |<1/\epsilon `$, $`|t|<\delta `$, $`\kappa U^n`$, such that $`g_0=g`$. Then there is a neighborhood $`U_1`$ of 0 in $`U`$ such that $`\left(g_\kappa \right)`$ is defined for $`\kappa U_1`$ and $`\left(g_\kappa \right)`$ depends smoothly on $`\kappa U_1`$.
###### Remark 14.4.
Since $`\left(g\right)`$ does not change when $`\epsilon `$ increases, it is clear that $`\sigma _+\left(\mu \right)`$ for $`|\mu |<1/\epsilon `$ can be written in terms of $``$. For example, if $`|\mu |<1`$, then $`\sigma _+\left(\mu \right)=\left(g(\frac{\lambda \mu }{\mu \lambda 1},t)\right)`$.
Similarly, one can calculate $`\sigma _{}\left(\mu \right)`$ by considering the inverse function for $`g(\lambda ,t)`$ in $`t`$ (with $`\lambda `$ being a parameter) instead of $`g(\lambda ,t)`$.
Consider now what changes if one takes the gluing functions $`g(\lambda ,t)`$ with $`g(\lambda ,0)0`$ and non-positive values of $`\mathrm{ind}\frac{g}{t}(\lambda ,0)`$ (as opposed to $`\mathrm{ind}=1`$). In such a case $`\mathrm{deg}𝒩\mathrm{\Sigma }=d`$ is non-negative, thus there is a $`\left(d+1\right)`$-parametric family of sections of $`\pi `$. By Proposition 12.8 we expect that a section is determined by its values at $`d+1`$ different points of $`^1`$.
Let us write the formula for the section in terms of $`g`$ and $``$. Use notations of Definition 0.6.
###### Proposition 14.5.
Suppose that $`k`$ numbers $`\lambda _1,\mathrm{},\lambda _k`$ satisfy $`0<|\lambda _l|<1`$, $`m`$ numbers $`\mu _1,\mathrm{},\mu _m`$satisfy $`|\lambda _l|>1`$. Consider a pair of functions satisfying
(14.1)
$$\sigma _{}\left(\lambda \right)=g(\lambda ,\sigma _+\left(\lambda \right)),|\sigma _+\left(\lambda \right)|<\delta \text{ for }|\lambda |<1/\epsilon ,$$
and conditions $`\sigma _+\left(\lambda _l\right)=a_l`$, $`l=1,\mathrm{},k`$, $`\sigma _{}\left(\mu _l\right)=b_l`$, $`l=1,\mathrm{},m`$. Suppose that $`\mathrm{ind}\frac{g}{t}(\lambda ,0)=1km`$. Then $`\sigma _+\left(0\right)=\left(𝒢_{\mathrm{\Lambda }\text{M},\left\{a_i\right\}\left\{b_i\right\}}\right)`$.
###### Proof.
Indeed, one can write
$$\sigma _+\left(\lambda \right)=\stackrel{~}{\sigma }_+\left(\lambda \right)F_+\left(\lambda \right)+\underset{l=1}{\overset{k}{}}a_lF_{+,l}\left(\lambda \right),\sigma _{}\left(\lambda \right)=\stackrel{~}{\sigma }_{}\left(\lambda \right)F_{}\left(\lambda \right)+\underset{l=1}{\overset{m}{}}b_lF_{,l}\left(\lambda \right).$$
Then $`\sigma _{}\left(\lambda \right)=g(\lambda ,\sigma _+\left(\lambda \right))`$ can be rewritten as $`\stackrel{~}{\sigma }_{}\left(\lambda \right)=𝒢_{\mathrm{\Lambda }\text{M},\left\{a_i\right\}\left\{b_i\right\}}(\lambda ,\stackrel{~}{\sigma }_+\left(\lambda \right))`$, and $`\sigma _+\left(0\right)=\stackrel{~}{\sigma }_+\left(0\right)`$. The only thing one needs to prove is that $`\mathrm{ind}\frac{𝒢_{\mathrm{\Lambda }\text{M},\left\{a_i\right\}\left\{b_i\right\}}}{t}=1`$, which follows from $`\mathrm{ind}F_+\left(\lambda \right)=k`$, $`\mathrm{ind}F_{}\left(\lambda \right)=m`$. ∎
###### Proof of Theorem 0.7 .
Correctness follows from Proposition 14.5. Show that $`w`$ is non-degenerate. Suppose that $`w/x=0`$ for some value of $`(x,y,z)`$. Recall that $`w\left(m\right)`$, $`m=(x,y,z)`$, is the value of $`\sigma (0,m)`$; here $`\sigma (\lambda ,m)`$ is a Kodaira–Spencer family of sections of $`𝔗`$, and $`x`$, $`y`$, $`z`$ are $`\sigma (\lambda _{1,2,3},m)`$. If $`w/x=0`$, this would mean that there is a one-parametric family of sections such that the infinitesimal family is non-vanishing, but infinitesimal family vanishes for $`\lambda \{0,\lambda _2,\lambda _3\}`$. However, by Lemma 13.13, the infinitesimal family is a section of $`𝒪\left(2\right)`$, thus cannot vanish at 3 distinct points.
Show that $`w`$ is $`(\lambda _1,\lambda _2,\lambda _3,0)`$-admissible. Glue domains $`𝔹_{1/\epsilon }^1\times 𝔹_\delta `$ and $`\left(^1\overline{𝔹}_\epsilon ^1\right)\times `$ together by gluing $`(\lambda ,t_+)𝔹_{1/\epsilon }^1\times 𝔹_\delta `$ to $`(\lambda ,t_{})=(\lambda ,g(\lambda ,t_+))\left(^1\overline{𝔹}_\epsilon ^1\right)\times `$ for $`\epsilon <|\lambda |<1/\epsilon `$, $`|t_+|<\delta `$. Call the resulting $`2`$-dimensional manifold $`𝔗`$. It is equipped with a projection $`\pi `$ to $`^1`$ and a section $`S=\left\{(\lambda ,0)\right\}`$ of this projection. As in the proof of Lemma 14.1, one can show that degree of $`𝒩S`$ is 2. Since $`\mathrm{deg}𝒩S0`$, $`𝒩S`$ is cohomologically trivial, and fibers of $`𝒩S`$ are generated by global sections. By Proposition 13.5, a neighborhood $`U`$ of $`S`$ in $`𝔗`$ is a twistor transform of a web of codimension 1 on a manifold $`M`$. Since $`dim\mathrm{\Gamma }(S,𝒩S)=3`$, $`dimM=3`$. Again, $`\mathrm{deg}𝒩S=2`$ implies that $`\mathrm{n}_m\left(\lambda \right)`$ of this web spans a quadratic cone in $`𝒯_m^{}M`$, thus this web is a Veronese web.
Taking a point $`uU`$, $`\pi \left(u\right)=\lambda ^1`$, gives a leaf of the foliation $`_\lambda `$ on $`M`$. Since fibers of $`U`$ over $`\lambda _1`$, $`\lambda _2`$, $`\lambda _3`$ and $`\lambda _4=0`$ are identified with subsets of $``$ by the construction of $`𝔗`$, this gives 4 functions $`x`$, $`y`$, $`z`$, $`W`$ on $`M`$, each constant on leaves of $`_{\lambda _{1,2,3,4}}`$. We may assume that $`W=W(x,y,z)`$ for an appropriate function $`W`$ defined in a neighborhood of $`0^3`$. By definition, the latter function is $`(\lambda _1,\lambda _2,\lambda _3,0)`$-admissible.
On the other hand, a point $`mM`$ induces a section of $`\pi `$. A section of $`\pi `$ which is close to $`S`$ is determined by two functions $`\sigma _+`$ and $`\sigma _{}`$ which satisfy (14.1). By Proposition 14.5 this section is determined by $`x=\sigma _+\left(\lambda _1\right)`$, $`y=\sigma _+\left(\lambda _2\right)`$, and $`z=\sigma _{}\left(\lambda _3\right)`$, moreover, $`\sigma _+\left(0\right)=w(x,y,z)`$. We conclude that $`w(x,y,z)=W(x,y,z)`$.
This implies the first statement of the theorem. The second statement is a direct corollary of the first one.
Given $`(A,B,C)`$, find $`\lambda _{1,2,3,4}^1`$ as in Remark 3.9. By a projective transform of $`^1`$ one can make $`\lambda _4=0`$, and $`\lambda _3=\mathrm{}`$. By transformations $`\lambda c\lambda `$ one can make $`\lambda _{1,2}`$ arbitrarily small. After this a transformation $`\lambda \frac{\lambda }{1+\lambda /N}`$ with $`N0`$ would produce a triple $`\lambda _{1,2,3}`$ with desired properties.
Given a non-degenerate solution $`\widehat{w}(x,y,z)`$ of the $`(A,B,C)`$-equation and $`\lambda _{1,2,3}`$ as found above, consider the $`3`$-dimensional Veronese web defined by Theorem 3.8. Let $`\stackrel{~}{𝔗}`$ is the twistor transform of this web. Then $`\stackrel{~}{\pi }:\stackrel{~}{𝔗}^1`$ is (locally near $`\mathrm{\Sigma }_{(0,0,0)}`$) isomorphic to $`\pi :𝔗^1`$; here $`𝔗`$ is glued using the function $`g(\lambda ,t)`$ defined in Theorem 10.1. It is clear that $`\mathrm{ind}g=\mathrm{deg}𝒩\mathrm{\Sigma }_{(0,0,0)}=2`$.
Functions $`x`$, $`y`$, $`z`$, $`\widehat{w}(x,y,z)`$ on the manifold of this Veronese web define local coordinates on the fibers $`\stackrel{~}{\pi }^1\left(\lambda \right)`$, $`\lambda \{\lambda _1,\lambda _2,\lambda _3,0\}`$, thus on $`\pi ^1\left(\lambda \right)`$. Denote these coordinates by the same symbols $`x`$, $`y`$, $`z`$, $`\widehat{w}`$. Investigate how these coordinates are related to coordinates $`t_+`$, $`t_{}`$ on two pieces of $`𝔗`$ it is glued of.
Recall that Theorem 10.1 defined the coordinate $`t_{}`$ on $`\left(^1\overline{𝔹}_\epsilon ^1\right)\times `$ by taking $`z`$-coordinates of the intersection point of leaves of the foliations with $`\gamma _2`$, which is the $`z`$-axis. The leaves of the foliation $`_{\lambda _3}`$ are $`z=\mathrm{const}`$, thus the coordinates $`z`$ and $`t_{}`$ on $`\pi ^1\left\{\lambda _3\right\}`$ coincide, and there is no translation of the argument $`z`$ of the function $`\widehat{w}`$. The coordinate $`t_+`$ over $`𝔹_{1/\epsilon }^1`$ is induced by taking $`x`$-coordinate of the intersection point of leaves with $`\gamma _1`$. Thus the coordinates $`x`$ and $`t_+`$ on $`\pi ^1\left(\lambda _1\right)`$ coincide, and there is no translation of the argument $`x`$. Similarly, the coordinates $`y`$ and $`t_+`$ on $`\pi ^1\left(\lambda _2\right)`$ differ by the transformation $`y=Y\left(t_+\right)`$; here $`y=Y\left(x\right)`$, $`z=0`$ are the equations of the curve $`\gamma _1`$.
Finally, the coordinate $`\widehat{w}`$ on $`\pi ^1\left(0\right)`$ is given by taking the value of $`\widehat{w}(x,y,z)`$ on the leaf of $`_0`$. The leaf which corresponds to a given value of $`t_+`$ passes through the point $`(t_+,Y\left(t_+\right),0)`$, thus the corresponding value of $`\widehat{w}`$ is $`\widehat{w}(t_+,Y\left(t_+\right),0)`$. ∎
###### Remark 14.6.
Consider the case when $`w(x,y,z)`$ is real for real values of $`x,y,z`$, and $`\lambda _{1,2,3}`$ are real. In such a case it is a meaningful question to reconstruct $`w`$ basing on the Cauchy data on a hypersurface w.r.t. which the linearization is hyperbolic. As we have seen in Remark 10.3, the nonlinear Riemann problem with gluing data provided on a neighborhood of a circle is not enough to treat such a problem. One should be able to treat gluing data on more general regions.
However, the results of suggest that providing the gluing data of some kind on the real axis alone should provide enough information. Note that this gluing data should be more general than one we consider here, since a literal application of our arguments leads to a function $`\frac{g}{t}`$ with zeros and/or poles on the real axis.
## 15. Appendix on transversal sections of webs
Some statements of this section are stated in complex-analytic case only. To restate them in real-analytic case is straightforward. Additionally, there is a $`C^{\mathrm{}}`$-treatment of some of these statements as well, see .
###### Definition 15.1.
Say that a submanifold $`NM`$ is transversal to the web $`_{}`$ on $`M`$ if it is transversal to any leaf of any foliation of the web.
If $`\mathrm{codim}_{}=r`$, then the usual count of dimensions shows that there are transversal varieties with dimensions down to $`r+dim\mathrm{\Lambda }`$. However, one should not expect them to exist is smaller dimensions, for example, Lemma 9.5 shows that there are no curves transversal to a complex-analytic Veronese web.
Obviously, a web $`_{}`$ on $`M`$ cuts out a smooth web $`_{}^{\left[N\right]}`$ on a transversal submanifold $`NM`$. Call this web the transversal section of $`_{}`$ by $`N`$. By definition of transversality, the germs of twistor transforms of $`_{}`$ and of $`_{}^{\left[N\right]}`$ near $`mN`$ coincide. This implies
###### Theorem 15.2.
Consider a separating airy web $`_{}`$ on $`M`$ and a transversal to $`_{}`$ submanifold $`NM`$. Then the transversal section web $`_{}^{\left[N\right]}`$ on $`N`$ determines the germ of $`M`$ near $`N`$ and the web $`_{}`$ on this germ (uniquely up to diffeomorphisms $`M\stackrel{~}{M}`$ which preserve $`N`$).
###### Remark 15.3.
Note that Theorem 10.1 and taken together with Theorem 13.12 imply a particular case of this statement: the Veronese web is locally determined by its restriction on the surface $`N`$ given by $`y=Y\left(x\right)`$ (in terms of Theorem 10.1).
On the other hand, classification of $`_{}^{\left[N\right]}`$ up to diffeomorphism can be much easier than classification of $`_{}`$, since leaves of $`_{}^{\left[N\right]}`$ have a smaller dimension. If $`dimN=r+dim\mathrm{\Lambda }`$, and $`dim\mathrm{\Lambda }=1`$, then leaves of $`_{}^{\left[N\right]}`$ have dimension 1. But to specify a foliation on $`N`$ of codimension $`dimN1`$ is exactly the same as to specify a direction $`\mathrm{d}_n𝒯_nN`$ in a tangent space at every point $`nN`$, there is no integrability condition involved (as, for example, one in Lemma 3.3). A family of foliations induces a family of directions $`\mathrm{d}_n\left(\lambda \right)`$. (Note that $`\mathrm{d}_n\left(\lambda \right)`$ is a direction in a tangent space, not in a cotangent space, as is $`\mathrm{n}_n\left(\lambda \right)`$ in the case of webs of codimension 1.)
If $`\mathrm{\Lambda }`$ is a compact complex curve, then a mapping $`\mathrm{d}:\mathrm{\Lambda }𝒯_nN`$ of given degree in general position is uniquely determined by images of $`P`$ points on $`\mathrm{\Lambda }`$ for an appropriate $`P>0`$. The standard arguments of algebraic geometry of curves show that for $`\mathrm{d}_n`$ one should expect this for $`Pr+2+\frac{g+\left(r+1\right)\left(dr1\right)}{r}`$; here $`g`$ is the genus of $`\mathrm{\Lambda }`$, $`d=dimM`$. Additionally, if this inequality is an equality, then the images of these $`P`$ points in $`𝒯_nN`$ should be expected to be arbitrary. Moreover, in the case $`g=0`$, $`r=1`$ these expectations can be easily checked to be true, as far as among these $`P`$ images no more than $`d1`$ glue into any point of $`^1`$. Additionally, the condition of general position can be removed, if one allows the degree of the mapping $`\mathrm{d}:^1^1`$ to drop. This leads to
###### Theorem 15.4.
Consider a complex-analytic Veronese web $`_{}`$ on $`M`$ and a transversal surface $`NM`$, $`dimN=2`$, $`dimM=d`$. Consider $`2d1`$ distinct points $`\lambda _1,\mathrm{},\lambda _{2d1}^1`$. Then the germ of $`_{}`$ near $`N`$ is determined (uniquely up to a diffeomorphism preserving $`N`$) by $`2d1`$ foliations $`^{\left(k\right)}=_{\lambda _k}^{\left[N\right]}`$, $`k=1,\mathrm{},2d1`$, of codimension 1 on $`N`$. The foliations $`^{\left(k\right)}`$ on $`N`$ can be taken arbitrarily with the restriction that at any point of $`N`$ no more than $`d1`$ foliation have any given tangent direction, and there is no mapping $`\mathrm{f}:^1^1`$ of degree less than $`d1`$ such that $`\mathrm{f}\left(\lambda _k\right)=𝒯_n^{\left(k\right)}`$, $`k=1,\mathrm{},2d1`$.
By a choice of coordinates on $`N`$ one can take last two of these foliations to be $`\left\{x=\mathrm{const}\right\}`$, $`\left\{y=\mathrm{const}\right\}`$, the rest to be $`\left\{w_k(x,y)=\mathrm{const}\right\}`$. Thus the collection $`(M,N,_{})`$ (up to the same transformations as in the theorem) is determined by $`2d3`$ functions on $`N`$.
The next step is to use the freedom in the choice of $`N`$ to reduce the number of parameters. Accidentally, a proper choice of $`N`$ also allows to ensure that no mapping like $`\mathrm{f}`$ exists.
Recall Lemma 9.5: given a Veronese web on $`M`$, a choice of a subset $`T^1`$ with multiplicities and the total count $`dimM1`$ determines a direction at each point of $`M`$. In particular, given $`T`$ and $`m_0M`$, there is a canonically defined curve $`\gamma _{m_0,T}m_0`$ on $`M`$ (one with the prescribed directions). Taking another subset $`T^{}`$, one can put a curve $`\gamma _{m,T^{}}`$ through every point $`m`$ of $`\gamma `$. Taken together, these curves $`\gamma _{m,T^{}}`$, $`m\gamma _{m_0,T}`$ sweep a surface $`N_{m_0,T,T^{}}`$ in $`M`$.
One can check that if $`TT^{}=\mathrm{}`$, then $`N_{m,T,T^{}}`$ is transversal to $`_{}`$ at $`m`$, thus in a neighborhood of $`m`$. A proof of the following statement is straightforward:
###### Lemma 15.5.
Suppose that $`N`$ is transversal to a Veronese web $`_{}`$ and contains a curve $`\gamma _{m,T}`$. Then $`\gamma `$ is a leaf of $`_\lambda ^{\left[N\right]}`$ for any $`\lambda T`$.
In particular, for $`N=N_{m,T,T^{}}`$ the foliations $`_\lambda ^{\left[N\right]}`$, $`\lambda T^{}`$, coincide. Since $`T^{}`$ contains $`d1`$ points, this condition ensures that no mapping $`\mathrm{f}`$ of degree smaller than $`d1`$ can exist. Additionally, the leaves of the foliations $`_\lambda ^{\left[N\right]}`$, $`\lambda T`$, which pass through $`m`$ coincide. This leads to the following
###### Corollary 15.6.
Consider a complex-analytic surface $`N`$, a point $`nN`$, $`2d1`$ distinct points $`\lambda _k`$ on $`^1`$, and $`2d1`$ foliations $`\stackrel{~}{}_k`$, $`k=1,\mathrm{},P`$, on $`N`$. Let $`\gamma _k`$ be the leaf of $`\stackrel{~}{}_k`$ through $`n`$. Suppose that $`\stackrel{~}{}_k=\stackrel{~}{}_k^{}`$ if $`1k,k^{}d1`$, $`\gamma _k=\gamma _k^{}`$ if $`dk,k^{}2d2`$, that for any fixed $`k`$, $`dk2d1`$, the foliations $`\stackrel{~}{}_1`$ and $`\stackrel{~}{}_k`$, $`dk2d1`$, have distinct directions at any point of $`N`$, and that directions of the foliations $`\stackrel{~}{}_k`$, $`dk2d1`$, are not all the same at any point of $`N`$. Then there is a complex-analytic Veronese web $`(M,_{})`$ of dimension $`d`$ and an embedding $`f:NM`$ such that $`\mathrm{Im}f=N_{f\left(n\right),T,T^{}}`$, and that $`f`$ identifies the foliations $`\stackrel{~}{}_k`$ on $`N`$ with foliations $`_{\lambda _k}^{\left[\mathrm{Im}f\right]}`$ on $`\mathrm{Im}f`$; here $`T=\{\lambda _d,\mathrm{},\lambda _{2d2}\}`$, $`T^{}=\{\lambda _1,\mathrm{},\lambda _{d1}\}`$. The germ of $`(M,_{})`$ near $`\mathrm{Im}f`$ is determined uniquely up to isomorphism.
This is a geometric local classification of complex-analytic Veronese webs: given $`Mn`$ and a Veronese web on $`M`$, $`N=N_{n,T,T^{}}`$ is canonically defined, thus foliations $`_{\lambda _k}^{\left[N\right]}`$ are canonically defined. These foliations satisfy the conditions of the corollary, and allow reconstruction of the web on $`M`$. In addition to the restriction that $`M`$ is defined only as a germ near $`N`$, there is another direction of locality in this result: $`N`$ can be embedded into $`M`$, not included into $`M`$.
One can make appropriate modifications to this statements if some of the points $`\lambda _k`$, $`1kd1`$ or $`dk2d2`$ can collide. In any case, count the number of parameters in this representation. It is enough to specify $`\stackrel{~}{}_1`$ and $`\stackrel{~}{}_k`$, $`dk2d1`$, one can suppose that $`\stackrel{~}{}_1`$ is $`\left\{x=\mathrm{const}\right\}`$, $`\stackrel{~}{}_d`$ is $`\left\{y=\mathrm{const}\right\}`$. Then one can write $`\stackrel{~}{}_k`$ as $`\left\{w_k(x,y)=\mathrm{const}\right\}`$, $`d+1k2d1`$. Assume that the point $`n`$ is given by $`x=y=0`$. As in , one can normalize $`w_{2d1}(x,y)=x+y+xyu_{2d1}(x,y)`$. Additionally, one can normalize $`w_k(x,y)`$ by $`w_k(0,y)=y`$. Since $`\frac{dw_k(x,0)}{dx}=0`$, one can write $`w_k(x,y)=y+xyu_k(x,y)`$. The functions $`u_k(x,y)`$ are defined uniquely up to a transformation
$$\stackrel{~}{u}_k(x,y)=C^1u_k(Cx,Cy),d+1k2d1.$$
Thus an analytic $`d`$-dimensional Veronese web on $`M`$ near a point $`mM`$ is locally uniquely determined by $`d1`$ functions $`u_k(x,y)`$ up the transformation above.
Note the similarity of this description with the Turiel classification of Veronese webs . In fact what we did above is just a geometric reformulation of this result. Unfortunately, our approach works in an analytic situation only, and does not imply the $`C^{\mathrm{}}`$-case of the Turiel classification.
###### Remark 15.7.
In the case of arbitrary separating airy webs and $`dimN=r+dim\mathrm{\Lambda }`$, here $`r=\mathrm{codim}_{}`$, it is not feasible to describe transversal sections of webs by specifying several foliations on $`N`$, since these foliations should satisfy too many conditions. However, if $`\mathrm{codim}_{}=1`$ (so there are no integrability conditions on $`_{}^{\left[N\right]}`$) one can make a substitution. The mapping $`\mathrm{n}_n:\mathrm{\Lambda }𝒯_n^{}N`$ induces a line bundle $`_n=\mathrm{n}_n^{}𝒪\left(1\right)`$ over $`\mathrm{\Lambda }`$, and an inclusion $`\iota _n:𝒯_nN\mathrm{\Gamma }(\mathrm{\Lambda },_n^{})`$.
Suppose that $`\mathrm{\Lambda }`$ is a compact curve, then the latter space is finite-dimensional. Since $`\iota _n`$ (up to multiplication by a constant) determines $`\mathrm{n}_n`$, it is enough to provide enough information to describe $`_n`$ and $`\iota _n`$. Since we are free to multiply $`\iota _n`$ by a constant, it is enough to know $`_n`$ up to isomorphism. We can see that to describe $`_{}^{\left[N\right]}`$, it is enough to describe the degree $`\delta `$ of $`_{}`$, provide the mapping $`l_{}:N\mathrm{Pic}^\delta \left(\mathrm{\Lambda }\right)`$ which sends $`n`$ to the class of $`_n`$ inside the Picard variety, the mapping $`\tau :N\mathrm{Gr}_2\left(\mathrm{\Gamma }(\mathrm{\Lambda },l_n)\right):n\mathrm{Im}\iota _n\mathrm{\Gamma }(\mathrm{\Lambda },l_n)`$, and an identification of $`𝒯_nN`$ with the $`2`$-dimensional vector subspace described by $`\tau \left(n\right)`$ up to a constant. In general position, given $`l`$ and $`\tau `$, one needs to know $`_\lambda ^{\left[N\right]}`$ for 3 values of $`\lambda `$ to provide such an identification.
If $`\mathrm{\Lambda }=^1`$, then $`\delta `$ is a number, and $`\mathrm{Pic}^\delta \left(\mathrm{\Lambda }\right)`$ has one point only. It is clear that $`_n𝒪\left(dimM1\right)`$, so it is enough to provide 3 foliations on $`N`$, and a mapping $`N\mathrm{Gr}_2\left(𝕍^{dimM}\right)`$. It is easy to see that this data is equivalent to the data of Theorem 15.4.
Consider now a transversal submanifold $`N`$ to a web $`_{}`$, and a submanifold $`\gamma N`$ of codimension $`\mathrm{codim}_{}`$. Let $`m\gamma `$. If $`𝒯_m\gamma 𝒯_mN`$ is in general position, then $`𝒯_m\gamma `$ is transversal to $`𝒯_m_\lambda `$ for $`\lambda \mathrm{\Lambda }Z`$, here $`Z`$ is a proper analytic subset of $`\mathrm{\Lambda }`$. Consequently, a small neighborhood of $`m`$ in $`\gamma `$ is transversal to $`_\lambda `$ for $`\lambda `$ in an open subset $`U\mathrm{\Lambda }`$. Reducing $`M`$ to a neighborhood of $`m`$, we obtain the corresponding sectional coordinate system on $`𝔗`$. Given two such submanifolds $`\gamma _1`$, $`\gamma _2`$, $`m\gamma _1\gamma _2`$ we obtain two subsets $`U_{1,2}\mathrm{\Lambda }`$, and the corresponding local identifications $`g_\lambda :\gamma _1\gamma _2`$, $`\lambda U_1U_2`$. This identification are obtained in the same way as in Section 9, the principal difference being that the whole construction is performed on $`N`$ instead of $`M`$.
Taking enough $`\gamma _k`$ to cover $`\mathrm{\Lambda }`$, the corresponding pairwise gluing functions determine the germ of $`𝔗`$ near $`\mathrm{\Sigma }_m`$ up to isomorphism, thus the germ of $`_{}`$ near $`m`$ up to isomorphism (assuming $`_{}`$ is airy). Note that to construct $`g_\lambda `$, we need to find a leaf of $`_\lambda ^{\left[N\right]}`$ which passes through a given point of $`\gamma _1`$, and find the intersection of this leaf with $`\gamma _2`$. Obviously, to do this it is enough to solve some ordinary differential equations.
Consequently, the construction of Theorem 10.1 can be generalized to arbitrary webs.
## 16. Appendix on computational complexity of the nonlinear Riemann transform
Continue using notations of Section 14. Consider not the mapping $`:g(\lambda ,t)\sigma _+\left(0\right)`$, but a more general mappings $`\stackrel{~}{}_\pm :g(\lambda ,t)\sigma _\pm \left(\lambda \right)`$. Let us introduce operators solving the linear Riemann problem: given a function $`\phi \left(\lambda \right)`$ defined for $`\epsilon <|\lambda |<1/\epsilon `$, define functions $`_+\phi `$ and $`_{}\phi `$ by the conditions $`\phi \left(\lambda \right)=\lambda _+\phi \left(\lambda \right)+_{}\phi \left(\lambda \right)`$ and the conditions that $`_+\phi \left(\lambda \right)`$ and $`_{}\phi \left(\lambda \right)`$ can be holomorphically extended on $`|\lambda |<1/\epsilon `$ and $`|\lambda |>\epsilon `$ correspondingly. Similarly, if $`\phi \left(\lambda \right)`$ is nowhere 0, and $`\mathrm{ind}\phi =0`$, define $`𝕄_+\phi `$ and $`𝕄_{}\phi `$ by $`\phi \left(\lambda \right)=𝕄_+\phi \left(\lambda \right)/𝕄_{}\phi \left(\lambda \right)`$ and the conditions that $`𝕄_+\phi \left(\lambda \right)`$ and $`𝕄_{}\phi \left(\lambda \right)`$ can be holomorphically extended on $`|\lambda |<1/\epsilon `$ and $`|\lambda |>\epsilon `$ correspondingly, these extensions are nowhere 0, and $`𝕄_+\phi \left(0\right)=1`$.
Uniqueness of $`𝕄_+\phi `$ and $`𝕄_{}\phi `$ is obvious, existence follows from the theory of the linear Riemann problem—or, what is the same, classification of line bundles over $`^1`$. Uniqueness of $`_+\phi `$ and $`_{}\phi `$ is obvious, existence follows from existence of $`\mathrm{log}𝕄_+e^\phi `$ and $`\mathrm{log}𝕄_{}e^\phi `$.
###### Lemma 16.1.
Denote by $`S^1`$ the circle $`|\lambda |=1`$. Then $`_\pm \phi |_{S^1}`$ is uniquely determined by $`\phi |_{S^1}`$. This induces two linear operators on real-analytic complex-valued functions on $`S^1`$. These operators can be extended to continuous linear operators in Sobolev spaces $`H^s\left(S^1\right)`$ for any $`s`$.
###### Proof.
The first statement is obvious, since $`\phi |_{S^1}`$ uniquely determines $`\phi `$. The second statement follows from $`_+\lambda ^k=c_k\lambda ^{k1}`$ and $`_{}\lambda ^k=c_k^{}\lambda ^{k1}`$ with $`c_k`$ and $`c_k^{}`$ being 0 or 1, and from the fact that $`\left(\lambda ^k\right)_k`$ is an orthogonal basis in the Sobolev spaces $`H^s\left(S^1\right)`$. ∎
Denote the continuations of operators $`_\pm `$ into $`H^s\left(S^1\right)`$ by the same symbols. Similarly, if $`s>1/2`$, then the mappings $`𝕄_\pm `$ can be considered as continuous mappings from an open subset of $`H^s\left(S^1\right)`$ into $`H^s\left(S^1\right)`$. Indeed, if $`s>1/2`$, then $`\phi e^\phi `$ is a continuously differentiable mapping $`H^s\left(S^1\right)H^s\left(S^1\right)`$ with an open image.
###### Lemma 16.2.
Consider $`\epsilon ,\delta >0`$ and a family $`g_\kappa (\lambda ,t)`$, $`\kappa K`$, of functions such that $`_{\epsilon \delta }\left(g_\kappa \right)`$ is well-defined for any $`\kappa K`$. Let $`\sigma _{\pm ,\kappa }\left(\lambda \right)=\stackrel{~}{}_\pm \left(g_\kappa \right)`$. Then
$$\frac{}{\kappa }\stackrel{~}{}_\pm \left(g_\kappa \right)=a_{\pm ,\kappa }^1\frac{b_{\pm ,\kappa }}{\kappa },$$
here
$$a_{\pm ,\kappa }\left(\lambda \right)=𝕄_\pm \left(\lambda ^1\frac{g_\kappa }{t}(\lambda ,\sigma _{+,\kappa }\left(\lambda \right))\right),b_{\pm ,\kappa }\left(\lambda \right)=_\pm \left(a_{,\kappa }\left(\lambda \right)g_\kappa (\lambda ,\sigma _{+,\kappa }\left(\lambda \right))\right).$$
###### Proof.
Fix $`\kappa _0K`$. We may assume that $`dimK=1`$, for example, $`K=𝔹_r^1`$. Let $`\sigma _\pm =\stackrel{~}{}_\pm \left(g_{\kappa _0}\right)`$, $`\delta _\pm =\frac{d}{d\kappa }\stackrel{~}{}_\pm \left(g_\kappa \right)|_{\kappa _0}`$. Then
$$\frac{g_{\kappa _0}}{t}(\lambda ,\sigma _+\left(\lambda \right))\delta _+\left(\lambda \right)+\frac{g_\kappa }{\kappa }|_{\kappa _0}(\lambda ,\sigma _+\left(\lambda \right))=\delta _{}\left(\lambda \right).$$
Since $`\mathrm{ind}\frac{g}{t}=1`$, one can write $`\frac{g_{\kappa _0}}{t}(\lambda ,\sigma _+\left(\lambda \right))`$ as $`\lambda a_+\left(\lambda \right)/a_{}\left(\lambda \right)`$, here $`a_+\left(\lambda \right)`$ and $`a_{}\left(\lambda \right)`$ have invertible holomorphic continuations into $`|\lambda |<1/\epsilon `$ and $`|\lambda |\epsilon `$ correspondingly. Similarly, write
$$a_{}\left(\lambda \right)\frac{g_\kappa }{\kappa }|_{\kappa _0}(\lambda ,\sigma _+\left(\lambda \right))=\lambda B_+\left(\lambda \right)+B_{}\left(\lambda \right),$$
here $`B_+\left(\lambda \right)`$ and $`B_{}\left(\lambda \right)`$ have holomorphic continuations into $`|\lambda |<1/\epsilon `$ and $`|\lambda |\epsilon `$ correspondingly. Then $`\sigma _+\left(\lambda \right)=a_+^1B_+`$, $`\sigma _{}\left(\lambda \right)=a_{}^1B_{}`$. Since operators $`_\pm `$ are linear and continuous in an appropriate topology, it is easy to see that $`B_\pm =\frac{db_{\pm ,\kappa }}{d\kappa }|_{\kappa _0}`$. ∎
Consider a vector space $`V=H^s\left(S^1\right)\times H^s\left(S^1\right)`$, $`s1/2`$. Denote the element of $`V`$ by $`(\sigma _+,\sigma _{})`$. In the conditions of Lemma 16.2 suppose that $`dimK=1`$. Define a mapping $`v_\kappa :UV:(\sigma _+,\sigma _{})(\delta _+,\delta _{})`$, here $`U`$ is an appropriate open subset of $`V`$, and $`\delta _\pm \left(\lambda \right)=a_{\pm ,\kappa }\left(\lambda \right)^1B_{\pm ,\kappa }\left(\lambda \right)`$,
$$a_{\pm ,\kappa }\left(\lambda \right)=𝕄_\pm \left(\lambda ^1\frac{g_\kappa }{t}(\lambda ,\sigma _+\left(\lambda \right))\right),B_{\pm ,\kappa }\left(\lambda \right)=_\pm \left(a_{}\left(\lambda \right)\frac{dg_\kappa }{d\kappa }(\lambda ,\sigma _+\left(\lambda \right))\right).$$
Since one can take value of elements of $`H^s\left(S^1\right)`$, $`s>1/2`$, at points, it makes sense to require that $`|\sigma _+\left(\lambda \right)|<\delta `$ for any $`\lambda `$, thus $`v_\kappa (\sigma _+,\sigma _{})`$ is indeed well-defined on an open subset of $`H^s\left(S^1\right)`$. Moreover, $`v_\kappa `$ is Lipschitz on $`K\times 𝔹`$, here $`𝔹`$ is any ball in $`H^s\left(S^1\right)`$.
###### Corollary 16.3.
Given a family $`g_\kappa (\lambda ,t)`$, $`\kappa K`$, $`dimK=1`$, of functions as in Definition 0.5, one can define a Lipschitz family $`v_\kappa `$ of vector fields on an open subset $`U`$ of $`H^s\left(S^1\right)\times H^s\left(S^1\right)`$ such that if $`\stackrel{~}{}_\pm \left(g_\kappa \right)`$ makes sense for any $`\kappa K`$, then the curve $`(\stackrel{~}{}_+\left(g_\kappa \right),\stackrel{~}{}_{}\left(g_\kappa \right))`$ is an integral curve of the ODE $`\frac{d\mathrm{\Phi }\left(\kappa \right)}{d\kappa }=v_\kappa \left(\mathrm{\Phi }\right)`$.
Now Lipschitz ODEs in Banach spaces enjoy most of the properties of finite-dimensional ODEs, and are not harder to solve. We conclude that in the setting of Corollary 14.3 one can calculate $`\left(g_\kappa \right)`$ by solving a Lipschitz ODE in a Hilbert space. In particular, Theorem 0.7 reduces solution of Cauchy problem for the nonlinear wave equation to solution of such an ODE.
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# For submission to Monthly Notices Evolution of the probability distribution function of galaxies in redshift-space
## 1 Introduction
One of the central goals of modern cosmology is to gain a quantitative understanding of the nature of the observed galaxy distribution. In order that such an understanding be complete it must incorporate an explanation of both the evolution of large-scale structure in the universe and of the relationship between galaxies and the underlying dark mass.
In the standard paradigm the origin of structure can be traced back to small perturbations in the mass density imprinted during an inflationary phase. These perturbations subsequently grow in amplitude through gravitational instability. This prescription is well tested in N-body simulations and can be used to make useful predictions about the clustering history of dark matter. Many of the details, however, are not well understood.
The picture becomes more complex when one considers the galaxy distribution. Firstly, the astrophysical processes that govern the formation of the galaxies themselves introduce a nontrivial relationship between the galaxies and the dark matter. This relationship is generally termed galaxy bias. The biasing relation may prove to be simple to describe, depending only upon the local density in dark matter at a given point (eg Coles 1993). Conversely it could be extremely complex, depending on the larger scale environment, tidal fields, feedback and entire history of the galaxy formation site. This could lead to non-local (Bower et al. 1993) or non-deterministic forms of biasing (Pen 1998, Dekel & Lahav 1999). Secondly, the observed galaxy distribution is distorted in 3D redshift catalogues due to the line of sight components of the galaxy peculiar motions adding to the general Hubble expansion (Kaiser 1987). These redshift distortions, if not properly modeled, may lead to misinterpretation when comparing theory with observation.
A complete specification of galaxy clustering may only truly be given by the full set of galaxy N-point correlation functions. Such functions form the solution to the so-called BBGKY system of dynamical equations. This approach, pioneered in the 1970s by Peebles and co workers (Davis & Peebles 1977, Groth & Peebles 1977, Fry & Peebles 1978, Peebles 1980), has met with little success in practice. From a theoretical perspective, a closed solution to the BBGKY hierarchy has never been found and remains a thorny problem. Observationally, measurements of the correlation functions have been restricted to the lowest few orders.
An alternative description, and the focus of this work, may be given by the Probability Distribution Function (PDF) of a random field. Strictly, we investigate the PDF at a single spatial location (the 1-point PDF), but in general the PDFs are a family of N-point distribution functions for the probability at N positions in space. The random fields of interest are the density fields of both dark matter and galaxies.
PDFs are useful in cosmology because, in principle, they encode much of the information contained within the full hierarchy of correlation functions, thus providing valuable information about gravitational evolution and initial conditions. Furthermore the discrete analogy of the 1-point PDF, the counts in cells statistic, is a relatively straightforward quantity to measure from galaxy surveys (Hamilton 1985, Alimi et al. 1990, Szapudi et al. 1992, 1996, Gaztañaga 1992, 1994, Bouchet et al. 1993, Kim & Strauss 1998), allowing for easy comparison between theory and observation.
The shape of the PDF is strongly influenced by the non-linear effects of gravitational instability. A cosmic field that is initially Gaussian random, the generic prediction of most inflationary models, will remain so only when its evolution is linear. When the evolution becomes nonlinear, the higher-order moments of the field (more correctly the cumulants) become non-zero for the first time resulting in a PDF that may be strongly skewed about its mean and cuspy around its peak.
Attempts to describe theoretically this evolution have been numerous (Fry 1985, Coles & Jones 1991, Bernardeau 1992, 1994, 1996, Bernardeau & Kofman 1995, Juszkiewicz et al. 1995, Colombi et al. 1997, Gaztanaga et al. 1999, Taylor & Watts 2000) with the result that, at least for quasi-linear evolution of the matter field, the behavior of the PDF is quite well understood. Since the density field of galaxies inferred from a galaxy redshift survey may not be a good representation of the underlying dark matter it is unlikely that the PDF of the dark matter density field is a realistic approximation for that of galaxies. Just as nonlinear gravitational evolution drives the PDF away from Gaussianity, so too will nonlinear bias and redshift distortions. The challenge is to separate out each of the effects and quantify the evolution of the PDF in terms of a few essential parameters.
There have been significantly fewer attempts to calculate the evolution of the PDF in redshift space than in real space. Most notably, Hui et al. (2000) extended the work of Bernardeau & Kofman (1995) to calculate the evolution using the Zel’dovich approximation. Despite being an excellent approximation for the nonlinear gravitational dynamics Zel’dovich’s solution does not produce a good fit to the PDF of N-body simulations. This is principally due to the formation of caustics in the density field at shell crossing. The result for the Zel’dovich PDF is an asymptotic high density tail that falls off like $`\delta ^3`$ in both real and redshift space. Tails like these are not observed in simulations.
Another way of incorporating the effect of redshift distortions and bias is to use the Edgeworth expansion (Juszkiewicz et al. 1995, Bernardeau & Kofman 1995). In this approximation, the PDF is reconstructed from the moments of the field which must be calculated separately using perturbation theory. Such a calculation was done by Juszkiewicz et al. (1993; but see Hivon et al. 1995 for more details), who used second order Lagrangian perturbation theory (Bouchet et al. 1995) to map the skewness of an unbiased field into redshift space. They did not, however, try to use the skewness to evolve the PDF via the Edgeworth expansion.
A problem with using the Edgeworth expansion is its tendency to produce unphysical features in the PDF. In particular is the appearance of negative probabilities at low densities and the formation of unphysical “wiggles” in the high density tail of the distribution when the variance grows. These problems can be overcome by reconstructing the PDF using a Gamma expansion (Gaztañaga et al. 1999) which shows a markedly better agreement with N-body data in real space.
In this paper we extend the formalism developed in a earlier work (Taylor & Watts 2000, henceforth TW2000) where we showed how the PDF transforms when the matter density field is propagated to second-order. This method is based on an exact propagation of Gaussian initial probabilities using the Chapman–Kolmogorov equation from statistical physics (e.g van Kampen 1992). In this paper we develop this approach to include the transformation to a local, second-order biased galaxy distribution, and apply second-order Eulerian perturbation theory to describe the mapping into redshift-space. We show how the method can be naturally extended to incorporate a stochastic (or hidden variable) bias, although we leave detailed calculations for a later paper (Watts & Taylor in preparation). Our calculation also provides a new analytic solution for the skewness in redshift-space. This result is different to that of Hivon et al. (1995) who estimated the skewness by inverse transform of the bispectrum. However they where unable to find a closed-form solution, and did not include the effects of nonlinear galaxy biasing. Here we derive a closed-form solution for the skewness, including second-order bias, and find that the nonlinear bias has a large effect.
The layout of this paper is as follows. In §2 we discuss the derivation of the 1-point Probability Distribution Function. The new result for the skewness, and a discussion of stochastic biasing may also be found here. In §3 we illustrate the various dependencies on cosmological parameters of the shape of the PDF. Our results are compared with N-body simulations and other approximations in §4, conclusions are given in §5.
## 2 The galaxy distribution function in redshift-space
### 2.1 Redshift distortions and biasing in second-order perturbation theory.
#### 2.1.1 Second-order perturbation theory
In Eulerian perturbation theory the density field $`\delta (𝒙,t)`$ at real space comoving position $`𝒙`$ and time $`t`$ in a flat Universe can be expanded into a series of separable functions;
$$\delta (𝒙,t)=\underset{n=1}{\overset{\mathrm{}}{}}\delta _n(𝒙,t)=\underset{n=1}{\overset{\mathrm{}}{}}D_n(t)\epsilon _n(𝒙),$$
(1)
In curved space this expansion is only separable to second-order. The evolution of $`\delta (𝒙,t)`$ can therefore be traced into the nonlinear regime by solving the fluid and Poisson’s equations for each order in the perturbation expansion (Peebles 1980, Juszkiewicz 1981, Vishniac 1983, Fry 1984, Bouchet et al. 1992, Bouchet et al. 1995);
$`\delta `$ $`=`$ $`\delta _1+\delta _2`$ (2)
$`=`$ $`\delta _1+{\displaystyle \frac{1}{3}}(2\kappa )\delta _1^2𝜼.𝒈+{\displaystyle \frac{1}{2}}(1+\kappa )E^2,`$
where
$$𝜼(𝒙,t)=\mathbf{}\delta _1(𝒙,t),$$
(3)
is the gradient of the linear density field and
$$𝒈(𝒙,t)=\mathbf{}^2\delta _1(𝒙,t),$$
(4)
is the linear peculiar gravity field<sup>1</sup><sup>1</sup>1In this paper we define $`4\pi G\rho _0=3/2\mathrm{\Omega }H^2=1`$ and the expansion parameter $`a(t)=1`$, since our final distribution function will be dimensionless. where $`^2`$ is the inverse Laplacian. The trace-free tidal tensor is given by
$$E_{ij}(𝒙,t)=(_i_j^2\frac{1}{3}\delta _{ij})\delta _1(𝒙,t).$$
(5)
#### 2.1.2 Redshift space distortions
To determine the redshift space density field one must consider the mapping from the real space comoving position to redshift space comoving position, $`𝒔`$, given by (see the excellent review by Hamilton, 1998, for a full account)
$$𝒔=𝒓+\widehat{𝒓}u$$
(6)
where $`u=\widehat{𝒓}𝒗(𝒓,t)/H_0`$ is the projection of the velocity field along the line of sight. In the distant observer approximation (Kaiser 1987) the result can be written as a series in $`u`$ and $`\delta `$,
$$\delta ^s(𝒔,t)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(u)^n}{n!}\frac{^n}{s^n}\left(\frac{\delta (𝒔,t)u^{}}{1+u^{}}\right),$$
(7)
where $`{}_{}{}^{}/r`$ and all quantities are evaluated in redshift-space coordinates.
Expanding this to second order we find
$$\delta ^s=\delta u^{}[u(\delta u^{})]^{}.$$
(8)
#### 2.1.3 Galaxy bias
Next we assume that the density field of galaxies, smoothed on some scale $`R`$, is some local function of the underlying smoothed field of dark mass. Following Fry & Gaztañaga (1993) we expand the galaxy density in powers of $`\delta `$;
$$\delta _g=\underset{n}{}\frac{b_n}{n!}\delta ^n.$$
(9)
The coefficients in this expansion, $`b_n`$, are the bias parameters. We make no assumptions about the biasing function other than that it is local and that it may be expanded in a Taylor series. This is a deterministic Eulerian biasing scheme, but can be generalised to a stochastic Eulerian biasing scheme to allow for the hidden effects of galaxy formation (Pen 1998, Dekel & Lahav 1999). In Section 2.4 we show how our intrinsically probabilistic approach can be used to incorporate a stochastic Eularian bias. Other alternatives for bias are nonlocal Eulerian biases (Bower et al. 1993) and Lagrangian biasing schemes, which are intrinsically nonlocal in Eulerian coordinates. The latter is possibly the most natural to arise from following halos in the Press-Schechter approach to galaxy formation (Press & Schechter 1974). Our approach is more phenomenological and we do not consider these possibilities here.
Combining equations (7) and (9) gives to second order (Heavens et al. 1998)
$$\delta _g^s=b_1\delta u^{}+\frac{b_2}{2}(\delta ^2\sigma _0^2)+u^2b_1\delta u^{}b_1\delta ^{}u+uu^{\prime \prime }$$
(10)
where all quantities are evaluated at $`𝒔`$.
In order that this expression has the correct expectation value, $`\delta _g^s=0`$, we have set $`b_0=b_2\sigma _0^2/2`$. Finally we rewrite equation (10) in terms of the linear density, gravity and tidal fields we shall use in later calculations;
$$\delta _g^s=\delta _{g,1}^s+\delta _{g,2}^s+\mathrm{\Delta }_{g,1}^s+\mathrm{\Delta }_{g,2}^s,$$
(11)
where
$$\delta _{g,1}^s=b_1\delta _1,$$
(12)
$$\delta _{g,2}^s=b_1\delta _2+\frac{b_2}{2}(\delta _1^2\sigma _0^2),$$
(13)
$$\mathrm{\Delta }_{g,1}^s=f_1\left(E_{zz}+\frac{\delta _1}{3}\right),$$
(14)
$`\mathrm{\Delta }_{g,2}^s`$ $`=`$ $`\mathrm{\Delta }_{g,1}^{s2}+\delta _{g,1}^s\mathrm{\Delta }_{g,1}^s{\displaystyle \frac{f_2}{3}}\kappa \left(\mathrm{\Pi }_{zz}^\delta +{\displaystyle \frac{\delta _1^2}{3}}\right)`$ (15)
$`+`$ $`{\displaystyle \frac{f_2}{2}}\kappa \left(\mathrm{\Pi }_{zz}^E+{\displaystyle \frac{E^2}{3}}\right)f_1g_z\left(b_1\eta _z+f_1F_{zzz}\right).`$
The first two terms in equation (11) deal with the non-linear evolution and second order bias, the last two terms bring in the effects of redshift distortions.
The new dynamical variables required for redshift-space are the gradient of the tidal field
$$F_{ijk}(𝒓,t)=_i_j_k^2\delta _1(𝒓,t),$$
(16)
and the second-order tidal fields
$$\mathrm{\Pi }_{ij}^\delta (𝒓,t)=(_i_j^2\frac{1}{3}\delta _{ij})\delta _1^2(𝒓,t),$$
(17)
and
$$\mathrm{\Pi }_{ij}^E(𝒓,t)=(_i_j^2\frac{1}{3}\delta _{ij})E^2(𝒓,t).$$
(18)
Note that in writing equations (11) to (15) we have made the plane parallel approximation, evaluating all of the redshift space contributions along a single cartesian axis.
The remaining parameters in the second-order model are (Bouchet et al. 1992, see also Catelan et al. 1995)
$`\kappa {\displaystyle \frac{D_2}{D_1^2}}`$ $``$ $`3/7\mathrm{\Omega }^{2/63}\mathrm{open}\mathrm{universe}`$ (19)
$``$ $`3/7\mathrm{\Omega }_m^{1/143}\mathrm{flat}\mathrm{universe}`$
where $`\delta _1(𝒙,t)=D_1(t)\epsilon _1(𝒙)`$, and $`D_1(t)(1+z)^{\mathrm{\Omega }^{0.6}}`$ is the linear growth function for an open model (Peebles 1980). The redshift-space parameters are (Peebles 1980, Lahav et al. 1991, Martel 1991)
$`f_1{\displaystyle \frac{d\mathrm{ln}D_1}{d\mathrm{ln}a}}`$ $``$ $`\mathrm{\Omega }_m^{3/5}\mathrm{open}\mathrm{universe}`$ (20)
$``$ $`\mathrm{\Omega }_m^{5/9}\mathrm{flat}\mathrm{universe}`$
and (Hivon et al. 1995)
$`f_2{\displaystyle \frac{d\mathrm{ln}D_2}{d\mathrm{ln}a}}`$ $``$ $`2\mathrm{\Omega }^{4/7}\mathrm{open}\mathrm{universe}`$ (21)
$``$ $`2\mathrm{\Omega }_m^{6/11}\mathrm{flat}\mathrm{universe}`$
Note that our expression for the redshift distorted and biased galaxy distribution, equation (11), does not agree with Hivon et al. (1995), even taking into account their slightly different projection into redshift space. Hivons et al.’s expression does not appear to have the correct expectation value, i.e. their quantity $`\delta ^s0`$.
### 2.2 The distribution of initial fields
Following the analysis of TW2000 we aim to find the joint probability for the fields in equation (10). Defining the parameter vector $`𝒚=(\delta _1,𝜼,𝒈,𝑬,\mathrm{\Pi }_{zz}^\delta ,\mathrm{\Pi }_{zz}^E,F_{zzz})`$, the initial distribution function is given by
$$P(𝒚)=\frac{1}{((2\pi )^ndet𝑪)^{1/2}}\mathrm{exp}\left(\frac{1}{2}𝒚^t𝑪^1𝒚\right),$$
(22)
where the covariance matrix is
$$𝑪=𝒚𝒚^t.$$
(23)
Taking the plane-parallel approximation for redshift-space distortions represents a significant simplification since any dependence of $`\delta _g^s`$ on the off-diagonal parts of $`\mathrm{\Pi }_{ij}^\delta `$, $`\mathrm{\Pi }_{ij}^E`$ and $`F_{ijk}`$ is removed, reducing the size of the required covariance matrix from $`34\times 34`$ to $`15\times 15`$.
Equation (22) is an approximation. The variables $`\mathrm{\Pi }^\delta `$ and $`\mathrm{\Pi }^E`$ are not strictly Gaussian random fields as they are generated from the square of linear fields. However, our definition of these quantities as trace-free makes them uncorrelated with all of the other fields. If we then assume that their distribution is Gaussian, as we may expect them to be as a result of the Central Limit Theorem, they become statistically independent Gaussian fields. To test our Gaussian approximation of the distribution of the $`\mathrm{\Pi }`$-fields we generated random Gaussian $`\delta `$-fields and calculated $`\mathrm{\Pi }_{zz}^\delta `$. From this we estimated its distribution and found that a Gaussian with variance $`\sigma ^2(\mathrm{\Pi }_{zz}^\delta )=8\sigma _0^4/45`$ (see equation 24) was indeed a good approximation. Later on, in Section 2.5, we shall see that these terms do not contribute to the skewness of the final distribution to lowest order. Hence the effects of this approximation will only be apparent in the higher moments and at higher order.
The non-zero elements of the covariance matrix are given by.
$`\delta ^2`$ $`=`$ $`\sigma _0^2,\eta _i\eta _j={\displaystyle \frac{1}{3}}\sigma _1^2\delta _{ij},`$
$`g_ig_j`$ $`=`$ $`{\displaystyle \frac{1}{3}}\sigma _1^2\delta _{ij},\eta _ig_j={\displaystyle \frac{1}{3}}\sigma _0^2\delta _{ij},`$
$`E_{ij}E_{kl}`$ $`=`$ $`{\displaystyle \frac{1}{15}}\sigma _0^2\left(\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}{\displaystyle \frac{2}{3}}\delta _{ij}\delta _{kl}\right)`$
$`\mathrm{\Pi }_{zz}^\delta \mathrm{\Pi }_{zz}^\delta `$ $`=`$ $`{\displaystyle \frac{8}{45}}\sigma _0^4\mathrm{\Pi }_{zz}^E\mathrm{\Pi }_{zz}^E={\displaystyle \frac{22}{135}}\sigma _0^4`$
$`g_iF_{jkl}`$ $`=`$ $`{\displaystyle \frac{1}{15}}\sigma _0^2(\delta _{ij}\delta _{kl}+\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk})`$
$`F_{zzz}\eta _z`$ $`=`$ $`{\displaystyle \frac{1}{5}}\sigma _0^2F_{zzz}F_{zzz}={\displaystyle \frac{1}{7}}\sigma _1^2`$ (24)
Evaluating the action of equation (22) we find
$`𝒚^t𝑪^1𝒚`$ $`=`$ $`{\displaystyle \frac{\delta ^2}{\sigma _0^2}}+{\displaystyle \frac{3}{(1\gamma _\nu ^2)}}\left({\displaystyle \frac{\eta ^2}{\sigma _1^2}}+{\displaystyle \frac{g^2}{\sigma _1^2}}2\gamma _\nu {\displaystyle \frac{𝜼.𝒈}{\sigma _1\sigma _1}}\right)`$ (25)
$`+`$ $`15{\displaystyle \frac{E^2}{\sigma _0^2}}+{\displaystyle \frac{63}{4\sigma _1^2}}\left(\eta _z{\displaystyle \frac{5}{3}}F_{zzz}\right)^2`$
$`+`$ $`{\displaystyle \frac{45}{8\sigma _0^4}}\left(\mathrm{\Pi }_{zz}^\delta \right)^2+{\displaystyle \frac{135}{22\sigma _0^4}}\left(\mathrm{\Pi }_{zz}^E\right)^2,`$
with determinant
$$det𝑪=\frac{22}{7}\frac{2^5}{3^{15}5^9}\sigma _0^{32}\sigma _1^2\frac{(1\gamma _\nu ^2)^3}{\gamma _\nu ^6}.$$
(26)
The correlation parameter is defined by
$$\gamma _\nu \frac{\sigma _0^2}{\sigma _1\sigma _1}.$$
(27)
which is the correlation coefficient of the velocity and gradient density fields. The variances are defined as
$$\sigma _n^2(k)=D_1^2(t)_0^{\mathrm{}}\frac{dk}{2\pi ^2}k^{2+2n}P(k)e^{k^2R_s^2},$$
(28)
where $`R_s`$ is the smoothing scale. For linear CDM power spectra $`\gamma _\nu `$ is found to lie in the range $`0.50<\gamma _\nu <0.65`$ (TW2000) for a wide range of scales. TW2000 discuss the effect of varying $`\gamma _\nu `$ and find that it is very small in the given range. We henceforth adopt the value $`\gamma _\nu =0.55`$ for all subsequent plots.
### 2.3 Propagation of the density distribution function
The multivariate Gaussian distribution for initial fields is propagated to later times using the Chapman-Kolmogorov equation;
$$P(𝜷)=𝑑𝜶W(𝜶|𝜷)P(𝜶).$$
(29)
where $`W(𝜶|𝜷)`$ is the transition probability from the state $`\alpha `$ to $`\beta `$. For a deterministic process, such as the gravitational evolution of a system from the linear to non-linear regimes, the transition probability is given by a Dirac delta function
$$W(\delta _g^s|𝒚)=\delta _D\left[\delta _g^s\delta _1^s\delta _2^s\mathrm{\Delta }_1^s\mathrm{\Delta }_2^s\right].$$
(30)
The advantage of using the Chapman-Kolmogorov approach is that is generalizes the change of variables to allow for stochasticity in the transformation between initial and final states. This will be further discussed in §2.4.
The probability distribution function for $`\delta _g^s`$ can be written as an expectation over the stochastic variables $`𝒚`$
$$P(\delta _g^s)=\delta _D(\delta _g^s\delta _1^s\delta _2^s\mathrm{\Delta }_1^s\mathrm{\Delta }_2^s)_y$$
(31)
The Characteristic Function for $`P(\delta _g^s)`$, defined by
$$𝒢(J)_{\mathrm{}}^{\mathrm{}}𝑑\delta _g^sP(\delta _g^s)\mathrm{exp}\left(iJ\delta _g^s\right)$$
(32)
and with inverse
$$P(\delta _g^s)=_{\mathrm{}}^{\mathrm{}}\frac{dJ}{2\pi }𝒢(J)\mathrm{exp}\left(iJ\delta _g^s\right)$$
(33)
can be expressed as the expectation value
$$𝒢(J)=\mathrm{exp}iJ(\delta _1^s+\delta _2^s+\mathrm{\Delta }_1^s+\mathrm{\Delta }_2^s)_y.$$
(34)
As in TW2000, equation (34) reduces to a set of Gaussian type integrals which yield the probability generating function, $`𝒢(J)`$, for a smoothed galaxy density field in redshift space;
$`𝒢(J)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{5}{}}}\theta _i^{1/2}\mathrm{exp}\left[{\displaystyle \frac{1}{2}}iJ\sigma _0^2b_2\right]\mathrm{exp}\left[{\displaystyle \frac{J^2\sigma _0^2N^2}{2\theta _1}}\right]`$ (35)
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \frac{2}{45}}{\displaystyle \frac{J^2\sigma _0^2f_1^2}{\theta _3}}\left(1{\displaystyle \frac{iJ\sigma _0^2NH}{f\theta _1}}\right)^2\right]`$
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{2}}J^2\sigma _0^8f_2^2\kappa ^2{\displaystyle \frac{1379}{1080}}\right]`$
where
$`\theta _1`$ $`=`$ $`1+2iJ\sigma _0^2A`$
$`\theta _2`$ $`=`$ $`\left(1+{\displaystyle \frac{4}{15}}iJ\sigma _0^2B\right)^4`$
$`\theta _3`$ $`=`$ $`1+{\displaystyle \frac{4}{15}}iJ\sigma _0^2B+{\displaystyle \frac{8}{45}}iJ\sigma _0^2f_1^2+{\displaystyle \frac{8}{45}}{\displaystyle \frac{J^2\sigma _0^4H^2}{2\theta _1}}`$
$`\theta _4`$ $`=`$ $`\left(1{\displaystyle \frac{2}{3}}iJ\sigma _0^2b_1+{\displaystyle \frac{1}{9}}J^2\sigma _0^4b_1^2{\displaystyle \frac{(1\gamma _\nu ^2)}{\gamma _\nu ^2}}\right)^2`$
$`\theta _5`$ $`=`$ $`1i{\displaystyle \frac{1}{5}}J\sigma _0^2\left(f_1^2{\displaystyle \frac{10}{3}}M\right)`$
$`+`$ $`J^2\sigma _0^4\left[M{\displaystyle \frac{(1\gamma _\nu ^2)}{9\gamma _\nu ^2}}\left(M+{\displaystyle \frac{18}{15}}f_1^2\right)+{\displaystyle \frac{1}{21}}{\displaystyle \frac{f_1^4}{\gamma _\nu ^2}}\left(1{\displaystyle \frac{21}{25}}\gamma _\nu ^2\right)\right]`$
and
$`A`$ $`=`$ $`{\displaystyle \frac{1}{3}}(2\kappa )b_1++{\displaystyle \frac{1}{9}}f_1^2+{\displaystyle \frac{1}{3}}b_1f_1{\displaystyle \frac{1}{9}}f_2\kappa +{\displaystyle \frac{1}{2}}b_2`$
$`B`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+\kappa )b_1+{\displaystyle \frac{1}{6}}f_2\kappa `$
$`H`$ $`=`$ $`f_1\left(b_1+{\displaystyle \frac{2}{3}}f_1\right)`$
$`M`$ $`=`$ $`b_1(1+f_1)`$
$`N`$ $`=`$ $`b_1+{\displaystyle \frac{1}{3}}f_1`$ (37)
The 1-pt PDF can then be found by numerical integration of equation (33). Equations (35) to (37) represent the main analytical result of the this paper.
### 2.4 Stochastic Bias Schemes
Instead of the local deterministic prescription for galaxy biasing described by equation (9) we may choose to implement a stochastic Eularian biasing scheme. In stochastic biasing (Pen 1998, Dekel & Lahav 1999) the relation between the underlying density field of dark matter and the galaxy density field is a random process with joint distribution $`P(\delta _g,\delta )`$.
If we assume that the biasing transition distribution, $`P(\delta _g,\delta )`$, is a bivariate Gaussian of the form
$`P(\delta _g,\delta )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi [\sigma _b^2+b_1^2\sigma _0^2(1r^2)]}}`$ (38)
$`\times `$ $`\mathrm{exp}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta _g^2+\delta ^2b_1^2(1+\sigma _b^2/b_1^2\sigma _0^2)2rb_1\delta _g\delta }{\sigma _b^2+b_1^2\sigma _0^2(1r^2)}},`$
where the covariances between the fields are given to first order by
$`\delta _g^2`$ $`=`$ $`b_1^2\sigma _0^2+\sigma _b^2,`$ (39)
$`\delta _g\delta `$ $`=`$ $`b_1r\sigma _0^2`$ (40)
where $`\sigma _b`$ is the variance of the random component determining galaxy formation and $`r`$ is the correlation coefficient between the mass and galaxy distribution, then the galaxy distribution function can be written
$$P(\delta _g)=P(\delta _g,f(𝒚))_y,$$
(41)
where $`f(𝒚)=\delta _1+\delta _2`$ in real space. The characteristic function for galaxies is then
$$G_{\delta _g}(J)=𝑑J^{}G_\delta (J^{})G_{\delta _g,\delta }(J,J^{})$$
(42)
where
$$\mathrm{ln}G_{\delta _g,\delta }(J,J^{})=\frac{1}{2}(J^2\sigma _0^2+J^2(b_1^2\sigma _0^2+\sigma _b^2)+2JJ^{}b_1r\sigma _0^2)$$
(43)
Equation (42) can be evaluated numerically.
The distribution function of galaxies can also be evaluated by numerically integrating
$$P(\delta _g)=_{\mathrm{}}^{\mathrm{}}\frac{dJ}{2\pi }𝒢_\delta (J)P(\delta _g,J)$$
(44)
where
$$P(\delta _g,J)=𝑑\delta e^{iJ\delta }P(\delta _g,\delta )$$
(45)
is the partially transformed characteristic function. We shall explore the effects of stochastic bias elsewhere, but note that it can be naturally incorporated into our probabilistic scheme. For the remainder of the paper we shall only discuss the deterministic scheme.
### 2.5 Skewness of a biased density field in redshift space
The skewness of the biased redshift-space galaxy distribution can be found from the characteristic function, $`𝒢(J)`$, by differentiation. The moment parameters are defined as
$$S_n=\frac{\delta ^n_c}{\sigma _0^{2n2}},$$
(46)
where the cumulants (or connected moments) of the field are given by
$$\delta ^n_c=\frac{^n}{[iJ^n]}\mathrm{ln}𝒢(J=0).$$
(47)
In redshift space the skewness parameter becomes
$$S_3^s=\frac{(\delta _g^s)^3_c}{\sigma _s^4},$$
(48)
where $`\sigma _s^2`$ is the linear, redshift space variance of the galaxy density field, $`\delta _g^s`$.
Taking the third derivative and setting $`J=0`$ we find the redshift-space skewness parameter;
$`S_3^sb_1^4F^2`$ $`=`$ $`4b_1^3\left(1+{\displaystyle \frac{2}{3}}\beta _1+{\displaystyle \frac{3}{25}}\beta _1^2\right)+3b_2b_1^2\left(1+{\displaystyle \frac{\beta _1}{3}}\right)^2`$ (49)
$``$ $`2\kappa b_1^3(1+{\displaystyle \frac{1}{3}}\beta _2)\left(1+{\displaystyle \frac{2}{3}}\beta _1+{\displaystyle \frac{7}{75}}\beta _1^2\right)`$
$`+`$ $`2b_1^4\beta _1\left(1+{\displaystyle \frac{19}{15}}\beta _1+{\displaystyle \frac{3}{5}}\beta _1^2\right)+{\displaystyle \frac{6}{25}}b_1^4\beta _1^4`$
where $`F(\beta _1)`$ is the linear redshift space enhancement factor for the redshifted variance given to lowest order by (Kaiser 1987)
$$F=1+\frac{2}{3}\beta _1+\frac{1}{5}\beta _1^2.$$
(50)
and
$$\beta _1\frac{f_1}{b_1},\beta _2\frac{f_2}{b_1}.$$
(51)
The variance in redshift space is just
$$\sigma _s^2=b_1^2\left(1+\frac{2}{3}\beta _1+\frac{1}{5}\beta _1^2\right)\sigma _0^2.$$
(52)
Equation (49) for the galaxy skewness in redshift space is the second major result of this paper. It is interesting to compare this with the results of Hivon et al. (1995) who derived the density skewness in redshift space. Hivon et al. derived their result by transforming the redshifted bispectrum but did not find a closed form solution and were left with a term containing an integral that was to be determined numerically. Our result for the skewness should coincide with that of Hivon et al. for the special case of $`b_1=1`$ and $`b_2=0`$, but does not. The differences between our result and that of Hivon et al. may result from the different approaches, or from the difference in initial expressions for $`\delta _g^s`$. We have not not been able to trace the origin of these differences. We have verified equation (49) by a direct calculation of the redshifted galaxy skewness, using equation (11) along with the correlators in equation (24), to evaluate the quantity $`(\delta _g^s)^3`$.
Equation (49) displays the hierarchical scaling associated with galaxy clustering. It has already been shown that a local bias will retain a hierarchical form if the underlying field displays hierarchical behavior (Fry & Gastañaga 1993), and that the lowest order of perturbation theory reproduces the hierarchical structure. Our result would seem to indicate that this structure holds for galaxies in redshift-space. This is an important result as it strengthens the claim that if galaxies display a hierarchical scaling then it is consistent that the underlying density field was initially Gaussian distributed and bias is local. Our results suggest that this can be extended as a test of galaxy distributions in redshift-space.
In Figure 1 we plot the skewness parameter $`S_3^s`$ as a function of linear bias, $`b_1`$, for different values of $`\mathrm{\Omega }_m`$ and a quadratic bias $`b_2`$. An immediate result is that in all the models the redshift-space skewness parameter, $`S_3^s`$, is less than the real-space skewness parameter, $`S_3`$, until $`b_1>1.5`$. This may seem surprising as one expects the effects of redshift space distortions will tend to make the density field look more evolved along the line of sight. This does happen, but the second-order skewness is dominated by the first-order increase in the variance, so that the ratio $`(\delta ^s)^3/(\delta ^s)^2^2`$ is smaller than its real space counterpart. This effect was previously observed by Hoyle et al. (2000).
Beyond $`b_11.5`$ the redshift skewness parameter rises above its real space value for all the models. This is because for large values of the linear bias parameter the second order skewness becomes comparable in magnitude to the redshifted variance.
For $`\beta _1\beta _21`$ and $`b_11`$ the redshifted skewness can be approximated by
$$S_3^s=\frac{1539}{343}b_1^{0.61}\beta _1^{0.02}+\frac{75}{49}\frac{b_2}{b_1^2}\beta _1^{9/14}$$
(53)
which is accurate to a few percent. This can be compared with the undistorted skewness (Fry & Gaztañaga 1993)
$$S_3=\frac{34}{7}\frac{1}{b_1}+3\frac{b_2}{b_1^2}.$$
(54)
Without second-order bias, $`S_3^s`$ is only weakly dependent on $`\beta _1`$. The main effect of redshift-space distortions in this case is to change the dependence on $`b_1`$. The second-order bias term, however, has a much stronger dependence on $`\beta _1`$, and has a lower amplitude than the undistorted term when $`\beta _1=1`$.
## 3 The shape of the PDF
The solution for the generating function, $`𝒢(J)`$, given by equations (35) – (37) may be integrated numerically without too much difficulty to yield the full distribution function $`P(\delta )`$. In figure 2 we show the shape of the resulting PDF for a range of values of the parameters $`b_1,b_2`$ and $`\mathrm{\Omega }_m`$.
To second order there are two important quantities working to distort the shape of the PDF from a Gaussian with linear mass variance $`\sigma _0`$. Firstly is the variance of the resulting field which acts to broaden the distribution. Secondly is skewness, which produces asymmetry in the PDF about its mean at $`\delta =0`$. These two quantities are effected to varying degrees by non-linear evolution, by non-linear bias and by the distorting effect of galaxy peculiar motions.
For the plots in figure 2 we chose a moderate value for the linear variance of the underlying dark matter density field, $`\sigma _0=0.20`$. The correlation parameter, $`\gamma _\nu `$, was 0.55. In the top panel we show the effect of a purely linear bias with values $`b_1=`$0.7, 1.0 and 1.2 and with $`b_2=0`$. For the lower panel the linear $`b_1`$’s match those above but in this case $`b_2=1.0`$ for every plot. Each plot shows three different PDFs: one (solid lines) in real space, one in redshift space for the case $`\mathrm{\Omega }_0=1.0`$ (dotted lines) and a third in redshift space but with $`\mathrm{\Omega }_0=0.3`$ (dashed lines). We have inset the same plots on logarithmic axis in order to emphasize the tails of the distribution.
In real space and with linear bias only, the shape of the PDF is dominated by the $`b_1^2`$ boost to the variance. Where $`b_1>1`$ the change to $`S_3^s`$ is small and any effect on the shape of the PDF caused by the skewness is masked by the enhanced variance When $`b_1<1`$ there is a more substantial change to $`S_3^s`$ but in this case the variance becomes small enough that the PDF looks relatively Gaussian regardless of what is happening to the skewness.
To second order the non-linear bias parameter, $`b_2`$, does not alter the variance of the underlying mass density field. The changes to the shape of the PDF when $`b_20`$ then reflect only the corresponding changes to $`S_3^s`$ We observe that the real space PDFs appear to peak sharply around $`\delta =1\sigma _s`$ with an abrupt drop off of the low density tail to the left of the peak. The high density tail is extended beyond the corresponding case where $`b_2=0`$. With low $`b_1`$ the effect of $`b_2`$ is most pronounced, causing an increasingly abrupt drop off in $`P(\delta )`$ on the low density side of the peak. This pronounced effect on the PDF would appear to be an important signature of quadratic bias, which cannot be reproduced by any combination of other parameters.
In redshift space the PDF is dominated by the enhancement of the variance due to linear redshift space distortions. As discussed in §2.5 the second-order change in the skewness is less than the linear change in variance resulting in a lower skewness parameter, $`S_3^s`$, in redshift-space than in real space. The result for the PDFs is a broader and less asymmetric appearance. This is most clearly seen in the PDFs where $`b_20`$. For these plots, the sharp peak in $`P(\delta )`$ around $`\delta =1\sigma _s`$ is significantly damped, and both the low and high density tails are much shallower than their real space counterparts.
## 4 Comparison of results with n-body simulations
### 4.1 The simulations
In this section we illustrate how our theoretical PDF compares with that measured from cosmological N-body simulations. For the comparison we chose to use Hugh Couchman’s Adaptive P<sup>3</sup>M N-body code (Couchman 1991) to model the evolution of the dark matter. The simulation volume was a cube of comoving side 200 $`h^1`$Mpc with periodic boundary conditions. We chose a CDM initial power spectrum (Bardeen at al. 1986) that was normalized to match the present day abundance of clusters. The variance on a scale of 8 $`h^1`$Mpc when the particle distribution was smoothed with top hat filters was $`\sigma _8=0.5\mathrm{\Omega }_m^{0.5}`$ in the final time-step of the simulation. The shape parameter used was $`\mathrm{\Gamma }=0.25`$, representing the best fit of the CDM power spectrum to galaxy clustering data. The simulation was performed on a 128<sup>3</sup> Fourier mesh with $`100^3`$ particles.
We investigated the numerical PDFs in three sets of circumstances:
* in real space with bias,
* in redshift space without bias,
* in redshift space with bias.
The first two scenarios were trivial to construct from the simulations. For the case of bias in real space, we binned the data and estimated the overdensity from $`\delta =(n_p\overline{n}_p)/\overline{n}_p`$, where $`n_p`$ is the number of particles in a cell and $`\overline{n}_p`$ is the mean number of particles per cell. The biased distribution was then estimated by the transformation
$$\delta _g=b_1\delta +\frac{b_2}{2}(\delta ^2\sigma _0^2)$$
(55)
where $`\sigma _0`$ is the variance on the scale of the binning.
Redshift distortions, in the absence of bias, were calculated using the peculiar velocities of the simulation particles. The distortions were made plane parallel in order to match our approximation from §2.1. The resulting particle distribution was smoothed with Gaussian filters of radius $`R_s`$. The PDF was then evaluated from the relative abundance of $`\delta `$ across the grid.
Measuring the combined effect of redshift distortions and bias was more difficult. A simple combination of the above methods could not be used, as we needed to identify biased galaxies before making the transformation to redshift space. Our basic biasing method calculates the biased density field on the grid so that information about individual galaxies was lost. In order to identify individual biased galaxies in the simulation we sampled a population of galaxies from the simulation particles so that the number of galaxies in a cell, $`n_g`$, was
$$n_g=\overline{n}_g\left[1+b_1\delta +\frac{b_2}{2}(\delta ^2\sigma _0^2)\right].$$
(56)
The resulting distribution could then be transformed to redshift-space by using the galaxies velocity.
The size of the grid was dictated by the need for $`n_p`$ to be sufficiently larger than $`n_g`$ so that underdense regions could be enlarged by the bias. Given this constraint we chose to smooth the particle distribution on a coarse mesh of between 9 and 11 cells, with a cell size of $`l=22h^1\mathrm{Mpc}`$ and $`18h^1\mathrm{Mpc}`$, respectively. As we calculate the linear variance using a Gaussian filter (equation 28) this corresponds to a smoothing scale of $`Rs=l/\sqrt{12}=6.5`$ and $`5h^1\mathrm{Mpc}`$. The total number of galaxies was constrained by the requirement that $`n_g`$ was not larger than $`n_p`$ in areas of high density, so we set the total number in the simulation volume $`N_g=5\times 10^5`$.
The errors on the numerical PDFs for the biased and distorted only simulations were taken from the standard deviation over 5 independent realisations of the simulated volume. Errors due to shot noise were incorporated into the theoretical PDF (TW2000) based on the mean density of particles in a cubical cell of side $`l=\sqrt{12}R_s`$. In practice the shot noise contribution was small for simulations due to the high particle density.
In the case of a combined biased and redshift-space distorted PDFs the need to sample the galaxies on a coarser mesh meant that the resulting PDF from an individual simulation had a larger scatter from simulation to simulation for the same value of $`\sigma _s^2`$. To avoid this we used a simulation with a smaller clustering variance. A relatively large number of particles could then to be found in each cell. For the case $`R_s5h^1`$Mpc there numbered on average 375 galaxies per cell sampled from around 750 simulation particles. When the numerical PDFs were averaged over 7 independent realisations the results were reasonably smooth, though the error bars (again the standard deviation over the 7 realisations) do reflect the scatter.
### 4.2 The Edgeworth Expansion
In addition to the PDF calculated from this work we also examine the PDF reconstructed from the Edgeworth expansion (Juszkiewicz et al. 1995, Bernardeau & Kofman 1995) using the redshift-space skewness parameter derived in Section 2.5. The distribution function from the Edgeworth expansion is, to second order,
$$P(\nu )=\frac{1}{\sqrt{2\pi }}\left[1+\frac{1}{3!}H_3(\nu )S_3^s\right]\mathrm{exp}\left[\frac{1}{2}\nu ^2\right]$$
(57)
where $`\sigma _s^2`$ is the redshift-space galaxy variance given by equation (52), $`S_3^s`$ is the skewness parameter from equation (49), $`H_n(\nu )`$ is an $`n^{th}`$-order Hermite polynomial and $`\nu =\delta ^s/\sigma _s`$ is the scaled redshift-space density field.
The main difference between the Edgeworth expansion and our own model for the PDF is that the Edgeworth relies upon two levels of approximation: firstly in the moments, which come from — in our case — a second order perturbative calculation; and secondly in the Edgeworth series itself, which reconstructs the non-linear PDF by expanding about a Gaussian distribution. Our PDF comes directly from the perturbation theory without any further constraints other than those implicit in the theory itself. In this way our model is perhaps more representative of the perturbation theory and its predictions and limitations. Another dissimilarity between the models, and a problem with trying to derive the PDF to second order, lies in the fact that second order perturbation theory does not get the higher moments of the distribution correct. In our calculation all of the higher moments exist but have terms missing that would come in in a higher order calculation. This may strongly effect the shape of our PDF, particularly around the peak. Such problems are not an issue for the Edgeworth and other similar approximations because in these models the cumulants are either put into the series explicitly, or else are zero.
### 4.3 Results
In comparing the models with simulation data we constructed the PDF so that the variance of the distorted, biased density field, $`\sigma _s^2`$, was constant for all choices of the parameters $`\mathrm{\Omega }`$, $`b_1`$ and $`b_2`$. This was arranged by using an appropriate value for $`\sigma _0`$ for each plot. For the case of bias in real space and of no bias in redshift space we chose $`\sigma _s=0.2`$. This relatively moderate variance was necessary because for the cases where $`b_1<1`$ the underlying $`\sigma _0`$ had to start quite large. For the combination of bias and redshift distortions we were limited to starting with a slightly higher linear mass variance because the simulation PDFs became dominated by sampling variance when the cell size was large. For these plots $`\sigma _s=0.25`$.
Figures 3 and 4 show the behavior of the PDF with various combinations of the bias parameters and with no redshift space distortions. The density parameter for these simulations was $`\mathrm{\Omega }=1.0`$, though in real space the quasi-linear evolution is not sensitive to this quantity (Bouchet et al. 1992, Martel & Freudling 1991).
The first set of plots (figure 3) show the effect of a linear bias only, $`b_1=0.7`$ in the top panel and $`b_1=1.5`$ in the bottom. The linear mass variance was $`\sigma _0=0.29`$ and $`0.13`$ respectively. The dominant effect on the shape is through the variance which tends to broaden and lower the amplitude PDF for $`b_1>1`$ and conversely makes the PDF more Gaussian when $`b_1.<1`$. The fit to the data (points) of the PDF from this work is good in both cases, though when $`\sigma _0`$ is relatively large our model misses the peak of the distribution and appears to be slightly too strongly skewed. The Edgeworth approximation (dotted lines) also shows excellent agreement with the simulations, particularly in the vicinity of the peak. The plots show the absolute value of the Edgeworth PDF. The “lobes” in the low density tail of the distribution represent negative probabilities in the Edgeworth PDF.
In figure 4 we show the fit to simulations when the second order bias term, $`b_2`$, is 1.0. The linear bias was set to $`b_1=0.7`$ (top) and $`b_1=1.0`$ (bottom). The linear mass variance was $`\sigma _0=0.29`$ (top) and $`\sigma _0=0.20`$ (bottom). To the order we are interested in $`b_2`$ only contributes to the PDF via $`S_3`$ and the higher moments. The effect of the second order bias is to fill in the void regions and enhance the peaks. The result for the PDF is a sharp drop-off of the low density tail and a dramatic amplification of the peak just beyond that point. The high density tail is also extended to account for the increased regions of very high density. The fit to the simulations for second order biasing is good so long as $`\sigma _0`$ is reasonably small, though still in the quasi-linear regime. The breakdown in the fit comes mainly in the low density tail which for the theoretical curves drops away too slowly when the variance is high. The Edgeworth PDF breaks down quite badly in the tails when $`b_20`$. This suggests that the Edgeworth approximation becomes rather unstable when the skewness parameter is large, even when the variance is low. The peak of the Edgeworth PDF is not so badly affected, however, and fits the data nicely so long as $`\sigma _0`$ is not too high.
Figure 5 shows the fit of the PDF to an unbiased field of galaxies in redshift space. We ran two sets of simulations with $`\mathrm{\Omega }_0=`$ 1.0 and 0.3. The linear variance used for the theoretical PDFs was $`\sigma _0=0.15`$ and $`0.17`$. The fit to both our model and the Edgeworth PDF is very good, although the linear variance in each case was quite small to ensure that $`\sigma _s=0.2`$. We found that our approximation began to break down when the redshift space variance was $`\sigma _s=0.4`$. Although the Edgeworth followed the peak and low density part of the distribution well to this variance and higher, its high density tail became unstable at a much lower $`\sigma _s`$. In both cases the PDFs look more Gaussian than in real space because of the reduction in $`S_3^s`$. The theoretical plots are almost identical in each case, as expected since we have arranged for the variance of the redshifted field to be identical in each case and anyway the skewness in redshift space is only very weakly dependant upon $`\mathrm{\Omega }_m`$ when $`b_1`$ is close to unity. The numerical PDFs are marginally different, although we suspect that this is a numerical effect as suggested by the large error bars around the peak of the distribution for $`\mathrm{\Omega }_m=0.3`$.
On much smaller scales we would expect there to be a real difference between the two PDFs as fingers of god, due to the pairwise velocities in clusters, become important. Our second-order analysis does not allow for these strongly nonlinear effects, and we restrict our analysis to scales large enough that these effects are not important. The agreement between theory and the simulations suggests that this is true. The effect we would expect to see would be a decrease in the variance and skewness of the PDFs in redshift-space.
Our final plots (figure 6) show the combined effect of both bias and redshift distortions. In the top row $`b_1=1.2`$ with $`b_2=0.0`$ and on the bottom we have set $`b_1=1.0`$ and made $`b_2=1.0`$. For both cases the density parameter was $`\mathrm{\Omega }_m=1.0`$. In each of the plots $`\sigma _s=0.25`$ — due to problems with creating the bias/redshift simulations we were unable to satisfactorily measure the PDFs for $`\sigma _s=0.2`$. The linear mass variance was therefore $`\sigma _0=0.18`$ in the top plot and $`0.19`$ in the bottom. The fit in both cases is very good — although the Edgeworth approximation becomes slightly unstable when $`b_2=1.0`$ because of the relatively high skewness. Although the redshift distortions do wash out some of the effect of the bias parameters, there are dependencies which are reproduced by the models to a reasonable degree of accuracy. This is encouraging for the purposes of constraining cosmological models using the PDF and data from galaxy redshift surveys.
In general there is a good similarity between all of the numerical PDFs and those found from either our analysis based on the Chapman–Kolmogorov equation or the Edgeworth approximation with the redshift space skewness parameter $`S_3^s`$. The Edgeworth expansion clearly breaks down sooner in the tails of the distribution, and for large values of the skewness parameter. Considering that both models use second order perturbation theory and have the same skewness and variance it is perhaps confusing that they behave so differently with various combinations of the parameters. However we must not forget that the Edgeworth approximation and our model rely on very different approaches to obtain the PDF as discussed in §4.2.
### 4.4 Constant $`\beta `$ PDFs
An interesting and potentially useful feature of modeling the PDF in redshift space comes when one considers parameter combinations that are degenerate to linear order. The combination
$$\beta _1=\frac{\mathrm{\Omega }_m^{0.6}}{b_1},$$
(58)
for example, is a common quantity to measure in linear analysis of galaxy redshift surveys (see e.g. Tadros et al. 1999 for a recent measurement from the PSCz redshift survey).
We constructed the PDFs of two fields sharing the same $`\beta _1`$ but with different combinations of the cosmological density and linear bias parameters, neglecting the effects of second order bias. We chose $`\beta _1=0.5`$ with $`\mathrm{\Omega }_m=0.3`$ and $`b_1=1.0`$ for one model and $`\mathrm{\Omega }_m=1.0`$ with $`b_1=2.0`$ for the other. The linear mass variance for each was inferred from the observed abundance of clusters from (Vianna & Liddle 1996)
$$\sigma (R_{th})=\sigma _8\left(\frac{R_{th}}{8h^1\mathrm{Mpc}}\right)^{\gamma (R)}$$
(59)
where
$$\gamma =(0.3\mathrm{\Gamma }+0.2)\left[2.92+\mathrm{log}\left(\frac{R_{th}}{8h^1\mathrm{Mpc}}\right)\right]$$
(60)
and with $`\mathrm{\Gamma }`$ the CDM shape function. The abundance of clusters was parameterised by
$$\sigma _8=0.6\mathrm{\Omega }_0^C,$$
(61)
where $`C0.49`$ for the closed and $`C0.43`$ for the open models. The underlying linear mass variance was found for each at constant top–hat filter radius, $`R_{th}`$. The resulting PDFs are shown in figure 6, the solid curve being the high-$`\mathrm{\Omega }`$ model and the dashed line the low-$`\mathrm{\Omega }`$ model. Clearly there is a marked difference between the two PDFs, a result of the different ways in which $`b_1`$ and $`\mathrm{\Omega }`$ contribute to the variance and the skewness. We investigate elsewhere whether this difference can be used to constrain cosmological parameters.
## 5 Summary
In this paper we have derived an expression for the nonlinear 1-point probability distribution of galaxies in redshift space. This function is useful in the analysis of redshift surveys as it is a quantity that can be directly observed. We have treated the nonlinearity of the density field with second-order Eulerian perturbation theory, and transformed from the mass-density distribution to a galaxy distribution using a second-order local bias prescription. This is then mapped to redshift space, again using second-order perturbation theory. The transformation of the initial probability distribution to the evolved distribution is carried out using the Chapman-Kolmogorov equation. This allows us to derive an exact expression for the second-order characteristic function for the galaxy distribution. We show that the Chapman-Kolmogorov equation is more general than just a transformation of variables, as we can use it to calculate the effects of Stochastic Bias Schemes. We shall investigate this in more detail elsewhere.
Taking the derivative of the galaxy characteristic function we derive a new expression for the skewness parameter of galaxies in redshift space, $`S_3^s`$. Unlike other derivations, we find a closed-form expression and include the effects of a quadratic bias term. We find that in general for values of the linear bias parameter $`b_1`$ below $`1.5`$, $`S_3^s`$ is smaller than its real space value, $`S_3`$. We also find that, while the first-order bias terms are largely independent of the linear distortion parameter $`\beta _1`$, or $`\mathrm{\Omega }_m`$, for values of $`b_1`$ around unity, as suggested by previous studies, the quadratic bias terms introduce a strong dependency on cosmological parameters. An analysis of the full 1-pt PDF in real and redshift space confirms these findings. In addition we find that quadratic bias produces a distinct sharp cut-of in the PDF for low-density regions.
Comparing our PDF with those measured from N-body simulations we find good agreement for all combinations of parameters. We also find that the Edgeworth series, using our expression for the redshift-space skewness parameter, fits the numerical PDFs rather well when the series is truncated to second order and the redshift space skewness and variance are used. The models, particularly the Edgeworth approximation, fit the data most poorly when the linear mass variance and the skewness are both high ($`\sigma _0>0.4`$). This occurs on small scales or where there is a large degree of non-linear bias.
Our analysis leaves two problems unresolved. The first is the issue of fingers-of-god, which we do not model. This restricts the scales on which we can model the PDF, but does not seem to be a problem when we compare our results to simulations. We will test elsewhere if this is a problem when we come to use the PDF to extract cosmological parameters. The second issue is that of smoothing the final density field. We neglect this in our analysis, although it is included in the analysis of other PDFs (see e.g. Bernardeau 1994). Again, when comparing with the results of simulations we do not see a significant effect. This may be a result of the type of power spectra we have tested our model against. The effect of smoothing is to transfer power across scales and for some power spectra this transfer will be minimal. However for the range of spectra we have investigated we only require the variance to be correctly calibrated against simulations and we find good agreement. An advantage of the PDF over other statistical measures is that it is straightforward to see if the fit is poor.
In future work the PDF we have derived will be compared with the PDF of galaxies measured from the PSCz galaxy survey (Watts & Taylor in preparation). This will be useful both in testing the gravitational instability hypothesis and in constraining cosmological models, using data from the nonlinear regime of the galaxy distribution. This step is vital for maximising the amount of information available from galaxy redshift surveys. Although we find that redshift distortions do dampen some of the effects of bias, there are dependencies that may be exploited. We have shown how using the PDF in conjunction with the abundance of clusters may provide a way to distinguish between cosmological models that are degenerate in linear analysis.
## 6 Acknowledgements
PIRW acknowledges the PPARC for a postgraduate studentship, ANT acknowledges the PPARC for a postdoctoral fellowship.
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# Hierarchical Radiative Quark Mass Matrices with an 𝑈(1)_𝑋 Horizontal Symmetry Model1footnote 11footnote 1Published:Rev. Mex. Fis. 48-1, 32 (2002)
## I INTRODUCTION
Although for several years a great effort has been done to shed some light on the mystery of the fermion masses it still is one of the outstanding puzzles of particle physics. There have been different approaches to explain the mass hierarchy, the fermion mixing and their possible relation to new physics. A good review covering widely these topics has been presented in fritzsch . In this work we will restrict our study to that of a new horizontal symmetry and the derived radiative corrections.
A possible answer to why the masses of the light quarks are so small compared with the electroweak scale is that they arise from radiative correctionsma , while the mass of the top quark and possibly those of the bottom quark and of the tau lepton are generated at tree level. This may be understood as a consequence of the breaking of a symmetry among families (a horizontal symmetry). This symmetry may be discretediscrete or continuouscontinuous . Here we consider the case of a continuous local horizontal symmetry, a $`U(1)_X`$ gauge group broken spontaneously. We limit our calculation of masses to those of the quark sector, insisting that at tree level only the top and bottom quarks acquire mass. Instead of assuming a texture for the quark masses from the beginning, we carry out a one loop and a two loop calculation of the mass matrices in terms of some parameters which are the tree level top and bottom quark masses, one Yukawa coupling and the entries in the mass matrices of the scalar bosons that participate in the loop diagrams which contribute to the quark masses.
This paper is organized in the following way: In Section II we describe explicitly the model, Section III contains the analytical calculations, while Section IV is devoted to a numerical fit of our equations. The conclusions are presented in section V.
## II The Model
We assume only three families, the Standard Model (SM) families, and we do not introduce exotic fermions to cancel anomalies. The fermions are classified as in the SM in five sectors f = q,u,d,l and e, where q and l are the $`SU(2)_L`$ quark and lepton doublets respectively and u,d and e are the singlets, in an obvious notation. In order to reduce the number of parameters and to make the model free of anomalies, we demand that the values X of the horizontal charge satisfy the traceless conditionZEP1
$$X(f_i)=0,\pm \delta _f,$$
(1)
where $`i=1,2,3`$ is a family index, with the constraint
$$\delta _q^22\delta _u^2+\delta _d^2=\delta _l^2\delta _e^2.$$
(2)
Eq. (1) guarantees the cancellation of the \[U(1)<sub>H</sub>\]<sup>3</sup> anomaly as well as those which are linear in the U(1)<sub>H</sub> hypercharge ($`[SU(3)_c]^2U(1)_H,[SU(2)_L]^2U(1)_H,[Grav]^2U(1)_H`$ and $`[U(1)_Y]^2U(1)_H`$). Eq. (2) is the condition for the cancellation of the U(1)<sub>Y</sub>\[U(1)<sub>H</sub>\]<sup>2</sup> anomaly. A solution of Eq. (2) which guarantees that only the top and bottom quarks get masses at tree level is given by (“doublets independent of singlets”, see Ref. ZEP1 )
$$\delta _l=\delta _q=\pm \mathrm{\Delta }\delta _u=\delta _d=\delta _e=\pm \delta .$$
(3)
To avoid tree level flavor changing neutral currents, we do not allow mixing between the standard model Z boson and its horizontal counterpart. Consequently the SM Higgs scalar should have zero horizontal charge. As a consequence, and since we insist in having a non-zero tree-level mass for the top and bottom quarks, the horizontal charges of these quarks should satisfy
$`X(q_3)+X(u_3)=0,`$ $`X(q_3)+X(d_3)=0`$ (4)
in order for the Yukawa couplings in Eq. (6) to be invariant, but then Eqs. (1) and (3) demand that they vanish,
$$X(u_3)=X(q_3)=X(d_3)=0,$$
(5)
which in turn implies $`X(l_3)=X(e_3)=0`$ (this defines the third family). The assignment of horizontal charges to the fermions is then as given in Table 1. The $`SU(3)_CSU(2)_LU(1)_Y`$ quantum numbers of the fermions are the same as in the Standard Model.
To generate the first and second family quark masses radiatively we must introduce new irreducible representations (irreps) of scalar fields, since the gauge bosons of $`G=G_{SM}U(1)_X`$ do not perform transitions between different families. Families are of course distinguishable (non degenerated) only below the scale of the SM symmetry breaking, when they become massive.
Looking for scalars which make possible the generation of fermion masses in a hierarchical manner, we divide the irreps of scalar fields into two classes. Class I (II) contains scalar fields which get (do not get) vacuum expectation value (VEV).
A proper choice of scalars should be such that no VEVs are induced, through couplings in the potential, for scalars in class II. In the model considered below scalars of class II have no electrically neutral components, so they never get out of its class. In our model we introduce two irreps of scalars of class I and six irreps of scalars of class II, with the quantum numbers specified in Table 2. Notice that we introduce just the minimum number of scalars of class I; that is, only one Higgs doublet of weak isospin to achieve the Spontaneous Symmetry Breaking (SSB) of the electroweak group down to the electromagnetic $`U(1)_Q`$, and one $`SU(2)_L`$ singlet $`\varphi _2`$ used to break $`U(1)_X`$. In this way the horizontal interactions affect the $`\rho `$ parameter only at higher orders.
With the above quantum numbers the quark Yukawa couplings that can be written may be divided into two classes, those of the D type which are defined by Fig 1a, and those of the M type which are defined in Fig 1b. The Yukawa couplings can thus be written as $`L_Y=L_{Y_D}+L_{Y_M}`$, where the D Yukawa couplings are
$$L_{Y_D}=Y^u\overline{q}_{L_3}\stackrel{~}{\varphi }_1u_{R_3}+Y^d\overline{q}_{L_3}\varphi _1d_{R_3}+h.c.,$$
(6)
with $`\stackrel{~}{\varphi }i\sigma _2\varphi ^{}`$, while the M couplings compatible with the symmetries of the model are
$`L_{Y_M}=Y_I[q_{1L}^{\alpha T}C\varphi _{3\{\alpha \beta \}}q_{2L}^\beta +q_{3L}^{\alpha T}C\varphi _{3\{\alpha \beta \}}q_{3L}^\beta +q_{2L}^{\alpha T}C\varphi _{4\{\alpha \beta \}}q_{3L}^\beta `$
$`+d_{2R}^TC\varphi _5d_{1R}+d_{3R}^TC\varphi _5d_{3R}+d_{3R}^TC\varphi _6d_{2R}`$
$`+u_{2R}^TC\varphi _7u_{1R}+u_{3R}^TC\varphi _7u_{3R}+u_{3R}^TC\varphi _8u_{2R}]+h.c.`$ (7)
In these couplings C represents the charge conjugation matrix and $`\alpha `$ and $`\beta `$ are weak isospin indices. Color indices have not been written explicitly. By simplicity and economy we have assumed only one Yukawa constant $`Y_I`$ for all the M couplings. Notice that $`\varphi _{3\{\alpha \beta \}}`$ is represented as
$$\varphi _3=\left(\begin{array}{cc}\varphi ^{4/3}& \varphi ^{1/3}\\ \varphi ^{1/3}& \varphi ^{2/3}\end{array}\right)$$
(8)
where the superscript denotes the electric charge of the field. The same applies for $`\varphi _4`$.
Scalar fields which are not SU(2)<sub>L</sub> doublets do not participate in D type Yukawa terms, they however contribute to the mass matrix of the scalar sector and in turn determine the magnitude of the radiatively generated masses of fermions, as we shall see below.
The most general scalar potential of dimension $`4`$ that can be written is
$`V(\varphi _i)\text{ }={\displaystyle \underset{i}{}}\mu _i^2\varphi _i^2+{\displaystyle \underset{i,j}{}}\lambda _{ij}\varphi _i^2\varphi _j^2+\eta _{31}\varphi _1^{}\varphi _3^{}\varphi _3\varphi _1+\stackrel{~}{\eta _{31}}\stackrel{~}{\varphi _1}^{}\varphi _3^{}\varphi _3\stackrel{~}{\varphi _1}`$
$`+\eta _{41}\varphi _1^{}\varphi _4^{}\varphi _4\varphi _1+\stackrel{~}{\eta _{41}}\stackrel{~}{\varphi _1}^{}\varphi _4^{}\varphi _4\stackrel{~}{\varphi _1}+{\displaystyle \underset{\begin{array}{c}ij\\ i,j1,2\end{array}}{}}\eta _{ij}\varphi _i^{}\varphi _j^2+(\rho _1\varphi _5^{}\varphi _6\varphi _2+`$
$`\rho _2\varphi _7^{}\varphi _8\varphi _2+\lambda _1\varphi _5^{}\varphi _1^\alpha \varphi _{3\{\alpha \beta \}}\varphi _1^\beta +\lambda _2\varphi _7^{}\stackrel{~}{\varphi _1}^\alpha \varphi _{3\{\alpha \beta \}}\stackrel{~}{\varphi _1}^\beta +`$
$`\lambda _3Tr(\varphi _3^{}\varphi _4)\varphi _2^2+\lambda _4\varphi _5\varphi _6\varphi _7\varphi _2+\lambda _5\varphi _5\varphi _6^{}\varphi _7^{}\varphi _8+h.c.),`$ (9)
where $`Tr`$ means trace and in $`\varphi _i^2\varphi _i^{}\varphi _i`$ an appropriate contraction of the $`SU(2)_L`$ and $`SU(3)_C`$ indices is understood. The gauge invariance of this potential requires the relation $`\mathrm{\Delta }=2\delta `$.
Now we proceed to describe the mechanism that produces the quark masses. In general we could have contributions of two types as depicted in Fig 1. In the present model however, we have only the diagrams of Fig 2 for the charge $`1/3`$ quark mass matrix elements and similar ones for the charge $`2/3`$ sector (these type of diagrams were first introduced in ma ); in those diagrams of Fig 2 the cross means tree level mixing and the black circle means one loop mixing. The diagrams in Fig 3a y 3b should be added to the matrix elements (1,3) and (3,1), respectively.
In the one loop contribution to the mass matrices for the different quark sectors only the third family of quarks appears in the internal lines. This generates a rank 2 matrix, which once diagonalized gives the physical states at this approximation. Then using these mass eigenstates we compute the next order contribution, obtaining a matrix of rank 3. After the diagonalization of this matrix we get the mass eigenvalues and eigenstates (A quark mass mechanism with some similar features to the one proposed here is given in ma ).
The VEVs of the class I scalar fields are
$`<\varphi _1>={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}0\\ v_1\end{array}\right),`$ $`<\varphi _2>=v_2,`$ (12)
and they achieve the breaking
$$G_{SM}U(1)_X\stackrel{<\varphi _2>}{}G_{SM}\stackrel{<\varphi _1>}{}SU(3)_CU(1)_Q.$$
(13)
The scalar field mixing arises after SSB from the terms in the potential that couple two different class II fields to one of class I. After SSB the mass matrices for the scalar fields of charge $`2/3`$ ($`\varphi _4`$,$`\varphi _3`$,$`\varphi _5`$,$`\varphi _6`$) and $`4/3`$ ($`\varphi _4`$,$`\varphi _3`$,$`\varphi _7`$,$`\varphi _8`$) are written, respectively, as
$`M_{2/3}^2=\left(\begin{array}{cccc}s_4^2& \lambda _3^{}v_2^2& 0& 0\\ \lambda _3v_2^2& s_3^2& \frac{\lambda _1^{}v_1^2}{2}& 0\\ 0& \frac{\lambda _1v_1^2}{2}& u_5^2& \rho _1v_2\\ 0& 0& \rho _1^{}v_2& u_6^2\end{array}\right)`$ and $`M_{4/3}^2=\left(\begin{array}{cccc}t_4^2& \lambda _3^{}v_2^2& 0& 0\\ \lambda _3v_2^2& t_3^2& \frac{\lambda _2v_1^2}{2}& 0\\ 0& \frac{\lambda _2^{}v_1^2}{2}& t_7^2& \rho _2v_2\\ 0& 0& \rho _2^{}v_2& t_8^2\end{array}\right),`$ (22)
where from Eq. (II) $`t_i^2=u_i^2=\mu _i^2+\lambda _{i1}v_1^2+\lambda _{i2}v_2^2`$ and $`s_i^2=t_i^2+\eta _{i1}v_1^2`$.
Notice that due to the scalar mixing in all the loop diagrams of Fig 2 and 3, the divergences in each one of these diagrams cancel as is physically expected, giving rise to finite contributions to the quark mass matrices.
Explicitly, the non vanishing contributions from the diagrams of Fig 2 to the mass terms $`\overline{d}_{iR}d_{jL}\mathrm{\Sigma }_{ij}^{(1)}+h.c.`$ read at one loop
$$\mathrm{\Sigma }_{22}^{(1)}=3m_b^{(0)}\frac{Y_I^2}{16\pi ^2}\underset{k}{}U_{1k}U_{4k}f(M_k,m_b^{(0)}),$$
(23)
$$\mathrm{\Sigma }_{23}^{(1)}=3m_b^{(0)}\frac{Y_I^2}{16\pi ^2}\underset{k}{}U_{2k}U_{4k}f(M_k,m_b^{(0)}),$$
(24)
$$\mathrm{\Sigma }_{32}^{(1)}=3m_b^{(0)}\frac{Y_I^2}{16\pi ^2}\underset{k}{}U_{1k}U_{3k}f(M_k,m_b^{(0)}),$$
(25)
where $`m_b^{(0)}`$ is the tree level contribution to the b quark mass, the 3 is a color factor, U is the orthogonal matrix which diagonalizes the mass matrix of the charge $`2/3`$ scalars,
$$(\varphi _4,\varphi _3,\varphi _5,\varphi _6)^T=U(\sigma _1,\sigma _2,\sigma _3,\sigma _4)^T,$$
where $`\sigma _i`$ are the eigenfields with eigenvalues $`M_i`$, and
$$f(a,b)\frac{1}{a^2b^2}[a^2ln\frac{a^2}{b^2}],$$
which is just a logarithmic contribution when $`a^2b^2`$. The resulting second rank mass matrix at this level is thus
$$M_d^{(1)}=\left(\begin{array}{ccc}0& 0& 0\\ 0& \mathrm{\Sigma }_{22}^{(1)}& \mathrm{\Sigma }_{23}^{(1)}\\ 0& \mathrm{\Sigma }_{32}^{(1)}& m_b^{(0)}\end{array}\right).$$
(26)
At effective two loops we obtain the following expressions:
$$\mathrm{\Sigma }_{11}^{(2)}=3\frac{Y_I^2}{16\pi ^2}\underset{k,i}{}m_i^{(1)}(V_{dL}^{(1)})_{2i}(V_{dR}^{(1)})_{2i}U_{2k}U_{3k}f(M_k,m_i^{(1)}),$$
(27)
$$\mathrm{\Sigma }_{12}^{(2)}=3\frac{Y_I^2}{16\pi ^2}\underset{k,i}{}m_i^{(1)}(V_{dL}^{(1)})_{3i}(V_{dR}^{(1)})_{2i}U_{1k}U_{3k}f(M_k,m_i^{(1)}),$$
(28)
$`\mathrm{\Sigma }_{13}^{(2)}=3{\displaystyle \frac{Y_I^2}{16\pi ^2}}{\displaystyle \underset{k,i}{}}m_i^{(1)}(V_{dL}^{(1)})_{3i}(V_{dR}^{(1)})_{2i}U_{2k}U_{3k}f(M_k,m_i^{(1)})`$ (29)
$`+3{\displaystyle \frac{Y_I^2}{16\pi ^2}}{\displaystyle \underset{k,i}{}}m_i^{(1)}(V_{dL}^{(1)})_{2i}(V_{dR}^{(1)})_{2i}U_{1k}U_{3k}f(M_k,m_i^{(1)}),`$
$$\mathrm{\Sigma }_{21}^{(2)}=3\frac{Y_I^2}{16\pi ^2}\underset{k,i}{}m_i^{(1)}(V_{dL}^{(1)})_{2i}(V_{dR}^{(1)})_{3i}U_{2k}U_{4k}f(M_k,m_i^{(1)}),$$
(30)
$`\mathrm{\Sigma }_{31}^{(2)}=3{\displaystyle \frac{Y_I^2}{16\pi ^2}}{\displaystyle \underset{k,i}{}}m_i^{(1)}(V_{dL}^{(1)})_{2i}(V_{dR}^{(1)})_{3i}U_{2k}U_{3k}f(M_k,m_i^{(1)})`$ (31)
$`+3{\displaystyle \frac{Y_I^2}{16\pi ^2}}{\displaystyle \underset{k,i}{}}m_i^{(1)}(V_{dL}^{(1)})_{2i}(V_{dR}^{(1)})_{2i}U_{2k}U_{4k}f(M_k,m_i^{(1)}),`$
where the k (i) index goes from 1 to 4 (from 2 to 3), $`V_{dL}^{(1)}`$ and $`V_{dR}^{(1)}`$ are the unitary matrices which diagonalize $`M_d^{(1)}`$ of equation (26) and $`m_i^{(1)}`$ are the eigenvalues. Therefore at two loops the mass matrix for d quarks becomes:
$$M_d^{(2)}=\left(\begin{array}{ccc}\mathrm{\Sigma }_{11}^{(2)}& \mathrm{\Sigma }_{12}^{(2)}& \mathrm{\Sigma }_{13}^{(2)}\\ \mathrm{\Sigma }_{21}^{(2)}& m_2^{(1)}& 0\\ \mathrm{\Sigma }_{31}^{(2)}& 0& m_3^{(1)}\end{array}\right).$$
(32)
For the up sector the procedure to obtain the masses is completely analogous. That is, the mass terms for the up sector come from graphs like those in Fig. 2 and 3, but replacing the $`\varphi _4`$,$`\varphi _3`$, $`\varphi _5`$ and $`\varphi _6`$ scalar fields by $`\varphi _4`$,$`\varphi _3`$, $`\varphi _7`$ and $`\varphi _8`$ and the quarks $`d_i`$ by the quarks $`u_i`$.
The CKM matrix takes the form
$$V_{CKM}=(V_{uL}^{(2)}V_{uL}^{(1)})^{}V_{dL}^{(2)}V_{dL}^{(1)},$$
(33)
where the unitary matrices $`V_{uL}^{(1)}`$ and $`V_{uR}^{(1)}`$ diagonalize $`M_u^{(1)}`$, and $`V_{uL}^{(2)}`$ and $`V_{uR}^{(2)}`$ diagonalize $`M_u^{(2)}`$, with an analogous notation used for the down sector.
It is important to mention here that the textures, particularly the zeros in the scalar and quark mass matrices (Eqs. 22 and 32) are not accidental neither imposed; they are just a direct consequence of the mass mechanism that we are introducing and of the gauge symmetry of the model.
## III Numerical Evaluation
### III.1 Experimental values
Since quarks are confined inside hadrons, their masses can not be directly measured. So, the quark mass parameters in the SM Lagrangian depend both on the renormalization scheme adopted to define the theory and on the scale parameter $`\mu `$ where the theory is being tested. In the limit where all quark masses are zero, the SM has an $`SU(3)_LSU(3)_R`$ chiral symmetry which is broken at an scale $`\mathrm{\Lambda }_\chi 1`$GeV. To determine the quark mass values one must use SM perturbation theory at an energy scale $`\mu >>\mathrm{\Lambda }_\chi `$ where non perturbative effects are negligible.
For illustration, the allowed ranges of quark massespdg in the modified minimal subtraction scheme $`(\overline{MS})`$ arenari :
$`m_u(1.GeV)`$ $`=`$ $`26.8MeV.`$
$`m_d(1.GeV)`$ $`=`$ $`412MeV.`$
$`m_s(1.GeV)`$ $`=`$ $`81230MeV.`$
$`m_c(m_c)`$ $`=`$ $`1.11.4GeV.`$
$`m_b(m_b)`$ $`=`$ $`4.14.4GeV.`$
$`m_t(Exp)`$ $`=`$ $`173.8\pm 5.2GeV`$
To get the relative magnitude of different quark masses in a meaningful way, one has to describe all quark masses in the same scheme and at the same scale. In our analysis we are calculating the quark masses at an energy scale $`\mu _m`$ such that $`M_Z<\mu _m<M_Xv_2`$, where $`M_X`$ is the mass scale where $`U(1)_X`$ is spontaneously broken. Since in our model there is no mixing between the Standard Model Z boson and its horizontal counterpart, we can have $`v_2`$ as low as the electroweak breaking scale. For simplicity, let us assume that our calculations are meaningful at the electroweak breaking scale and from the former values for the quark masses let us calculate, in the $`\overline{MS}`$ scheme, the quark masses at the $`m_t`$ scalefusaoka and at the $`M_Z`$ scalefritzsch . Those values calculated in the references cited, are presented in tables 3 and 4 respectively.
On the other hand, the CKM matrix elements are not ill defined and they can be directly measured from the charged weak current in the SM. For simplicity we assume that they are real, and as discussed in Ref.fusaoka , they are almost constant in the interval $`M_Z<\mu <`$ a few TeV. Their current experimental valuepdg are given in the Tables 3 and 4.
### III.2 Evaluation of the parameters
In order to test the model using the least possible number of free parameters, let us write the scalar mass matrices in the following form:
$`M_{2/3}^2=\left(\begin{array}{cccc}a_+& b& 0& 0\\ b& a_+& c_+& 0\\ 0& c_+& a_+& d_+\\ 0& 0& d_+& a_+\end{array}\right)`$ and $`M_{4/3}^2=\left(\begin{array}{cccc}a_{}& b& 0& 0\\ b& a_{}& c_{}& 0\\ 0& c_{}& a_{}& d_{}\\ 0& 0& d_{}& a_{}\end{array}\right).`$ (42)
Using the central value of the CKM elements in the PDG bookpdg and the central values of the six quark masses at the top mass scalefusaoka , we build the $`\chi ^2`$ function in the ten parameter space defined by $`(a_+,a_{},b,c_+,c_{},d_+,d_{},Y_I,m_b^{(0)},m_t^{(0)})`$, where $`m_b^{(0)}`$ and $`m_t^{(0)}`$ are the tree level quark masses for the bottom and top quarks respectively. Expressions for the eigenvectors and eigenvalues of the mass matrices involved in the numerical evaluation were obtained using MATHEMATICA, and the $`\chi ^2`$ function was minimized using MINUIT from the CERNLIB packagesminuit ; both Monte Carlo and standard routines were used in the minimization process. The tree level masses of the top ($`m_t^{(0)}`$) and bottom ($`m_b^{(0)}`$) quarks were restricted to be around the central values $`\pm `$ 10 % in order to assure consistency with the assumption that radiative corrections are small. The $`\chi ^2`$ function presents an even symmetry with respect to 5 of the parameters of the matrices in Eq. (42); we find that there are 32 parameter domains where the $`\chi ^2`$ function takes small values. For the extremal points this even symmetry is not an exact one, but all the zones have a Yukawa constant of the order of 10 and give masses and CKM matrix elements in good agreement with the available experimental values. The numerical results on one of the 32 minima of the ten parameter space are shown in Table 5. We use those values (which minimize $`\chi ^2`$) to calculate, in the context of our model, the fifteen predictions for $`m_q(m_t)`$ for $`q=u,d,c,s,t,b`$ and $`(CKM)_{ij}`$ for $`i,j=1,2,3`$. The numerical results are shown in Table 3.
For the sake of comparison, we repeat the same calculations but now using the central values of the six quark masses at the $`M_Z`$ scalefritzsch . The numerical results are shown in Table 4.
Let us make two comments: First, the values for the parameters in the scalar field square mass matrices are of order $`10^{17}(\text{MeV})^2`$ (see Table 5), so, the scalar physical masses are of order $`10^3`$ TeV. Second, the rounding errors allow us to take safely up to five significative figures in the masses and in the CKM matrix elements
As can be seen from Tables 3 and 4, even under the assumption that the CKM matrix elements are real, the numerical values are in good agreement with the allowed experimental results.
## IV CONCLUSIONS
By introducing a $`U(1)_X`$ gauge flavor symmetry and enlarging the scalar sector, we have presented a mechanism and an explicit model able to generate radiatively the hierarchical spectrum of quarks masses and CKM mixing angles. The horizontal charge assignment to particles is such that we do not need to go beyond the known three generations of quarks and leptons. Also, at tree level only the t and b quarks get masses. To generate radiatively the masses for the light families we have introduced some new exotic scalars. All of these new scalars are charged and color non-singlets, so they can not get VEV as is required in the loop graphs.
Our numerical results are presented in Tables 3 and 4. Even though we are guessing the $`U(1)_X`$ mass scale, the two sets of results do not differ by much and they agree fairly well with the experimental values, meaning that the mass scale associated with the horizontal symmetry may be in the range $`100GeV<M_X<1.0TeV`$. A closer look to our analysis shows that we are translating the quark mass hierarchy to the quotient $`v_1/v_2`$ which is the hierarchy between the electroweak mass scale and the Horizontal $`U(1)_X`$ mass scale. In this way we demonstrate the viability that new physics at the electroweak mass scale, or just above it, may help to explain the long-lasting puzzle of the enormous range of quark masses and mixing angles.
Since quarks carry baryon number $`B=1/3`$, the color sextet scalars we have introduced must have $`B=2/3`$ (the scalar singlets $`\varphi _1`$ and $`\varphi _2`$ have B=0); in this way $`L_{Y_M}`$ is not only $`U(1)_X`$ invariant but conserves color and baryon number as well. On the other hand, $`V(\varphi )`$ does not conserve baryon number; as a matter of fact, the term $`\lambda _4\varphi _5\varphi _6\varphi _7\varphi _2`$ violates baryon number by two units ($`\mathrm{\Delta }B=\pm 2`$) and could induce neutron-antineutron oscillations. A roughly estimate of such oscillations shows that they are proportional to $`v_2(M_{\varphi _5}^2M_{\varphi _6}^2M_{\varphi _7}^2)^110^{19}GeV`$ which is negligible in principle. Any way, in the worse of the situations, since the offending term does not enter on the mass matrix for the Higgs scalars (it is just there), it may be removed in more realistic models by the introduction of a discrete symmetry.
Our results are encouraging; even under the assumption that the CKM matrix is real, and without knowing exactly the $`U(1)_X`$ mass scale, the numerical predictions are in the ballpark, implying also a value of order 10 for the Yukawa coupling $`Y_I`$, and masses for the exotic scalars being of order $`10^3`$ TeV. Our model presents thus a clear mechanism able to explain the mass hierarchy and mixing of the quarks.
Finally let us mention that in the work presented here, the Higgs scalar used to produce the SSB of the SM gauge group down to $`SU(3)_C\times U(1)_Q`$ has zero horizontal charge, and as a consequence the standard $`Z`$ boson does not mix with the horizontal counterpart. However, due to recent interest langacker on the phenomenology of a $`Z^{}`$, it is worth to study the possibility of allowing this mixing in future work.
###### Acknowledgements.
This work was partially supported by Conacyt in Mexico and Colciencias and BID in Colombia. One of us (A. Z.) acknowledges the hospitality of the theory group at CERN and useful conversations with Marcela Carena and Jean Pestieau.
## V BIBLIOGRAPHY
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# The Surroundings of Disturbed, Active Galaxies
## 1 Introduction
The neighborhoods of the most active ULIRG’s, Mark273, Mark231, Arp 220 and NGC6240 plus the radio galaxy 3C31$`/`$NGC383 are examined for associated objects.
## 2 Markarian 273
A strong jet emerges S from the disturbed central regions of this ULIRG. Fig.1 shows that both Mark 273 (z = .038) and the compact object Mark 273x (z = .458) are strong X-ray sources. The deep optical image in Fig. 2 illustrates the conclusion by Xia et al. (1999) that ”Mark 273x is at the tip of the plume” reaching NE from Mark 273.
Fig.3 addresses the question of whether there is any X-ray connection between Mark 273 and the AGN to the NE. The outer isophotes do not appear as rounded as one would expect from accidental adjacency of unrelated, symmetrical sources. (See for example the contouring of two adjacent, unrelated images in Arp 1998a, Fig. 1-12). In addition the bright central regions of Mark 273 are conspicuosly extended toward Mark 273x.
In fact what we see is a compact optical object of R = 19.6 mag. and redshift z = .458 connected back to a disturbed Seyfert 2 of V = 14.9 mag. by both optical and X-ray emitting material. Mark 273x is the same kind of quasar-like object which was physically associated with somewhat brighter Seyfert galaxies at a 7.5 sigma level (Radecke 1997; Arp 1997). Also there is a recent precedent for a quasar of z = 2.15 at the tip of an X-ray/optical filament emerging along the minor axis of the active galaxy NGC3628 (Flesch and Arp 1999).
The luminosity of the quasar-like Mark 273x as derived by Xia et al. (1998) when they believed it to be connected to Mark 273 was M<sub>R</sub> =-17.6 mag. But we will argue later that Mark 273 itself is quasar-like in that it has an intrinsic redshift and should be moved of the order of 5 magnitudes (in modulus) closer.
Smoothing and contouring the lower resolution, PSPC X-ray measures produces Fig. 4. It is seen that the outer contours of Mark 273x verify the elongation to the NE, along a line back toward the central galaxy. Regardless of whether this x-ray connection is due to a jet or to some merger related activity it does establish the fact that the two objects are at the same distance. There is, however a clear X-ray jet visible in the PSPC measures. It comes out of Mark 273 to the SE and is delineated even in the 3,4 and 5 sigma contours (also more weakly in the opposite direction).
There is no reason not to have jets coming out of active galaxies in different directions. In fact there is considerable evidence for lines of X-ray sources emerging from Seyfert and other active galaxies (Arp 1996;1997). Why those apparent ejections take place in roughly orthogonal directions, as here in Mark 273, there seems to be no ready explanation, but it does seem to be an empirical result. As we shall see there is evidence for X-ray material and objects coming out further along these same principal directions.
As for the relation of Mark 273x to Mark 273, the ejection of higher redshift quasars and quasar-like objects has been argued for more than 30 years (see Arp 1983; 1998a). But the interesting additional information given by these disturbed, active galaxies is that when quasars do not come out along relatively unimpeded along the minor axis, that they then interact with the material of the ejecting galaxy and disrupt it, giving rise to disturbed morphologies, entrainments and fragmentation (Arp 1999a).
## 3 Markarian 273 Radio Features
It is most interesting to consider the information in other wavelengths, not only to confirm the features we have already seen but also to try to understand the physcial processes involved in the ejections, if that is the mechanism for producing these associations.
In Fig.5 we see the map of the Mark 273 region in continuum radio wavelengths. The most prominent feature is a broad ejection coming out NW and SE, agreeing closely with the direction of the X-ray jet. (Close to the nucleus, agreement between X-ray and radio jets is characteristic of very active radio galaxies like Virgo A and Centaurus A). In the counter jet direction the radio extension is longer but weaker. Overall, there seems to be confirmation that both X-ray and radio emitting material are coming out of Mark 273 at position angles of about 130 and 310 degrees.
The next most conspicuous feature in Fig. 5 is the very strong radio source about 5.5 arcmin SW of Mark 273. This radio source is a close double (about 1.1 arcsec separation and p.a. about 20 deg when viewed on the VLA, high resolution mode, FIRST). Doubleness of companion radio and X-ray sources will be noted as we explore other active galaxy neighborhoods. But of considerable significance for the ejection hypothesis for Mark 273x is the fact that it is itself a strong radio source, and that it is aligned, as exactly as can be measured, across Mark 273 with the strong radio source to the SW.
Fainter radio extensions than contoured in the preceding Figure are shown here in Fig. 6. They extend both NW and SE generally along the direction of the radio and X-ray jets. (The arrow in Fig. 6 marks the X-ray jet pictured in Fig. 4). Surprisingly there are two low surface brightness radio filaments extending from Mark 273 and ending on two very strong X-ray sources (in Fig. 7 they are marked as Nos. 21N and 38 and are at 6.2 and 6.6 arcmin respectively from Mark 273). They are marked with x’s in Fig. 6. One is a quasar of z = .941 and the other is a blue stellar object (BSO) with a similar apparent magnitude, of 18.1. The latter will almost certainly turn out to be a second quasar in a pair. In the deep red exposure of Fig. 2 there are twin, low surface brightness optical filaments coming out of the S end of Mark 273. (Seen better in deeper image in J.Hibbard’s home page.) They appear to lead into the radio filaments which are observed in Fig. 6. In summary, it appears that in all the wavelengths that can be checked, radio, optical and X-ray, there is a connection of these two strong X-ray sources back to the center of Mark 273.
Because of the apparent quasar nature and consequent high redshift of these two strong X-ray sources it is customary to ask whether they could be accidentally projected background sources. We can answer this suggestion by noting that the strength of these X-ray sources: Mark 273x = 8.3 cts/ks, z = .941 quasar = 37.7 cts/ks and 18.1 BSO = 21.2 cts/ks indicates that they would have a chance of being background sources falling as close as they do to Mark 273 of p = .003, .01 and .04. Altogether getting just these three objects this close to an arbitrary active galaxy would only be about one chance in a million. Xia et al. (1999) note in passing that it is interesting that ”… the X-ray companions of the three nearest ULIG’s (Arp 220, Mrk 273 and Mrk 231) are all background sources…”.
## 4 X-ray field around Mark 273
Since X-ray sources and quasars have been associated with Seyfert galaxies out to distances of nearly a degree it is of interest to see what the full PSPC field of ROSAT around mark 273 looks like.
The strong sources that we have been discussing are seen close to the center of Fig. 7. The two probable quasars are seen to be exceptionally strong X-ray sources at 21 and 38 cts/ks, much stronger even than Mark273 and 273x which have 13 and 8 cts/ks. As we go off-axis the point spread of the images enlarges so only the brighter sources are registered. It is clear, however, that sources of all brightness are distributed in an elongated pattern, roughly NE to SW through Mark 273. This line is close to the direction of Mark 273x and the strong, double radio source pictured in Fig. 5 (p.a. $``$ 40 deg.).
Of course the radio and X-ray ejections which lead toward the two strong X-ray sources (21N and 38) are in an almost orthogonal direction to this (p.a $``$ 140 deg.). They are suggested to be connected with the the strong X-ray sources 16 and 30 which lie further out in this direction. As Table 1 shows, both the 16 and 30 cts/ks X-ray sources are optically identified as blue stellar objects (BSO’s) and therefore are most probably a quasar pair like so many of the X-ray sources paired across Seyfert galaxies.
Ejection in more than one direction from active galaxies is commonly observed (see Arp 1996; 1997a; Komossa and Schulz 1998). The most important aspect of Fig. 7, however, is that numerous X-ray sources in this field are aligned preferentially across Mark 273, confirming that they are predominantly not background.
TABLE 1
Bright X-ray sources in Mark 273 field
| Cts | R.A. (2000) | Dec. | E | O-E | ID | Remarks |
| --- | --- | --- | --- | --- | --- | --- |
| 49 | 13 41 36.2 | 55 14 37 | 16.93 | 1.99 | c.g. | |
| 26 | 13 42 06.3 | 56 39 15 | 16.52 | 1.20 | BSO | 47”off pos |
| 16 | 13 42 21.6 | 56 14 50 | 19.88 | .73 | BSO | |
| 18 | 13 44 13.5 | 56 30 34 | 18.30 | .34 | BSO | |
| 21S | 13 44 57.3 | 55 28 22 | 18.9 | .85 | ? | extend? |
| 38 | 13 44 47.3 | 55 46 56 | 17.76 | .77 | QSO | z = .941 |
| 13 | 13 44 42.1 | 55 53 13 | 14.9\* | .77\* | S2 | Mark 273 |
| 8 | 13 44 47.1 | 55 54 10 | 19.25 | 1.78 | c.g. | Mark 273x |
| 21N | 13 45 12.0 | 55 47 59 | 18.1 | .57 | BSO |
| 29 | 13 45 33.4 | 55 24 06 | 19.26 | .66 | BSO |
| 8 | 13 46 06.7 | 56 04 29 | 19.08 | .28 | BSO |
| 30 | 13 46 38.9 | 55 27 09 | 19.44 | .32 | BSO |
| 48 | 13 46 59.1 | 56 07 04 | 18.21 | 1.49 | BSO |
* B and B-V magnitudes
## 5 Infrared sources aligned with Mark 273
In order to investigate the Mark 273 surroundings in longer wavelengths the Simbad on-line catalog was consulted. It was noticed immediately that there was a conspicuous string of what turned out to be infrared sources coming generally SW from the ULIRG galaxy. Fig. 8 shows a plot of nearly the same sized area as pictured in the preceding PSPC X-ray map. Only infrared point sources with 100 micron fluxes greater than .5 Jy have been plotted. The sources are as listed in the IRAS faint source catalog of 173,044 sources. It is clear that there is a general line of deep infrared sources coming from the center of the Mark 273 region in roughly the same direction (p.a. = 235 deg.) as the Mark 273x/radio double and X-ray material ( 220 deg.).
## 6 Radio sources around Mark 273
Among the three closest radio sources to Mark 273 at 4.85 GHz (Becker et al. 1991), there is a 42 mJy extended source 18 arcmin NE and a 27 mJy source 35 arcmin SW. A little further in this direction (46 arcmin) there is a radio source of 31 mJy which appears to coincide with the strong X-ray source of 49 cts/ks mapped in Fig. 7. This same source is identified in Table 1 as a compact galaxy. These outer radio sources lie closely along the NE-SW direction of the inner radio sources Mark 273, 273x and dbl source SW (see Fig. 5).
Radio ejections from active galaxy nuclei are well accepted. X-ray ejections from Mark 273 have been shown here. Physical association of X-ray sources has been been shown in the past - for example six Seyferts, including Mark 273, were shown to have a factor 8 higher X-ray source density around them (Turner et al. 1993). A sample of 24 of the brightest Seyferts were shown to have paired X-ray sources across them, the great majority of which must be quasars or quasar-like (Radecke 1997; Arp 1997). The alignment of radio sources and X-ray quasars from these conspicuously ejecting ULIRG’s would suggest, therefore, that material in various different forms is ejected.
We now see that for Mark 273 the infrared sources in Fig. 8 are distributed around a principal line of ejection. It raises the question of how they are related to the ejection process and the other ejected material.
## 7 Markarian 231
Fig. 9 shows that flanking Mark 231 there is a double radio source about 7.5 arcmin W and another double 2.2 arcmin to the E. (The latter has a faint companion to the NE making it actually a triple). Fig. 9b aligned directly below Fig. 9, and to the same scale, shows that radio material is being ejected toward the E from Mark 231 and that presumably the sources to the W represent counter ejection in the other direction. The X-ray source to the east (left) is optically identified with a BSO and extended in the NE direction as are also its radio contours.
It is important to note that there is a faint but definite X-ray jet extending to the N from Mark 231. This is the same direction as the radio extensions to the N and some small spots of radio emission as shown in Fig. 9b. Closely to the S of Mark 231 there is en extended radio source (Condon and Broderick 1998 ) and in the inner 60 milliarcsec there is a strong N-S triple source (Ulvestad et al. 1999; Taylor et al. 1999). All of this is evidence for ejection and outflow of material in this direction, a direction which only shifts by 5 deg. in p.a. from the innermost to the outermost aligned sources (over 3 deg. distant for the infrared sources that we shall see later).
As we have seen in other cases here there is evidence for ejection in more than one direction from an active galaxy. In the case of Mark 231 the strong double radio source 7.6 arcmin slightly S of W is an obvious counter ejection to the radio sources emerging from Mark 231 on the E. The W source is conspicuous in that that it is a close double. (The sources on the E are also closely multiple).
The most striking object of all, however, is the double X-ray source which is aligned back through the double radio source to the nucleus of the galaxy. An obvious interpretation is that as the X-ray source travelled outward from Mark 231 it had its radio plasma stripped from it. This could occur in transit of a homogeneous medium surrounding Mark 231, in conjunction with a discrete event of generation of radio plasma. Or it could be simply an encounter with a discrete density enhancement in the medium.
It is apparent in Fig. 9 that at a slightly greater angle to the SW, there is another X-ray source with a radio source trailing at about the same distance behind it. In this and the previous case it is noticeable that the radio sources are slightly above the line from the X-ray source back to Mark 231. This suggests that the medium which which stripped the plasma is rotating relatively slowly counterclockwise.
Proof of such a mechanism would illuminate several puzzles about ejection of matter from active galaxies. To this end we point to Fig. 9c which shows that the double X-ray source is identified with a double quasar. (The brighter X-ray source is catalogued as z = 1.272 and the optical image 52 arcsec E is appreciably bluer implying a highly probable QSO of redshift to be determined). The properties which support the relation of the double X-ray source to the double radio source are then: 1) The relation of the stronger components of the double to the weaker. 2) The similar orientation of the radio double, p.a. = 65 deg., and the X-ray double, p.a. = 75 deg. The separation of 52 arcsec of the images in the X-ray and optical compared to the 24 arcsec separation of the radio images would then imply a continuing separation of the X-ray images since the event of the stripping of the radio plasma.
Thanks to the detailed observations of Alan Stockton (astro-ph/9801056) we can support this picture with observations of the radio galaxy 3C212. Fig. 10 shows that the optical objects f and g have passed out beyond the radio material. The peculiar, identical shape of the radio contours and the optical material leaves no doubt that the radio plasma has been removed suddenly from the optical objects as they progressed away from the ejecting galaxy. It is also revealing to note that the redshift of the material in the galaxy is z = 1.05 whereas the optical object, f, has z = .93. (Ejection velocities relative to the active objects, and presumably their surrounding medium, have been shown to be of the order of .1c. (Arp 1998a)).
If this suggested mechanism of plasma stripping in fact operates, it furnishes us with a general explanation of why the radio emission and the X-ray/optical objects can be generally but not exactly in the same positions. In fact we will see later that in the famous radio galaxy, 3C31/NGC383, that strong radio filaments trail toward the position of some high redshift galaxies and an even higher redshift quasar. But the radio track is rotated somewhat from the X-ray/optical objects.
If the plasma generating events in quasars are intermittent it might help us understand why, among similar appearing quasars, some are radio sources and some are not. The higher percentage of quasars identified with X-ray sources than with radio sources could be explained by the fact that X-rays and optical synchrotron radiation would decay relatively quickly and therefore be much closer to the denser core of the quasar itself. Periodic gas ejection events could also supply multiple clouds to account for those quasars with multiple absorption line systems.
Markarian 231 turns out to be particulary well suited to test these suggestions because the lines of radio and X-ray ejection are so well marked. Proceeding from the milliarcsec interior they appear fairly constant in direction out to almost a degree on the sky. Perhaps this is due to the lack of rotational or precessional angular momentum imparted by these particular ULIRG’s, as suggested by their irregular morphologies.
## 8 Surroundings of Mark 231 in the Infrared
In the infrared field of Fig. 11 the long arrow shows the direction and distance of the 136 cts/ks X-ray source, 3C277.1, plus the spread of the radio ejections around this direction. The short arrow shows the extent of the radio and X-ray ejection N from Mark 231. The IR source is in the direction of the radio and X-ray ejection at p.a. = 243 deg. (The probable identification of this source is Y Uma but there is a 1RXS, X-ray BSO within 23 arcmin.) There are also three sources aligned to the N at about p.a. = -10 deg., which is close to the inner radio and X-ray ejection to the N at about p.a. = -5 deg.
If these cold sources (they are perhaps even more prominent at 100 microns) are ejected from Mark 231 they confirm the result obtained in Fig. 8 for Mark 273. The large angular separation of around 3 degrees on the sky, however would imply that Mark 231 was much closer than the conventional redshift distance. Cold infrared sources would also pose again the question of whether they were proto quasars or by-products of the ejection.
## 9 The X-ray field around Mark 231
Fig. 12 shows the distribution of the brightest X-ray sources around Mark 231 (14.3 cts/ks). This is the same elongation of radio sources and X-ray sources as discussed in connection with Fig. 9. The strong, 136 cts/ks X-ray source which is 36 arcmin SW of Mark231 is optically identified with a relatively bright BSO at E = 16.5 mag. as shown in Table 2. There is a moderately strong X-ray source, 22.3 cts/ks at 50 arcmin to the NE of Mark 231 (not registered in Fig. 12) which is identified with a bright BSO of E = 17.4 mag. This is the most promising pair of AGN’s to be identified with the major ejection event of which we see the traces in the interior. The 17.5 cts/ks X-ray source is the candidate double quasar shown in Fig.9.
TABLE 2
Bright X-ray sources in Mark 231 field
| Cts | R.A. (2000) | Dec. | E | O-E | ID | Remarks |
| --- | --- | --- | --- | --- | --- | --- |
| 136 | 12 52 26.2 | 56 34 20 | 16.47 | 1.17 | BSO | 3C277.1 |
| 15.4 | 12 54 49.1 | 57 04 52 | 19.72 | .47 | BSO | |
| 16.1 | 12 54 51.5 | 56 44 30 | 17.03 | 1.43 | Bcg | |
| 17.5 | 12 54 56.7 | 56 49 42 | 19.21 | 1.53 | QSO | z = 1.272 |
| | 12 55 02.9 | 56 49 54 | 19.24 | .65 | BSO | see Fig. 11 |
| 8.1 | 12 55 24.7 | 56 56 14 | 19.53 | .23 | BSO | |
| 4.9 | 12 55 28.2 | 56 46 40 | - | - | - | no candidate |
| 14.3 | 12 56 14.2 | 56 52 25 | 13.84\* | .84\* | S1 | Mark 231 z = .041 |
| 4.2 | 12 56 30.8 | 56 52 19 | 19.06 | .44 | BSO | |
| 4.8 | 12 56 48.3 | 57 03 45 | 19.77 | .63 | BSO | |
| 22.3 | 13 00 33.5 | 57 28 35 | 17.35 | .62 | BSO | |
| 78 | 13 00 43.3 | 56 21 28 | 19 | blue | BSO | 30”S of brt star |
| 117 | 13 00 52.1 | 56 41 05 | 16.27 | .88 | BSO | possibly one or |
| | 13 00 54.5 | 56 41 11 | 19.03 | 1.12 | BSO | both are QSO’s |
* V and B-V
Optical identifications will be discussed after Arp 220.
## 10 Surroundings of Arp 220
This ULIRG displays a chaotic morphology with strong dust absorption. Fig. 14 shows that neutral hydrogen extends from the active galaxy down to companions situated about 2 arcmin to the SW (Hibbard, Vacca & Yun, 2000).
Fig. 15 shows that the brighter X-ray sources form a N-S string in the shape of a shallow ”S”. These X-ray sources seem to be ejected oppositely from the ULIRG with a small amount of rotation. The material which formed the galaxies to the SW has apparently been ejected along this same path but seems to have encountered resistance which has slowed its exit and allowed the development into fairly normal galaxies very close to the ejecting ULIRG. In Fig. 14 the stream of hydrogen drawn out of the parent galaxy and ending exactly on the most active companion (an X-ray and radio source) is evidence for the interaction which braked the normal escape of this material.
Redshifts of ejected compact material are, of course, initially high. We have seen ejected quasars and quasar-like objects in the preceding examples of the active ULIRG’s, Mark 273 and Mark 231. Similar candidates are apparent around Arp 220. Much previous evidence has indicated that these intrinsic redshifts decline as the objects evolve into more normal galaxies ( Arp 1998a,b). The only question is where they evolve, close to the galaxy due to interaction, or further out if the proto galaxy came out, for example, along the minor axis of an ordered galaxy. Empirically there has been evidence starting in 1970 that companion galaxies have systematically higher redshifts than their parent galaxies and this could only be reasonably accounted for by time dependent, intrinsic redshifts associated with their more recent creation.
We will see another example, almost identical to the companions just discussed for Arp 220, when we later consider briefly the active radio source 3C31/NGC383.
## 11 Secondary ejection from the z = .09 galaxies
In order to measure the redshifts of the galaxies immediately SW of Arp 220 Ohyama et al. (1999) placed a one arcsec wide slit at an angle that passed through the three largest companions. As Fig. 16 shows, this slit serendipitously intercepted 4 much fainter galaxies on a line SE of the brightest companion. The redshifts of three of these objects could be measured from their emission lines. They turned out to be z = .528, .529 and .523.
Pursuant to conventional practice the investigators concluded that all of these objects were accidental projections of objects at different distances. The picture was that of galaxies of z = .09 at 5x greater distance than Arp 220, with the higher z objects at a further 6x greater distance, all accidentally aligned closely in the same line of sight. However, if one computes the density of just the three faint objects in the 1 x 16.5 arcsec area of the sky in which they were found, one comes up with the astonishing figure of more than 2 million per sq. deg. Out to a radius of 37 arcsec from the galaxy B/C there should be over 800 such objects unless they are preferentially distributed along just the alignment of the spectrograph slit back to the large galaxy. It is instead attractive to consider that these z = .5 objects were ejected from the active X-ray and radio galaxy at z = .09. They themselves are emission line objects which makes the chance of their being accidental background galaxies less (Arp 1982,p.62).
The brighter galaxies at z =.09, in turn, are much too close to Arp 220 to be accidental. (At their apparent brightness of $`m_R`$ = $``$ 15.5 mag. about .03 galaxies would be expected this close to Arp 220 and, of course, very few would be X-ray sources). Moreover these otherwise morphologically normal galaxies would be of supposed quasar luminosities if placed at their redshift distances (Ohyama et al. 1999 Table 2). Overall, considering the HI and X-ray connections, it would be difficult to avoid interpreting these z = .09 galaxies as having been ejected from Arp 220 and then having ejected even higher redshift objects. Empirically it seems to be another case of associated companion galaxies of various higher redshifts (Arp 1998a).
## 12 The larger neighborhood of Arp 220
Fig. 17 shows a larger field around Arp 220 in X-rays. It is clear that the outermost X-ray sources, Nos. 1 and 10, continue the ”S” shape from the more interior regions pictured in Fig. 15. It is important to note that sources Nos. 2 and 9 (RSO and BSO in Fig. 15) are exactly aligned and almost exactly spaced across the central ULIRG. The outer pair, Nos. 1 and 10 in Fig. 17, are also exactly aligned and almost exactly spaced.
The inner and outer pairs of X-ray sources therefore support the picture of ejection with a slight counterclockwise rotation. The fainter sources interior to the BSO source in Fig. 15 lead directly in along this spiral path, through the companions at z = .09, and then into the ULIRG. The overall picture is then one of material ejected from the active galaxy some of which develops into smaller objects of variously higher redshift. We next look at the outermost regions surrounding Arp 220 in an attempt to define over how large an area on the sky the association extends. To do this we process the full field of the archived PSPC observations as was done earlier for Mark 273 and Mark 231. The results are shown in Fig. 18.
The most striking aspect of Fig. 18 is that the X-ray sources, particularly the faint ones, are distributed in long extensions to the NW and SE. The inner ”S” shape has turned over into an alignment across Arp 220 with a relatively small rotation of the inner sources, the same as observed in Figs. 15 and 17. One can see the slightly smoothed background in Fig. 17 extending away on either side in this large distribution of X-ray radiation.
Of considerable interest is the radio source 3C321. It is optically identified as a S2 galaxy of V = 16.0 mag. Although one of the more intense radio galaxies in the sky it was not detected in X-rays by the Einstein IPC (Fabbiano et al. 1984). In the ROSAT PSPC observation reduced here it is clearly identifiable but too near the edge of the field (56 arcmin) to obtain a reliable flux measure. Note that 3C321 at z = .096 could represent material ejected unimpeded at the same epoch as the z = .090 galaxies which were stopped close to the SW edge of Arp 220.
As to the radio map of 3C321, two strong lobes are well aligned back toward Arp 220. Since there is no ready explanation for such an alignment one would ordinarily assume it was accidental. We have seen, however, near alignments in double radio sources near Mark 273, Mark 231 and along the X-ray ejection axis from NGC2639 and NGC5985. (The latter two are not discussed here).
That 3C321 is physically associated with the apparent ejection alignment from Arp 220 is supported by the map in Fig. 19. There all the bright radio sources and all the bright infrared point sources have been plotted for a degree around Arp 220. The line through the ULIRG comprises 6 out of the 8 radio sources in the area and 5 out of the 7 IRAS sources. Moreover that line coincides closely with the line of X-ray sources, as can be seen by referring to Fig. 18 directly above it.
## 13 Bright X-ray sources in the Arp 220 field
It is very important to note the optical identifications of the X-ray sources which make up this extended spiral alignment in the Arp220 field. Table 3 shows that the strong source 20.3N is identified with a blue compact galaxy of relatively bright apparent magnitude. The other strong source, 20.3S is identified with a blue stellar object which is also unusually bright in apparent magnitude. These two strong sources form a conspicuous pair across the active central galaxy.
TABLE 3
Bright X-ray sources in the Arp 220 field
| Cts | R.A. (2000) | Dec. | E | O-E | ID | Remarks |
| --- | --- | --- | --- | --- | --- | --- |
| 20.3N | 15 33 54.7 | 23 56 15 | 16.34 | .93 | c.g. | BSO’s near |
| 5.0 | 15 34 07.7 | 23 29 39 | 18.28 | 1.19 | BSO | |
| 8.9 | 15 34 51.0 | 23 46 05 | 20.65 | .73 | BSO | possibly |
| ” | 15 34 51.8 | 23 46 03 | 21.59 | $`<`$1.59 | BSO | dbl QSO |
| 9.6 | 15 34 55.0 | 23 28 43 | - | - | gal | SW plume |
| 7.3 | 15 34 57.3 | 23 20 12 | 13.88\* | .88\* | Sey | Arp 220 |
| 2.8 | 15 35 02.3 | 23 15 29 | 15.93 | 2.98 | RSO | |
| 6.9 | 15 35 06.1 | 23 36 56 | 19.61 | 2.18 | RSO | |
| 4.2 | 15 32 08.2 | 23 28 21 | 17.27 | .99 | BSO | |
| 9.4 | 15 37 09.6 | 23 28 36 | 19.88 | 1.22 | BSO: | |
| 20.3S | 15 37 14.5 | 23 00 40 | 17.74 | .76 | BSO | |
* V and B-V
As experience has shown, the outer members of ejected pairs tend to be brighter, lower redshift, transitions between quasars and galaxies. (See particularly Chu et al. 1998; Arp 1998b; 1999). It would therefore be predicted that 20.3N would spectroscopically turn out to be an active, moderately high redshift, compact galaxy and 20.3S a moderately low red shift quasar. The inner quasar candidates would be predicted to turn out to be fainter, higher redshift quasars.
In this respect it is interesting to note that in Table 3 the source at 6.9 cts/ks is identified with a very red optical object. A ROSAT HRI position agrees with the PSPC in identifying a very faint optical candidate. Together with its redness it would suggest that the quasar candidate was highly dust obscured. This in turn would be highly interesting because it would suggest that the famously heavy dust absorption in the central ULIRG extended to at least 8 arcmin angular distance from Arp 220.
Of course there are a number of further quasar candidates identified in Table 3. A line consisting of 9.4, 4.2 and 5.0 cts/ks sources goes roughly E-W. This is consonant with ejection in more than one direction as observed in the cases of Mark 273 and Mark 231 and other cases mentioned earlier. The pairs of quasar candidates indentified in Tables 1 through 3 should be studied with the aim of obtaining further empirical data concerning the speed, orientation and evolutionary behavior of the ejected matter.
## 14 Soft X-ray absorption from Arp 220
When a column of hydrogen ($`N_H`$) containing X-ray absorbing metals intervenes, the soft X-rays will be absorbed and the hardness ratio will approach HR1 = 1. As Fig. 20 shows, the sources around Arp 220 are conspicuously shifted toward this high hardness ratio. The dashed line represents a control field around a bright star, HR8905, which is at a comparable galactic latitude. The estimated excess in hardness ratio from Fig. 20 is about 0.4 in HR1. This translates into a difference in visual absorption of $`A_V`$ = .22 mag.
In the detailed discussion in Appendix A we will argue that the actual absorption is much greater than this minimum value.This means that in addition to absorption from our own galaxy at this galactic latitude there is extra reddening and absorption from the material in the environs of the ULIRG. The surprising conclusion is that the absorption extends out at least as far as 20 arc minutes and probably beyond 30 arcmin radius! Thus the gaseous component of the system extends over a degree in diameter on the sky. This angle is of the order subtended on the sky by larger Local Group Galaxies such as M33!
We have seen individual, bright X-ray sources associated out to a diameter approaching 2 degrees around Arp 220. But they may not be bound. The absorbing gas, however, is much more likely to be travelling at less than the escape velocity and this raises the question of what the system will evolve into - a large low redshift galaxy or a low redshift cluster of glaxies? In either case the empirical evidence on the apparent diameter requires the system to be much closer than the redshift distance of the central, ejecting galaxy.
The same situation applies to the ULIRG’s Mark 273 and Mark 231. Although we do not show the HR1 off - axis plots here, it is readily apparent that they have 5 -9 sources which fall significantly above field sources at the same galactic latitude and extend to radii of 20 arcmin and beyond. Such sources mark these systems as also extending to large angular diameters on the sky and thus closer than their redshift distances.
## 15 Spectroscopic observations in Progress around Arp 220
A joint observing project on optically identified X-ray sources around active galaxies is currently being carried out by E.M. Burbidge, Y. Chu and H. Arp. Measures by Chu with the 2.2 meter Beiging telescope and E.M. Burbidge with the 3 meter Lick Observatory reflector have yielded results some of which are in the process of being reduced and reported.
In the case of Arp 220 it can be said that the two brightest X-ray sources in the field are 20.3 cts/ks (20.3N and 20.3S in Table 3). The physical association of this pair of X-ray sources is suggested by the fact that they are so much brighter in X-rays than the remaining sources in the field and that they are almost exactly equal in X-ray flux. These two have now been confirmed as quasars. they are of low and rather similar redshift which would support the inference from their diametric positions that they had been ejected in opposite directions from the active nucleus.
An even more exciting result is in the process of final reduction. This is the pair of X-ray sources labeled RSO and BSO in Fig. 15. These two sources are aligned as exactly as can be measured across the nucleus of Arp 220 and are evenly spaced, each at about 7.3 arcmin distance. It is also apparent from Fig. 15 that there is a line of 4 or 5 X-ray sources which curves through the z = .09 X-ray galaxies and lead directly to the southern member of the pair of sources, designated BSO.
The candidate labeled RSO (red stellar object) is exceptionally red for a quasar (O - E = 2.18 mag.) which makes it about 21.8 mag in the blue - very difficult for a 3 meter telescope located on Mt. Hamilton above the lights of San Jose. Nevertheless a 50 minute spectroscopic exposure recorded strong emission lines which unmistakably signaled a quasar! Immediately thereafter the southern member of the pair (BSO) was observed and it turned out to have less strong lines but at very closely the same wavelength as the RSO. As a result we have a pair of quasars extremely well aligned across Arp 220 with closely the same redshift.
Added to the pairing properties, of course, is the line of X-ray sources connecting BSO back to Arp 220. The unusual redness of the northern member of the pair may be related to the huge amounts of dust and absorption in the ULIRG. If entrained material is drawn out in the process of ejection, the quasar could be involved in a dense cloud of dust. The southern member is not so red (O - E = 1.07) but this is still somewhat red for the usual quasar. X-ray material, however, seems to have been shed behind it in its outward track so more of the obscuring material may have been shed than for the RSO ejected in the other direction.
## 16 Note in appendix which deals with absorption over the field associated with ULIRG’s and disturbed galaxies
Please see Appendix A for detailed discussion of the X-ray hardness ratios as a measure of absorption and also some consequences of depressing the counts of more distant, background sources.
## 17 Note added on quasar near Mrk 273
E.M. Burbidge has kindly allowed her spectroscopically determined redshift of z = 1.168 for the 18.1 mag. BSO SE of Mrk 273 to be quoted. This is a particularly important redshift because of its similar magnitude, strong X-ray flux and close proximity to the previously known z = .941 quasar. When corrected to the reference frame of Mrk 273 ( z= .038) the new redshift goes to z = 1.089 and the .941 redshift goes to z = .870. This means the redshifts fall very close to the quantized redshift peak at z = .96 for quasars in general (Arp et al. 1990). In fact one redshift differs by only +.066 and the other by -.055 from the z = .96 peak. These are characteristic peculiar velocities relative to the parent galaxy (Narlikar and Arp in preparation).
## 18 Field of NGC6204
There is only a 5.2 ks PSPC exposure available for this fourth, bright ULIRG. Consequently only the brightest X-ray sources can be examined (above 3.7 cts/ks). Fig. 22 shows that these X-ray sources are distributed principally along directions NE and SW from NGC6240.
Inspecting the sources within 1 deg. radius around NGC6240 as compiled by SIMBAD it is clear that also radio sources run approximately out to the NE along this same line. In Fig. 21 we show in a gray scale map the disposition of radio sources as recorded in the VLA continuum survey (nvss). Proceeding outward from NGC6240 there is an extended radio source 5.4 arcmin to the east of the ULIRG. It is low surface brightness and not readily identifiable with any optical object. Probably it represents some radio plasma ejected in this direction. Further out along a line to the NE is a string of four radio sources. Somewhat above this line are four X-ray sources. The arrangement and spacing of the outer three sources is similar - and reminiscent of the correlation between radio and X-ray sources which we saw around Mark 231 (Fig. 9).
Going in the other direction, we see three strong radio sources in a line to the W. The striking feature here is that two of the sources are very close doubles. We have seen a strong tendency in all of the previously mentioned active objects for the associated sources to be double. At such small separations they stand out from the rest of the radio sources in the field of Fig. 21 and are therefore additionally indicated to be physically associated with NGC6240. The doubleness might be attributed to their (recent) origin in a typical opposite ejection from an ejected denser body which may or may not be in the near vicinity.
The radio sources pictured in Fig. 21 are mostly confirmed in the 4850 MHz surveys (Griffith et al. 1995) and the 365 MHz survey (Douglas et al. 1996). But those fields extend further to the S than pictured in Fig. 21 and show a strong radio source to the SW which strengthens the appearance of a straight NE - SW line through NGC6240 (as in the map only of X-rays shown in Fig. 22).The only three galaxies shown by SIMBAD in this field also fall approximately along this line. The central of these three falls about 12 arcmin NE of NGC6240, is 15.7 mag. and an IRAS source. Optical identifications should be made for these X-ray and radio sources and spectroscopic measures made also on the galaxies in order to study the details of their association.
## 19 The bright radio galaxy 3C31/NGC383
Although NGC383 is not classified as an ULIRG (it only has a flux of 1.2 Jy at 100 microns), it is 1.5 mags brighter in the blue than Arp 220 and exhibits some of the same characteristics as the galaxies we have just been investigating. In particular it is a strong X-ray source and has strong, ejected radio filaments which lead toward companion galaxies which are also X-ray sources. Curiously, it has an almost identical redshift to Arp 220, z = .017 for NGC383 compared to z = .018 for the extreme ULIRG. NGC383 is, however, the 13th strongest radio E galaxy in the sky at 1420 MHz, and among that group, the 9th brightest apparent magnitude E galaxy (Arp 1968a). It has an emission lines which are ionized from the nucleus and probably shows a weak, broad emission component (Owen, O’Dea and Keel 1990).
The most important feature for our purposes, however, is shown in Figs. 23 and 24 where it is seen that there is a tight group of galaxies which are strong in X-rays (1E 0104) and exhibit an X-ray tail pointing back toward NGC383 . The resemblance to Arp 220, as it was pictured in Figs. 13 -15, is remarkable . The startling fact then emerges that the redshift of the group associated with Arp 220 is z = .09, very similar to the redshift of the group associated with NGC383 which is z = .11
There is strong additional evidence for ejection of X-ray objects in Fig. 23. Most of the discrete X-ray sources surrounding NGC 383 lie on opposite ends of diameters passing close to the central galaxy. Most prominent is a line of sources at p.a. = 23 deg., including the 1E 0104 X-ray galaxy group. Both outer sources are elongated along this diameter. Fig. 23 also shows a pair of strong sources at p.a. = 100 deg. with material connecting back to the central X-ray mass. A weaker pair is seen at p.a. = -45 deg. with the NW component a double elongated back to the center. The innermost pair is particularly well aligned at p.a. = 60 deg. with the SW component exhibiting isophotes which conspicuously connect back to the central galaxy and the NE component being a close double source.
The most significant aspect of these observations must be the X-ray filaments and isophotal extensions which lead back to the central source. Just examining the X-ray map of the NGC383 surroundings, where the extensions fill a region out to 33 arc min radius, would seem to offer the most direct demonstration possible of the ejection origin of these sources. (The same kind of evidence is available from the galaxy cluster Abell 754 - see Arp 1998a, Fig. 7-17 and also Mark205 on the cover and Fig 1-7). The individual X-ray sources around NGC383 of course should be optically identified as in Tables 1- 3 around the ULIRG’s. Presumably they would be mostly quasars whose properties could be usefully compared to those around the ULIRG’s.
## 20 The connection between radio and X-rays
The radio map superposed on the X-ray map in Fig. 25 gives perhaps the clearest evidence for the physical relation between the radio material which is accepted as being ejected, and X-ray material for which similar evidence has been accumulating. We notice that the radio filaments lead in the general direction of the X-ray sources, but they deviate somewhat, principally toward the end of the track. This is essentially the same result as gleaned from all the preceding ULIRG’s dicussed in this paper. In NGC383 the radio track is more continuous and well marked but it is suggested that it is be due to one of the same two causes, or combination thereof:
1) The X-ray sources and radio material are ejected along the same track but the motion of the local medium separates the radio plasma from the denser X-ray emitting bodies.
2) Successive ejections are rotated and the radio track is an older remnant of a preceding ejection.
In either case the radio and X-ray material should be superposable by simple rotations and translations. In NGC383 the southern extensions seem to be well fitted with a small rotation. The northern radio extension seems to drift, except for a small spot, considerably westward at its end, perhaps due to some perturbation in the medium or local event. (For a purely translational drift note 3C212 in Fig. 10 of this paper).
A corollary to this mechanism is that if the radio and X-ray material try to penetrate through appreciable material in the ejecting galaxy, the radio plasma will be stripped without ever getting out. This would be a natural explanation of why X-ray ejections seem to come out in several directions but radio ejections tend more to occupy one main channel.
## 21 Quasars associated with NGC383
There are some bright radio quasars and quasar like objects catalogued within a degree or two of NGC383 ( z = .603, 1.71 and a Sey1 with z = .015 which has a companion of Z = .287). But the one for which it is possible to calculate a very high probability of association is B in Fig. 24. It has z = 2.027 and falls only 10 arcsec from A, the main galaxy in the group which has z = .11.
How probable is this to be an accident? In an objective prism survey Weedman (1985) found quasars with 2 $``$ z $``$ 2.5 to have a density of .25/sq. deg. at a magnitude of m<sub>4500</sub> = 19 mag. Quasar B has R = 18.9 mag. and therefore would have a chance of approximately 6$`x10^6`$ of falling accidentally within 10 arcsec of galaxy A. Moreover the analysis of Komossa and Boringher (1999) make it seem likely that galaxy A and its group are the source of the strong X-rays. The association of an active X-ray galaxy with the quasar by chance is then vastly less likely. It seems difficult to escape the conclusion that we have another example of a physical association of a high redshift quasar with a low redshift galaxy (Burbidge 1996; Arp 1998a;1999). It is to be noted that, although the quasar is within the optical boundaries of galaxy A it shows no absorption lines at z = .11 (Hewitt and Burbidge 1993). It shows a series of absorption systems from z = 1.97 down to z = 1.75 but the absence of absorption lines from A would conventionally argue that it was not behind A.
Of course we have argued that the origin of the group A galaxies was by ejection from the central NGC383. This could make the quasar a secondary ejection from an active, evolving companion or a later epoch proto galaxy either entrained or traveling the same track. What then is the probability that A has come from NGC 383, or for that matter that the quasar has come directly from NGC383?
Fig. 26 shows the X-ray contours obtained from the Einstein IPC. These observations are with a less sensitive detector than used for the ROSAT maps in Figs. 23 and 25 but at a somewhat higher energy, namely the .3 to 3.5 keV band. It is clear that the isophotes of 3C31 are extended toward 1E 0104 and the contours of the X-ray source are extended back toward 3C31. Allowing a generous $`\pm `$ 10 degrees in pointing coincidence, the chances of accidental mutual alignment are about 1 in 100. But now we have to ask what are the chances a 3C radio source would be encountered in an area of 16.4 arcmin radius. About 400 3C sources in about 30,000 sq deg. of sky gives another factor of about 1 in 100. So the total chance of finding a galaxy as active as NGC383 with mutual X-ray alignments to 1E 0104 is about 1 in 10,000.
But perhaps the most compelling evidence for association is the similarity with Arp 220, where a tight group of X-ray galaxies of about the same redshift are linked directly back to the ULIRG by X-ray and HI material.
## 22 Distances and quantizations of redshifts
The angular distance from Arp 220 to its associated (z = .09) galaxy group is 2 arcmin whereas it is 16 arcmin in the (z = .11) NGC383 case. The brightness of the NGC383 system makes it tempting to acount for this by arguing that NGC383 is nearer the observer. But we should remember that the 3C321 sytem at z = .096 was also associated with Arp 220 but at 56 arcmin radial distance. We argued there that the retarding interaction with the ejecting galaxy might determine the different distances traveled by the ejecta at their observed evolutionary stage.
The distances of these parent galaxies is difficult to estimate because their redshifts appear to be an intrinsic property related to their age rather than a measure of their distance. There are several ways of seeing this. One is through the evident quantization of their redshifts. For example the redshifts of Mark 273 and Mark 231 are z = .038 and z = .041 respectively. The redshifts of Arp 220 and NGC383 are z = .018 and z = .017 respectively. The latter redshifts are very close to the 5000 km/sec peak in redshifts over the whole sky but particularly of the Perseus - Pisces filament that stretches almost 90 deg. across the sky. (see Arp 1987 p. 129ff; Arp 1998a p. 149ff). There is an obvious problem with interpreting this as a shell of galaxies expanding away from our own in every direction. As an alternative the above references indicate that, empirically, galaxies at these redshifts give evidence of being associated with brighter, more nearby galaxies. In the case of 3C31/NGC383 it falls only 3.2 deg. away from the Local Group Seyfert 3, NGC404.
It is also evident that quasars in these systems which have measured redshifts come close to the global redshift peaks of z = .06, .30, .60, .96, 1.41, 1.96, 2.64 …etc. For example the quasar B, SSW of NGC383 has z = 2.027. (In the reference frame of NGC383 z<sub>Q</sub> = 1.976). The intrinsic redshifts seem to decline in discrete steps as they evolve toward more normal galaxies. It would be natural to reason that the ULIRG’s and other active galaxies are part of this evolutionary sequence and still have a dominant component of the age related redshift which is quantized. For this reason it might be better to estimate the distance from the apparent magnitude of the quasars. In the systems we have been examining, the quasars are quite bright in apparent magnitude. The quasar B again, for example, is R = 18.9 mag. This is bright for a quasar of this redshift and is similar to the properties of the line of quasars coming SW from M33, another Local Group galaxy near NGC404 and NGC383. (See Arp 1987 p.71ff).
An other example of an object like those we have been studying is shown in Fig. 27. It is the disturbed spiral IC1767. Amazingly it has a redshift of z = .0175, just between that the .018 of Arp 220 and the .017 of NGC383. It has been argued to be a southern extension of the Perseus-Pisces Cluster (Maurogordato, Proust and Balkowski 1991). Like NGC383 it is only a modest infrared, IRAS source of flux 1.0 Jy at 100 microns (but rising steeply toward longer wavelengths).
Like Arp 220 and NGC383, IC1767 shows evidence for ejecting radio sources and higher redshift objects. Its most outstanding feature is a pair of very strong radio sources aligned across it at 1.3 deg separation on the sky. This pair was so strong it was identified 31 years ago (Arp 1968b). Subsequently the radio sources were measured to be quasars, both very close to z = .6. (If one refers the redshifts to the central galaxy they become z = .588 and .640, which, after allowing for about cz = .02 toward and away ejection velocity, gives both quasars very close to the major quantization peak of z = .60. The combination of these properties with the strength of the radio sources and closeness of alignment gives negligible chance of accidental association of these quasars.
The argument now centers on the properties of these z $``$ .6 quasars. they are unusually bright in apparent magnitude, V = 17.09 and 16.40 mag., exceptionally bright in radio wavelengths, and unusually widely spread across the central galaxy. All these properties argue for an unusually closeby system. Thus it becomes another argument for the active objects like Arp 220, NGC383 and IC1767 to belong to the Perseus-Pisces redshift peak, but for all these galaxies to be much closer than their conventional redshift distance.
## 23 Summary: Empirical Properties of Active Galaxies
The operational definition of an ”Ultra Luminous” galaxy is one which deviates strongly above the Hubble relation between redshift and apparent magnitude. Since the class of galaxies we are investigating here differs in spectroscopic and morphological characteristics from the one that defines the Hubble relation, the latter cannot be used to define the distances or luminosities of the ULIRG’s. In fact the class of active and high surface brightness galaxies - Seyferts, compacts and quasars - defines a conspicuously steeper slope than the Hubble slope (Arp 1968a; 1998a, Fig. 2 - 19), ruling out expansion velocity redshifts unless the luminosities are progressively adjusted to compensate.
We have investigated here a sample of active galaxies with the most extreme deviation from the Hubble relation. We have found them to be strong X-ray sources, show vigorous ejection of radio material and generally disturbed morphologies. They do not look like the relaxed, symmetric galaxy forms of the most luminous nearby galaxies - whose distances do not depend on redshift. Characteristically the central active galaxy has physically associated companions which are even more active and of various degrees of higher redshift. The activity and morphological disturbances seem to be associated with lower rather than higher luminosity and there are many new cases presented here of apparently younger companions with higher redshifts than their galaxy of origin.
Mark 273 at z = .038 shows X-ray and optical connections to an active companion object at z = .458. Optical and radio connections lead toward a pair of optical and X-ray bright quasars one of which has z = .941. Many strong X-ray pairs across Mark 273 have been optically identified as almost certain quasar candidates. A line of 100 micron infrared sources extends over 20 arcmin from Mark 273. Their nature and relationship to Mark 273 is a mystery to be investigated.
Mark 231 at z = .041 shows well marked radio ejections in two, roughly orthogonal directions. X-ray sources are associated with these radio patches in such a way as to suggest the radio plasma is in varying degrees of being stripped from them. The close double nature of many ejected radio sources is again seen and becomes a general characteristic of ejecta from these active systems. A string of infrared sources stretches over 3 degrees from Mark 231 in the direction of a major interior radio ejection from this active ULIRG! Eleven quasar candidates are optically identified as BSO or Bcg candidates, many in pairs.
Arp 220 at z = .018 has a group of X-ray galaxies at z = .09 linked to it by neutral hydrogen extensions. Moreover they are a part of a chain of X-ray sources emanating from the central ULIRG in a shallow spiral shape which has knots paired closely across the central object. The interpretation would seem to require ejection with rotation. The 50 x 50 arcmin field shows the inner ”S” shape bending over into an extended spiral distribution of X-ray sources of the order of 2 degrees in diameter! Radio and X-ray sources are shown to extend along this same line and for the same distance. Heavy absorption over the Arp 220 area shows up in the hardening of the X-ray sources over an area of the order of a degree in diameter.
NGC6240 at z = .024 shows in a short exposure, X-ray sources extending on a line either side of the the ULIRG. Radio surveys show lines of radio sources extending out farther than 40 arcmin. Double sources are again featured.
The bright radio galaxy NGC383 at z = .017 shows pairs of X-ray sources connected diametrically across it. The strongest pair includes a group of X-ray galaxies at z = .11 . There is a quasar of redshift z = 2.027 only 10 arcsec from the central galaxy in this group and they are both apparently associated with the central radio galaxy (3C31). The radio extensions from 3C31 are the strongest of any of the cases and appear to be only slightly displaced from the main X-ray ejections.
Note the similarity of the X-ray galaxies of z = .11 associated with NGC383 to the X-ray galaxies of z = .09 connected to Arp 220. In one of the best marked ejection lines from a Seyfert galaxy, NGC3516 (Chu et al. 1998), the last X-ray galaxy in the line has z = .089. The latter confirms both the ejection origin of these galaxies and the quantized nature of their redshifts. It is suggested that only the distances from the ejecting central galaxy change in accordance with the amount of their interaction as they exit.
An important final point should be that the preceding paper analyzes a more or less complete sample of the brightest ULIRG’s. Fainter examples would be expected to follow somewhat the same patterns. In recent sample of faint ULIRG’s (Stanford et al. 2000) there is a marked tendency for the central objects, pictured in the K band, to have diametric companions or two or three condensations emerging in a jet like configuration.
## 24 What are the distances of these active central glaxies?
If we cannot use the conventional redshift distances, then what are the real luminosities of the systems we have been investigating in this paper? I see two possibilities of estimating non-redshift distances:
1) The systems all seem to be ejecting high redshift quasars of the kind that have been physically associated with nearby galaxies. If we view the quasar redshift as signaling an age related phase in its evolution toward a normal galaxy then it may be that at this phase all have similar luminosities (for example if ejected masses are similar). Then we could judge by the apparent magnitude of the associated quasars what the distance was. A preliminary assessment seems to indicate the quasars around the active galaxies investigated here are somewhat brighter than those associated with the sample of Seyerts in Arp 1997. Those Seyferts were judged primarily to be in the Local Supercluster so the active galaxies investigated here would be indicated to be somewhat closer.
2) The angular size of the associations of material we have seen around the ULIRG’s in this paper are of the order of a degree radius or more. This would class them with Local Group galaxies. One might argue that the outer sources are escaping into the general field from the central galaxy and only within a much smaller radius does the material fall back in and then evolve into a smaller, more relaxed galaxy or group of galaxies. (This may finally offer a legitimate use for merger calculations).
The intriguing question then arises: Will Mark 273 and Mark 231 evolve into systems like Arp 220 and NGC383 and what will the latter two systems evolve into? NGC383 is firmly a member of the Perseus-Pisces chain. As such it consists of redshifts between 4000 and 6000 km/sec and exhibits extensions of several degrees on the sky. If the redshifts continue to lower and the luminosities increase as all these objects age they may develop into what we consider more ”local” groups and clouds of galaxies.
There are two kinds of observations which could shed light on this question:
1) There was an infrared, NICMOS observation (Scoville et al. 1998) of the central regions of Arp 220 which discovered 8 possible star clusters (unresolved). The authors calculated that at the redshift distance of Arp 220 these clusters would have an absolute K magnitude of M<sub>K</sub> = -13.5 mag. This turned out to be brighter by 1.5 magnitudes than, for example, any in NGC5128, the giant E galaxy also know as Centaurus A. Now if we were to move Arp 220 a factor of 10 closer, that would lower the luminosity of these objects to M<sub>K</sub> = -8.5, about that of the brightest stars in a galaxy. Have stars been resolved in Arp 220? A factor of 10 closer than its redshift distance would place the ULIRG about half way between our Local Group and our Local Supercluster.
2) Quasar redshifts, as mentioned earlier, are ”quantized”. Empirically the spread around the preferred values is less for quasars associated with low redshift galaxies (Arp et al. 1990). This has been interpreted as meaning that the $`\pm `$ .1c ejection velocities carry some quasars out into the field, but those that are captured by the ejecting galaxy must slow down, lose this .1c ejection velocity and assume more closely their intrinsic, sharply quantized redshift values. When such data is available for the systems discussed here it may be possible to identify the gravitational diameter of the system and hence judge its distance by this angular size criterion
In general it could be remarked that there is a quite well established, empirical sequence of evolution from high redshift, compact, low luminosity quasars to brighter medium redshift quasars to compact, active and disturbed galaxies and finally to more relaxed, normal glaxies (Arp 1998b). The intrinsic redshift drops in steps along this sequence and it is implied that the ULIRG’s considered here are in a phase of still evolving toward lower redshift - they are evolving from quasars, and like quasars, they are much closer, and less luminous than their conventional redshift distances. If this is true, one piece of terminology which might be salvaged from previous work is the name ULIRG - except that it now would have to mean ”Under Luminous” galaxy.
## 25 Acknowledgements
I would like to thank the many observers who allowed use of their results, some of them unpublished, which enabled a comprehensive view of the properties of these bright, very active galaxies.
## 26 Appendix A
The ratio of hard to soft X-rays is HR1 = (A - B)/(A + B), where A is the flux between .4 to 2.4. keV and B is between .07 and .4 keV. Testing the behavior of the hardness ratio on a five assumed homogeneous star fields showed that the hardness ratio varies with off-axis distance, crowding in the center of the field and source strength. The control field HR8905 was chosen as:
1) the closest match to Arp 220 in all these characterisitics. (N<sub>H</sub> = .43 for Arp 220 and .50 x $`10^{21}cm^2`$ for the star field).
2) Only sources greater than 1.3 cts/ks were used and
3) appreciably extended sources were excluded.
Correcting for a difference of .07 in HR1 due to the different $`N_H`$ in the two fields gives an estimated excess in hardness ratio from Fig. 20 of about 0.4 in HR1. This translates into a difference in N$`{}_{H}{}^{}=.4x10^{21}`$ in the hydrogen column and, using the relation A<sub>V</sub> = .56N<sub>H</sub>(10$`{}_{}{}^{21}cm^{}2)`$ developed by Predehl and Schmitt (1995), a difference in visual absorption of A<sub>V</sub> = .22 mag. is obtained. Doubling this for the absorption through the entire cluster gives $`E_{BV}`$ = .15 and A<sub>B</sub> = .6 mag.
This must be a minimum absorption for the material associated with Arp 220. Some reasons are:
1) When the hardness ratio approaches 1 it becomes an insensitive measure of the absorption (see Pietsch, Trinchieri and Vogler 1998, Fig. 1).
2) If there are any foreground sources in Fig. 20 they would lower the excess hardening attributed to Arp 220.
3) Some objects associated with the ULIRG may be obscured by dense clouds and be below the detection level of the exposure.
In general, however, the points in Fig. 20 spread around the mean by roughly $`\pm `$ .4 as if we were seeing some on the near side and some on the far side of the group.
An absorption value from completely different type of measurement, the Balmer decrement of the emission line galaxy B in Fig. 13 gives $`A_V`$ = .51 $`\pm `$.1mag. (Ohyama et al. 1999). Galaxy B, attached to Arp 220 by X-ray material is surely near the center of the group which implies the absorption A<sub>B</sub> through the whole group of 1.4 mag. Since Arp 220 with its enormous absorption and reddening (Becklin and Wynn Williams 1987 estimate 50 magnitudes of visual absorption ot the center) lies also in Fig. 20 very near HR1 = 1.0, it implies that that HR1’s near this value signal much larger absorptions than the minimums we have estimated. The important result of Fig. 20 then becomes the fact that these high absorptions extend out past 30 arcmin radius.
It should also be noted that recent measures at 200 microns (Alton et al. 1998) show colder dust extending to larger radii than 100 micron dust in resolved galaxies. If this is true of Arp 220 its diameter would grow even larger and raise the question of what is the role of this dust in the evolution of this system.
One point of general interest is that with a minimum of $`A_B`$ = .6 mag. absorption, the normal background count of X-ray sources must be appreciably depressed. But with the larger absorptions indicated in the further discussion, background sources could be almost completely supressed and therefore almost all the sources actually belong to Arp 220. This is a point which has not been considered when analyzing presumably distant quasars which show absorption lines from intervening absorption clouds.
For example, in 32 total fields searched for quasars behind elliptical galaxies and clusters of galaxies Knezek and Bregman (1998) found 13 quasars out to distances of r $``$ 13 arcmin with X-ray counts $``$ 10 cts/sec. The density of these quasars reached from 5 to 8 per sq. deg. between 8 and 2 arcmin from the center. The expected background density of X-ray quasars using their detection success rate was about 1.4 quasars per sq. deg. This concentration of quasars toward the E galaxies and clusters is illustrated in Fig. 28.
The above result addresses a long standing claim that the observation of lower redshift absorption lines in high redshift quasar spectra proved that the quasars were at their large, emission line redshift distances. Of course the only thing that the observation proved was that some quasars were not in front of the clusters. They could be inside the cluster or, if they did not show absorption lines they could be in front, or behind but seen between clear patches. The only decisive observation would be to see if they were more numerous in the direction of the galaxies than the surrounding background. The Knezek and Bregman observations seem to have inadvertently established that. The point that the Arp 220 observations now make, however, is that because of the absorption of the galaxy or cluster (as evidenced by the absorption lines in the quasar spectra) the background count of quasar density should be to some extent lower in the direction of the galaxy or cluster. Thus the result quoted above for an excess number in the direction of the galaxies is actually a lower limit because an undepressed background was used for comparison.
In the direction of clusters of galaxies ”…the significant overdensity of background bright quasars … on a scale of 10 arcmin…” was noted by Wu and Fang (1996). Later these authors reexamined the possibility that the excess could be due to gravitational lensing and concluded ”…the quasar galaxy association remains an unsolved puzzle in today’s astronomy…” (Zhu, Wu and Fang 1997). It should be emphasized that evidence for absorption in the cluster can only strengthen this conclusion.
Figure Captions:
Fig. 1 X-ray contours on a POSSII, R image of Markarian 273 and 273x from Xia et al. 1998
Fig. 2 Deep R image from the Univ. of Hawaii 88-inch telescope by John Hibbard. Arrow points to Mrk 273x.
Fig. 3 Smoothed and contoured HRI X-ray image of Mark 273 and Mark273x (72 arcsec NE) by Thomas Boller.
Fig. 4 Deeper, lower resolution PSPC X-ray image smoothed and contoured by Thomas Boller.
Fig. 5 VLA, 20cm continuum map of Mark 273 by Min Su Yun.
Fig. 6 VLA radio continuum map from survey (nvss). Strong X-ray sources are marked with x’s, including a quasar of z = .941 and a BSO of E = 18.1 mag.
Fig. 7 X-ray photons in a 50 arcmin radius PSPC field. Brighter sources from the standard detection algorithm are marked in counts per kilosecond. Mark 273 is at center with 13 cts/ks and Mark 273x has 8. The two strong, candidate X-ray quasars in Fig. 6 are marked 21, referred to in text as 21N for North, and 38. The data on the labeled sources is given in Table 1. Note elongation of source distribution NE to SW through Mark 273 and also pairs of sources SE to NW. The picture is 50 arcmin on a side (1 sky pixel = $`1/2`$ arcsec)
Fig. 8 A plot of all infrared sources of greater than .5 Jy at 100 microns within a one deg. radius of Mark 273.
Fig. 9 (a) High resolution radio map of 10.4 x 20.7 arcmin field around Mark 231. From 20cm, VLA, FIRST survey. Notice the pairs of double radio sources flanking the central galaxy. (b) At the same scale as above, a PSPC X-ray exposure (sources are dark) overlayed with contours of radio sources from the low resolution VLA survey. (c) Optical map of double X-ray source consists of z = 1.272 quasar and O = 19.9 mag. BSO at p.a. = 75 deg. (nvss).
Fig.10 Radio contours are here superposed on an HST image of the radio galaxy 3C212 (Stockton 1998) Images marked f and g are optical objects
Fig. 11 At an infrared wavelength of 60 microns, a 6.1 x 6.1 deg. IRAS survey field is pictured
Fig. 12 X-ray photons within about 50 arcmain from a PSPC exposure. Cts/ks are labeled for brighter sources. Mark 231 has 14.3 cts/ks. Data for sources are given in Table 2
Fig. 13 Optical identification of the four largest galaxies around Arp 220 which has z = .018. Ohyama et al. 1999
Fig. 14 HI contours at the redshift of Arp 220 (z = .018) leading down to a group of galaxies at z = .09. By J. Hibbard
Fig. 15 Hard X-ray band, .5 to 2.4 keV, smoothed, showing curved string of sources leading down from Arp 220
Fig. 16 The z = .09 companions SW of Arp 220 with 4 small high redshift objects circled. Ohyama et al. 1999
Fig. 17 Standard detection X-ray sources in a 43 x 43 arcmin field. No. 3 is Arp 220. No. 7 is a star
Fig. 18 X-ray photons (broad band .11 to 2.4 keV, slightly smoothed) within about a 50 arcmin radius with brighter sources labeled in cts/ks. Arp 220 is in the center, unlabeled, but with a count of 7.3 cts/ks. The z = .09 companion galaxies are an extended source of 9.6 cts/ks merged with Arp 220 to the SW at his scale (also unlabeled). The 8.9 and 6.9 sources are Nos. 1 and 2 respectively in Fig. 17, the 2.8 source is No. 10 in Fig. 17 and the source just above the 2.8 (unlabeled in this Figure) is No. 9 in Fig.17. The radio source 3C321 and its two lobes are sketched in to scale.
The enlarged view in Fig. 17 illustrates how the sources and background radiation in the central regions turn over from an approximately vertical ”S” shape to a spiral distribution elongated in the NW to SE direction in this map
Fig. 19 Plot of all radio sources $`>`$ 70 mJy at 4.85 Ghz (triangles) and all infrared sources $`>`$ .7 Jy at 100 microns (filled circles) within 1 degree of Arp 220. From surveys by Becker et al 1991; Moshir et al 1989
Fig. 20 Hardness ratio, HR1, of X-ray sources as a function of angular distance from Arp 220. The dashed line indicates averages from a control field at similar galactic latitude
Fig. 21 The dark spots show strong radio sources from the VLA low resolution survey (nvss). Small x’s show X-ray sources from Fig. 22. Two of the radio sources on a line W of the ULIRG are just discernible as very close doubles. Dashed outline represents an extended radio source. It is 42.5 arcmin from NGC6240 to the westernmost double
Fig. 22 A PSPC X-ray map around NGC6240 (radius $``$ 1 deg.) showing standard detection sources with liklihood greater than 10. From S. Komossa, private communication
Fig. 23 A PSPC X-ray map at .5 to 2.0 keV of 3C31/NGC383. Sources pairing across the center are indicated by their position angles. Adapted fom Komossa and Boehringer 1999
Fig. 24 The galaxy group 16.6 arc min SSW of NGC383 (the + marks the centroid of the the X-ray source 1E 0104). Galaxies A,C,E,F have redshift z = .11. B is a quasar with z = 2.027. Red image by Gioia et al. 1986. Insert at higher contrast shows the quasar embedded in the image of the dominant galaxy
Fig. 25 The same X-ray map as in Fig. 23 but now with the radio jet contours. Overlay of a .6 GHz radio map from Strom et al. 1983 by Komossa and Böhringer
Fig. 26 The X-ray map of NGC383/3C31 from the Einstein Laboratory IPC. Two stars have been removed from the figure published by Gioia et al 1986
Fig. 27 Very strong radio sources form a pair across the disturbed spiral IC1767. Their redshifts are close to the z = .60 quasar peak and the galaxy’s redshift is at the 5000 km/sec galaxy peak
Fig. 28 Density of quasars found near E galaxies and the centers of galaxy clusters by Knezek and Bregman (1998). The background is calculated from the density of X-ray sources with their quasar detection rates but with no background count suppression from absorption by the central object
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# Quark combinatorics for production ratios in hadronic 𝑍⁰ decays
## Abstract
Model independent verification of the quark combinatorics rules, which govern the ratios of the yields of secondaries, is presented for jet processes. Because of a large number of produced resonances in the hadron jets, a test of the quark combinatorics rules is hardly possible in the central region, $`x_{hadron}0.2`$. However, a model-independent verification is plausible at $`x_{hadron}1`$. It is shown that for the large-$`x_{hadron}`$ kinematical region the quark combinatorial relations are in a reasonable agreement with data for $`\rho ^0/\pi ^0`$ and $`p/\pi ^+`$ ratios.
Initially the rules of quark combinatorics were suggested in and checked in multiple production processes, see and references therein. Henceforth these rules were applied to the decay processes of resonances .
Concerning multiple hadron production processes, it is necessary to underline that quark combinatorics rules manifest important features of the process in which quarks form hadrons, that is, soft colour neutralization and soft hadronization .
There exist two types of predictions of quark combinatorics for multiple production processes. The first one is the prediction of the yield ratios for hadrons belonging to the same quark multiplet. For example, we have for the lightest meson nonet:
$$\rho ^0/\pi ^0=3,K^0/K^0=3,K^\pm /K^\pm =3.$$
(1)
Two statements form the basis for these relations: $`(i)`$ wave functions of hadrons from the same multiplet are equal, and $`(ii)`$ quarks which are fused into hadrons are spin- and flavour-non-correlated.
The second type of prediction is related to the yields of mesons and baryons. For the hadronization of the quark $`q_i`$ in $`(q,\overline{q})`$-sea the production rule reads:
$$q_i+(q,\overline{q})_{\text{sea}}\frac{1}{3}B_i+\frac{2}{3}M_i+\frac{1}{3}M+(M,B,\overline{B})_{\text{sea}},$$
(2)
where $`B_i`$ and $`M_i`$ are baryons and mesons containing the quark $`q_i`$ and $`M`$, $`B`$ and $`\overline{B}`$ are mesons and baryons of the sea.
The prompt verification of Eqs. (1) and (2) is rather difficult because in multiple production processes a number of resonances is produced. Generally one can write
$$M=\underset{L}{}\mu _LM_L\text{ and }B=\underset{L}{}\beta _LM_L,$$
(3)
where $`L=0,1,2,\mathrm{}`$ define the multiplet, while $`\mu _L`$ and $`\beta _L`$ are production probabilities of mesons and baryons of the given multiplet in the process of quark hadronization. The probability $`\mu _L`$ is determined by characteristic relative momenta of the fused quarks.
A straightforward way to overcome the ambiguities related to the resonance production is as follows: one can take into account all existing resonances and their decay into all possible channels. This is the scenario suggested in Refs. . However, in this way one faces a set of problems. The matter is that the number of resonances which are observed and cited in the compilation is a comparatively small fraction of the whole set of existing states. The basis for this statement is provided, for example, by recent investigations of meson production data where a large number of new meson states with masses in the region 1950–2350 MeV is reported. One should keep in mind that, naturally, those resonances are first discovered which can be easily detected. And one should also recall that all observed resonances have multiplet partners which are also produced with approximately equal probabilities and the decay of those form a background which prevents the direct investigation of Eqs. (1) and (2). There exists one more effect which does not allow the realization of the program discussed in at the time being. This is the effect of the accumulation of widths of overlapping resonances by one of them; it has been observed for scalar/isoscalar mesons in the region 1200–1600 MeV . As a result of width accumulation, a broad state $`(\mathrm{\Gamma }/2400`$ MeV) has been formed; similar states can exist, as was stressed in , in other waves, in other mass regions. Currently it does not seem possible to take into account the productions and decays of such broad states.
Still, investigation of jet processes in the region of large $`x`$ opens a way to test the relations (1) and (2): this article is devoted to a presentation of this possibility.
The resonance decays increase the contribution of lighter hadrons, such as pions and kaons in case of mesons, and nucleons in case of baryons. However, in jet processes $`Z^0q\overline{q}hadrons`$ which were studied in , the spectra, being maximal at $`x0`$, decrease rapidly when $`x1`$. It results in a rapid increase of contribution of the promptly produced particles with the increase of $`x`$, since the decay product carries only a fraction of $`x`$ of the initial resonance. It is just this feature which enables us to estimate the ratio of probabilities for promptly produced hadrons.
The jets of light hadrons are formed in the processes $`Z^0u\overline{u}`$, $`Z^0d\overline{d}`$ and $`Z^0s\overline{s}`$. According to
$$\mathrm{\Gamma }_{u\overline{u}}/(\mathrm{\Gamma }_{u\overline{u}}+\mathrm{\Gamma }_{d\overline{d}}+\mathrm{\Gamma }_{s\overline{s}})=0.258\pm 0.031\pm 0.032,$$
(4)
$$\mathrm{\Gamma }_{d\overline{d},s\overline{s}}/(\mathrm{\Gamma }_{u\overline{u}}+\mathrm{\Gamma }_{d\overline{d}}+\mathrm{\Gamma }_{s\overline{s}})=0.371\pm 0.016\pm 0.016.$$
The leading quarks, the probabilities of which are given by Eq. (4), determine the hadron content at $`x1`$. To find the ratio $`V/P`$ at large $`x`$ we have fitted the spectra $`(1/\sigma _{tot})d\sigma /dx`$ of Ref. to the sum of exponents $`\mathrm{\Sigma }C_ie^{b_ix}`$; the results of the fit are presented in Fig. 1 for the spectra of $`\pi ^\pm `$, $`\pi ^0`$, $`\rho ^0`$ and $`(p,\overline{p})`$. The ratio of the fited curves, together with calculation errors, for $`\rho ^0/\pi ^0`$ is shown in Fig. 2a as a shaded area. We see that in the region $`0.6<x<0.8`$ the data are in an reasonable agreement with quark combinatorics prediction $`\rho ^0/\pi ^0=3`$.
The same fitting procedure has been done for the $`K^0`$, $`K^0`$, $`K^\pm `$ and $`K^\pm `$ spectra. The figures 2b and 2c demonstrate the ratios $`K^0/K^0`$ and $`K^\pm /K^\pm `$: the data give systematically lower values compared to the prediction (1), although there is no strong contradiction.
The verification of the barion-to-meson ratio given by Eq. (2) is of great interest and principal meaning as well. Quark combinatorics predicts for the proton-to-pion ratio:
$$p/\pi ^+0.20$$
(5)
In Fig. 2d one can see the ratio $`p/\pi ^+`$ given by the fit to the data (shaded area) and the prediction of quark combinatorics (5): the agreement at $`x>0.2`$ is quite good.
Let us give a brief comment to the calculation result $`p/\pi ^+0.20`$ for leading particles in jets. In the jet created by a quark the leading hadrons are produced in the proportions as follows:
$$q_i\frac{1}{4}B_i+\frac{1}{2}M_i+\frac{1}{4}M.$$
(6)
We consider only the production of hadrons belonging to the lowest (baryon and meson) multiplets, and, hence, keep only the terms with $`L=0`$ in Eq. (3). In our estimations we assume $`\beta _0\mu _0`$, and therefore we substitute $`B_iB_i(0)`$, $`M_iM_i(0)`$ and $`MM(0)`$. The precise content of $`B_i(0)`$, $`M_i(0)`$ and $`M(0)`$ depends on the proportions in which the sea quarks are produced. We assume flavour symmetry breaking for sea quarks, $`u\overline{u}:d\overline{d}:s\overline{s}=1:1:\lambda `$, with $`0\lambda 1`$. For the sake of simplicity, we put first $`\lambda =0`$ (actually the ratio $`p/\pi `$ depends weakly on $`\lambda `$). Then for the $`u`$-quark initiated jet we have:
$`B_u(0)={\displaystyle \frac{2}{15}}p+{\displaystyle \frac{1}{15}}n+(\mathrm{\Delta }\text{resonances }),`$ (7)
$`M_u(0)={\displaystyle \frac{1}{8}}\pi ^++{\displaystyle \frac{1}{16}}\pi ^0+{\displaystyle \frac{1}{16}}(\eta +\eta ^{})+(\text{ vector mesons }),`$ (8)
$`M(0)={\displaystyle \frac{1}{16}}\pi ^++{\displaystyle \frac{1}{16}}\pi ^0+{\displaystyle \frac{1}{16}}\pi ^{}+{\displaystyle \frac{1}{16}}(\eta +\eta ^{})+(\text{ vector mesons }).`$ (9)
The hadron content of the $`d`$-quark initiated jet is determined by isotopic conjugation $`pn`$, $`np`$, $`\pi ^+\pi ^{}`$, and the content of antiquark jets is governed by charge conjugation; in jets of strange quarks only sea mesons ($`M`$) contribute to the ratio $`p/\pi ^+`$.
Taking into account Eq. (5) and probabilities for the production of different jets (4), we obtain $`p/\pi ^+0.21`$ for $`\lambda =0`$. We can easily get the ratio $`p/\pi ^+`$ for arbitrary $`\lambda `$: the decomposition of the ensembles $`B_i(0)`$, $`M_i(0)`$, $`M(0)`$ with respect to hadron states has been performed in Ref. , see Appendix D (Tables D.1 and D.2). But, as was stressed above, this ratio is a weakly dependent function of $`\lambda `$: at $`\lambda =1`$ we have $`p/\pi ^+0.20`$.
In conclusion, jet-induced processes where the spectra $`(1/\sigma _{tot})d\sigma /dx`$ are rapidly decreasing with the growth of $`x`$ enable us to perform a model-independent verification of the rules of quark combinatorics. Experimental data on particle yields for the decays $`Z^0q\overline{q}hadrons`$ are in qualitative agreement with the predictions of quark combinatorics. Still, the available experimental data do not allow to conclude decisively about the accuracy of the quark combinatorics predictions.
The authors are indebted to V.V. Anisovich for useful discussions. The paper was partly supported by RFBR grant 98-02-17236.
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# 1 Introduction
## 1 Introduction
In superstring theories it turns out to be possible to lower the string scale without lowering the Planck scale . Most notably, Arkani-Hamed, Dimopoulos and Dvali have proposed the radical possibility that the fundamental scale of quantum gravity might no longer be associated with the Planck mass $`M_\mathrm{P}=1.2\times 10^{19}`$ GeV, but the true scale of quantum gravity, $`M_\mathrm{F}`$, could be many orders of magnitude smaller than $`M_\mathrm{P}`$, close to TeV energies. In such a novel theoretical framework, the standard-model (SM) particles can only live in a $`(1+3)`$-dimensional Minkowski subspace that constitutes our observable world, whereas gravity may freely propagate to a number $`n`$ of large extra dimensions. Furthermore, the ordinary Planck mass $`M_\mathrm{P}`$ would be related to the genuinely fundamental scale $`M_\mathrm{F}`$ through
$$M_\mathrm{P}M_F(RM_\mathrm{F})^{n/2},$$
(1.1)
where $`R`$ denotes the compactification radius, which is considered to be common for all extra compact dimensions. The case $`n=1`$ and $`M_\mathrm{F}`$ of order TeV leads to a visible macroscopic compactification radius and is therefore not viable. Moreover, astrophysical and cosmological considerations give rise to a lower limit on $`M_\mathrm{F}`$ of order 100 TeV, for the scenario with $`n=2`$ extra dimensions , while $`M_\mathrm{F}`$ can be as low as 1 TeV for theories with $`n>2`$ dimensions.
In addition to gravity, one might think that fields which are singlets under the Standard-Model gauge group could also propagate in the $`[1+(3+n)]`$-dimensional space. As such, one might consider isosinglet neutrinos or axion fields . In fact, within the context of theories of TeV-scale quantum gravity, the latter realization is theoretically compelling for the solution of the strong CP problem through the Peccei-Quinn (PQ) mechanism. According to this idea, the strong CP-odd parameter $`\theta `$ may be dynamically eliminated by means of the spontaneous breakdown of a global U(1) symmetry. On the other hand, phenomenological and astrophysical considerations place lower and upper limits on the breaking scale $`v_{\mathrm{PQ}}`$ of the PQ-U(1) symmetry, which has to be many orders of magnitude larger than the TeV scale of quantum gravity. Therefore, in order to account for this large mass scale, one inevitably has to introduce a singlet higher-dimensional axion field into the QCD Lagrangian, with a higher-dimensional PQ-breaking scale $`\overline{v}_{\mathrm{PQ}}`$ that could even be much smaller than 1 TeV. As we will see below, as a result of the compactification of the large extra dimensions, the effective four-dimensional PQ-breaking scale $`v_{\mathrm{PQ}}`$ can be obtained from $`\overline{v}_{\mathrm{PQ}}`$, after multiplying the latter by the huge higher-dimensional volume factor $`(M_\mathrm{F}R)^{n/2}M_\mathrm{P}/M_\mathrm{F}`$. In this way, the PQ-breaking scale $`v_{\mathrm{PQ}}`$ may reside in the phenomenologically allowed region. Another feature of the higher-dimensional axionic theories is that their mass spectrum consists of a tower of Kaluza-Klein (KK) excitations, which have an almost equidistant mass-spacing of order $`1/R`$. The lowest KK excitation may be identified with the ordinary PQ axion and specifies the strength of each KK state to matter.
This tower of axionic modes has two phenomenological consequences. First, for a fixed value of the axion coupling constant to matter or photons, a given source such as the Sun will emit axions of each mode up to the kinematic limit. The high multiplicity of the KK axion modes thus leads to much larger flux than would be otherwise expected. Second, the large mass of the KK modes compared to usual axions dramatically increases the width of the decay process $`a\gamma \gamma `$ by opening up phase space. Therefore, one may plausibly search for the decay photons of the solar KK axion flux in a laboratory experiment.
In this paper, we shall analyze the potential of a terrestrial axion detector to observe the radiative decay of solar KK axion modes. In particular, such a detector proves inexpensive and may run in parallel with the CERN Axion Solar Telescope (CAST) which will be built from a decommissioned LHC test magnet . We will find that the suggested terrestrial detector may reach the unprecedented sensitivity of the $`10^2`$-eV level ($`g_{a\gamma \gamma }10^{12}`$ GeV<sup>-1</sup>) to the fundamental PQ axion mass ($`m_{\mathrm{PQ}}`$).
The paper is organized as follows: in Section 2 we briefly describe the basic low-energy structure of a generic theory that includes higher-dimensional axions. In Section 3 we compute the solar flux of massive KK axions. In Section 4 we estimate the event rates of photons due to axion decays as seen by a terrestrial detector. Section 5 summarizes our conclusions.
## 2 Axions in large extra dimensions
Before discussing the higher-dimensional case, let us first recall the main phenomenological predictions of the axion theories in four dimensions. The axionic sector of the effective Lagrangian which is of interest to us has the generic form
$$_{\mathrm{eff}}=\frac{1}{2}(_\mu a)(^\mu a)\frac{1}{2}m_{\mathrm{PQ}}^2a^2+\frac{g_{a\gamma \gamma }}{4}aF_{\mu \nu }\stackrel{~}{F}^{\mu \nu },$$
(2.1)
where $`a`$ is the PQ axion, $`F_{\mu \nu }`$ and $`\stackrel{~}{F}^{\mu \nu }`$ are the electromagnetic field-strength tensor and its associate dual tensor, and
$$g_{a\gamma \gamma }=\frac{\xi \alpha _{\mathrm{em}}}{\pi }\frac{1}{v_{\mathrm{PQ}}}$$
(2.2)
is the effective axion-photon-photon coupling. The multiplicative parameter $`\xi `$ in Eq. (2.2) is generally of order unity, and crucially depends on the axion model under study . Furthermore, the PQ axion mass $`m_{\mathrm{PQ}}`$ is related to the breaking scale $`v_{\mathrm{PQ}}`$ of the PQ U(1) symmetry through
$$m_{\mathrm{PQ}}\frac{m_\pi ^2}{v_{\mathrm{PQ}}},$$
(2.3)
where $`m_\pi 135`$ MeV is the pion mass. Astrophysical and cosmological limits indicate that
$$10^9\mathrm{GeV}\stackrel{<}{_{}}v_{\mathrm{PQ}}\stackrel{<}{_{}}10^{12}\mathrm{GeV},$$
(2.4)
which in turn by virtue of Eq. (2.3) implies that
$$10^2\mathrm{eV}\stackrel{>}{_{}}m_{\mathrm{PQ}}\stackrel{>}{_{}}10^5\mathrm{eV},$$
(2.5)
respectively. The lifetime of the PQ axion is easily calculated to be
$$\tau (a\gamma \gamma )=\frac{64\pi }{g_{a\gamma \gamma }^2m_{\mathrm{PQ}}^3}10^{48}\mathrm{days}\times \left(\frac{10^{15}\mathrm{GeV}^1}{g_{a\gamma \gamma }}\right)^2\left(\frac{10^5\mathrm{eV}}{m_{\mathrm{PQ}}}\right)^3.$$
(2.6)
For $`g_{a\gamma \gamma }=10^{15}`$ GeV<sup>-1</sup>, which corresponds to $`m_{\mathrm{PQ}}=10^5`$ eV, the axion lifetime turns out to be much larger than the age of the universe. The prospect of detecting photonic axion decays would have remained hopeless, even if one had considered larger axion masses. For instance, for $`m_{\mathrm{PQ}}=10^1`$ eV ($`g_{a\gamma \gamma }=10^{11}`$ GeV<sup>-1</sup>), the axion decay is still undetectable with a lifetime $`\tau (a\gamma \gamma )10^{27}`$ days.
We shall now focus on the higher-dimensional case. Following Refs. , we introduce one singlet axion field $`a(x^\mu ,𝐲)`$ which feels the presence of a number $`\delta n`$ of large extra dimensions, denoted by $`𝐲=(y_1,y_2,\mathrm{},y_\delta )`$. The relevant axionic sector may then be determined by the effective Lagrangian
$$_{\mathrm{eff}}=d^\delta 𝐲\left[\frac{1}{2}M_F^\delta (_\mu a)(^\mu a)+\frac{1}{2}M_F^\delta (_\delta a)(^\delta a)+\delta ^{(\delta )}(𝐲)\frac{\xi \alpha _{\mathrm{em}}}{\pi }\frac{a}{\overline{v}_{\mathrm{PQ}}}F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }\right],$$
(2.7)
where $`\overline{v}_{\mathrm{PQ}}`$ denotes the original higher-dimensional PQ-breaking scale. In Eq. (2.7), the axion field is compactified on a Z<sub>2</sub> orbifold with an orbifold action : $`𝐲𝐲`$, i.e. the axion field satisfies the properties: $`a(x^\mu ,𝐲)=a(x^\mu ,𝐲+2\pi R)`$ and $`a(x^\mu ,𝐲)=a(x^\mu ,𝐲)`$. The latter gives rise to the KK decomposition:
$$a(x^\mu ,𝐲)=\underset{𝐧=\mathrm{𝟎}}{\overset{\mathrm{}}{}}a_𝐧(x^\mu )\mathrm{cos}\left(\frac{𝐧𝐲}{R}\right),$$
(2.8)
where $`𝐧=(n_1,n_2,\mathrm{},n_\delta )`$ is a $`\delta `$-dimensional vector that labels the individual KK excitations, and $`_{𝐧=\mathrm{𝟎}}^{\mathrm{}}_{n_1=0}^{\mathrm{}}_{n_2=0}^{\mathrm{}}\mathrm{}_{n_\delta =0}^{\mathrm{}}`$. Substituting Eq. (2.8) into Eq. (2.7) and taking the PQ mechanism into consideration, we arrive at the effective Lagrangian
$$_{\mathrm{eff}}=\frac{1}{2}\underset{𝐧=\mathrm{𝟎}}{\overset{\mathrm{}}{}}(_\mu a_𝐧)(^\mu a_𝐧)\frac{1}{2}m_{\mathrm{PQ}}^2a_0^2\frac{1}{2}\underset{𝐧\mathrm{𝟎}}{\overset{\mathrm{}}{}}\frac{𝐧^2}{R^2}a_𝐧^2+\frac{\xi \alpha _{\mathrm{em}}}{\pi }\underset{𝐧=\mathrm{𝟎}}{\overset{\mathrm{}}{}}\frac{r_𝐧a_𝐧}{v_{\mathrm{PQ}}}F_{\mu \nu }\stackrel{~}{F}^{\mu \nu },$$
(2.9)
with $`r_0=1`$ and $`r_{𝐧\mathrm{𝟎}}=\sqrt{2}`$. From Eq. (2.9), it is easy to read off the effective couplings of the KK axions to photons,
$$g_{a_𝐧\gamma \gamma }=\frac{r_𝐧\xi \alpha _{\mathrm{em}}}{\pi }\frac{1}{v_{\mathrm{PQ}}}g_{a\gamma \gamma }.$$
(2.10)
Instead of a Z<sub>2</sub> orbifold compactification, one could have equally considered the compactification on a $`\delta `$-dimensional torus , leading to modified coupling constants by factors of order unity. For the sake of simplicity we will always assume that the KK axion modes couple to photons with the usual PQ coupling $`g_{a\gamma \gamma }`$; it is trivial to insert model-dependent factors in the final result.
Few comments are now in order in connection with the effective KK Lagrangian (2.9). First, we should remark that the higher-dimensional PQ-breaking scale $`\overline{v}_{\mathrm{PQ}}`$ may be very low at the TeV scale, when compared to the usual four-dimensional one $`v_{\mathrm{PQ}}`$, i.e.
$$\overline{v}_{\mathrm{PQ}}\left(\frac{M_\mathrm{F}}{M_\mathrm{P}}\right)^{\delta /n}v_{\mathrm{PQ}}.$$
(2.11)
The suppression mechanism is very analogous to the case, in which the fundamental scale of quantum gravity can be reduced to the electroweak scale in the presence of large extra dimensions (cf. Eq. (1.1)). In Eq. (2.11), the simplest setting is to consider that both gravity and axions live within the same higher-dimensional space, i.e. $`\delta =n`$. Second, one notices that the lowest KK state constitutes the PQ axion of the theory which determines the size of the coupling of the KK axions to photons. Finally, the KK-axion masses are given by
$$m_{a_0}=m_{\mathrm{PQ}}\frac{1}{R},m_{a_𝐧}\frac{n}{R},$$
(2.12)
with $`n=|𝐧|=\sqrt{n_1^2+\mathrm{}+n_\delta ^2}>0`$. It is interesting to observe that for the higher-dimensional scenarios under discussion, the mass-spacing of the KK axions is always larger than PQ masses lying in the phenomenologically favoured region, with $`m_{\mathrm{PQ}}\stackrel{<}{_{}}0.01`$ eV. For example, for $`\delta =2`$ and $`M_\mathrm{F}100`$ TeV , one obtains $`1/R1`$ eV, while for $`\delta =3`$ and $`M_\mathrm{F}1`$ TeV, the inverse of the compactification radius reaches a much higher value, i.e. $`1/R10`$ eV.
The lifetime of an individual axionic KK state $`a_𝐧`$ may easily be computed from Eq. (2.6). In this way, we find
$$\tau (a_𝐧\gamma \gamma )\left(\frac{m_{\mathrm{PQ}}}{m_{a_𝐧}}\right)^3\tau (a_0\gamma \gamma ).$$
(2.13)
We observe that the lifetime of the KK axion $`a_𝐧`$ decreases rapidly with the third power of its mass. For example, the lifetime of one single (solar) KK-axion mode with $`m_{a_𝐧}=10`$ keV and $`g_{a\gamma \gamma }=10^{11}`$ GeV<sup>-1</sup> (corresponding to $`m_{\mathrm{PQ}}=10^1`$ eV) is $`\tau (a_𝐧\gamma \gamma )10^{12}`$ days, which is 15 orders of magnitude smaller than the respective one obtained in usual four-dimensional theories of PQ axions.
## 3 Solar flux of Kaluza-Klein axions
### 3.1 Primakoff process
In order to calculate the solar flux of KK axion modes we restrict ourselves to hadronic axion models where these particles do not couple to electrons at tree level. The dominant production processes will thus involve the axion-photon interaction; the axion-nucleon coupling will not be important in the Sun. The usual PQ axions are primarily produced by the Primakoff process $`\gamma +ZeZe+a`$ where a thermal photon in the solar interior converts into an axion in the Coulomb fields of nuclei and electrons in the solar plasma. In addition, the KK modes can be produced by the photon coalescence process $`\gamma \gamma a`$. For PQ axions, this process is suppressed by the small mass and actually is kinematically forbidden in the solar plasma because the effective photon mass (plasma frequency) is about 0.3 keV. However, with a temperature in the Sun of around 1.3 keV, the solar KK axions will be produced with masses up to several keV, rendering the coalescence process an important contribution.
Beginning with the Primakoff process, the production cross section on a target with charge $`Ze`$ in a nonrelativistic plasma is found to be
$$\frac{d\sigma _{\gamma a}}{d\mathrm{\Omega }}=\frac{g_{a\gamma \gamma }^2Z^2\alpha }{8\pi }\frac{|𝐤\times 𝐩|^2}{𝐪^4}\frac{𝐪^2}{𝐪^2+\kappa ^2},$$
(3.1)
where $`𝐤`$ is the photon momentum, $`𝐩`$ the axion momentum, and $`𝐪=𝐤𝐩`$ the momentum transfer. The last factor takes account of screening effects where the Debye-Hückel screening scale is given by
$$\kappa ^2=\frac{4\pi \alpha }{T}\frac{\rho }{m_u}\left(Y_e+\underset{j}{}Z_j^2Y_j\right).$$
(3.2)
In Eq. (3.2), $`\rho `$ is the mass density, $`m_u`$ the atomic mass unit (approximately the proton mass), $`Y_e`$ the number of electrons per baryon in the medium, and $`Y_j`$ the number of various nuclear species $`j`$ per baryon with nuclear charge $`Z_j`$. The medium is assumed to be nonrelativistic, and recoil effects by the targets have been neglected since typical photon energies of a few keV are much smaller than even the electron mass. It turns out that we have the approximate relation, $`\kappa 7T`$, between the screening scale $`\kappa `$ and the temperature $`T`$ in the relevant regions of the Sun.
Summing over all target species of the medium, the photon-axion transition rate is finally
$$\mathrm{\Gamma }_{\gamma a}=\frac{g_{a\gamma \gamma }^2T\kappa ^2}{32\pi ^2}\frac{|𝐤|}{\omega }𝑑\mathrm{\Omega }\frac{|𝐤\times 𝐩|^2}{𝐪^2(𝐪^2+\kappa ^2)},$$
(3.3)
where $`\omega `$ is the photon energy and the factor $`|𝐤|/\omega `$ is the relative velocity between photons and target particles. The angular integration can be performed explicitly, leading to
$`\mathrm{\Gamma }_{\gamma a}`$ $`=`$ $`{\displaystyle \frac{g_{a\gamma \gamma }^2T\kappa ^2}{32\pi }}{\displaystyle \frac{k}{\omega }}\{{\displaystyle \frac{\left[(k+p)^2+\kappa ^2\right]\left[(kp)^2+\kappa ^2\right]}{4kp\kappa ^2}}\mathrm{ln}\left[{\displaystyle \frac{(k+p)^2+\kappa ^2}{(kp)^2+\kappa ^2}}\right]`$ (3.4)
$`{\displaystyle \frac{(k^2p^2)^2}{4kp\kappa ^2}}\mathrm{ln}\left[{\displaystyle \frac{(k+p)^2}{(kp)^2}}\right]1\},`$
where $`k=|𝐤|`$ and $`p=|𝐩|`$.
The effective “photon mass” in the medium, the plasma frequency, is small in the Sun, typically about 0.3 keV, while the temperature near the solar center is $`T=1.3\mathrm{keV}`$ and typical photon energies are $`3T4\mathrm{keV}`$. Therefore, we ignore the plasma frequency and treat photons as strictly massless. In a photon-axion transition the energy is conserved because we ignore recoil effects. Therefore, we use $`k=E`$ with $`E`$ the axion energy and $`p=\sqrt{E^2m^2}`$ so that finally
$$\mathrm{\Gamma }_{\gamma a}=\frac{g_{a\gamma \gamma }^2T\kappa ^2}{32\pi }\left\{\frac{(m^2\kappa ^2)^2+4E^2\kappa ^2}{4Ep\kappa ^2}\mathrm{ln}\left[\frac{(E+p)^2+\kappa ^2}{(Ep)^2+\kappa ^2}\right]\frac{m^4}{4Ep\kappa ^2}\mathrm{ln}\left[\frac{(E+p)^2}{(Ep)^2}\right]1\right\}.$$
(3.5)
Note that the expression in curly brackets expands for small momenta as
$$\{\mathrm{}\}=\frac{8p^2}{3(\kappa ^2+m^2)}+𝒪(p^4),$$
(3.6)
so that the emission of slow-moving axions is suppressed.
The axion flux at Earth, differential with regard to the axion energy $`E`$, is then found by multiplying the transition rate with the blackbody photon flux in the Sun, and integrating over a standard solar model,
$$\mathrm{\Phi }_a=\frac{dF_a}{dE}=\frac{1}{4\pi d_{}^2}_{\mathrm{sun}}d^3𝐫\mathrm{\Gamma }_{\gamma a}\frac{1}{\pi ^2}\frac{E^2}{e^{E/T}1}.$$
(3.7)
Here $`T`$ and $`\kappa ^2`$ depend on the location in the Sun and $`d_{}=1.50\times 10^{13}\mathrm{cm}`$ is the distance to the Sun. We stress that no velocity factor appears for massive axions because in a stationary situation all axions produced per second must traverse a spherical shell around the Sun within one second.
In Ref. an approximation formula for the axion flux at Earth was given which we slightly modify and extend to the case of massive KK axions,
$$\mathrm{\Phi }_a=4.20\times 10^{10}\mathrm{cm}^2\mathrm{s}^1\mathrm{keV}^1\left(\frac{g_{a\gamma \gamma }}{10^{10}\mathrm{GeV}^1}\right)^2\frac{Ep^2}{e^{E/1.1}0.7}(1+0.02m),$$
(3.8)
where $`E`$, $`p`$ and $`m`$ are to be measured in keV. This approximation formula is typically good to better than $`\pm 15\%`$ for all relevant conditions, and even better than a few percent for the most relevant case of axion masses of larger than a few keV. In Fig. 1 we show the energy dependence of the flux of massive KK axions at Earth for three typical choices of the axion mass: $`m=5`$, 10 and 15 keV.
### 3.2 Photon Coalescence
In order to calculate the production rate of axions from the process $`\gamma \gamma a`$ in a thermal medium, we approximate the Bose-Einstein photon distribution by a Maxwell-Boltzmann one, i.e. we use $`e^{\omega /T}`$ instead of $`1/(e^{\omega /T}1)`$ for the photon occupation number. This approximation is justified since we are interested only in axion masses and thus axion energies of order the temperature or larger. The production rate of axions of energy $`E`$ per unit volume and unit energy interval is then found to be
$$\frac{dN_a}{dE}=\frac{g_{a\gamma \gamma }^2m^4}{128\pi ^3}pe^{E/T},$$
(3.9)
where again $`p=\sqrt{E^2m^2}`$ is the axion momentum. Integrating this expression over a standard solar model we find the axion flux at Earth. It is approximately represented by
$$\mathrm{\Phi }_a=1.68\times 10^9\mathrm{cm}^2\mathrm{s}^1\mathrm{keV}^1\left(\frac{g_{a\gamma \gamma }}{10^{10}\mathrm{GeV}^1}\right)^2m^4p\left(\frac{10}{0.2+E^2}+1+0.0006E^3\right)e^E,$$
(3.10)
where again $`m`$, $`E`$ and $`p`$ are to be taken in keV. For $`1\mathrm{keV}<E<16\mathrm{keV}`$ the quality of the approximation is better than 5%. Both lower and higher energies are irrelevant for our purposes. In Fig. 2 we display numerical estimates of the flux of KK axions at Earth as a function of their energy for three selected values of axion mass: $`m=5`$, 10 and 15 keV. A direct comparison of Fig. 1 with Fig. 2 reveals that the photon coalescence process becomes more important than the Primakoff one for the heavier KK-axion modes.
### 3.3 Axion limit from solar energy loss
As a next step we consider the energy loss of the Sun as a function of $`g_{a\gamma \gamma }`$. To this end we first calculate the solar axion luminosity as a function of the KK axion mass
$$L_a(m)=4\pi d_{}^2_m^{\mathrm{}}𝑑EE\mathrm{\Phi }_a(E)$$
(3.11)
for the two processes. Then we need to sum over all KK modes with their different masses. Instead, we integrate over the density of modes which is $`R^\delta `$, where $`R`$ is the compactification radius and $`\delta `$ the number of compactified dimensions. Therefore, the axion luminosity is
$$L_a=\frac{2\pi ^{\delta /2}}{\mathrm{\Gamma }(\delta /2)}R^\delta _0^{\mathrm{}}𝑑mm^{\delta 1}L_a(m),$$
(3.12)
where the first factor is the surface of the $`\delta `$ dimensional unit sphere, i.e. 2 for $`\delta =1`$, $`2\pi `$ for $`\delta =2`$ and $`4\pi `$ for $`\delta =3`$. Numerically, we write the result in the form
$$L_a=AL_{}\left(\frac{g_{a\gamma \gamma }}{10^{10}\mathrm{GeV}^1}\right)^2\left(\frac{R}{\mathrm{keV}^1}\right)^\delta $$
(3.13)
where $`L_{}`$ is the luminosity of the Sun and the values of the coefficients $`A`$ for the two processes and different dimensions $`\delta `$ are given in Table 1. It depends on $`\delta `$ which of the processes is more important.
Helioseismology implies that a novel energy-loss mechanism of the Sun must not exceed something like $`0.2L_{}`$ . This limit translates into the constraint
$$\left(\frac{g_{a\gamma \gamma }}{10^{10}\mathrm{GeV}^1}\right)\left(\frac{R}{\mathrm{keV}^1}\right)^{\delta /2}<\{\begin{array}{cc}3.3\hfill & \text{for }\delta =1\text{,}\hfill \\ 1.0\hfill & \text{for }\delta =2\text{,}\hfill \\ 0.31\hfill & \text{for }\delta =3\text{.}\hfill \end{array}$$
(3.14)
As an example we use the simplest setting of $`\delta =n=2`$ large extra dimensions, with $`M_\mathrm{F}=100`$ TeV and $`R=10^3`$ keV<sup>-1</sup>, leading to $`g_{a\gamma \gamma }<10^{13}\mathrm{GeV}^1`$. For $`\delta =n=3`$ large extra dimensions, with $`M_\mathrm{F}=1`$ TeV and $`R=10^2`$ keV<sup>-1</sup>, we get an even better limit of $`g_{a\gamma \gamma }<0.3\times 10^{13}\mathrm{GeV}^1`$. This is to be compared with the solar PQ axion limit of $`g_{a\gamma \gamma }<10^9\mathrm{GeV}^1`$ . Of course, the KK limits could have been estimated by simply scaling the standard limit with the multiplicity of KK modes and observing that the maximum allowed mass is a few keV before the solar flux gets suppressed by the kinematic threshold.
## 4 Flux of decay photons
### 4.1 Numerical estimates
The KK axions emerging from the Sun are neither nonrelativistic nor strongly relativistic. The average speed is 0.95 (in units of the speed of light) for $`m=1\mathrm{keV}`$, 0.79 for 3 keV, 0.66 for 5 keV, 0.57 for 7 keV, and 0.51 for 9 keV. Therefore, the decay photons will have a considerable angular spread relative to the direction of the Sun. The event rate in a detector thus depends crucially on its geometry. For our simple estimate we will assume that the detector consists of a volume $`V`$, and that any x-ray produced within this volume will be detected with unit efficiency, independently of its direction.
In view of the solar energy-loss limits derived above we further note that even keV-mass axions are long-lived relative to the Sun-Earth distance so that the axion flux on its way to Earth is not significantly diminished by radiative decays. Therefore, at any given time the total number of solar axions of mass $`m`$ per unit energy interval in the detector is
$$\frac{dN_a}{dE}=\frac{V\mathrm{\Phi }_a}{v}$$
(4.1)
where $`v=p/E`$ is the axion velocity. In the laboratory frame they decay with a rate $`(m/E)\mathrm{\Gamma }_{a\gamma \gamma }=(g_{a\gamma \gamma }^2/64\pi )m^4/E`$, each decay producing 2 photons with energies which are uniformly distributed in the range $`(Ep)/2\omega (E+p)/2`$. This implies that in order to get a decay photon of energy $`\omega `$ the parent axion must have $`E\omega +m^2/4\omega `$ and that the photon energies from a given axion decay are spread over an interval of length $`p`$. Altogether, then, we find for the differential event rate of decay photons from axions of mass $`m`$
$$\frac{dN_\gamma (m,\omega )}{d\omega }=\mathrm{\Gamma }_{a\gamma \gamma }mV_{\omega +m^2/4\omega }^{\mathrm{}}𝑑E\frac{2\mathrm{\Phi }_a}{p^2}.$$
(4.2)
Finally, in order to obtain the total event rate due to all modes of the tower of KK modes we proceed as before by integrating over the density of modes so that
$$\frac{dN_\gamma (\omega )}{d\omega }=\frac{2\pi ^{\delta /2}}{\mathrm{\Gamma }(\delta /2)}R^\delta V_0^{\mathrm{}}𝑑mm^\delta \mathrm{\Gamma }_{a\gamma \gamma }_{\omega +m^2/4\omega }^{\mathrm{}}𝑑E\frac{2\mathrm{\Phi }_a}{p^2}.$$
(4.3)
Numerically, we write this in the form
$$\frac{dN_\gamma (\omega )}{d\omega }=A_\delta \left(\frac{g_{a\gamma \gamma }}{10^{10}\mathrm{GeV}^1}\right)^4\left(\frac{R}{\mathrm{keV}^1}\right)^\delta \left(\frac{V}{\mathrm{m}^3}\right)f_\delta (\omega )$$
(4.4)
where $`A_\delta `$ is a rate given in Table 2 and $`f_\delta (\omega )`$ is a spectrum with its integral normalized to unity. These normalized functions are surprisingly well approximated by the simple analytic form
$$f_\delta (\omega )=a\omega ^be^{0.9\omega }$$
(4.5)
where $`a`$ and $`b`$ are given in Table 2 for each $`\delta `$. Of course, $`\omega `$ is understood in keV. In Table 2 we also give the average photon energies. In particular, as can be seen from Fig. 3, the energy distributions of the decay photons are shifted towards to the few keV energy range.
### 4.2 Experimental sensitivity
On the experimental side, we assume a 1 m<sup>3</sup> cubic detector of the Micromegas type. This is a new kind of gas detector which can be used to measure photon interactions with good space and energy resolution . A small detector of this kind, with a surface of $`15\times 15\mathrm{cm}^2`$, was used in Saclay (on the surface) and measured 1.2 neutral particles per second in a 1 keV wide energy interval centred at 1 keV. At these energies, practically all photons entering the chamber interact in the gas, so we have a measurement of the neutral particle flux through a surface of $`15\times 15\mathrm{cm}^2`$, which is about $`53\mathrm{neutral}\mathrm{particles}/\mathrm{m}^2/\mathrm{sec}`$ .
In the search for axion decays into two gammas, the background originates from two neutral particles interacting in the gas within the resolving time of the chamber. Therefore, one can choose the gas so that the mean absorption length of 1 keV photons is 0.3 cm . As a result, the interaction points of the two photons from axion decay will be very close to each other, in a cell with volume $`\mathrm{\Delta }x\mathrm{\Delta }y\mathrm{\Delta }z=1\mathrm{cm}^3`$. In the Micromegas chamber $`\mathrm{\Delta }x`$ and $`\mathrm{\Delta }y`$ are measured directly and $`\mathrm{\Delta }z`$ is measured from the time interval between the two signals. For $`\mathrm{\Delta }z=1\mathrm{cm}`$, the time interval is $`2\times 10^7`$ sec. Thus, in a small cell of 1 cm<sup>3</sup> volume, the rate of events from two uncorrelated neutral particles is $`5.6\times 10^{12}`$ events/sec. As there are $`10^6`$ cells when going from 1 cm<sup>3</sup> cell size to 1 m<sup>3</sup> size of the detector, the background rate becomes 0.5 events per day.
At this point, we should remark that we have not used two additional criteria to reduce the background:
* Real photons in the keV region entering the detector from outside will interact very close to the detector walls. If one requires that the events occur in a fiducial volume at some distance (few cm) from the walls, only photons generated inside the detector volume are important.
* If axions are non-relativistic, the two photons will have approximately equal energies.
Of course, a more precise estimate of the background requires a measurement with a realistic detector in the environment where the experiment is going to be performed.
Applying now Eq. (4.4) to the simplest setting of $`\delta =n=2`$ large extra dimensions, with $`M_\mathrm{F}=100`$ TeV and $`R=10^3`$ keV<sup>-1</sup>, we find the rate
$$R_\gamma 0.05\mathrm{events}\mathrm{day}^1\mathrm{m}^3\left(\frac{g_{a\gamma \gamma }}{10^{12}\mathrm{GeV}^1}\right)^4.$$
(4.6)
Consequently, the suggested terrestrial detector outlined above will be sensitive to an effective $`a\gamma \gamma `$-coupling $`g_{a\gamma \gamma }\stackrel{<}{_{}}2.\times 10^{12}`$ GeV<sup>-1</sup>, corresponding to a fundamental PQ mass $`m_{\mathrm{PQ}}10^2`$ eV. In particular, for $`\delta =n=3`$ large extra dimensions, with $`M_\mathrm{F}=1`$ TeV and $`R=10^2`$ keV<sup>-1</sup>, we obtain an estimate for the rate
$$R_\gamma 1.0\mathrm{events}\mathrm{day}^1\mathrm{m}^3\left(\frac{g_{a\gamma \gamma }}{10^{12}\mathrm{GeV}^1}\right)^4.$$
(4.7)
From this last result, one can readily see that the axion detector will be maximally sensitive to an effective $`a\gamma \gamma `$-coupling $`g_{a\gamma \gamma }\stackrel{<}{_{}}6.\times 10^{13}`$ GeV<sup>-1</sup>, corresponding to a fundamental PQ mass $`m_{\mathrm{PQ}}3.\times 10^3`$ eV.
Finally, it would be interesting to know whether measurements of $`\gamma `$-rays coming from the Sun could impose severe constraints on the $`2\gamma `$-decay mode of axions and hence on the parameters of the higher-dimensional axionic models under consideration . According to recent analyses , the solar x-ray luminosity in the range of interest to us, i.e. above 0.4 keV, is
$$L_{\mathrm{x}\mathrm{rays}}10^9\mathrm{events}/\mathrm{cm}^2/\mathrm{sec}10^{17}\mathrm{events}/\mathrm{day}/\mathrm{m}^2.$$
(4.8)
As the decay path available for solar axions is the distance to the Sun of $`1.5\times 10^{11}\mathrm{m}`$, the x-ray luminosity is by many orders of magnitude larger than the one expected from the decays of the KK axions.
### 4.3 Laboratory limits on solar axions
Lowest level underground experiments searching for weakly interacting massive particles (WIMPs) and other particles offer independent limits on the effective axion-to-photon coupling $`g_{a\gamma \gamma }`$. Specifically, these experiments report the following lower limit on the integrated event rate in the energy range below 10 keV:
$$R_\gamma ^{\mathrm{exp}}\stackrel{<}{_{}}20000\mathrm{events}\mathrm{day}^1\mathrm{m}^3.$$
(4.9)
These highly sensitive experiments measure the deposited energy but they are unable to distinguish between 1-prong and 2-prong events.
Applying Eq. (4.9) to Eqs. (4.6) and (4.7), we are able to derive for the first time experimental limits on $`g_{a\gamma \gamma }`$ in theories with KK axions. In this way, we find the upper limits
$$g_{a\gamma \gamma }\stackrel{<}{_{}}2.5\times 10^{11}\mathrm{GeV}^1,$$
(4.10)
for $`\delta =2`$ and $`M_F=100`$ TeV, and
$$g_{a\gamma \gamma }\stackrel{<}{_{}}1.2\times 10^{11}\mathrm{GeV}^1,$$
(4.11)
for $`\delta =3`$ and $`M_F=1`$ TeV. Evidently, our suggested underground detector will improve at least by one order of magnitude the present experimental limits which we derived in Eqs. (4.10) and (4.11). The latter upper limits should also be contrasted with the weaker upper bound: $`g_{a\gamma \gamma }\stackrel{<}{_{}}6.\times 10^{10}`$ GeV<sup>-1</sup>, which is obtained from recent experimental searches for conventional PQ axions coming from the Sun .
## 5 Conclusions
We have examined the potential of an underground detector shielded from cosmic-ray backgrounds for detecting KK axions coming from the Sun. The solar KK axions may be produced via the Primakoff process $`\gamma +ZeZe+a`$ or via the photon coalescence process $`\gamma \gamma a`$. In either case, we have calculated the expected flux of the KK axions, as well as estimated possible limits derived from helioseismology. We find that solar KK axions might lead to observable signatures in terrestrial experiments. In fact, the characteristic 2$`\gamma `$-decay mode of the KK axions offers a unique possibility to drastically reduce the cosmic background by coincidental triggering both of the emitted photons. Our elaborate estimates have shown that a terrestrial detector of 1 m<sup>3</sup> size may be sensitive to a fundamental PQ-axion mass up to $`10^2`$ eV, which amounts to having an effective axion-photon coupling $`g_{a\gamma \gamma }2.\times 10^{12}`$ GeV<sup>-1</sup>, in theories with 2 large extra dimensions and a fundamental quantum-gravity scale $`M_\mathrm{F}=100`$ TeV. In particular, in theories with 3 large compact dimensions with $`M_\mathrm{F}=1`$ TeV, the suggested detector is capable of probing PQ-axion masses up to $`3.\times 10^3`$ eV, corresponding to an effective axion-photon coupling $`g_{a\gamma \gamma }6.\times 10^{13}`$ GeV<sup>-1</sup>. Most importantly, the experimental detector under discussion will considerably improve, at least by one order of magnitude, the corresponding experimental limits on $`g_{a\gamma \gamma }`$ in theories with KK axions, which we derived in Eqs. (4.10) and (4.11) based on present data obtained from underground experiments.
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# Localization of eigenstates in a modified Tomonaga-Luttinger model
## I Introduction
The concept of Anderson localization in the Fock space of many-body systems has been recently adopted for the studies of interacting electrons in finite-size conductors. The traditional field of localization consists of the investigation of electronic motion in disordered solids. These studies were initiated by Anderson who considered tight-binding Hamiltonians on real-space lattices with random on-site energies. The behavior of the system was found to depend on the parameter $`Z=W/V`$, where $`W`$ is the typical variance of the random potential and $`V`$ is the hopping matrix element. For $`Z`$ above some critical value the eigenfunctions of the Hamiltonian are exponentially localized with a characteristic scale $`\xi `$, called the localization length, which leads to the absence of diffusion of electrons. The Anderson problem inspired numerous studies (for reviews see ). The scaling theory of localization predicted the strong dependence of the localization properties on the dimension of the lattice. It was found, for instance, that in one spatial dimension the states are always localized for any strength of the hopping.
The ideas of localization can be extended to the investigation of the properties of interacting many-body Hamiltonians in the basis of eigenstates of the corresponding free system. The central problem considered there is the mixing of a particular state representing one particle excited above the Fermi surface with many-particle states responsible for emergence of a quasi-particle described by the usual Landau Fermi liquid theory. For this purpose the analogy between the one-particle and the many-body problems was demonstrated in by mapping the Hilbert space of interacting electrons onto an effective tight-binding model. Each site of the effective lattice represents the eigenstate of the non-interacting system and the bonds represent the interaction matrix elements. After several approximations, which consist mainly in neglecting the possibility of closed loops on this effective lattice, the resulting model was identified with the Cayley tree with branching number $`K`$ depending on the energy $`ϵ`$ of the initial excitation, representing a particle excited above the Fermi sea. The Cayley tree model was previously studied and was found to undergo a localization transition for the value of parameter $`Z=Z_cK\mathrm{ln}K`$. In the context of the original many-electron problem, the system is in the localized regime of the many-particle states on the states of the non-interacting system for energy below the lower threshold $`ϵ^{}=\mathrm{\Delta }\sqrt{g/\mathrm{ln}g}`$, where $`\mathrm{\Delta }`$ is the one-particle level-spacing and $`g`$ is dimensionless conductance, i.e. conductance in units of $`e^2/\mathrm{}`$. It was also argued that until the excitation energy $`ϵ`$ reaches another critical value $`ϵ^{}=\mathrm{\Delta }\sqrt{g}>ϵ^{}`$ the states are not completely delocalized, i.e. they do not mix large number of non-interacting states, but rather a small portion of them, corresponding to a subtree of the whole lattice. The last result was not confirmed in the subsequent studies of this model by super-symmetry methods. It was found there that the states are strongly correlated above the threshold $`ϵ>ϵ^{}`$.
In more recent numerical studies of the interacting models it was observed that the choice of the type of the lattice affects the localization-delocalization transition point. The crossover between the spectral statistics and the properties of eigenstates such as inverse participation ratio were shown to be sensitive to the type of the lattice. It was doubted that the actual transition and the properties of states in the delocalized phase can be accounted for properly using the pure tree-like structure of the lattice. It was claimed in that the infinite tree-like lattice used in may be inadequate for description of the splitting of the quasi-particle peak into many-body constituents. The finiteness of the lattice used in this work results in the absence of the sharp transition from localized to delocalized states. In the numerical work , where the possibility of closed loops on the lattice was taken into account using the random-matrix theory no delocalization transition was observed. Rather the participation ratio and spectral width of the eigenstates of the interacting system exhibited a smooth crossover from almost localized to delocalized states.
The exact solution for the delocalized phase (which corresponds to the strongly interacting regime) is not available in these studies in any limit . It would be useful to have a model which can clarify the nature of the extended states in the strongly interacting regime. The candidates for such a system are of course one-dimensional interacting system whose analytical solution is available for any strength of the interactions. The present work was inspired by the observation that the eigenstates of the simplest interacting problem, namely the chiral Tomonaga-Luttinger model are extended in the effective lattice of the non-interacting Slater determinants. This model is usually solved by the bosonization technique, which can be viewed as a generalization of the solutions by Fourier transform of tight-binding models without disorder. This analogy can be traced more explicitly if one considers the Hilbert space spanned by the Slater determinants as an effective lattice where each site is labeled by partition of some integer. Due to the conservation of total momentum the Hamiltonian of the chiral Tomonaga-Luttinger model is block-diagonal in the total momentum index $`N`$ and the linear dispersion implies that energies of all the states inside one block are degenerate and equal $`N\mathrm{\Delta }`$. The interactions lift this degeneracy mixing the Slater determinants into coherent superpositions which are described by the bosonic quantum numbers. Thus the bosonic eigenstates of the full interacting model can be viewed as a generalization of the extended Bloch states in the tight-binding models with degenerate on-site energies.
The linearization of the dispersion is the approximation which allows to solve the problem of interacting electrons in one dimension. The deviations from the linearity mixes the bosonic states, so they cannot any longer provide the solution to the model. In this work we introduce the random dispersion into the Tomonaga-Luttinger model and study the properties of its eigenstates. It is expected to model a system where the single particle spectrum is complicated. The question we address is whether this model exhibits localization, since the randomness introduced into the dispersion leads to randomness of the diagonal matrix elements of the Hamiltonian analogous to the diagonal disorder introduced into tight-binding Anderson models. We have found that the states remain localized in the basis of Slater determinants, i.e. resemble the states of the non-interacting systems.
The outline of this paper is as follows. In section 2 we review some known facts about the standard Tomonaga-Luttinger model. We describe the structure of the Hilbert space as an effective lattice and discuss the bosonization solution in this context. In section 3 we introduce randomness into the dispersion law of the Tomonaga-Luttinger model and simplify the interaction matrix. Then we solve the resulting model with assumption of independence of the on-site energies. Section 4 is devoted to the study of relevance of the correlations between the energies. In the last section the results are discussed and further investigations are proposed.
## II Hilbert space of Tomonaga-Luttinger model
In this section we review a model exactly solvable for any interaction strength. The simplest model of interacting fermions is the one-branch (chiral) Tomonaga-Luttinger model . Let $`n`$ label the one-particle orbital with convention that in the ground state $`|0`$ of the non-interacting model, $`n=1`$ denotes the first unoccupied state. Since a very large number of states is occupied (the number of fermions is large) the number $`n`$ can be allowed to run from $`\mathrm{}`$ to $`+\mathrm{}`$. Another consequence is that the one-particle spectrum is assumed to be equidistant: $`E_n=n\mathrm{\Delta }`$. Both the linear dispersion and the fact that the spectrum is unbound are crucial for the solubility of the model. Introducing creation and annihilation operators $`c_n^{}`$, $`c_n`$ for a fermion on the orbital $`n`$, with usual anti-commutation relations $`\{c_n,c_n^{}^{}\}=\delta _{n,n^{}}`$ (other anti-commutators vanish) the Hamiltonian of the problem can be written as
$$\widehat{H}=\widehat{H}_0+\widehat{U}=E_0+\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}E_n:c_n^{}c_n:+\underset{m=1}{\overset{+\mathrm{}}{}}\underset{n,n^{}=\mathrm{}}{\overset{+\mathrm{}}{}}V_mc_{n+m}^{}c_nc_{n^{}m}^{}c_n^{}$$
(1)
The normal ordering of operators having infinite expectation values is defined as $`:\widehat{O}:=\widehat{O}0|\widehat{O}|0`$. As a result the constant $`E_0`$ absorbs the (infinite) free ground state energy together with the Hartree term:
$$E_0=0|\widehat{H}|0=\underset{n0}{}E_n0|c_n^{}c_n|0+\frac{V_0}{2}\underset{n,n^{}0}{}0|c_n^{}c_nc_n^{}^{}c_n^{}|0$$
(2)
The operator of “total momentum”
$$\widehat{P}=\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}n:c_n^{}c_n:$$
(3)
commutes with the full Hamiltonian. Due to the linear dispersion this operator is proportional to the free part of the Hamiltonian, i.e. the first term on the RHS of (1).
Instead of solving this model by the standard method of bosonization using identities among second quantized operators (see for example ) we prefer to reformulate our interacting problem in order to relate it to the problems usually discussed within the framework of Anderson localization. We start by labeling the eigenstates of non-interacting problem, which are the familiar Slater determinants, by shifts of particles with respect to the non-interacting ground state as shown in the Fig. 1: the uppermost particle is shifted up $`\lambda _1`$ levels, the second one is shifted up $`\lambda _2`$ levels etc. The numbers $`\lambda _i`$ are defined to satisfy $`\lambda _i\lambda _{i+1}`$ for each $`i`$. If a particle is not shifted it is not recorded. The set $`(\lambda )`$ of positive integers satisfying $`\lambda _i=N`$ forms the set of partitions of $`N`$ and is conveniently displayed by a Young diagram — array with $`\lambda _i`$ boxes in the $`i`$-th row. Following the standard notations in the mathematical literature the numbers $`\lambda _i`$ will be referred to as *parts* in what follows. For example, the partition of $`N=28`$, $`(\lambda )=(7,6,5,5,3,2)`$ sometimes written as $`(7,6,5^2,3,2)`$ that is displayed in Fig. 1 is represented graphically as
$$(\lambda )\text{ }$$
This pictorial representation of Slater determinants has nice properties under particle-hole interchange: recording shifts of the holes (which are moved downwards in energy) rather than particles in Fig. 1 yields the partition $`\stackrel{~}{(\lambda )}=(6,6,5,4,4,2,1)`$ whose graphical representation is just the transpose of the Young diagram for $`(\lambda )`$:
$$\stackrel{~}{(\lambda )}\text{ }$$
Partition $`\stackrel{~}{(\lambda )}`$ is said to be conjugate or dual to $`(\lambda )`$. This duality results from the fact that each move of a particle upwards is accompanied by a move of a hole downwards.
The eigenstates of the non-interacting Hamiltonian labeled by partitions of $`N`$ are degenerate and their energy is given by
$$\underset{i}{}\left(E_{1i+\lambda _i}E_{1i}\right)=\mathrm{\Delta }\underset{i}{}\lambda _i=N\mathrm{\Delta }$$
(4)
The interactions conserve this number. The Hamiltonian (1) is therefore block-diagonal and does not mix the partitions of different $`N`$ so we can restrict ourselves to the subspace of the Hilbert space spanned by partitions of $`N`$.
The dimensionality of this subspace is given by $`p(N)`$ — the number of partitions of $`N`$. There is no closed analytical expression for this number but the Euler expression for the generating function
$$p(q)=\underset{N=1}{\overset{\mathrm{}}{}}p(N)q^N=\underset{l=1}{\overset{\mathrm{}}{}}\frac{1}{1q^l}$$
(5)
allows to obtain the asymptotic behavior of $`p(N)`$ for large $`N`$ and it is given by the formula of Hardy-Ramunajan:
$$p(N)\frac{e^{\pi \sqrt{2N/3}}}{4\sqrt{3}N}$$
(6)
It demonstrates that the number of “sites” in the subspace increases rapidly with $`N`$.
Since for a given subspace $`N`$ the free part of the Hamiltonian is proportional to the unit matrix the solution consists in diagonalization of the interacting part $`U`$. The interaction term in (1) is conveniently represented in our picture as effective binding between the sites. In order to present the results in a compact and transparent way we drop the requirement on a partition labeling the basis state to be represented by a non-increasing sequence. Instead we extend the notion to any finite sequence of integers (either positive or negative or zero) and use the following rules to relate this sequence to the standard (non-increasing) partition:
1. In any sequence $`(\lambda _1,\lambda _2,\mathrm{})`$ two consecutive parts may be interchanged provided that the preceding part is decreased by unity and the succeeding part increased by unity, the state vector which corresponds to this partition acquires a minus sign, i.e.
$$|(\lambda _1,\mathrm{},\lambda _i,\lambda _{i+1},\mathrm{})=|(\lambda _1,\mathrm{},\lambda _{i+1}1,\lambda _i+1,\mathrm{})$$
(7)
2. If any part exceeds by unity the preceding part the partition corresponds to zero vector, i.e.
$$|(\lambda _1,\mathrm{},\lambda _i,\lambda _i+1,\mathrm{})=0$$
(8)
3. If the last part of the sequence is negative the partition corresponds to the zero vector.
These rules, used in the theory of symmetric functions are in fact nothing but a standard properties of Slater determinants under action of second quantized operators. Rule 1 results from the fact that $`i`$-th particle is excited to the level $`n_i`$ such that
$$n_i=\lambda _i+1i$$
with the energy $`E_i=n_i\mathrm{\Delta }`$ while the $`(i+1)`$-th particle is excited to the level
$$n_{i+1}=\lambda _{i+1}i$$
and interchanging consecutive rows in a (Slater) determinant results in a sign change. Rule 2 results from the fact that $`\lambda _{i+1}=\lambda _i+1`$ implies $`n_{i+1}=n_i`$ which contradicts the Pauli principle. Rule 3 reflects the fact that negative $`\lambda _i`$ describe (forbidden) transitions to occupied states inside the Fermi sea.
Let us denote by $`\eta _{(\lambda )}`$ the factors $`1,0`$ or $`+1`$ acquired when bringing a partition $`(\lambda )`$ to its standard form. In these notations the matrix element of the interactions between the partitions $`(\lambda )`$ and $`(\mu )`$ can be represented as
$$U_{(\lambda ),(\mu )}=(\lambda )|\widehat{U}|(\mu )=\underset{m}{}V_m\underset{i,j}{}(\mathrm{}\lambda _im\mathrm{})|(\mathrm{}\mu _jm\mathrm{}).$$
(9)
Since the basis of Slater determinants is orthonormal the following expression is obtained
$$U_{(\lambda ),(\mu )}=\underset{m}{}V_m\underset{i,j}{}\eta _{(\mathrm{}\lambda _im\mathrm{})}\eta _{(\mathrm{}\mu _jm\mathrm{})}\delta _{(\mathrm{}\lambda _im\mathrm{}),(\mathrm{}\mu _jm\mathrm{})}.$$
(10)
The eigenstates of $`U`$ are
$$|(l)=\underset{(\lambda )}{}\chi _{(l)}^{(\lambda )}|(\lambda )$$
(11)
where $`\chi _{(l)}^{(\lambda )}`$ are the characters of the irreducible representation $`(\lambda )`$ of an element of the symmetric group $`S_N`$ belonging to the class $`(l)`$. For this purpose it is shown in Appendix A that
$`{\displaystyle \underset{i,j}{}}\eta _{(\mathrm{}\lambda _im\mathrm{})}\eta _{(\mathrm{}\lambda _im\mathrm{}\lambda _j+m\mathrm{})}\chi _{(l)}^{(\mathrm{}\lambda _im\mathrm{}\lambda _j+m\mathrm{})}=mn_m\chi _{(l)}^{(\lambda )}`$ (12)
where $`n_m`$ is the number of times the cycle $`m`$ appears in the class $`(l)`$. This leads to
$$\widehat{U}|(l)=\underset{m}{}V_mmn_m\underset{(\lambda )}{}\chi _{(l)}^{(\lambda )}|(\lambda )=\left(\underset{m}{}V_mmn_m\right)|(l).$$
(13)
The state $`|(l)`$ is therefore an eigenstate of the Hamiltonian (1) with the eigenvalue
$$E_{(l)}=E_0+N\mathrm{\Delta }+\underset{m}{}V_mmn_m.$$
(14)
The orthogonality of the characters implies that the states $`|(l)`$ form an orthogonal basis, namely:
$$(l)|(l^{})=\frac{g}{g_{(l)}}\delta _{(l),(l^{})}$$
(15)
where $`g=N!`$ is the total number of elements in $`S_N`$ and the number of elements in the class $`(l)`$ is
$$g_{(l)}=N!\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m^{n_m}n_m!}.$$
(16)
We normalize the states $`|(l)`$ as
$$|(l)=\sqrt{\frac{g_{(l)}}{g}}\underset{(\lambda )}{}\chi _{(l)}^{(\lambda )}|(\lambda )$$
(17)
so that they form an orthonormal basis. These are standard bosonic states usually labeled by the occupation numbers $`n_m`$ that are the eigenstates of the interacting Hamiltonian (1). The transformation with the matrix
$$(\lambda )|(l)=\sqrt{\frac{g_{(l)}}{g}}\chi _{(l)}^{(\lambda )}$$
(18)
is the unitary transformation of the basis vectors in our Hilbert space, which diagonalizes the interacting Hamiltonian (1). When the free part of the Hamiltonian is degenerate in the subspace $`N`$ the states $`|(l)`$ play the role of the extended states (in the Hilbert space). In order to verify this statement we calculated numerically the inverse participation ratio $`P`$, defined as
$$P=\underset{(\lambda )}{}|(\lambda )|(l)|^4$$
(19)
for some arbitrary states $`(l)`$. Although no closed expression is available for the general amplitude (18) of the state $`|(l)`$ on the site $`(\lambda )`$, the beautiful recurrent relation of the characters of the symmetric group (A14) (based on the Frobenius formula) was implemented numerically to calculate these amplitudes. The inverse participation ratio averaged over the bosonic states $`(l)`$ is shown in Fig. 2 for different values of $`N`$. As the number of sites $`N_{site}=p(N)`$ increases the inverse participation ratio $`P`$ was found to decay as $`P=1/N_{site}^\alpha `$, where $`\alpha 0.8`$. This allows us to conclude that the bosonic states $`(l)`$ are indeed extended in the basis of free Slater determinants $`(\lambda )`$.
## III Simplified Interactions
The existence of extended states in the chiral Tomonaga-Luttinger model can be attributed to the degeneracy of the diagonal elements of the Hamiltonian (1), since the interactions (10) have short-range in the space of partitions. In order to define more precisely what we mean by short-range in our Hilbert space we have to order the partitions in some order. Then the interaction can be shown to connect only blocks of elements that are not too far one from the other.
It is natural to order the partitions in the lexicographic order, i.e. if $`\lambda _i=1,2,3\mathrm{}`$ numbers the letters in an alphabet the lexicographic order of partitions corresponds to the order of words in Arabic or Hebrew dictionary. For example, the 11 partitions of $`N=6`$ in the lexicographic order are:
$$\begin{array}{cc}k=1& \hfill (1,1,1,1,1,1)\\ & \\ k=2& \hfill (2,1,1,1,1)\\ & \\ k=3& \hfill (3,1,1,1)\\ & \\ k=4& \hfill (2,2,1,1)\\ & \hfill (4,1,1)\\ & \\ k=5& \hfill (3,2,1)\\ & \hfill (5,1)\\ & \\ k=6& \hfill (2,2,2)\\ & \hfill (4,2)\\ & \hfill (3,3)\\ & \hfill (6)\end{array}$$
(20)
The partitions are grouped by $`k`$—the number of 1’s (that appear always on the right). The $`k`$-th group contains partitions of the form
$$(\lambda _1,\lambda _2,\mathrm{}\lambda _l,1^{Nk})$$
(21)
with constraint $`\lambda _i2`$, with the only exception for the first group containing $`(1^N)`$. The typical Young diagram of a partition in the $`k`$-th group is of the form
(22)
$`Nk\{\begin{array}{c}\mathrm{}\hfill \\ \text{ }\hfill \end{array}`$ (25)
To calculate the number of partitions in the $`k`$-th group we note that exactly $`p(k)`$ partitions have at least $`Nk`$ ones. Among them $`p(k1)`$ have at least $`N(k1)`$ ones and the rest $`n(k)=p(k)p(k1)`$ have *exactly* $`Nk`$ ones. Therefore the $`k`$-th group contains $`n(k)`$ partitions.
Usually the matrix element of the interaction potential $`V_m`$ is a decreasing function of the momentum transfer $`m`$, with some effective range $`m_c`$. In what follows it will be convenient to analyze a model where the only non-zero matrix element corresponds to $`m=1`$ with $`V_1=V`$. It will not affect qualitatively our results, since we are primarily interested in the limit of large $`N`$ subspace. The matrix element of this interaction is given by the simplified version of (10)
$$U_{(\lambda ),(\mu )}=V\underset{i,j}{}\delta _{(\mathrm{}\lambda _i1\mathrm{}),(\mathrm{}\mu _j1\mathrm{})}.$$
(26)
The factors $`\eta `$ are omitted since the partitions $`(\mathrm{}\lambda _i1\mathrm{})`$ and $`(\mathrm{}\mu _j1\mathrm{})`$ are either already in the standard (non-increasing) form or corresponds to zero vector and do not contribute to the sum. Therefore the interaction couples partitions $`(\lambda )`$ and $`(\mu )`$ that differ in two of the $`\lambda _i`$’s, namely $`\lambda _i\mu _i=1`$ and $`\lambda _j\mu _j=1`$. The Young diagram of the partition $`(\lambda )`$ is obtained from the one of the partition $`(\mu )`$ by moving one square under the constraint that the new Young diagram is legitimate. Consider the $`k`$-th group. If $`\mu _i>1`$ and $`\mu _j>2`$ the coupling is between partitions inside this group. This corresponds to moving a square between $`\mu _i>1`$ rows of the Young diagram. In order to couple to neighboring groups a square corresponding to $`\mu _i=1`$ should be moved. Moving a square from/to any row with $`\mu _j>1`$ to/from the region where $`\mu _i=1`$ leads to a coupling to the neighboring group $`k\pm 1`$. The exception is if a square is moved from $`\mu _i=2`$ (that has only two squares and after the move only one is left) and is attached as the last square in the tail of $`\mu _j=1`$, or if the last square is moved to $`\mu _j=1`$ turning it into a $`\mu _j=2`$ row. In this case the length of the tail of $`\mu _j=1`$ is changed by two, i.e. the coupling is to the group $`k\pm 2`$. These are the only moves of one square, therefore only the groups $`k^{}`$ such that $`|kk^{}|2`$ are coupled to the $`k`$-th group and the interaction between the various blocks of partitions is of short range. Fig. 3 shows the structure of the interaction matrix for $`N=14`$.
It is therefore natural to expect that the diagonal disorder will destroy the extended states of the chiral Tomonaga-Luttinger model. In the original model in the absence of interactions the eigenenergies $`E_{(\lambda )}`$ are independent of the partition and take the constant value $`E_{(\lambda )}=E_0+N\mathrm{\Delta }`$. We introduce a modification and assume that the $`E_{(\lambda )}`$ are independent random variables. This assumption is made for the sake of simplicity. It cannot be correct since there are $`p(N)`$ values of $`E_{(\lambda )}`$ that depend on $`2N`$ level spacings. An assumption that is more physical will be introduced in the following section.
In order to be able to solve the model with random energies we simplify further the interaction (26). From the experience in research of Anderson localization we expect that the behavior of the eigenvectors will not be affected by the precise form of the short-range off-diagonal matrix elements, so we make the following simplifications. We assume that the effect of transitions inside the same group of partitions is to delocalize the states within the group. The localization within a group is expected only to enhance the overall localization. Next, we neglect the coupling to the $`k+2`$-th group. The reason is that for a given partition $`(\lambda )`$ there exists only one coupling term of this kind. It couples
$$((\mathrm{}),1^{Nk})\text{and}((\mathrm{})^{},2,1^{Nk2})$$
(27)
This should be compared it with the couplings of $`(\lambda )`$ to the next group:
$$((\mathrm{}),1^{Nk})((\mathrm{})^{},1^{Nk1})$$
(28)
The number of such couplings is given by the number of different $`\lambda _i2`$ in the partition $`(\lambda )`$, which is generally much larger than 1 for large enough $`k`$. It is equal to the number of rows with $`\lambda _i2`$ in the Young diagram of the partition. Consider for example the partition $`(6,4^2,3,1^4)`$ of $`N=21`$ belonging to he the 17-th group corresponding to the Young diagram (25). It has (apart of ones) three different parts: 6,4,3 and is coupled to 3 partitions in the 18-th group:
$$\begin{array}{ccc}\text{ }\hfill & \text{ }\hfill & \text{ }\hfill \\ & & \\ (7,4^2,3,1^3)\hfill & (6,5,4,3,1^3)\hfill & (6,4^3,1^3)\hfill \end{array}$$
(29)
Another way to visualize the number of coupling is to count the number of concave corners in the Young diagram that is equal to the number of parts namely lines of different length larger than unity.
In Appendix B we show that $`d(k)`$ — the number of different parts in all partitions of $`k`$ with $`\lambda _i2`$ is exactly $`p(k2)`$ — the number of (any) partitions of $`(k2)`$. Therefore the mean number of transitions from the $`k`$-th group to the $`k+1`$-th group is given by $`t(k)=d(k)/n(k)=p(k2)/\left(p(k)p(k1)\right)`$.
The Schrödinger equation for our simplified model can be written as a tight-binding form where the groups of partitions are the effective sites:
$$(EE_{(\lambda )})\psi (\lambda )_k=\underset{(\mu )}{}U_{(\lambda ),(\mu )}\psi (\mu )_{k1}+\underset{(\nu )}{}U_{(\lambda ),(\nu )}\psi (\nu )_{k+1}$$
(30)
where $`\psi (\lambda )_k`$ denotes the value of the wave function on site$`(\lambda )`$ in the $`k`$-th group. Let us assume that the wave function is extended within the $`k`$th group. We introduce a simpler model that is expected to exhibit similar localization properties . We replace the interaction $`U_{(\lambda ),(\mu )}`$ by an average interaction that couples each partition in the $`k`$-th group to all partitions in the group $`k\pm 1`$,
$$V_k=\frac{_{(\nu )}U_{(\lambda ),(\nu )}}{n(k+1)}$$
(31)
In this approximation (30) reduces to
$$(EE_{(\lambda )})\psi (\lambda )_k=V_{k1}\underset{(\mu )}{}\psi (\mu )_{k1}+V_k\underset{(\nu )}{}\psi (\nu )_{k+1}$$
(32)
The numerator in (31) can be estimated as $`V`$ times the mean number of couplings with partitions in $`(k+1)`$-th group:
$$V_k=V\frac{t(k)}{n(k+1)}=V\frac{d(k)}{n(k)n(k+1)}$$
(33)
Denoting the sum of the wave-functions over a group by $`C_k=_{(\lambda )}\psi (\lambda )_k`$ we find with the help of (32) that
$$_kC_k=V_{k1}C_{k1}+V_kC_{k+1}$$
(34)
where
$$\frac{1}{_k}=\underset{(\lambda )}{}\frac{1}{EE_{(\lambda )}}$$
(35)
and the sum is over all partitions in the group $`k`$.
For random independent $`E_{(\lambda )}`$ all the $`_k`$ are independent random variables and the tight-binding model (34) is an Anderson model at zero energy. In the spirit of transfer matrix method it is expected that also for $`_k`$ given by (35) similar localization properties will be found. The properties of the solution of this one-dimensional model can be analyzed within the framework of Anderson localization in one-dimensional systems .
One can take the uniform distribution of energies that turns out to be convenient for calculations. Remembering that each many-body energy is a sum of $`N`$ level spacings we assume the following distribution:
$$P(E)=\{\begin{array}{cc}1/\sqrt{N}\mathrm{\Delta },& \sqrt{N}\mathrm{\Delta }/2<E<\sqrt{N}\mathrm{\Delta }/2\\ 0,& \text{otherwise}\end{array}$$
(36)
For a variety of distributions the distribution of $`_k`$ tends to the Cauchy distribution:
$$P(_k)=\frac{1}{\pi }\frac{\delta _k(E)}{_k^2+\delta _k(E)^2}$$
(37)
where $`\delta _k(E)=\delta (E)/n(k)`$ and for (36) $`\delta (E)=\sqrt{N}\mathrm{\Delta }/\pi `$.
Multiplying the equation (34) by $`n(k)n(k+1)/(Vd(k))`$ and observing that in the limit $`k\mathrm{}`$ one finds $`n(k)/n(k+1)1`$ and $`d(k)/d(k1)1`$ so that (34) takes the form
$$_kC_k=C_{k1}+C_{k+1}$$
(38)
It is a Lloyd model with a $`k`$-dependent distribution of on-site energies characterized by parameter
$$\delta _k(E)=\frac{\delta (E)}{V}\frac{n(k+1)}{d(k)}$$
(39)
which in the limit $`k\mathrm{}`$ has the property
$$\delta _k(E)\sqrt{\frac{N}{6k}}\frac{\mathrm{\Delta }}{V}$$
(40)
We generalize the solution of the Lloyd model to the present case of site-dependent distribution of energies. We observe that the distribution parameter depends weakly on the site number $`k`$ and use the site-dependent localization length ansatz:
$$1/\xi _k=\mathrm{ln}\frac{|C_k|}{|C_{k+1}|}=\mathrm{cosh}^1\sqrt{1+\frac{\delta _k^2}{4}}\mathrm{ln}\delta _k=\mathrm{ln}\sqrt{\frac{N}{6k}}\frac{\mathrm{\Delta }}{V}$$
(41)
The last approximation is justified in the weak-coupling limit $`V/\mathrm{\Delta }1`$.
The validity of the site-dependent localization length assumption is checked numerically. The on-site energies $`_k`$ were generated from the Cauchy distribution (37) with parameter $`\delta _k`$ given by (40). The inverse localization length $`1/\xi `$ (Lyapunov exponent) was calculated using the standard technique of transfer matrices and it was found that the site-dependent ansatz (41) describes well (at least qualitatively) the behavior of the localization length. Thus for independently distributed energies the typical behavior of the wave functions is described by coefficients that decay as
$$C_ke^{k/\xi _k}=\mathrm{exp}\left(k\mathrm{log}\sqrt{\frac{N}{6k}}\frac{\mathrm{\Delta }}{V}\right).$$
(42)
In the strong coupling limit similar considerations lead to
$$C_ke^{k/\xi _k}=\mathrm{exp}\left(k\sqrt{\frac{N}{24k}}\frac{\mathrm{\Delta }}{V}\right),$$
(43)
but in this limit these require some further justification. The importance of this result is that numerical studies of more realistic many-particle energy distributions confirm the validity of (42). This will be the subject of the next section.
## IV Effects of energy correlations
In the last section we considered a model in which the diagonal matrix elements $`E_{(\lambda )}`$ are independent random variables. We did not specify the distribution law of these energies, since the solution does not depend on fine details, but rather on general properties of the distribution. We would like to stress that the assumption of statistical independence of the energies $`E_{(\lambda )}`$ is not realistic for a many-body system. To illustrate this statement a simple argument can be given. Suppose we are dealing with a disordered mesoscopic system and the one-particle spectrum is random, i.e. consists of levels with independently distributed spacings $`\mathrm{\Delta }_j=E_jE_{j1}`$. The generalization of the expression (4) for the energy of the partition $`(\lambda )`$ for this case is
$$E_{(\lambda )}=\underset{i}{}\left(E_{1i+\lambda _i}E_{1i}\right)=\underset{i}{}\underset{j=2i}{\overset{\lambda _i+1i}{}}\mathrm{\Delta }_j$$
(44)
In order to describe the energies $`E_{(\lambda )}`$ in the subspace $`N`$ we need $`2N`$ random quantities $`\mathrm{\Delta }_j`$. On the other hand the number of states in this subspace is given by $`p(N)`$, the number of partitions of $`N`$ which grows exponentially for $`N\mathrm{}`$. Therefore there exist much more many-body energies $`E_{(\lambda )}`$ than independent level spacings. Therefore the many-particle energies $`E_{(\lambda )}`$, which are the sums of the one-particle energies, are statistically dependent.
The analytical calculation of the (joint) distribution of the parameters $`_k`$ defined in (35) for the effective tight-binding model (34) seems to be hopeless. Nevertheless we expect some features of this distribution, based on the formula (35). For an eigenstate the eigenvalue $`E`$ is in the interval $`[1+E_{(\lambda )}^{min},+1+E_{(\lambda )}^{max}]`$, where $`E_{(\lambda )}^{min}`$ and $`E_{(\lambda )}^{max}`$ are the minimal and maximal values of $`E_{(\lambda )}`$. The value of $`E`$ can be very close one of the $`E_{(\lambda )}`$. Therefore the terms in the sum on the RHS of (35) can be positive or negative and they strongly fluctuate in magnitude. Therefore there is a finite probability that this sum will take a value in the vicinity of zero. Consequently it is expected that the distribution $`P(_k)`$ is broad with a diverging second moment or variance:
$$\mathrm{\Delta }_k^2=(_k_k)^2\mathrm{}$$
(45)
Another important issue is the effective statistical independence of $`_k`$ for different $`k`$-s. It is known for Anderson localization that the long-range correlations between the on-site energies change dramatically the criteria for the onset of the localization (see and references therein). The measure of statistical dependence of energies in our case is the correlation matrix defined as
$$C_{kl}=\frac{(_k_k)(_l_l)}{\mathrm{\Delta }_k\mathrm{\Delta }_l}$$
(46)
Because of the strong fluctuations between the terms of the sum on the RHS of (35), that are different for the various groups, it is reasonable that the off-diagonal terms of the correlation function are much smaller than the diagonal ones. Usually for localization only pair correlations are important . In the view of the diverging variance we assume that $`C_{kl}=\delta _{kl}`$, an assumption that will be tested numerically.
These facts were checked numerically and in the approximation of independent energies the effective distribution $`P(_k)`$ can be calculated. For this purpose $`M=1000`$ realizations of $`N=30`$ level spacings $`\mathrm{\Delta }_j`$ were generated from the exponential distribution:
$$P(\mathrm{\Delta })=e^\mathrm{\Delta }$$
(47)
with $`\overline{\mathrm{\Delta }}=1`$. Then we generated numerically the whole list of $`p(30)=5604`$ partitions of $`N=30`$ and for each realization of level spacings the energies $`E_{(\lambda )}`$ were calculated using (44). Out of these energies the effective energies $`_k`$ were obtained using the definition (35) for $`E=0`$. Several statistical tests were performed and results are shown in Fig. 4-5.
The correlation matrix of the energies $`_k`$ was computed. The $`\mathrm{}`$ averages were performed over realizations of the level spacings. The result is shown in Fig. 4. It is plausible that the correlations are of short range in the space of indices $`k`$ and only the diagonal elements of $`C_{kl}`$ are of appreciable magnitude. This behavior can be attributed to the diverging variance.
Due to the fact that the distribution of the energies $`_k`$ is expected to have a diverging variance the distribution $`P(_k)`$ cannot be calculated from a histogram. Instead its characteristic function
$$X_k(p)=e^{ip_k}$$
(48)
was computed and plotted for several values of $`k`$ in Fig. 5 as a function of $`p`$. The diverging second moment due to the fact that $`P(x)1/x^2`$ manifests itself in the discontinuity of the derivative at $`p=0`$ as can be seen from the Fig 6. Therefore we have strong evidence that the energies $`_k`$ have a broad distribution.
Observing the behavior of the characteristic function shown on the Fig. 5 and Fig. 6 we introduced a scaling ansatz for the distribution function:
$$P(_k)=P_k()=\frac{n(k)}{\mathrm{\Delta }}\stackrel{~}{P}\left(n(k)\frac{}{\mathrm{\Delta }}\right).$$
(49)
The universal function $`\stackrel{~}{P}(x)`$ is independent of $`k`$. The scaling of the distribution implies that the characteristic function (48) satisfies
$$X_k(p)=\stackrel{~}{X}\left(\frac{p\mathrm{\Delta }}{n(k)}\right),\stackrel{~}{X}(q)=_{\mathrm{}}^{\mathrm{}}𝑑x\stackrel{~}{P}(x)e^{iqx}$$
(50)
which is easy to check by plotting $`X_k(p)`$ as a function of $`q=p\mathrm{\Delta }/n(k)`$ for different values of $`k`$. These curves are displayed in Fig. 7 and are indeed found to be close to each other in spite of the fact that $`n(k)`$ ranges from $`n(1)=1`$ to $`n(30)=1039`$.
We found numerically that the universal characteristic function can be well approximated by the general form
$$\stackrel{~}{X}(q)=\mathrm{exp}\left(a|q|\frac{bq^2}{2}\right)$$
(51)
The value of parameters were found to be $`a1.15`$, $`b6.32`$. This is the characteristic function of a distribution of a sum of two random variables $`x_1`$ and $`x_2`$, where $`x_1`$ is drawn from Cauchy distribution with probability density
$$p_1(x)=\frac{1}{\pi }\frac{a}{x^2+a^2}$$
(52)
and the second is distributed normally according to
$$p_2(x)=\frac{1}{\sqrt{2\pi b}}e^{\frac{x^2}{2b}}.$$
(53)
The distribution of $`x=x_1+x_2`$ is thus given by
$$\stackrel{~}{P}(x)=_{\mathrm{}}^+\mathrm{}𝑑yp_1(xy)p_2(y)$$
(54)
and its characteristic function is (51). It has a diverging variance due to the fact that $`\stackrel{~}{P}(x)`$ behaves as $`\frac{1}{\pi }\frac{a}{x^2}`$ for large $`|x|`$.
The scaling factor $`\mathrm{\Delta }/n(k)`$ and the limiting behavior of the distribution function are the same as those found in the case of mutually independent energies. This allows us to conclude that the behavior of the coefficients $`C_k`$ given by (42), that was obtained in the model with uncorrelated $`E_{(\lambda )}`$, is satisfied also by the eigenfunctions of the more realistic model of the present section.
## V Discussion
The localization properties of the eigenstates of the modified interacting chiral Tomonaga-Luttinger model were studied. For total energy $`N\mathrm{\Delta }`$ subspace we defined an effective lattice of Slater determinants labeled by partitions of $`N`$ which represent the eigenstates of the system in the absence of interactions. We have observed that in the case of linear dispersion the eigenstates of the interacting model described by the bosons are extended over this lattice. When the random one-particle levels are introduced instead of the equidistant spectrum, the bosons are no more the eigenstates and the problem of diagonalizing of the Hamiltonian. In order to achieve this goal we have grouped the different sites (partitions) according to the number of ones, representing the distance of the most excited particle from the Fermi surface in the original Slater determinant. This is to be contrasted with the grouping of many-body states into generations according to the number of excited particles and holes as was firstly done in . Each resulting group of partitions was given a label $`k`$. The interactions were shown to be of short range in $`k`$. The interactions were further simplified by assuming that they interconnect only groups of partitions corresponding to the adjacent values of the index $`k`$ (and the interactions between groups $`k`$ and $`k\pm 2`$ were neglected). Assuming further that the states are extended within each group $`k`$ we obtained an effective tight-binding Anderson model where each group $`k`$ has an effective random energy and the hopping takes place between the nearest neighbors only. In the assumption of independently distributed on-site energies the model is found to be close to the Lloyd model characterized by a broad distribution of the on-site energies. The main feature of our effective model is that the distribution of energies is $`k`$-dependent. We employed the site-dependent localization length assumption, which consists in assuming the parameters of distribution change slowly, so locally the localization length is given by the same formula as for the Lloyd model with parameters of the local distribution instead of the uniform global one. Up to an unknown numerical factor this solution describes well the behavior of the localization length as a function of $`k`$ and it enabled us to show that the eigenstates are localized in the space of groups of partitions, $`k`$. The weak-coupling and the strong-coupling regimes can be identified according to the value of the parameter $`V/\mathrm{\Delta }`$. In the weak-coupling regime $`V\mathrm{\Delta }`$, the amplitudes decay as $`\mathrm{exp}\left(k\mathrm{log}\sqrt{\frac{N}{6k}}\frac{\mathrm{\Delta }}{V}\right)`$, namely nearly exponentially. In the strong coupling regime $`V\mathrm{\Delta }`$, this behavior changes to $`\mathrm{exp}\left(\sqrt{\frac{Nk}{24}}\frac{\mathrm{\Delta }}{V}\right)`$.
One important point which was discussed is the effect of correlations between the many-body energies of the non-interacting system. In the studies devoted to the localization in many-body systems the effective on-site energies corresponding to the eigenstates of the non-interacting system are assumed to be independent random variables, so their spectrum is described by a Poissonian sequence. This is not the case for the true many-body fermionic system, since these energies are sums of independent energies of one-particle excitations. The number of independent random variables is therefore much less than the number of different many-body states. This situation is similar in some sense to one faced in random walk problems, where the increments are distributed independently, while there is a strong correlation between the positions of the random walk at different times. In our work we considered this problem numerically and found that although the correlation between the exact energies corresponding to the eigenstates of the free model are not a priori negligible the random energies of the effective sites $`k`$ are found to be very weakly correlated. Moreover it was found that the parameters of the distribution of the effective energy of the group $`k`$ scale with the number of the partitions inside the group in the same way as they do in the case of independent many-body energies. This led us to the conclusion that our result for the localization of the eigenstates are valid for more realistic assumptions on the distribution of the many-body energies.
The main open problem arising from the present discussion is the study of the interacting electrons in dimensions higher than one by the methods of high-dimensional bosonization. In this technique the fermionic system is described as an infinite set of one-dimensional Luttinger models attached to a point on the Fermi surface. Therefore the interactions mix the states belonging to different points in addition to the coupling within the same point (represented by a Luttinger model), which can be treated within the framework developed in the present paper. This possibly can be a mechanism for the delocalization of the states and a resulting localization transition must be compared with the transition encountered in the recent studies of interacting electrons in finite systems.
Another question to be explored in the future is the nature of localization in the Hilbert space for single-particle dispersions that are more realistic than the random ones.
## Acknowledgements
The authors would like to thank Alex Kamenev for useful remarks and Eric Akkermans for stimulating discussions in the initial stages of the project. This research was supported in part by the U.S.–Israel Binational Science Foundation (BSF), by the Minerva Center for Non-linear Physics of Complex Systems, by the Israel Science Foundation, by the Niedersachsen Ministry of Science (Germany) and by the Fund for Promotion of Research at the Technion.
## A A property for characters of $`S_N`$
In this appendix we show that the property (12)
$`{\displaystyle \underset{i,j}{}}\eta _{(\mathrm{}\lambda _im\mathrm{})}\eta _{(\mathrm{}\lambda _im\mathrm{}\lambda _j+m\mathrm{})}\chi _{(l)}^{(\mathrm{}\lambda _im\mathrm{}\lambda _j+m\mathrm{})}=mn_m\chi _{(l)}^{(\lambda )}`$
is satisfied for the characters of the symmetric group $`S_N`$. We did not find this property in any textbook on the representations of groups, but its validity follows from the known properties of the characters of $`S_N`$ found, for example, in and . We present the proof of this property here for the sake of completeness of the discussion.
The Frobenius formula as it is written in the classic textbooks on group theory states the following identity between two antisymmetric functions of $`N`$-dimensional vector $`z=(z_1,z_2,\mathrm{},z_N)`$:
$$S_{(l)}(z)D(z)=\underset{(\lambda )}{}\chi _{(l)}^{(\lambda )}\mathrm{\Psi }_{(\lambda )}(z)$$
(A1)
where $`\chi _{(l)}^{(\lambda )}`$ is the character of the class $`(l)`$ of the symmetric group $`S_N`$ in the irreducible representation $`(\lambda )`$. With the partition $`(\lambda )`$ one associates the totally antisymmetric function:
$$\mathrm{\Psi }_{(\lambda )}(z)\underset{P}{}\text{sgn}(P)z_{P(1)}^{\lambda _1+N1}z_{P(2)}^{\lambda _2+N2}\mathrm{}z_{P(N)}^{\lambda _N}$$
(A2)
where the sum runs over all the permutations of $`1,2,\mathrm{},N`$, the sign of permutation being $`\pm `$ for an even/odd permutation. This is, up to a normalization constant, a Slater determinant constructed out of one-particle wave-functions $`z^{n_i}=z^{\lambda _i+Ni}`$, where $`z=e^{i\theta }`$. The factor $`D(z)`$ in (A1) is the Vandermonde determinant:
$$D(z)=\left|\begin{array}{ccccc}1& 1& \mathrm{}& 1& \\ z_1& z_2& \mathrm{}& z_N& \\ z_1^2& z_2^2& \mathrm{}& z_N^2& \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ z_1^{N1}& z_2^{N1}& \mathrm{}& z_N^{N1}& \end{array}\right|=\underset{i<j}{}(z_iz_j)$$
(A3)
representing the ground state of $`N`$ fermions. Let $`S_m`$ be a power sum of variables $`z_j`$ defined as follows:
$$S_m(z)=\underset{j}{}z_j^m$$
(A4)
The function $`S_m(z)`$ is totally symmetric function of the $`z_j`$-s. These functions are linearly independent as can be checked by calculating the Jacobian:
$$\left|\frac{S_m}{z_j}\right|=N!D(z).$$
(A5)
The class $`(l)`$ is characterized by $`n_1,n_2,\mathrm{},n_m,\mathrm{}`$ cycles of length $`1,2,\mathrm{},m,\mathrm{}`$. For this class the totally symmetric function $`S_{(l)}`$ is defined as a product of power sums
$$S_{(l)}(z)=S_1^{n_1}S_2^{n_2}\mathrm{}S_N^{n_N}=\underset{m=1}{\overset{\mathrm{}}{}}S_m^{n_m}$$
(A6)
Since for different $`m`$-s the power sums $`Sm`$ are linearly independent so are the monomials $`S_{(l)}`$ for different partitions $`(l)`$.
The Frobenius formula (A1) can be inverted with help of the orthogonality relation of the characters:
$$\mathrm{\Psi }_{(\lambda )}(z)=\frac{1}{g}\underset{(l)}{}g_{(l)}\chi _{(l)}^{(\lambda )}S_{(l)}(z)D(z)$$
(A7)
where $`g=N!`$ is the order of the $`S_N`$ and $`g_{(l)}`$ is the number of elements of $`S_N`$ in the class $`(l)`$ given by (16).
We turn now to develop a recurrence relation between the characters. For this purpose it is instructive to consider the product of the antisymmetric function $`\mathrm{\Psi }_{(\mu )}(z)`$ and a power sum $`S_m(z)`$. Since the power sum is totally symmetric the summands can be rearranged in the order set by each permutation and one obtains
$`S_m(z)\mathrm{\Psi }_{(\mu )}(z)`$ $`=`$ $`{\displaystyle \underset{P}{}}\text{sgn}(P)\left({\displaystyle \underset{i}{}}z_{P(i)}^m\right)z_{P(1)}^{\mu _1+N1}z_{P(2)}^{\mu _2+N2}\mathrm{}z_{P(N)}^{\mu _N}=`$ (A8)
$`=`$ $`{\displaystyle \underset{i}{}}\eta _{(\mathrm{}\mu _i+m\mathrm{})}\mathrm{\Psi }_{(\mathrm{}\mu _i+m\mathrm{})}(z)`$ (A9)
where $`\eta _{(\mathrm{}\mu _i+m\mathrm{})}=1,0`$ or $`+1`$ is determined accordingly to the rules of Section 2.
Having established the relation between the Slater determinants $`\mathrm{\Psi }_{(\mu )}`$ for various partitions we turn to develop a recurrence relation between the characters. Consider class $`(l)`$ with at least one $`m`$-cycle. Let us denote by $`(l_m^{})`$ the class obtained by removing one $`m`$-cycle from $`(l)`$. It is then clear that
$$S_{(l)}=S_mS_{(l_m^{})}=\left(\underset{j}{}z_j^m\right)S_{(l_m^{})}$$
(A10)
The partition $`(l_m^{})`$ defines a class in the symmetry group $`S_{Nm}`$. The Frobenius formula (A1) for this class reads
$$S_{(l_m^{})}D(z)=\underset{(\mu ^{})}{}\chi _{(l_m^{})}^{(\mu ^{})}\mathrm{\Psi }_{(\mu ^{})}$$
(A11)
while for $`(l)`$ it can be rewritten with the help of (A9) in the form
$`S_{(l)}D(z)`$ $`=`$ $`S_mS_{(l_m^{})}D(z)={\displaystyle \underset{(\mu ^{})}{}}\chi _{(l_m^{})}^{(\mu ^{})}S_m\mathrm{\Psi }_{(\mu ^{})}=`$ (A12)
$`=`$ $`{\displaystyle \underset{(\mu ^{})}{}}\chi _{(l_m^{})}^{(\mu ^{})}{\displaystyle \underset{i}{}}\eta _{(\mathrm{}\mu _i^{}+m\mathrm{})}\mathrm{\Psi }_{(\mathrm{}\mu _i^{}+m\mathrm{})}`$ (A13)
Comparing the coefficients of the $`\mathrm{\Psi }_{(\mu ^{})}`$ in (A13) with those in (A1) we find the recurrence relation of the characters
$$\chi _{(l)}^{(\lambda )}=\underset{i}{}\eta _{(\mathrm{}\lambda _im\mathrm{})}\chi _{(l_m^{})}^{(\mathrm{}\lambda _im\mathrm{})}$$
(A14)
This relation yields the characters of $`S_N`$ in terms of the characters of the symmetric group $`S_{Nm}`$.
In order to derive (12) take some particular $`m`$ and let each partition $`(l)`$ contain $`m`$ as its part $`n_m^{(l)}`$ times. For each $`(\lambda )`$ consider the following sum over partitions:
$$F_{(\lambda )}(z)=\frac{1}{g}\underset{(l)}{}g_{(l)}mn_m^{(l)}\chi _{(l)}^{(\lambda )}S_{(l)}(z)$$
(A15)
If $`(l)`$ has no parts equal to $`m`$, $`n_m^{(l)}=0`$ and $`(l)`$ does not contribute to this sum. We want to factor out $`S_m(z)`$ from the sum and for this purpose we can write $`S_{(l)}(z)=S_m(z)S_{(l_m^{})}(z)`$, so that the sum is over $`(l_m^{})`$. With the help of (16) we see that the expansion coefficients satisfy
$`g_{(l)}mn_m/g=mn_m{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k^{n_k}n_k!}}={\displaystyle \frac{mn_m}{m^{n_m}n_m!}}{\displaystyle \underset{km}{}}{\displaystyle \frac{1}{k^{n_k}n_k!}}`$ (A16)
$`={\displaystyle \frac{1}{m^{n_m1}\left(n_m1\right)!}}{\displaystyle \underset{km}{}}{\displaystyle \frac{1}{k^{n_k}n_k!}}=g_{(l_m^{})}/g^{}`$ (A17)
where $`g^{}=(Nm)!`$. The sum in (A15) takes the form
$$F_{(\lambda )}(z)=\left(\underset{j}{}z_j^m\right)\frac{1}{g^{}}\underset{(l_m^{})}{}g_{(l_m^{})}\chi _{(l)}^{(\lambda )}S_{(l_m^{})}(z)$$
(A18)
Application of the recurrence relation (A14) to the characters in the sum together with the inverse Frobenius formula (A7) yields
$`F_{(\lambda )}(z)D(z)`$ $`=`$ $`S_m(z){\displaystyle \underset{i}{}}\eta _{(\mathrm{}\lambda _im\mathrm{})}{\displaystyle \frac{1}{g^{}}}{\displaystyle \underset{(l_m^{})}{}}g_{(l_m^{})}\chi _{(l_m^{})}^{(\mathrm{}\lambda _im\mathrm{})}S_{(l_m^{})}(z)D(z)=`$ (A19)
$`=`$ $`{\displaystyle \underset{i}{}}\eta _{(\mathrm{}\lambda _im\mathrm{})}S_m(z)\mathrm{\Psi }_{(\mathrm{}\lambda _im\mathrm{})}(z)`$ (A20)
The relation (A9) is applied to show that
$$F_{(\lambda )}(z)D(z)=\underset{i,j}{}\eta _{(\mathrm{}\lambda _im\mathrm{})}\eta _{(\mathrm{}\lambda _im\mathrm{}\lambda _j+m\mathrm{})}\mathrm{\Psi }_{(\mathrm{}\lambda _im\mathrm{}\lambda _j+m\mathrm{})}(z)$$
(A21)
Application of the inverse Frobenius formula (A7) yields
$$F_{(\lambda )}(z)=\frac{1}{g}\underset{ij}{}\underset{(l)}{}\eta _{(\mathrm{}\lambda _im\mathrm{})}\eta _{(\mathrm{}\lambda _im\mathrm{}\lambda _j+m\mathrm{})}g_{(l)}\chi _{(l)}^{(\mathrm{}\lambda _im\mathrm{}\lambda _j+m\mathrm{})}S_{(l)}(z)$$
(A22)
Making use of the linear independence of $`S_{(l)}`$ for different partitions $`(l)`$ and equating the coefficients of the $`S_{(l)}(z)`$ with the corresponding coefficients in (A15) results in the required relation (12).
## B Number of distinct parts of partitions.
In this appendix we show that $`d(N)`$ – the number of distinct parts in the partitions of $`N`$ with parts $`\lambda _i2`$ equals $`p(N2)`$ – the number of unrestricted partitions of $`N2`$. Together with the generating function (5) for the number of partitions let us introduce the generating function
$$\stackrel{~}{p}(z,q)=\underset{N=1}{\overset{\mathrm{}}{}}\underset{d=1}{\overset{\mathrm{}}{}}\stackrel{~}{p}(N,d)q^Nz^d$$
(B1)
for the number $`\stackrel{~}{p}(N,d)`$ of partitions of $`N`$ having entries that satisfy $`\lambda _i2`$ and having exactly $`d`$ distinct parts. For example for $`N=6`$ and $`d=2`$ this number is $`\stackrel{~}{p}(N,d)=1`$, corresponding to the partition $`(4,2)`$ as can be seen from the list (20). The number of distinct parts $`d(N)`$ in all partitions of $`N`$ with parts $`2`$ is then given by
$$d(N)=\underset{d=1}{\overset{N}{}}d\stackrel{~}{p}(N,d)$$
(B2)
It is generated by the following function
$$d(q)=\underset{N=1}{\overset{\mathrm{}}{}}d(N)q^N=z\frac{d}{dz}\stackrel{~}{p}(z,q)|_{z=1}$$
(B3)
The explicit expression for $`\stackrel{~}{p}(z,q)`$ can be found easily summing directly over all the appropriate partitions:
$$\stackrel{~}{p}(z,q)=\underset{(\lambda ),\lambda _i2}{}q^{|(\lambda )|}z^{d(\lambda )}$$
(B4)
where $`d(\lambda )`$ is the number of distinct parts in $`(\lambda )`$ and $`|(\lambda )|=_i\lambda _i`$. It is convenient to use the sets of “occupation numbers” $`\{n_m\}`$, where $`n_m`$ – the number of times $`m`$ appears in the partition $`(\lambda )`$. Noting that $`|(\lambda )|=_i\lambda _i=_mmn_m`$ and $`d(\lambda )=_m\theta (n_m)`$, where $`\theta (x)`$ is the Heavyside step function we rewrite (B4) as a sum over all the configurations $`\{n_m\}`$, with $`m2`$
$`\stackrel{~}{p}(z,q)`$ $`=`$ $`{\displaystyle \underset{\{n_m\},m2}{}}q^{_{m=2}^{\mathrm{}}mn_m}z^{_{m=2}^{\mathrm{}}\theta (n_m)}={\displaystyle \underset{m=2}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_m=0}{\overset{\mathrm{}}{}}}q^{mn_m}z^{\theta (n_m)}=`$ (B5)
$`=`$ $`{\displaystyle \underset{m=2}{\overset{\mathrm{}}{}}}\left(1+{\displaystyle \frac{zq^m}{1q^m}}\right)`$ (B6)
Differentiating this expression with respect to $`z`$ at $`z=1`$ we get
$$d(q)=z\frac{d}{dz}\underset{m=2}{\overset{\mathrm{}}{}}\left(1+\frac{zq^m}{1q^m}\right)|_{z=1}=q^2\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{1q^m}=q^2p(q)$$
(B7)
where $`p(q)`$ is given by (5). The function $`q^2p(q)`$ is the generating function for $`p(N2)`$. The equality of the generating functions $`d(q)`$ and $`q^2p(q)`$ implies the equality of the coefficients: $`d(N)=p(N2)`$, i.e. the number of different parts in all partitions of $`N`$ with entries $`2`$ is given by the number of unrestricted partitions of $`N2`$. Applying this argumentation to all partitions ending with $`Nk`$ ones yields $`d(k)=p(k2)`$.
## C Distribution of harmonic sum of independent random variables.
We consider the harmonic mean $`X`$ defined as
$$\frac{1}{X}=\frac{1}{N}\underset{n=1}{\overset{N}{}}\frac{1}{x_n}$$
(C1)
of $`N`$ independent random variables $`x_1,x_2,\mathrm{}x_N`$ drawn from the common distribution with probability density $`P(x)`$. The only requirement we impose on $`P(x)`$ is that
$$\underset{x0}{lim}P(x)=C>0$$
(C2)
Consider the random variable $`y=1/x`$. Its distribution function $`Q(y)`$ can be calculated by the standard methods. For $`y\pm \mathrm{}`$ it has the asymptotic behavior
$$Q(y)\frac{C}{y^2}$$
(C3)
and its second moment diverges. The general theory of broad distributions predicts that in the limit $`N\mathrm{}`$ the probability density of the sum (C1) $`Y=1/X`$ is Lorenzian
$$Q_N(Y)=\frac{1}{\pi }\frac{\pi C}{(\pi C)^2+Y^2}$$
(C4)
and consequently the probability density $`P_N(X))`$ of variable $`X=1/Y`$ is again a Lorenzian:
$$P_N(X)=\frac{1}{\pi }\frac{\delta }{\delta ^2+X^2}$$
(C5)
characterized by the parameter $`\delta `$ related to the original distribution $`P(x)`$ by the formula:
$$\delta =\frac{1}{\pi C}$$
(C6)
In particular for the uniform distribution (36) it takes the value $`\delta =\mathrm{\Delta }/\pi `$.
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# Teeny, tiny Dirac neutrino masses: an unorthodox point of view 1footnote 11footnote 1Talk presented at The Second Tropical Workshop on Particle Physics and Cosmology, Neutrino and Flavor Physics,1-6 May 2000, San Juan, Puerto Rico
## Why should one bother with a teeny, tiny Dirac neutrino mass?
A subtitle to this talk should perhaps go like “A see-saw-like mechanism without a Majorana mass”. Here, I shall try to present arguments as to why it is interesting and worthwhile to study scenarios in which neutrinos possess a mass which is pure Dirac in nature. Along the way, I shall try to argue that one should perhaps try to separate the issue of a see-saw like mechanism from that of a Majorana mass. By ’see-saw-like mechanism”, it is meant that a “tiny” mass arises due to the presence of a very large scale.
The suggestions that neutrinos do indeed possess a mass came from three different sources, all of which involve oscillations of one type of neutrino into another type. They are the SuperKamiokande atmospheric neutrino oscillation, the solar neutrino results, and the LSND result kayser . The present status of these three oscillation experiments is well presented in this workshop. The future confirmation of all three will certainly have a profound impact on the understanding of the origin of neutrino masses. In particular, it is now generally agreed that if there were only three light, active (i.e. electroweak non-singlet) neutrinos, one would not be able to explain all three oscillation phenomena. The confirmation of all three results would most likeky involve the presence of a sterile neutrino.
Whatever the future experiments might indicate, one thing is probably true: If neutrinos do have a mass, it is certainly tiny compared with all known fermion masses. Typically, $`m_\nu O(10^{11})`$(Electroweak Scale). Why is it so small? Is it a Dirac or a Majorana mass? This last question presently has no answer from any known experiment. The nature of the mass will no doubt have very important physical implications. The route to a gauge unification will certainly be very different in the two cases. Whether or not the mass is Dirac or Majorana, there is probably some new physics which is responsible for making it so tiny. What is the scale of this new physics? What are the possible mechanisms which could give rise to the tiny mass? In trying to answer these questions, one cannot help but realize that there is something very special about neutrinos (specifically the right-handed ones) which make them different from all other known fermions. Do they carry some special symmetry?
One example of new physics which might be responsible for a small neutrino mass is the ever-popular and beautiful see-saw mechanism of Gell-Mann, Ramond and Slansky seesaw , in which a Majorana mass arises through a lepton number violating process. Generically, one would have $`m_\nu m_{D\nu }^2/`$, with $`m_{D\nu }`$ Electroweak Scale, and $``$ some typical GUT scale. Since one expects $`m_{D\nu }`$, one automatically obtains a tiny Majorana neutrino mass. The actual detail of the neutrino mass matrix is however quite involved and usually depends on some kind of ansatz. But that is the same old story with any fermion mass problem anyway. The crucial point is the fact that the very smallness of the neutrino mass comes from the presumed existence of a very large scale $``$ compared with the electroweak scale. This mechanism has practically become a standard one for generating neutrino mass. Why then does one bother to look for an alternative?
First of all, there is so far no evidence that, if neutrinos do have a mass, it should be of a Majorana type. If anything, the present absence of neutrinoless double beta decay might indicate the contrary. (Strictly speaking, what it does is to set an upper limit on a Majorana mass of approximately 0.2 eV, although actually it is a bound on $`_iU_{ei}^2m_i)`$. Therefore, this question is entirely open. In the meantime, it is appropriate and important to consider scenarios in which neutrinos are pure Dirac. The questions are: How can one construct a tiny Dirac mass for the neutrinos? How natural can it be? Can one learn something new? Are there consequences that can be tested?
## A model of teeny, tiny Dirac neutrino mass
The construction of the model reported in this talk was based on two papers hung1 ; hung2 . There exists several other works Dirac on Dirac neutrino masses which are very different from hung1 ; hung2 . The first one hung1 laid the foundation of the model. The second one hung2 is a vastly improved and much more detailed version, with new results not reported in hung1 . In constructing this model, we followed the following self-imposed requirements:
1) The smallness of the Dirac neutrino mass should arise in a more or less natural way.
2) The model should have testable phenomenological consequences, other than just merely reproducing the neutrino mass pattern for the oscillation data.
3) One should ask oneself if one can learn, from the construction of the model, something more than just neutrino masses. This also means that one should go beyond the neutrino sector to include the charged lepton and the quark sectors as well. This last sentence refers to work in progress and will not be reported here.
### Description of the model
Before describing our model, let us briefly mention a few facts. First of all, it is rather easy to obtain a Dirac mass for the neutrino by simply adding a right-handed neutrino to the Standard Model. This right-handed neutrino (one for each generation) is an electroweak singlet and, as a result, can have a gauge-invariant Yukawa coupling: $`g_\nu \overline{l}_L\varphi \nu _R+h.c.`$. The Dirac neutrino mass would then be $`m_\nu =g_\nu \varphi `$. With $`\varphi 173GeV`$, a neutrino mass of O(1 eV) would require a Yukawa coupling $`g_\nu 10^{11}`$. Although there is nothing wrong with it, a coupling of that magnitude is normally considered to be extremely fine-tuned, if it is put in by hand! Could $`g_\nu 10^{11}`$ be dynamical? Would the limit $`g_\nu 0`$ lead to some new symmetry? What would it be? This new symmetry would be the one that protects the neutrino mass from being “large”.
In choosing such a symmetry, we followed our self-imposed requirement # 3: One should learn something more from it than just merely providing a symmetry to protect the neutrino mass. First, in order to implement the symmetry protection, one should assume that this new symmetry is particular to the neutrinos, in particular the right-handed ones since left-handed neutrinos are weak interaction partners of standard charged leptons. Therefore, it will be assumed that all fermions other $`\nu _R`$’s are singlets under this new symmetry.
One of the reasons we adhere to our requirement #3 is the wish to work from the bottom up, instead of from the top down. As a result, we would try to make every increased step in energy as meaningful as possible.
The symmetry chosen in hung1 ; hung2 is a chiral gauge symmetry. It is $`SU(2)_{\nu _R}`$, where the subscript $`\nu _R`$ means that only $`\nu _R`$’s carry $`SU(2)_{\nu _R}`$ quantum numbers. Why $`SU(2)_{\nu _R}`$? Because it is a chiral gauge $`SU(2)`$ which has a very important property: For Weyl fermions transforming as doublets under such a group, there exists an argument due to Witten witten that says, in a nutshell, that, because of the presence of a non-perturbative global anomaly the number of such Weyl doublets has to be even in order for the theory to be well-defined. (In the language of quantum field theory, this means that the generating functional should be non-vanishing.) Amusingly enough, a long-forgotten fact about the SM is related to the Witten anomaly. The absence of such anomaly for $`SU(2)_L`$ require an even number of electroweak doublets per family, which is the case there: one lepton and three quark doublets. (From a historical prespective, one might say that, had this constraint be known in the early seventies, the SM, as we now know it, would have had an extra strong argument in its favor, prior to the SLAC experiment.) In our case, let $`\nu _R`$ transform as doublets under $`SU(2)_{\nu _R}`$, i.e. we now have $`\eta _R=(\nu _R,\stackrel{~}{\nu }_R)`$. (Cosmological issues concerning $`\stackrel{~}{\nu }_R`$’s are discussed in hung2 .) The absence of the Witten anomaly then requires the number of $`\eta _R`$’s to be even. If furthermore, $`\eta _R`$’s carry some family indices (if a family symmetry exists) then this constraint can have a profound implication on the issue of family replication. Further remarks can be found in hung2 .
We know that families do mix. In consequence, we need some kind of family symmetry. The family symmetry chosen in hung1 ; hung2 is a gauge symmetry. This choice is pure prejudice: We believe that gauge theories are better choices for a family symmetry because of the fact that they do provide strong constraints on matter representations and because one might want to mimic the vertical symmetry (the electroweak interactions). We choose $`SO(N_f)`$ as our family gauge group with all fermions transforming as vector representions in order to avoid the usual traingle anomaly. Our model is given by the following extension of the SM:
$$SU(3)_cSU(2)_LU(1)_YSO(N_f)SU(2)_{\nu _R}$$
(1)
In hung1 ; hung2 , I have discussed the various arguments used to constrain $`N_f`$. In this talk, I shall however restrict $`N_f`$ to be $`N_f=4`$. This means that this is a four-family model. Is a 4th generation ruled out by experiment as one often hears? The answer is: Not at all! For instance, the usual question is the following: What about the Z width which tells us that there are only three light neutrinos? This does not apply to the case when the 4th neutrino is more massive than half the Z mass. Why then would it be so heavy when the other three neutrinos are so light? Isn’t it unnatural? The answer is NO as we shall see below. Then, what about the 4th generation quarks and charged leptons? There exists a review physrep dealing extensively with this question. A quick summary of that review is the statement that there is plenty of room for the discovery of the 4th generation, either at the next upgraded collider experiments at the Tevatron, or at the LHC. I shall now turn to the basic results of the model.
### Basic Reults of the model
I shall describe the results which are based solely on the assumption that only two oscillation results are correct: The solar and atmospheric neutrino oscillation data. I shall mention at the end a possibility in case all three oscillation experiments are confirmed.
In a nutshell, here are the results obtained in hung2 . 1) We obtain three light, near degenerate neutrinos. 2) The “tiny” masses are obtained dynamically at one loop. One will see below the reason for the use of the term “see-saw-like mechanism with Dirac mass”. 3) The masses of the light neutrinos, $`m_{\nu _i}`$ ($`i=1,2,3`$), and $`\mathrm{\Delta }m^2`$ are correlated in an interesting way: a) If the MSW solution is chosen for the solar neutrino problem, the masses can be as large as O(few eV’s) and can provide enough mass for the Hot Dark Matter (HDM); b) If the vacuum solution is chosen instead, the masses are found to be at most $``$ 0.1 eV and, as a result, are too small to be relevant to the HDM. 4) There are a number of phenomenological consequences which can be tested: There is no neutrinoless double beta decay since the mass is Dirac; There is a possibility of detection of “light” (a couple of hundreds of GeV’s) vector-like fermions; etc… 5) There are a number of possible cosmological consequences: Baryon asymmetry through neutrinogenesis with a pure Dirac neutrino mass; Perhaps some of the very heavy vector-like fermions could be the source of Ultra High Energy Cosmic Rays.
In writing down the Lagrangian for our model, we take into account the fact that our point here is to obtain a pure Dirac mass. Therefore, B-L will be assumed. The particle content is listed in the Table
and the Lagrangian is given by
$`_{Lepton}^Y`$ $`=`$ $`g_E\overline{l}_L^\alpha \varphi e_{\alpha R}+G_1\overline{l}_L^\alpha \mathrm{\Omega }_\alpha F_R+G_{M_1}\overline{F}_L\varphi _{1R}+G_{M_2}\overline{F}_L\stackrel{~}{\varphi }_{2R}+G_2\overline{}_{1L}\mathrm{\Omega }_\alpha e_R^\alpha +`$ (2)
$`G_3\overline{}_{2L}\rho _m^\alpha \eta _{\alpha R}^m+M_F\overline{F}_LF_R+M_1\overline{}_{1L}_{1R}+M_2\overline{}_{2L}_{2R}+h.c.`$
After integrating out the $`F`$, $`_1`$, and $`_2`$ fields, the relevant part of the effective Lagrangian below $`M_{F,1,2}`$ reads
$`_{Lepton}^{Y,eff}`$ $`=`$ $`g_E\overline{l}_L^\alpha \varphi e_{\alpha R}+G_E\overline{l}_L^\alpha (\mathrm{\Omega }_\alpha \varphi \mathrm{\Omega }^\beta )e_{\beta R}+`$ (3)
$`G_N\overline{l}_L^\alpha (\mathrm{\Omega }_\alpha \stackrel{~}{\varphi }\rho _i^\beta )\eta _{\beta R}^i+H.c.,`$
where
$$G_E=\frac{G_1G_{M_1}G_2}{M_FM_1};G_N=\frac{G_1G_{M_2}G_3}{M_FM_2}.$$
(4)
As one can see, all neutrinos are massless when $`SO(4)SU(2)_{\nu _R}`$ is unbroken. Assume $`<\mathrm{\Omega }>=(0,0,0,V)`$ and $`<\rho >=(0,0,0,V^{}s_1)`$, the 4th neutrino gets a mass $`m_N=\stackrel{~}{G_N}\frac{v}{\sqrt{2}}`$ with $`\stackrel{~}{G}_N=G_1G_{M_2}G_3\frac{VV^{}}{M_FM_2}`$. One can arrange the masses and couplings in such a way that $`\stackrel{~}{G}_N`$ O(1), and $`m_N`$ O(100 GeV). There is nothing unnatural about such a choice. It is natural in this scenario to have the 4th neutrino having a mass of O(100 GeV). One point which is worth emphasizing again is the following: The breaking of $`SU(2)_{\nu _R}`$ (in addition to $`SO(4)`$) is essential for $`m_N`$ to be non-vanishing! At this stage (tree-level), there are three massless neutrinos.
It turns out, at one loop, that the three formely-massless neutrinos acquire a common mass, i.e. they are degenerate:
$$\frac{m_\nu }{m_N}=\frac{\mathrm{sin}(2\beta )}{32\pi ^2}I(\frac{M_F}{M_2},\frac{M_F}{M_G},\frac{M_F}{M_P})$$
(5)
where $`M_{F,2,G,P}`$ are masses of particles which participate in the loop diagram and where $`\mathrm{tan}\beta =V^{}/V`$. This is shown in Fig. 1.
Here $`I(\frac{M_F}{M_2},\frac{M_F}{M_G},\frac{M_F}{M_P})`$ is given by
$`I({\displaystyle \frac{M_F}{M_2}},{\displaystyle \frac{M_F}{M_G}},{\displaystyle \frac{M_F}{M_P}})`$ $`=`$ $`{\displaystyle \frac{1}{M_FM_2}}\{{\displaystyle \frac{M_F[M_F^2(M_G^2\mathrm{ln}(\frac{M_G^2}{M_F^2})M_P^2\mathrm{ln}(\frac{M_P^2}{M_F^2}))+M_G^2M_P^2\mathrm{ln}(\frac{M_P^2}{M_G^2})]}{(M_G^2M_F^2)(M_P^2M_F^2)}}`$ (6)
$`(M_FM_2)\}.`$
One important remark is in order here. From Eq. (5), one notices that the neutrino mass does not depend explicitely on the value of the masses $`M_{F,2,G,P}`$ but only on their ratios. If one takes $`M_F`$ as a “base” mass for example, it turns out that one can obtain quite a small mass for the light neutrinos, $`m_\nu O(10^{11})m_N`$, as long as one has, e.g., $`M_2M_F`$, with $`M_F`$ being an arbitrary number which can be as low as experimentally allowed i.e. O(200 GeV) physrep . (Remember that $`F`$ stands for $`F=(F^0,F^{})`$, where $`F^0`$ and $`F^{}`$ are degenerate in mass.) This can be seen from Fig. 2.
Just to illustrate this point, a numerical example would be helpful. Take $`m_N=100`$ GeV, $`M_2/M_F=10^9`$, $`M_P/M_F=5`$, $`M_G/M_F=10^4`$, we obtain $`m_\nu =1.4`$ eV. I wish to emphasize that this is not a prediction. It will have to come from some deeper theory which will fix the above mass ratios and hence $`m_\nu /m_N`$. This is only meant to illustrate the fact that, in our model, it is quite natural to get a teeny, tiny Dirac neutrino mass. To go further, one needs to lift the degeneracy of the light neutrinos. Before showing how one would go about doing it, let us explain what we meant by “see-saw-like mechanism without a Majorana mass”.
The function $`I`$ shown in Eq. (6) has the following limit: $`IM_F/M_2`$ for $`M_2M_GM_P>M_F`$. From Eq. (5), one can see that, in the above limit, $`m_\nu \frac{m_NM_F}{M_2}\frac{\mathrm{sin}(2\beta )}{32\pi ^2}`$. If $`M_FO(m_N)`$ ($`M_FM_2`$), one would obtain
$$m_\nu \frac{m_N^2}{M_2}\frac{\mathrm{sin}(2\beta )}{32\pi ^2}.$$
(7)
This is a typical see-saw-like relation! With a “low” mass ($``$ Electroweak scale) vector-like fermion, $`F`$, one can qualitatively see that $`m_\nu `$ can be very small when $`M_2M_F`$. This behaviour is very reminescent of the see-saw mechanism, except that, in our case, the mass is Dirac.
Now the next step is to introduce mixing among the neutrinos in order to lift the mass degeneracy. This can be acomplished, in our model, by involving mixing in the scalar sector. At this stage, the degeneracy among the three light neutrinos is due to a remaining global $`SO(3)`$ symmetry. This remaining global symmetry can be explicitely broken by the scalar sector as shown in hung2 . The offshoot of all this is, in the end, the fact that this explicit breaking depends on a parameter denoted by $`b`$ in hung2 . It turns out that, in a paper under preparation hung3 , $`b`$ itself will be severely constrained when our model is extended to the quark sector. This is because the same scalar sector is also involved in the quarks. It is satisfying to see the link between the quark and lepton sectors. This is however not the subject of this talk and I shall now return to the task at hand.
The first case which was investigated in hung2 is when the scalar sector is written down in such a way that there is no mixing between the 4th neutrino and the other three in the mass matrix. Something interesting happens here. It turns out that the mass splittings are quasi-degenerate, in the sense that $`|m_2^2m_1^2||m_3^2m_2^2|`$. If this model were to explain both solar and atmospheric oscillation data, this quasi-degeneracy of $`\mathrm{\Delta }m^2`$ has to be lifted. Also, the fact that the oscillation data appear to show $`\mathrm{\Delta }m_{Solar}^2\mathrm{\Delta }m_{Atmospheric}^2`$ implies, in the context of our model, that indirectly the data suggests the existence of a 4th neutrino whose mixing with the lighter three will lift the quasi-degeneracy of $`\mathrm{\Delta }m^2`$. Before showing how this could be done, let us see what these results imply. If the vacuum solution for solar neutrinos is preferred, i.e. $`\mathrm{\Delta }m^210^{10}eV^2`$, then it is found that the median mass value of three almost degenerate neutrinos is $`\overline{m}_\nu 0.1eV`$. As stated earlier, this is not enough for the HDM scenario. If the MSW solution is preferred, i.e. $`\mathrm{\Delta }m^210^5eV^2`$, the median mass value could be $`\overline{m}_\nu `$ O(few eV’s), a reasonable value for the HDM scenario.
The fact that $`\mathrm{\Delta }m_{Solar}^2\mathrm{\Delta }m_{Atmospheric}^2`$ indicates, in the context of this model, the existence of more than three light neutrinos. In hung2 where only the atmospheric and solar data were taken into account, this means that it indicates the existence of a 4th neutrino. Again through the scalar sector, one can construct e.g. a mixing between the 3rd and 4th neutrino (other possibilities exist). The size of the mixing determines the correct mass splittings. It was found that there are strong constraints on the some of the scalar masses when one requires that $`|\mathrm{\Delta }m_{12}^2|10^5eV^2`$ and $`|\mathrm{\Delta }m_{23}^2|10^3eV^2`$.
The next question concerns the oscillation angles. To find out what they are, one needs to know the leptonic “CKM” matrix: $`V_L=U_l^{}U_\nu `$. In dealing with the neutrino sector of our model, we have presented a case where $`U_\nu `$ can be computed. It is basically given by:
$$U_\nu ^{(3)}=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}& 0\\ \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}& 0\\ 0& 0& 1\end{array}\right)$$
(8)
As for $`U_l`$ which requires a detailed study of the charged lepton sector, a construction is in progress. In the meantime, just for the purpose of illustration, Reference yanagida has been used in which a simple ansatz for the charged lepton sector was given. The reason for using this reference is because it contains an ansatz for $`U_\nu `$ which is similar to ours. Therefore the results should be similar: a small angle MSW solution and a large angle atmospheric solution.
## Epilogue
I have presented in this talk a model which can “naturally” give rise to a teeny, tiny Dirac neutrino mass, without resorting to the concept of a Majorana mass. What was shown was the need to differentiate between the see-saw mechanism and the existence of a Majorana mass. In this model, the smallness of the light neutrino mass arises in a see-saw-like fashion, with the mass being purely Dirac. As we have argued in the Introduction, the reason for constructing such a model is twofold: a) One does not know experimentally whether the mass is Majorana or Dirac; b) The physics is very different in the two cases. At this stage of our knowledge, it is perhaps prudent to explore all different possibilities. Since so much has already been worked out with models using a Majorana mass, any new model which takes a different route should have a clear motivation and predictable consequences. For our model, we have presented clearly our motivation: naturally small Dirac neutrino mass, family replication, etc..; and predictions concerning the neutrino sector: Vacuum solution $`\iff ̸`$ HDM, MSW solution $``$ HDM, $`\mathrm{\Delta }m_{Atmospheric}^2\mathrm{\Delta }m_{Solar}^2`$ as an indirect indication of a 4th heavy neutrino. Other phenomenological implications include:
1) There is no neutrinoless double beta decay because of the fact that we have a Dirac neutrino here.
2) The existence of “long-lived” and “light” (i.e. $`100GeV`$) vector-like leptons ($`F`$) whose detection might be possible at the LHC. A study of this kind of search can be found in a comprehensive review physrep . The quark counterparts should also be detectable hung3 .
3) The existence of several scalars with masses of order TeV’s.
There are several other phenomelogical issues to be discussed. For lack of space, a few of those will be briefly mentioned. One is the S parameter for example. It is well-known that, to leading order, vector-like fermions which carry electroweak $`SU(2)_L`$ quatum numbers do not contribute to S if one has a degenerate $`SU(2)`$ doublet. The reason for this being so is because the right-handed contribution cancels exactly the left-handed contribution. Therefore, to leading order, there is no constraint from the S parameter on the mass of the $`F`$-fermions. This point and other issues concerning quarks and leptons beyond the third generation are discussed in physrep . Issues such as the decay of the heavy 4th neutrino can be consulted in hung2 . Also, another issue such as the magnitude of flavor-changing neutral currents, e.g. $`\mu e\gamma `$, will be discussed in an upcoming paper dealing with the charged lepton sector. However, a preliminary statement can be made. For example, in the case of $`\mu e\gamma `$, there are two kinds of contributions: One coming from the propagation of neutrinos with a non-zero mass inside the loop diagram for the process, and the other one coming from diagrams involving the new vector-like fermions. It turns out that both contributions are negligible: 1) In the first case, it is because $`m_\nu M_W`$; 2) In the second case, it is because of the cancellations of the type described in hung2 .
As far as the cosmological implications are concerned, there are:
1) Can the fermion $`_2`$ be the source of Ultra High Energy Cosmics Rays (E $`>10^{11}GeV`$)? For example if $`M_F200GeV`$ then $`M_22\times 10^{11}GeV`$? Would the decays of $`_2`$ (e.g. $`_{2R}W_L^\pm F^{}`$ high energy quarks and leptons) be responsible for UHECR? This deserves a closer look.
2) The possibility of Baryon Asymmetry from neutrinogenesis. This is a scenario of Ref. lindner . The ingredients needed for such a scenario to work are basically: 1) a tiny, pure Dirac neutrino mass; 2) A decay process from some superheavy particles at the GUT (or similar) scale into right-handed neutrinos such that there is an asymmetry between right-handed neutrinos and anti-neutrinos; 3) a B+L violating process from the electroweak sphaleron. Since the Dirac neutrinos have a tiny Yukawa coupling (which is dynamical in our case), the part of B+L which is stored in the right-handed neutrinos, $`(B+L)_R`$, survived the sphaleron “washout”. So if one starts out with B-L=0, this process can generate a net baryon number, $`n_B=n_Ln_{\nu _R}`$.
Last but not least, a Dirac neutrino mass would certainly imply a different route to unification, diffrent from the popular scenario such as $`SO(10)`$.
One last remark is in order here. If all three oscillation experiments were to be confirmed in the future, there seems to be a need for a sterile neutrino. How will it fit in our framework? It turns out that some modifications of the previous analysis will be needed but the basic framework is still the same. This work is in progress.
I would like to thank Jose Nieves, Terry leung, Art Halprin and Qaisar Shafi for a wonderful workshop. This work is supported in parts by the US Department of Energy under grant No. DE-A505-89ER40518.
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# Constraints on Cold H₂ Clouds from Gravitational Microlensing Searches.
## 1 Introduction
The nature of the dark matter dominating the mass of the Galaxy remains elusive. Although nonbaryonic forms of dark matter are one possible explanation, it is possible that much or most of the mass may be baryonic. Recent results from gravitational microlensing experiments, such as MACHO and EROS (Alcock et al. 2000; Lasserre et al. 2000), have put severe constraints on the abundances of compact baryonic objects in the Galaxy, - such as brown dwarfs, old white dwarfs, neutron stars, or black hole remnants - thus virtually removing stellar objects as candidates for being the major constituent of the dark matter (Alcock et al. 2000; Freese, Fields, & Graff 2000; Lasserre et al. 2000).
It has recently been proposed that the dark matter in the Galaxy could consist of small, cold clouds of H<sub>2</sub> (Pfenniger, Combes, & Martinet 1994; De Paolis et al. 1995a, 1995b, 1996; Gerhard & Silk 1996). Walker & Wardle (1998) pointed out that such clouds could be responsible for the extreme scattering events (ESEs; Fiedler et al. 1994), which would occur when the line of sight to an extragalactic point radio source is crossed by the photoionized cloud envelope. In a later paper, Wardle & Walker (1999) considered the thermal stability of such clouds and showed that these clouds could be stable against heating by cosmic rays. These clouds could have masses of the order of Jupiter’s mass ($`10^3M_{}`$) and radii $`R_{cl}10`$ AU. If they exist, these clouds could also naturally explain the core radii of galaxies (Walker 1999) as well as the $`\gamma `$-ray emission from the Galactic halo (Kalberla, Shchekinov, & Dettmar 1999; Sciama 2000). Cloud-cloud collisions would steadily resupply the Galactic disk with gas to sustain steady star formation.
The proposed clouds are not compact enough to produce any gravitational microlensing (Henriksen & Widrow 1995) but Draine (1998) demonstrated that they can produce magnification of background stellar sources. Refraction of light passing through the clouds would cause amplification of background stars in a way resembling gravitational microlensing, which provides us with the possibility of using searches for gravitational microlensing events to either detect such gas clouds or constrain the properties of the cloud population.
In this paper we calculate the lensing rate of the stars in the Large Magellanic Cloud (LMC) by clouds of molecular hydrogen, taking into account the spatial distribution of the clouds and all the relevant motions which contribute to this rate: motion of the observer on the Earth, proper motion of the LMC, the velocity distribution of stars in the Cloud, and the velocity distribution of lensing clouds in the Galactic halo. We suppose that the dark matter is composed of clouds of only one size and mass and that the clouds are transparent.
The paper is organized as follows: in §2 we review some aspects of the lensing by gaseous clouds. In §3 we derive basic formulae for the rate of lensing events produced by the gaseous clouds and for timescale distribution of this sort of lensing. In §4 we present our results in the form of parametric plots, covering a wide range of cloud models, and compare our results with those obtained by collaborations undertaking searches for gravitational microlensing events. We determine the region of parameter space which is not excluded by these experiments and other constraints. Finally, in §5 we compare our results with the limits placed on cloud models by other authors from different arguments.
## 2 Physics of gaseous lensing
The physics of lensing by gaseous clouds has been discussed by Draine (1998) and we repeat here only the salient points. The refractive index $`m`$ is related to the gas density $`\rho `$ by
$$m(\lambda )=1+\alpha (\lambda )\rho ;$$
(1)
$`\alpha (4400\mathrm{\AA })=1.243\mathrm{cm}^3\mathrm{g}^1`$ and $`\alpha (6700\mathrm{\AA })=1.214\mathrm{cm}^3\mathrm{g}^1`$ for H<sub>2</sub>/He gas with 24% He by mass (AIP Handbook 1972).
For small deflections, a light ray with impact parameter $`b`$ will be deflected through an angle
$$\varphi (b)=2\alpha b\underset{b}{\overset{\mathrm{}}{}}\frac{dr}{(r^2b^2)^{1/2}}\frac{d\rho }{dr}.$$
(2)
Let $`D_{\mathrm{sl}}`$ and $`D_{\mathrm{ol}}`$ be the distance from source to lens, and from observer to lens (see Figure 1). If $`b_0`$ is the distance of the lens center from the straight line from source to observer, then the apparent distance $`b`$ of the image from the lens is given by the lensing equation
$$bb_0=D\varphi (b),D\frac{D_{\mathrm{sl}}D_{\mathrm{ol}}}{D_{\mathrm{sl}}+D_{\mathrm{ol}}}.$$
(3)
For a point source, the image magnification is given by
$$M(b)=\frac{|b|}{b_0}\frac{1}{1D\varphi ^{}(b)},$$
(4)
$$\varphi ^{}(b)\frac{d\varphi (b)}{db}=2\alpha \underset{0}{\overset{\mathrm{}}{}}𝑑z\left[\frac{b^2}{r^2}\frac{d^2\rho }{dr^2}+\frac{z^2}{r^3}\frac{d\rho }{dr}\right].$$
(5)
where $`r^2=b^2+z^2`$. For a given $`b_0`$ there will be an odd number $`N(b_0)`$ of solutions $`b_i(b_0)`$, $`i=1,\mathrm{},N`$. The total amplification $`A(b_0)=_{i=1}^NM(b_i)`$. The “trajectory” of the lens relative to the source is characterized by a “source impact parameter” $`p`$ and a displacement $`x`$ along the trajectory; for any $`x`$ we have $`b_0=(p^2+x^2)^{1/2}`$, and the “light curve” is just $`A(b_0)`$ vs. $`x`$.
Following Draine (1998) we define a dimensionless “strength” parameter
$$S\frac{\alpha \rho D}{R_{cl}}=0.355\left(\frac{M_{cl}}{10^3M_{}}\right)\left(\frac{\mathrm{AU}}{R_{cl}}\right)^4\left(\frac{D}{10\mathrm{kpc}}\right),$$
(6)
where $`M_{cl},R_{cl}`$, and $`\rho 3M_{cl}/4\pi R_{cl}^3`$ are the cloud mass, radius, and mean density.
As the pressure density-relation in such clouds is uncertain, Draine (1998) considered polytropic equations of state for polytropic index $`1.5n<5`$. (For $`n=1.5`$ the cloud is isentropic, while for $`n5`$ the central density becomes infinite).
For each polytropic index $`n`$ there exists a specific value of $`S=S_{cr}`$ such that for $`S<S_{cr}`$ equation (3) has only one solution and magnification of the lens is always finite, even for $`b_0=0`$. For $`S>S_{cr}`$ equation (3) can have three solutions for sufficiently small $`b_0`$ and caustic lensing can occur. In this caustic regime the magnification becomes infinite for $`b_0=0`$ and also for some finite $`b_{0c}`$ given by the condition that at this $`b_{0c}`$ two of the roots of equation (3) coincide. At $`b_0=b_{0c}`$ light curves exhibit conspicuous caustics with magnification going to infinity. Cases of caustic and noncaustic lensing are illustrated in Figure 2 for $`n=1.5`$.
For $`S<S_{cr}`$, we can obtain a simple analytical formula describing the dependence of the central magnification $`M(0)`$ upon the strength parameter $`S`$. If we define dimensionless variables $`\widehat{\rho }\rho /\rho `$ and $`\widehat{r}r/R_{cl}`$, it is straightforward to show that as $`b_00`$
$$bb_02Sb\underset{0}{\overset{1}{}}\frac{d\widehat{r}}{\widehat{r}}.\frac{d\widehat{\rho }}{d\widehat{r}}$$
(7)
The integral in this expression depends only on the cloud density profile.
Substituting this result into (4) we get
$$M(0)=\frac{1}{\left(1S/S_{cr}\right)^2},$$
(8)
where
$$S_{cr}\left(2\underset{0}{\overset{1}{}}\frac{d\widehat{r}}{\widehat{r}}\frac{d\widehat{\rho }}{d\widehat{r}}\right)^1.$$
(9)
Thus $`M(0)\mathrm{}`$ as $`SS_{cr}`$, but it is finite for smaller $`S`$. At $`S=S_{cr}`$ the number of solutions of equation (3) changes from $`1`$ to $`3`$ (for $`b_0=0`$) and for $`SS_{cr}`$ we enter the caustic regime. For $`n=1.5`$ $`S_{cr}=0.026`$ and for $`n=4.5`$ $`S_{cr}=1.72\times 10^6`$.
## 3 Rate of gaseous lensing events
When a cloud crosses the line of sight between the distant star and the observer on the Earth, it amplifies or attenuates the brightness of the star \[in fact each gaseous lensing curve exhibits some demagnification before and after the peak magnification (Draine 1998), but we will be primarily interested in amplification in this paper\]. This effect can be observed by searches for gravitational microlensing in our Galaxy which monitor the brightness of a large number of stars.
Let $`D_S`$ be the distance to the source stars being lensed by gaseous clouds in the Galaxy. Consider a lens located at some distance $`D_{OL}=xD_S`$ from the observer; the distance from the lens to the source star is $`D_{LS}=(1x)D_S`$, and $`D=x(1x)D_S`$.
Both lensing regimes, caustic and noncaustic, are important. Since the lensing parameter
$$S=\frac{\alpha \rho D_S}{R_{cl}}x(1x),$$
(10)
maximum $`S`$ is attained when the lens is placed midway between the source and the observer, and is equal to $`S_{max}=\alpha \rho D_S/(4R_{cl})`$. As $`x`$ changes towards $`x=0`$ or $`x=1`$, $`S`$ declines to zero. Thus, if $`S_{max}>S_{cr}`$, there are critical values of $`x`$
$$x_{1,2}=\frac{1}{2}\left(1\sqrt{1\frac{S_{cr}}{S_{max}}}\right),$$
(11)
such that for $`x_1<x<x_2`$ lensing is caustic while for $`x<x_1`$ and $`x>x_2`$ lensing is in the noncaustic regime. If $`S_{max}<S_{cr}`$ lensing is always noncaustic.
For the observer to see a noncaustic lensing event with a magnification larger than some threshold magnification $`A_t`$, the lens has to pass near the source star with a sufficiently small (unlensed) impact parameter $`b_0`$. In other words, magnification $`M>A_t`$ if $`b_0<b_{0t}(A_t)`$, where $`b_{0t}(A_t)`$ is given by
$$M(b_{0t}(A_t))=A_t.$$
(12)
As the cloud moves through the sky, all the source stars lying in a strip on the sky along the lens trajectory with angular width $`2\delta `$ given by
$$\delta (A_t,x)=b_{0t}(A_t)/D_{OL},$$
(13)
are amplified with magnification $`M>A_t`$.
Let $`𝐯_S`$, $`𝐯_L`$, and $`𝐯_O`$ be the transverse velocities (i.e. velocities perpendicular to the line from observer to source star) of the source star, the lens, and the observer. Then the relative transverse velocity $`𝐯_{}`$ of the gaseous lens and the source star as seen by the observer is just
$$𝐯_{}=𝐯_L\left[x𝐯_S+(1x)𝐯_O\right]$$
(14)
(Han & Gould 1996).
It is important to distinguish events with different durations. Any real gravitational microlensing experiment has some finite probability $`<1`$ of detecting a lensing event which occurred during the monitoring campaign and this detection efficiency function $`\varphi `$ depends strongly upon the timescale of the observed event (Alcock 1997). We define the timescale $`\tau `$ of the lensing event as the time which the light curve spends above a threshold magnification $`A_t`$:
$$\tau =\frac{2D_{OL}\sqrt{\delta ^2\phi ^2}}{|𝐯_{}|},$$
(15)
where $`\phi `$ is the angular impact parameter of the lens’s trajectory relative to the source star; $`\phi <\delta `$ (otherwise $`M<A_t`$).
In calculations of the rate of lensing we must take into account the fact that lensing clouds and source stars have some distribution in velocity space; we average the lensing rate over these distributions to obtain the true expected rate. The general formula for the event rate per lensing cloud is
$$d\dot{N}=2\mathrm{\Sigma }_S\underset{0}{\overset{\delta }{}}𝑑\phi 𝑑𝐯_L𝑑𝐯_Sf_L(𝐯_L)f_S(𝐯_S)\frac{|𝐯_{}|}{D_{OL}}\varphi _\tau (\tau (\phi ,|𝐯_{}|)),$$
(16)
where $`f_L`$ and $`f_S`$ are the velocity distribution functions for lensing clouds and source stars, respectively, $`\mathrm{\Sigma }_S`$ is the number of source stars per unit solid angle on the sky, and $`\varphi _\tau `$ is the detection efficiency function for a lensing event with timescale $`\tau `$.
We assume for simplicity isotropic Maxwellian distribution functions for the velocities of both lenses and source stars with dispersions $`\sigma _L`$ and $`\sigma _S`$ respectively, taking into account the fact that the object containing stars is moving in general, i.e. there is some offset velocity $`𝐯_c`$ in the latter distribution:
$$f_L(𝐯_L)=\frac{1}{\pi ^{3/2}\sigma _L^{3/2}}e^{𝐯_L^2/\sigma _L^2},$$
(17)
and
$$f_S(𝐯_S)=\frac{1}{\pi ^{3/2}\sigma _S^{3/2}}e^{(𝐯_S𝐯_c)^2/\sigma _S^2}.$$
(18)
Our calculations are not very sensitive to this assumption in the sense that the order of magnitude result does not depend strongly upon the exact shape of the velocity distribution. The important assumption of isotropy of the distribution function permits analytical simplifications.
Consider a coordinate system with the $`z`$-axis lying along the line of sight from observer to the source, $`x`$-axis perpendicular to the line of sight so that observer’s velocity $`𝐯_O`$ lies in the $`xz`$-plane, and $`y`$-axis perpendicular to those two, so that $`v_{Oy}=0`$. Taking distributions (17) and (18) it is shown in Appendix A that the total event rate per source is
$`\dot{N}_{tot}=4D_S^2{\displaystyle \underset{0}{\overset{1}{}}}𝑑x{\displaystyle \frac{xn_L(𝐱)\delta (x,A_t)v_{ch}^3}{\sigma _L^2+x^2\sigma _S^2}}e^C{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}u^2\kappa (uv_{ch})e^{Au^2}I_0\left(Bu\right)𝑑u,`$ (19)
where $`I_0`$ is the modified Bessel function of order zero, $`n_L(𝐱)`$ is the number density of lensing clouds, $`v_{ch}(x)`$ is some characteristic velocity (the final result does not depend upon the choice of this velocity), function $`\kappa `$ is defined as
$$\kappa (u)=\underset{0}{\overset{1}{}}\varphi _\tau \left(\frac{2\delta (A_t,x)xD_S}{u}\sqrt{1s^2}\right)𝑑s,$$
(20)
and $`A(x),B(x)`$ and $`C(x)`$ are given by
$`A={\displaystyle \frac{v_{ch}^2}{\sigma _L^2+x^2\sigma _S^2}},B={\displaystyle \frac{2v_{ch}}{\sigma _L^2+x^2\sigma _S^2}}\sqrt{x^2v_{cy}^2+(xv_{cx}+(1x)v_{Ox})^2},`$
$`C={\displaystyle \frac{x^2v_{cy}^2+(xv_{cx}+(1x)v_{Ox})^2}{\sigma _L^2+x^2\sigma _S^2}}.`$ (21)
This formula for the lensing rate is quite general. For example, it is directly applicable to the case of gravitational microlensing \[with corresponding $`\delta (A_t,x)`$\], in which case it is identical to the formula for the event rate obtained by Griest (1991).
It is interesting to know the timescale distribution of the event rate to have some idea of where to look for the events produced by the gaseous lensing. The timescale distribution is given by (see Appendix A)
$`{\displaystyle \frac{d\dot{N}_{tot}}{d\tau }}=32D_S{\displaystyle \frac{\varphi _\tau (\tau )}{\tau ^4}}{\displaystyle \underset{0}{\overset{1}{}}}𝑑x{\displaystyle \frac{n_L(𝐱)b_{0t}^4(x,A_t)}{\sigma _L^2+x^2\sigma _S^2}}e^C{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{u^4}{\sqrt{1u^2}}}e^{Au^2}I_0\left(Bu\right)𝑑u,`$ (22)
where $`b_{0t}(x,A_t)`$ is defined in (12) and $`v_{ch}`$, entering the definitions of $`A`$, $`B`$, and $`C`$ in (21), is given by $`v_{ch}=2b_{0t}(x,A_t)/\tau `$ \[compare with a similar formula by Griest (1991) for the case of gravitational microlensing\].
Closely related to the timescale distribution $`d\dot{N}_{tot}/d\tau `$ is its integral from $`\tau _{min}`$ to infinity, that is the frequency of events with timescales exceeding $`\tau _{min}`$,
$`\dot{N}_{tot}(\tau >\tau _{min})={\displaystyle \underset{\tau _{min}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\dot{N}_{tot}}{d\tau }}𝑑\tau `$
$`=32{\displaystyle \frac{D_S}{\tau _{min}^3}}{\displaystyle \underset{0}{\overset{1}{}}}𝑑x{\displaystyle \frac{n_L(𝐱)b_{0t}^4(x,A_t)}{(\sigma _L^2+x^2\sigma _S^2)}}e^C{\displaystyle \underset{0}{\overset{1}{}}}u^2\stackrel{~}{\kappa }(u)e^{Au^2}I_0\left(Bu\right)𝑑u,`$ (23)
where function $`\stackrel{~}{\kappa }`$ is now
$`\stackrel{~}{\kappa }(u)={\displaystyle \underset{0}{\overset{\sqrt{1u^2}}{}}}\varphi _\tau \left({\displaystyle \frac{\tau _{min}}{u}}\sqrt{1s^2}\right)𝑑s.`$ (24)
If one wants to know the true event rate, for $`100\%`$ efficiency, one need only set $`\varphi _\tau =1`$ in formulae (22) and (23), as we will do in our consideration of the event rate distributions.
In the next section we show typical plots of $`\dot{N}_{tot}(\tau >\tau _{min})`$ for some clouds characterized by polytropes with $`n=1.5`$ and $`n=4.5`$.
## 4 Results
In this section we consider the lensing of stars located in the LMC and compare with the results obtained by the MACHO collaboration, which has monitored $`11.9`$ million LMC stars for $`5.7`$ yrs (Alcock et al. 2000).
### 4.1 Parameters of the event rate calculations
The LMC has a distance $`D_{LMC}55`$ kpc from the Sun, and its position in galactic coordinates is $`b=32.8^{}`$ and $`l=281^{}`$. Jones et al. (1994) finds the motion of the LMC to be
$`v_l=16\pm 60\mathrm{km}\mathrm{s}^1`$
$`v_b=328\pm 60\mathrm{km}\mathrm{s}^1`$ (25)
$`v_{rad}=250\pm 10\mathrm{km}\mathrm{s}^1.`$
We use these proper motion measurements for obtaining our event rates. The velocity dispersion of the lensing clouds was assumed to be $`\sigma _L=220\mathrm{km}\mathrm{s}^1`$ and that of the stars in the LMC was taken to be $`\sigma _S=10\mathrm{km}\mathrm{s}^1`$.
We performed our calculations both for the case of the true event rate (for $`100\%`$ detection efficiency function) and also for the the case of the MACHO detection efficiency function which is given in Alcock et al (1997). In the first case we assume a threshold magnification $`A_t=1.1`$ and in the second we take $`A_t=1.35`$, which is the current threshold magnification of the MACHO experiment (Alcock et al. 2000). Estimation of the MACHO detection efficiency function $`\varphi _\tau `$ is problematic because it is known in terms of the duration of a gravitational microlensing event, which differs significantly from our definition (15) for the gaseous lensing timescale. The method by which we estimate $`\varphi _\tau `$ is described in Appendix B, and our adopted $`\varphi _\tau `$ is shown on Figures 3 and 4.
We model the lensing clouds by self-gravitating H<sub>2</sub>-He (nonrotating) polytropes of radius $`R_{cl}`$. The polytropic index $`n`$ ($`T\rho ^{1/n}`$) must be in the range $`1.5n<5`$; for $`n<1.5`$ the cloud would be convectively unstable, while for $`n5`$ the central density is infinite. We do not expect $`T(r)`$ and $`\rho (r)`$ to be accurately described by a polytropic model, but a slight rise in temperature toward the interior may be reasonable since the interior will be heated by high energy cosmic rays, with the few cooling lines \[e.g., H<sub>2</sub> 0-0S(0) 28.28$`\mu \mathrm{m}`$, or HD 0-0R(1) 112$`\mu \mathrm{m}`$\] very optically thick. Furthermore, polytropes have $`T(R_{cl})=0`$, so the density structure near the surface is unphysical. Table 1 of Draine (1998) gives various properties for H<sub>2</sub>-He polytropes. In our current calculations we consider gaseous clouds with polytropic indices $`n=1.5`$ and $`n=4.5`$, to compare extremes of behavior<sup>1</sup><sup>1</sup>1 We assume that it is the density-radius relation of the clouds which is determined by the polytropic law, not the behavior of the gas on dynamical timescales (otherwise the $`n=4.5`$ polytrope would be dynamically unstable.) .
Following Widrow & Dubinski (1998) we adopt a Navarro, Frenk, & White (1996) model of a spherical halo, composed of gaseous clouds:
$$n_L(𝐫)=f_{cl}\frac{M_{\mathrm{MW}}}{4\pi M_{cl}}\frac{1}{r(r+a_s)^2}$$
(26)
where $`r`$ is the distance between the galactic center and the point of interest, $`a_s=15.9`$ kpc is the core radius, and $`0f_{cl}<1`$ is the fraction of the total mass of the Milky Way $`M_{\mathrm{MW}}=6.5\times 10^{11}\mathrm{M}_{}`$ contributed by the population of self-gravitating H<sub>2</sub> clouds.
### 4.2 Timescale distributions
On Figures 3 and 4 we plot the dependence of the integrated timescale distribution $`\dot{N}_{tot}(\tau >\tau _{min})`$, as a function of $`\tau _{min}`$, for some representative values of the cloud mass $`M_{cl}`$ and radius $`R_{cl}`$. The cloud parameters are given in Table 1 and we assume $`\varphi _\tau =1`$.
All timescale distribution curves exhibit a plateau for small timescales and then a rapid falloff at larger event durations. From these plots we see that the characteristic timescale of the lensing events is quite small, varying from several days to several tens of days. This is as expected, because cloud sizes are small, $`10`$ AU, sufficient amplification occurs only when $`b_0<(0.10.3)R_{cl}`$, and for typical transverse velocity $`200`$ km $`\mathrm{s}^1`$ we get just this range of event durations.
The time distribution seen by a real experiment would be modified by the detection efficiency, which may be small for short timescale events. In addition, events lasting $`50`$ days, which occur when the transverse velocity of the lens is quite small, will be affected by changes in the velocity of the Earth over this period (the so-called “parallax effect” (Paczyński 1996; Gould & Andronov 1999)).
### 4.3 Results for the event rate
Calculation of the event rate is quite straightforward in the case of noncaustic lensing because in this case foreground star crossing the cloud produces magnifications $`A>A_t`$ in one continuous span (see Figure 2 for the cases $`S<S_{cr}`$). In the caustic regime the situation is more complicated, since for clouds located at the middistance between observer and foreground star these magnifications $`A>A_t`$ will be not only in the region of the closest approach but also at the caustic spikes (see Figure 2 for the cases $`S>S_{cr}`$) and we want to account for this.
We adopted a simple approach for the caustic regime. One can see from Figure 2 (upper panel) that for a given polytropic index there exists a specific value of $`S_{}>S_{cr}`$ such that the inner minimum of the lightcurve (between the central part and one of the caustic spikes) has magnification $`A_t`$. For example, in the case $`n=1.5`$ $`S_{}=0.112`$ for $`A_t=1.1`$ or $`S_{}=0.101`$ for $`A_t=1.35`$. It is obvious that when $`S_{cr}<S<S_{}`$ the picture is similar to the noncaustic case since magnifications $`A>A_t`$ are reached in one continuous span. If $`S>S_{}`$ there are also caustic spikes producing events but one can easily see that the duration of events caused by these spikes decreases rapidly with growing $`S`$ (Draine 1998) and we simply ignore them in our calculations. That is we assume for $`S>S_{}`$ that $`A>A_t`$ only in the central part of the lightcurve. This gives us only a lower limit on the event rate but a useful one since the number of events produced by caustics is relatively small and the caustic spikes are of short duration.
In Figures 5 and 6 we show the true event rate (i.e. for efficiency $`\varphi _\tau =1`$) due to gaseous lensing for various cloud parameters and for two polytropic indices of the clouds: $`n=1.5`$ and and $`n=4.5`$. We characterize the clouds by mass $`M_{cl}`$ and $`S_{max}`$, the “strength” parameter for a cloud located midway between the Earth and the LMC: $`S_{max}=\alpha \rho D_{LMC}/4R_{cl}`$, where $`D_{LMC}=55`$ kpc. For each value of $`S_{max}`$ and $`M_{cl}`$ the cloud radius $`R_{cl}`$ is obtained from equation (6):
$`R_{cl}=\left({\displaystyle \frac{3\alpha M_{cl}D_{LMC}}{16\pi S_{max}}}\right)^{1/4}=1.25\times 10^{13}\mathrm{cm}\left({\displaystyle \frac{M_{cl}}{10^3M_{}}}\right)^{1/4}{\displaystyle \frac{1}{S_{max}^{1/4}}}.`$ (27)
Lines of constant cloud radius are plotted on each of Figures 5-6 as dot-dashed lines. We consider here events with all timescales and magnifications larger than $`A_t=1.1`$.
Our adopted bound on $`M_{cl}`$ comes from two considerations. First, clouds which are too massive would be similar to the objects known to be unstable to collapse to form low mass stars. Second, Wardle & Walker (1999) obtained an upper bound on $`M_{cl}<10^{1.7}M_{}`$ by considering a cooling mechanism which could stably radiate away the heat deposited by cosmic rays.
We restricted the parameter space of $`S_{max}`$ and $`M_{cl}`$ because of the following reasons. The upper bound in $`S_{max}`$ is just $`10`$ because this corresponds to a cloud with $`M_{cl}=10^{1.7}M_{}`$ and radius $`R_{cl}=1`$ AU. We do not consider smaller radii because extreme scattering events favor clouds with sizes of the order of several AU (Draine 1998).
The minimum value of $`S`$ for a cloud to be detectable by light curve monitoring with a threshold magnification $`A_t`$ can be obtained from equation (8):
$$S_{det}=S_{cr}(11/\sqrt{A_t}),$$
(28)
where $`S_{cr}`$ is given by equation (9). For $`S_{max}>S_{det}`$ the cloud population will produce a nonzero rate of events with $`A_{max}>A_t`$.
Dotted lines show the restriction imposed by the requirement that the gravitational binding energy of gas at the surface exceed the thermal energy for an estimated atmospheric temperature $`T`$ (Draine 1998):
$$M_{cl}>6\times 10^5\left(\frac{T}{10\mathrm{K}}\right)\left(\frac{R_{cl}}{\mathrm{AU}}\right)M_{}.$$
(29)
Two boundaries are shown, for atmospheric temperatures $`T=10`$ K and $`T=15`$ K. The shaded region below these curves is prohibited.
The thick solid lines are contours of constant event rate. The lensing rate is very high for all the cloud models appreciably to the right of the “rate $`=0`$” line.
Figures 7 and 8 show the results for the lensing rate for polytropes with $`n=1.5`$ and $`n=4.5`$ for a detection efficiency function approximately that of the MACHO experiment. All the notation is the same as in the case of $`\varphi _\tau =1`$. The “rate $`=0`$” boundary has moved (relative to Figures 5 and 6) because we now require a threshold magnification $`A_t=1.35`$. Small glitches on the contours are artifacts of the numerical procedure. We consider here the events with all timescales filtered through the detection efficiency function discussed in Appendix B.
The rate which would be seen by the MACHO experiment is significantly smaller than the rate for $`A_t=1.1`$ and $`\varphi _\tau =1`$, sometimes by two orders of magnitude. This is primarily a consequence of the reduced sensitivity to short timescale events of the MACHO detection efficiency function. Nevertheless the predicted rate is still quite high and in the next subsection we compare the predicted rate to that actually observed.
### 4.4 Comparison with the MACHO results
How shall we compare the results of the MACHO experiment with our simulated plot of what it would see if the dark matter is composed of molecular clouds? First of all, we need to estimate what fraction of the events actually observed by MACHO could be attributed to gaseous lensing. This can be done by noticing that the typical timescale of gaseous lensing events is very small, while the shortest timescale event reported by MACHO has $`\widehat{t}=34`$ days. For the types of timescale distribution seen in Figures 3 and 4 it is extremely unlikely that there are more than $`1`$ or $`2`$ gaseous lensing events among those observed by MACHO, since otherwise a large number of short timescale events (with durations less than $`10`$ days) would be seen, in conflict with the distribution of timescales observed by the MACHO and EROS (Alcock et al. 2000; Lasserre et al. 2000) collaborations. It would also be completely inconsistent with the EROS and MACHO combined limits on the rate of very short timescale events, from $`15`$ minutes to several days, for which they claim that the analysis of two years of data on $`8.6`$ million stars found no short-duration “spike” events (Alcock et al. 1998), implying an upper limit $`\dot{N}10^7`$ events yr<sup>-1</sup>.
The MACHO collaboration observed $`11.9`$ million stars in the LMC for $`5.7`$ years. This means that the rate which could possibly be due to gaseous lensing events can be at most $`(0.010.03)\times 10^6\mathrm{events}\mathrm{yr}^1`$. It is clear from our simulated observations of gaseous lensing on Figures 7 and 8 that if $`f_{cl}0.1`$ the contour describing such a rate will lie very close to the line $`S_{max}=S_{det}`$, where the event rate drops to $`0`$. Any cloud model with $`S_{max}<S_{det}`$ will of course be undetectable and thus be allowed. Models with $`S_{max}>1.1S_{det}`$ are prohibited, because otherwise MACHO would have seen an enormous number of gaseous lensing events.
These considerations and the restrictions described in §4.3 are combined together to produce an exclusion plot (Figure 9). One can see that an allowed region exists where clouds basically cannot be seen by the MACHO experiment and are not prohibited by the maximum mass and “thermal evaporation” constraints. For larger $`n`$ (softer equation of state) the allowed region shrinks. It is likely that for $`n5`$ this region disappears completely, though we did not run our calculations for $`n`$ larger than $`4.5`$.
Also in the case $`n=4.5`$ another allowed region appears for small cloud radii (see Figure 8). It is bounded by the level contour corresponding to the rate $`3\times 10^8\mathrm{events}\mathrm{yr}^1`$ and goes into the region of the small cloud radii (see Figure 9 b).
The allowed regions can be further narrowed by noting that clouds are highly unlikely to have high central temperatures. Detailed studies of their thermal structure are not available, but it seems reasonable to exclude models of clouds with central temperatures $`>(50100)`$ K. For $`n=1.5`$ this places another limit within the allowed region (see Figure 9a). In the case $`n=4.5`$ the allowed region for small radii virtually disappears because of this restriction (see Figure 9b).
In the absence of realistic models for the thermal structure of these clouds, we can only use $`n=1.5`$ polytropes as the most conservative model from the standpoint of gaseous lensing. For this case, we see that the allowed region is quite substantial, and contains the model with $`M_{cl}=10^3M_{}`$ and $`R_{cl}=10`$ AU which was favored by Draine (1998).
It is also possible to expand the allowed region by assuming that the clouds contribute only a small part $`f_{cl}1`$ of the dark matter (see §4.1), but to get any significant effect we need to suppose that this fraction is $`<10^2`$ which would rule out these clouds as a main constituent of the dark matter, and eliminate them as the explanation for the extreme scattering events.
It is clear that our lensing restriction represented on Figure 9 is sensitive to the adopted density profile of an individual cloud, but it is relatively insensitive to changes in the spatial and velocity distributions of the dark matter, since changes in these parameters are unlikely to produce variations in the event rate of $`2`$ orders of magnitude, which would be needed to move the restriction given by the observed lensing rate significantly away from the line “rate $`=0`$”.
## 5 Discussion
Wardle & Walker (1999) suggested that heating caused by cosmic rays in the cold molecular clouds can be balanced if particles of solid H<sub>2</sub> can form. They argue that these particles cool the cloud by thermal continuum radiation if the cloud mass lies in the range $`10^{7.5}10^{1.7}M_{}`$. Our calculations show that the lensing event rate would grow to extremely high values for masses below $`10^4M_{}`$ (see Figures 7 and 8 ).
Kalberla, Shchekinov, & Dettmar (1999) proposed that the $`\gamma `$-ray background emission from the Galactic halo seen in EGRET ($`E>100`$ MeV photons) data can be produced by interaction of high-energy cosmic rays with small dense clouds. They modeled this effect with different types of spatial distribution of dark matter in the form of these clouds and found that their best fit to the EGRET data for uniform density clouds occurred for $`M_{cl}/R_{cl}^210^3M_{}/(6AU)^2`$. Their conclusions about the cloud parameters are model dependent, but the quoted values fall within the allowed range on Figure 9 for $`n=1.5`$.
Though it is impossible to observe the gravitational microlensing by the molecular clouds in our Galaxy it becomes possible if one observes them in other galaxies. A search for such lenses in the Virgo cluster is now underway (Tadros, Warren, & Hewett 2000) and preliminary results indicate that galaxy halos are unlikely to primarily consist of object with masses smaller than $`10^5M_{}`$. We confirm this restriction but also place more rigorous constraints on the cloud properties since we can reject many models with higher cloud masses.
Another possible approach would be to look for periods of demagnification in the light curves of lensing events. For instance, the OGLE-II survey (Udalski et al. 1997) has collected a large sample of light curves for stars in the Galactic bulge, including $`300`$ lensing events (Woźniak 2000). Of these, two or three appear to show demagnification preceding or following the central peak in brightness, although the statistical significance of the demagnification has not yet been established (Woźniak, private communication). Demagnification cannot be produced by gravitational lensing by one or more point masses, but would be a characteristic signature of gaseous lensing.
## 6 Summary
In this paper we have considered in detail the lensing of stars in the LMC by small self-gravitating molecular clouds which have been proposed as a candidate for the dark matter in the Galaxy. Lensing would occur because of the refraction of optical light in the clouds, with resulting magnification of the source. We have developed a semianalytical formalism for calculation of the lensing rate including the spatial distribution of the dark matter, finite velocity dispersions of the lensing clouds and the source stars in LMC, and proper motion of the observer and LMC itself.
Our calculations were carried out for a single mass and size cloud population. One can easily extend the analysis to a distribution of cloud masses and sizes by simply convolving event rates obtained for a uniform cloud population with the appropriate distributions.
Lensing events might be detectable by searches for gravitational microlensing. This has allowed us to obtain constraints on the cloud properties by comparing our calculations with observational results obtained by the MACHO collaboration.
We found that almost the only possibility for the dark matter to be in the form of such molecular clouds is for the clouds to be sufficiently weak lenses so that their lensing effects are below the detection threshold of the MACHO experiment, since otherwise a very large gaseous lensing event rate would have been detected. This still leaves an allowed region in parameter space where these clouds could exist and not contradict the limitations posed by lensing experiments and simple physical considerations.
Though we have performed event rate calculations for only one halo model \[given by equation (26)\], we expect our results to be relatively insensitive to the particular form of the spatial distribution of these clouds, or the assumed shape of the velocity distribution of the clouds and the stars in the LMC.
Future microlensing experiments with a lower detection threshold magnification $`A_t`$ could either detect these clouds or strengthen the constraints on their properties.
## 7 Acknowledgements
We thank B.Paczyński and P.Woźniak for very helpful discussions, and R.H.Lupton for availability of the SM plotting package. This research has been supported in part by NSF grants AST-9619429 and AST-9988126.
## Appendix A Derivation of equations (19) - (23)
For the purposes of the lensing rate calculations the velocities of the lens, observer, and source along the line of sight are not important, because only transverse motions are significant. It means that we can immediately perform integrations in equation (16) over $`v_{Sz}`$ and $`v_{Lz}`$.
Taking into account (15) we can perform integration over $`\phi `$ in (16) and define the result as a function $`\kappa (𝐯_{})`$:
$$\kappa (𝐯_{})\frac{1}{\delta (x,A_t)}\underset{0}{\overset{\delta (x,A_t)}{}}\varphi _\tau (\tau (\phi ))𝑑\phi =\underset{0}{\overset{1}{}}\varphi _\tau \left(\frac{2\delta (x,A_t)D_{OL}}{|𝐯_{}|}\sqrt{1s^2}\right)𝑑s,$$
(A1)
so that now
$$d\dot{N}=2\delta \frac{\mathrm{\Sigma }_S}{D_{OL}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑v_{Lx}\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑v_{Ly}\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑v_{Sx}\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑v_{Sy}|𝐯_{}|f_L(v_{Lx},v_{Ly})f_S(v_{Sx},v_{Sy})\kappa (|𝐯_{}|).$$
(A2)
We now change variables from $`v_{Lx}`$ and $`v_{Ly}`$ to $`v_x`$ and $`v_y`$ via equation (14):
$`d\dot{N}=2\delta {\displaystyle \frac{d\mathrm{\Sigma }_L}{\pi ^2\sigma _L^2\sigma _S^2D_{OL}}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑v_x{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑v_y{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑v_{Sx}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑v_{Sy}\sqrt{v_x^2+v_y^2}\kappa (\sqrt{v_x^2+v_y^2})`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{(v_x(1x)v_{Ox}xv_{Sx})^2+(v_yxv_{Sy})^2}{\sigma _L^2}}{\displaystyle \frac{(v_{Sx}v_{cx})^2+(v_{Sy}v_{cy})^2}{\sigma _S^2}}\right\}.`$ (A3)
We integrate over $`v_{Sx}`$ and $`v_{Sy}`$ to obtain after lengthy but straightforward calculations
$`d\dot{N}=2\delta {\displaystyle \frac{\mathrm{\Sigma }_S}{\pi (\sigma _L^2+x^2\sigma _S^2)D_{OL}}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑v_x{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑v_y\sqrt{v_x^2+v_y^2}\kappa (\sqrt{v_x^2+v_y^2})`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{v_x^2+v_y^2}{\sigma _L^2+x^2\sigma _S^2}}+\zeta _1v_x+\zeta _2v_yC\right\},`$ (A4)
where
$$C=\frac{x^2v_{cy}^2+(xv_{cx}+(1x)v_{Ox})^2}{\sigma _L^2+x^2\sigma _S^2},$$
(A5)
$$\zeta _1=2\frac{xv_{cx}+(1x)v_{Ox}}{\sigma _L^2+x^2\sigma _S^2},\zeta _2=2\frac{xv_{cy}}{\sigma _L^2+x^2\sigma _S^2}.$$
(A6)
Introducing polar coordinates in velocity space $`v_x=v\mathrm{cos}\beta `$, $`v_y=v\mathrm{sin}\beta `$ and normalizing $`v`$ to some characteristic velocity $`v_{ch}`$ we get
$`d\dot{N}=2\delta {\displaystyle \frac{\mathrm{\Sigma }_Sv_{ch}^3}{\pi (\sigma _L^2+x^2\sigma _S^2)D_{OL}}}e^C{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑uu^2\kappa (uv_{ch})e^{Au^2}{\displaystyle \underset{0}{\overset{2\pi }{}}}e^{uv_{ch}(\zeta _1\mathrm{cos}\beta +\zeta _2\mathrm{sin}\beta )}𝑑\beta ,`$ (A7)
where
$$A=\frac{v_{ch}^2}{\sigma _L^2+x^2\sigma _S^2}.$$
(A8)
The last integral in equation (A7) can be reduced to
$`{\displaystyle \underset{0}{\overset{2\pi }{}}}e^{uv_{ch}\sqrt{\zeta _1^2+\zeta _2^2}cos(\beta +\beta _0)}𝑑\beta ={\displaystyle \underset{0}{\overset{2\pi }{}}}e^{uv_{ch}\sqrt{\zeta _1^2+\zeta _2^2}cos\beta }𝑑\beta =2\pi I_0\left(Bu\right),`$ (A9)
where $`I_0`$ is the modified Bessel function of order zero, and
$$B=v_{ch}\sqrt{\zeta _1^2+\zeta _2^2}=\frac{2v_{ch}}{\sigma _L^2+x^2\sigma _S^2}\sqrt{x^2v_{cy}^2+(xv_{cx}+(1x)v_{Ox})^2}=2\sqrt{AC}.$$
(A10)
So, finally we get for the event rate produced by one lens
$`d\dot{N}(x)=4\delta {\displaystyle \frac{\mathrm{\Sigma }_Sv_{ch}^3}{(\sigma _L^2+x^2\sigma _S^2)D_{OL}}}e^C{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}u^2\kappa (uv_{ch})e^{Au^2}I_0\left(Bu\right)𝑑u.`$ (A11)
Thus, the total event rate per source, produced by all lenses between us and LMC, is given by formula (19).
Now, we obtain the formula for the timescale distribution (22). Note that $`\tau `$ and $`\phi `$ – the angular distance from the lens’s trajectory to the source star in the perpendicular direction – are directly related by equation (15):
$$d\phi =\delta (x,A_t)\frac{d\tau }{\tau }\frac{u^2}{\sqrt{1u^2}},$$
(A12)
with $`u=v_{}/v_{ch}`$, $`v_{ch}=2b_{0t}(x,A_t)/\tau `$. To obtain $`d\dot{N}_{tot}/d\tau `$ we simply omit the integration over $`\phi `$ in (16). Also, the velocity $`v_{}`$ is limited by $`2b_{0t}(x,A_t)/\tau `$. All the other steps are analogous to those which we have done in derivation of (19) and we obtain equation (22).
If we are interested only in events whose durations are larger than some chosen value $`\tau _{min}`$, then to get the total rate we should integrate in (16) over $`\phi `$ not up to $`\delta (x,A_t)`$, but up to
$$\phi _{max}=\sqrt{\delta ^2\frac{v_{}^2\tau _{min}^2}{4D_{OL}}}=\delta \sqrt{1u^2},$$
(A13)
with $`u=v_{}\tau /2\delta xD`$. One can easily see that we will get in this case for $`\dot{N}_{tot}(\tau >\tau _{min})`$ the formula (23).
## Appendix B Detection efficiency function
The detection efficiency function $`\varphi `$ characterizes the sensitivity of the experiment to events with different durations. We will use the detection efficiency $`\varphi _{\widehat{t}}(A,\widehat{t})`$ assumed by MACHO experiment to estimate how it would affect gaseous lensing (see Figure 8 in Alcock et al. (1997)).
We first recall some simple facts about gravitational microlensing. The amplification $`M`$ in the case of gravitational microlensing is determined by the formula
$`M(u)={\displaystyle \frac{u^2+2}{u(u^2+4)^{1/2}}},u=(u_{min}^2+4t^2/\widehat{t}^2)^{1/2}`$
$`u_{min}={\displaystyle \frac{b_{min}}{r_E}},\widehat{t}=2{\displaystyle \frac{r_E}{v_{}}},r_E=\left[{\displaystyle \frac{4GmDx(1x)}{c^2}}\right]^{1/2},`$ (B1)
where $`b_{min}`$ is the impact parameter, $`m`$ is the mass of the lensing object, $`t`$ is the time of observation, and $`v_{}`$ is, again, the relative velocity of the lens and the source star.
Recall that we defined the timescale $`\tau `$ to be the time spent with magnification $`M>A_t`$. From equation (B1) we find that, for gravitational microlensing,
$`\tau _{gm}=\sqrt{2}\widehat{t}\left[{\displaystyle \frac{A_t}{\sqrt{A_t^21}}}{\displaystyle \frac{A_{max}}{\sqrt{A_{max}^21}}}\right]^{1/2}`$
$`=\widehat{t}\left[u_t^2u_{min}^2\right]^{1/2},`$ (B2)
where $`u_t`$ is the value of $`u`$ at which the magnification is equal to $`A_t`$ and $`A_{max}`$ is the maximum magnification for a given $`u_{min}`$.
We use a very simple algorithm for converting detection efficiency function from $`\widehat{t}`$ to $`\tau `$: since $`\tau \widehat{t}`$, their ratio $`\tau /\widehat{t}`$ is a function of $`A_{max}`$ (or $`u_{min}`$) for a given $`A_t`$. Thus we can average this ratio over the distribution of the maximum magnifications $`A_{max}`$ and assume that the value of known $`\varphi _{\widehat{t}}(\widehat{t})`$ is equal to the $`\varphi _\tau (<\tau >)`$, that is
$$\varphi _\tau (\tau )=\varphi _{\widehat{t}}\left(\tau \left(<\tau /\widehat{t}>\right)^1\right).$$
(B3)
The calculation of $`<\tau /\widehat{t}>`$ is quite straightforward and can be most easily done in terms of $`u_{min}`$ (see the last equality in equation (B2)). Since the distribution of lensing events over $`u_{min}`$ is just constant
$$<\tau /\widehat{t}>=\frac{\underset{0}{\overset{u_t}{}}\left[u_t^2u_{min}^2\right]^{1/2}𝑑u_{min}}{\underset{0}{\overset{u_t}{}}𝑑u_{min}}=\frac{\pi }{4}u_t=\frac{\pi }{2\sqrt{2}}\left(\frac{A_t}{\sqrt{A_t^21}}1\right)^{1/2}.$$
(B4)
For $`A_t=1.35`$, accepted by MACHO experiment, $`<\tau /\widehat{t}>=0.776`$, which means that $`\varphi _\tau (\tau )=\varphi _{\widehat{t}}(1.288\tau )`$.
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# Effective field theory of slowly-moving “extreme black holes”
## I introduction
Recently the properties of black holes in the Einstein-Maxwell-dilaton (EMD) theory have been studied by many authors. We are interested in the EMD theory since it can be regarded as an effective field theory of string and/or the gravitational theory with a dimensional reduction. An “extreme black hole”, which we consider in the present paper, is a kind of solitons with balanced three long-range forces (the gravitational force, the electric force, and the dilatonic force). <sup>*</sup><sup>*</sup>*If a dilaton field is coupled, the extreme limit of the black hole solution has a singularity in general. Therefore, we use the “quotation marks” around the extreme black hole in this paper. For the balance, the charge and mass of the “black hole” must satisfy the certain condition.
For the static case, the system with an arbitrary number of “extreme black holes” is known to be described by the Papapetrou-Majumdar solution and its descendants. In the present paper, we consider slowly-moving “extreme black holes”. Since they can be regarded as point particles, we are able to describe the collective behavior of them by a scalar field. Though the similar system has been studied before, our approach is different from previous works and may give new insights of the collective phenomena of such objects. The paper is arranged as follows.
In Sec. II, we derive a theoretical model of a non-relastivistic scalar field incorporating the low-energy interaction among “extreme black holes”. As an application, we examine an isothermal sphere of “extreme black holes” by using the technique of finite-temperature field theory in Sec. III. As a result, we will show that the gas of “extreme black holes” lumps at high temperature. In Sec. IV, we derive the expression for the energy of the system after eliminating the potentials by the field equations. We also obtain the total energy of the classical many-body system. Section V is devoted to conclusion.
The derivation of the effective lagrangian for the model in $`(N+1)`$ dimensions is exhibited in Appendix.
## II Effective Lagrangian
In this section, we derive an effective lagrangian for “extreme black holes” of which a typical velocity $`v`$ is small. In the classical theory, the action of particles with mass $`m`$ and charge $`e_0`$, coupled to the gravitational field, the electromagnetic field $`A_\mu `$, and the dilaton field $`\varphi `$, is written as
$$I=𝑑s\left[me^{a\varphi }+e_0A_\mu \frac{dx^\mu }{ds}\right],$$
(1)
where $`a`$ is the coupling constant for the dilaton field. Then, the four-momentum of the particle is
$$P_\mu =me^{a\varphi }g_{\mu \nu }\frac{dx^\nu }{ds}e_0A_\mu .$$
(2)
This four-momentum satisfies the following equation:
$$g^{\mu \nu }(P_\mu +e_0A_\mu )(P_\nu +e_0A_\nu )+m^2e^{2a\varphi }=0.$$
(3)
In quantum theory, we can rewrite this equation (3) into a wave equation for a wave function $`\phi `$:
$$\left[g^{\mu \nu }(P_\mu +e_0A_\mu )(P_\nu +e_0A_\nu )+m^2e^{2a\varphi }\right]\phi =0,$$
(4)
where $`\phi `$ possesses information on the “dynamics” of the collective behavior. So the action which yields this wave equation is
$`S_m`$ $`=`$ $`{\displaystyle }d^4x\sqrt{g}[\phi ^{}e^{a\varphi }g^{\mu \nu }(P_\mu +e_0A_\mu )(P_\nu +e_0A_\nu )\phi m^2e^{a\varphi }\phi ^{}\phi ],`$ (5)
where we regard $`\phi `$ as a “field” from now on. Therefore the total action including long-range interactions is
$`S`$ $`=`$ $`{\displaystyle }d^4x{\displaystyle \frac{\sqrt{g}}{16\pi }}[R2_\mu \varphi ^\mu \varphi e^{2a\varphi }F_{\mu \nu }F^{\mu \nu }]+S_m,`$ (6)
where $`R`$ is the scalar curvature and $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$. The Newton constant is set to unity.
This action leads to the following field equations:
$`^2\varphi +{\displaystyle \frac{a}{2}}e^{2a\varphi }F^2+\mathrm{\hspace{0.17em}4}\pi a\left[e^{a\varphi }\phi ^{}(P+e_0A)^2\phi e^{a\varphi }m^2\phi ^{}\phi \right]=0,`$ (7)
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }R=2\left[_\mu \varphi _\nu \varphi {\displaystyle \frac{1}{2}}g_{\mu \nu }(\varphi )^2\right]+e^{2a\varphi }\left[\mathrm{\hspace{0.17em}2}F_{\mu \nu }^2{\displaystyle \frac{1}{2}}g_{\mu \nu }F^2\right]`$ (8)
$`+16\pi \{e^{a\varphi }Re[\phi ^{}(P_\mu +e_0A_\mu )(P_\nu +e_0A_\nu )\phi `$ (9)
$`{\displaystyle \frac{1}{2}}g_{\mu \nu }\phi ^{}(P+e_0A)^2\phi ]{\displaystyle \frac{1}{2}}g_{\mu \nu }e^{a\varphi }m^2\phi ^{}\phi \},`$ (10)
$`_\mu \left[e^{2a\varphi }F^{\mu \nu }\right]=8\pi e_0e^{a\varphi }\phi ^{}g^{\nu \lambda }(P_\lambda +e_0A_\lambda )\phi ,`$ (11)
where $`F^2=F_{\mu \nu }F^{\mu \nu }`$ and $`F_{\mu \nu }^2=F_{\mu \lambda }F_\nu ^\lambda `$. $`R_{\mu \nu }`$ is the Ricci tensor.
Now, in order to consider only the lowest order in a typical velocity $`v`$ of “extreme black holes”, we assume the following ansätze: These ansätse are the same as those in .
$`ds^2`$ $`=`$ $`U^2\left(dt+B_idx^i\right)^2+U^2d𝒙^2,`$ (12)
$`U(𝒙)`$ $`=`$ $`V(𝒙)^{\frac{1}{1+a^2}},`$ (13)
$`e^{2a\varphi }`$ $`=`$ $`V^{\frac{2a^2}{1+a^2}},`$ (14)
$`A_0(𝒙)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+a^2}}}\left(1{\displaystyle \frac{1}{V}}\right),`$ (15)
$`A_i(𝒙)B_i(𝒙)`$ $`=`$ $`O(v),`$ (16)
where $`i,j,\mathrm{}`$ denotes the spatial indices.
If there is no matter source and the vacuum is static ($`A_i=B_i=0`$), the ansätze together with the mass-charge relation
$$\frac{e_0}{m}=\sqrt{1+a^2},$$
(17)
reduce the field equations to
$$^2V=0,$$
(18)
which implies that the vacuum solution represents an arbitrary number of “extreme black holes”. One can find that the relation between mass $`m`$ and electric charge $`e_0`$ of a particle corresponds to “extreme black holes”. Thus we take this relation (17) hereafter.
Since we consider the low energy limit, $`P_0m=Emm`$, we find
$`P_0+e_0A_0`$ $`=`$ $`P_0+{\displaystyle \frac{e_0}{\sqrt{1+a^2}}}\left(1{\displaystyle \frac{1}{V}}\right)=P_0+m\left(1{\displaystyle \frac{1}{V}}\right)`$ (19)
$``$ $`m{\displaystyle \frac{1}{V}},`$ (20)
$`P_i+e_0A_iB_i(P_0+e_0A_0)`$ $``$ $`P_i+e_0\left(A_i+{\displaystyle \frac{1}{\sqrt{1+a^2}}}{\displaystyle \frac{1}{V}}B_i\right)`$ (21)
$``$ $`P_i+e_0\widehat{A}_i,`$ (22)
where
$$\widehat{A}_iA_i+\frac{1}{\sqrt{1+a^2}}\frac{1}{V}B_i,$$
(23)
and we define
$`\widehat{F}_{ij}`$ $``$ $`_i\widehat{A}_j_j\widehat{A}_i`$ (24)
$`=`$ $`\overline{F}_{ij}+{\displaystyle \frac{1}{\sqrt{1+a^2}}}{\displaystyle \frac{1}{V}}G_{ij},`$ (25)
where
$`\overline{F}_{ij}`$ $``$ $`F_{ij}+B_iF_{j0}B_jF_{i0},`$ (26)
$`G_{ij}`$ $``$ $`_iB_j_jB_i.`$ (27)
Taking the low energy or non-relativistic limit $`P_0m=Emm`$, $`|P_i+e\widehat{A}_i|^2m^2v^2m^2`$, we simplify the field equations (7), (10) and (11) explicitly. From the dilaton field equation (7), the time-time component of the gravitational field equation (10), and the temporal component of the electromagnetic field equation (11), in the lowest order, we obtain
$$^2V+8\pi (1+a^2)m^2U^3|\phi |^2=0.$$
(28)
From the time-space component of the gravitational field equation (10), we get
$`_{\mathrm{}}\left[V^{\frac{2(a^21)}{1+a^2}}{\displaystyle \frac{1}{V^2}}G_\mathrm{}i\right]`$ $`=`$ $`4\sqrt{{\displaystyle \frac{1}{1+a^2}}}_{\mathrm{}}\left[V^{\frac{2(a^21)}{1+a^2}}{\displaystyle \frac{1}{V}}\overline{F}_\mathrm{}i\right]+\mathrm{\hspace{0.17em}4}\sqrt{{\displaystyle \frac{1}{1+a^2}}}{\displaystyle \frac{1}{V}}_{\mathrm{}}\left[V^{\frac{2(a^21)}{1+a^2}}\overline{F}_\mathrm{}i\right]`$ (30)
$`\mathrm{\hspace{0.25em}32}\pi {\displaystyle \frac{m}{V}}e^{a\varphi }\phi ^{}(P_i+e_0A_i)\phi ,`$
On the other hand, the spatial component of the electromagnetic field equation (11) reads
$$_{\mathrm{}}\left[V^{\frac{2(a^21)}{1+a^2}}\overline{F}_\mathrm{}i\right]=8\pi e_0e^{a\varphi }\phi ^{}(P_i+e_0A_i)\phi .$$
(31)
Here we define an antisymmetric tensor field $`H_{ij}`$ as
$$H_{ij}4\sqrt{\frac{1}{1+a^2}}\overline{F}_{ij}+\frac{1}{V}G_{ij}.$$
(32)
From Eqs. (30) and (31), the equation for the antisymmetric tensor field is found to be
$$_{\mathrm{}}\left[V^{\frac{2(a^21)}{1+a^2}}\frac{1}{V}H_\mathrm{}i\right]=0.$$
(33)
Note that from Eqs. (25) and (32), we can obtain the following relation:
$$\frac{1}{4}\frac{1}{V^2}G^2\overline{F}^2=\frac{1+a^2}{3a^2}\left(\widehat{F}^2\frac{1}{4}H^2\right),$$
(34)
where we assumed $`a^2/=3`$.
Now we consider the effective lagrangian up to the lowest order. Taking the above estimations into the total action (6), we find:
$$S=d^4x,$$
(35)
where the effective lagrangian density is
$``$ $`=`$ $`\sqrt{g}\{{\displaystyle \frac{1}{16\pi }}[R2(\varphi )^2e^{2a\varphi }F^2]`$ (37)
$`+[\phi ^{}e^{a\varphi }g^{\mu \nu }(P_\mu +e_0A_\mu )(P_\nu +e_0A_\nu )\phi m^2e^{a\varphi }\phi ^{}\phi \left]\right\}`$
$``$ $`{\displaystyle \frac{1}{16\pi }}V^{\frac{2(a^21)}{1+a^2}}\left[{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{V^2}}G^2\overline{F}^2\right]`$ (39)
$`+\left[VU^3\phi ^{}\left(P_0+e_0A_0\right)^2\phi {\displaystyle \frac{1}{V}}m^2U^3\phi ^{}\phi {\displaystyle \frac{1}{V}}V^{\frac{2(a^21)}{1+a^2}}U^3\phi ^{}\left(P_i+e_0\widehat{A}_i\right)^2\phi \right]`$
$`=`$ $`{\displaystyle \frac{1}{16\pi }}V^{\frac{2(a^21)}{1+a^2}}\left[{\displaystyle \frac{1+a^2}{3a^2}}\left(\widehat{F}^2{\displaystyle \frac{1}{4}}H^2\right)\right]`$ (41)
$`+\left[VU^3\phi ^{}\left(P_0+e_0A_0\right)^2\phi {\displaystyle \frac{1}{V}}m^2U^3\phi ^{}\phi {\displaystyle \frac{1}{V}}V^{\frac{2(a^21)}{1+a^2}}U^3\phi ^{}\left(P_i+e_0\widehat{A}_i\right)^2\phi \right].`$
In the first term of this effective lagrangian, $`H_{ij}`$ can be regarded as an independent field, for $`H_{ij}`$ does not couple to the scalar field $`\phi `$. Because we consider only the interactions among black holes, we can set $`H_{ij}0`$.
Then Eq. (30), which comes from the time-space component of the gravitational field equation (10), and Eq. (31), from the spatial component of the electromagnetic field equation (11), can be read as
$$\frac{1+a^2}{3a^2}_{\mathrm{}}\left[V^{\frac{2(a^21)}{1+a^2}}\widehat{F}_\mathrm{}i\right]=8\pi e_0e^{a\varphi }\phi ^{}(P_i+e_0\widehat{A}_i)\phi ,$$
(42)
where we assumed $`a^2/=3`$. For $`a^2=3`$, the scalar field does not couple to the force field. For a while, we take $`a^2/=3`$.
To proceed further, we introduce a non-relativistic field $`\psi `$:
$$\psi \sqrt{2m}U^{3/2}\phi ,$$
(43)
where for the spatial measure $`(g^{(3)})^{1/4}=U^{3/2}`$, we obtain a correct measure for a usual spatial volume.
Finally, using the approximation
$`\left(P_0+e_0A_0\right)^2`$ $`=`$ $`\left(P_0m+{\displaystyle \frac{m}{V}}\right)^2`$ (44)
$``$ $`{\displaystyle \frac{m^2}{V^2}}+2{\displaystyle \frac{m}{V}}\left(P_0m\right),`$ (45)
together, we get the effective lagrangian density in the low energy limit:
$``$ $`=`$ $`\psi ^{}\left(P_0m\right)\psi {\displaystyle \frac{1}{2mV^{(3a^2)/(1+a^2)}}}\psi ^{}\left(𝑷+e_0\widehat{𝑨}\right)^2\psi `$ (47)
$`+{\displaystyle \frac{1}{16\pi }}{\displaystyle \frac{1+a^2}{3a^2}}{\displaystyle \frac{1}{V^{2(1a^2)/(1+a^2)}}}\widehat{F}^2(a^2/=3),`$
where $`V`$ satisfies the following equation:
$$^2V+4\pi (1+a^2)m|\psi |^2=0.$$
(48)
As a check, varying this effective lagrangian (47) with respect to $`\widehat{𝑨}`$, we can derive again the field equation (42) in the low energy approximation:
$`{\displaystyle \frac{1+a^2}{3a^2}}_{\mathrm{}}\left[V^{\frac{2(a^21)}{1+a^2}}\widehat{F}_\mathrm{}i\right]`$ $`=`$ $`4\pi {\displaystyle \frac{e_0}{m}}V^{\frac{a^23}{1+a^2}}\psi ^{}(P_i+e_0\widehat{A}_i)\psi `$ (49)
$`=`$ $`8\pi e_0e^{a\varphi }\phi ^{}(P_i+e_0\widehat{A}_i)\phi (a^2/=3).`$ (50)
For $`a^2=3`$, since the scalar field does not couple to the vector potential, the effective lagrangian density at the lowest order is
$$=\psi ^{}\left(P_0m\right)\psi \frac{1}{2m}\psi ^{}𝑷^2\psi (a^2=3),$$
(51)
which seems to describe a free field in the low energy limit.
## III Gas of “Extreme Black Holes” at Finite Temperature
In this section, we apply the effective theory to the study of the thermal system of “extreme black holes”. First we rewrite the effective lagrangian (47) derived in the preceding section as follows: We assume the case of $`a^2/=3`$, unless particularly indicated.
$$=\psi ^{}i\frac{}{t}\psi \frac{1}{2\stackrel{~}{m}}\psi ^{}\left(𝑷+\stackrel{~}{e}\widehat{𝑨}\right)^2\psi +\frac{1}{16\pi }\widehat{F}^2,$$
(52)
where we have used the notations:
$`\stackrel{~}{m}`$ $``$ $`mV^{\frac{3a^2}{1+a^2}},`$ (53)
$`\stackrel{~}{e}^2`$ $``$ $`{\displaystyle \frac{3a^2}{1+a^2}}V^{\frac{2(1a^2)}{1+a^2}}e_0^2`$ (54)
$`=`$ $`(3a^2)V^{\frac{2(1a^2)}{1+a^2}}m^2.`$ (55)
Here we have used the relation between mass $`m`$ and electric charge $`e_0`$ is again
$$\frac{e_0}{m}=\sqrt{1+a^2}.$$
(56)
Next we consider the field theory at finite temperature. In this approach, we assume that $`V`$ is nearly constant and ignore the problem of operator ordering. We define the propagator of the vector field $`\widehat{𝑨}`$:
$$\nu _{ij}(𝒒)=\nu (𝒒)\left(\delta _{ij}\frac{q_iq_j}{q^2}\right).$$
(57)
In the lowest order, the self-energy of the scalar field, here assumed to follow Bose-Einstein statistics, is <sup>§</sup><sup>§</sup>§There is no tadopole contribution, because of the derivative coupling.
$`\mathrm{\Sigma }(𝒌)`$ $`=`$ $`{\displaystyle \frac{1}{\beta 𝒱}}{\displaystyle \frac{1}{4\stackrel{~}{m}^2}}{\displaystyle \underset{q,\mathrm{}}{}}{\displaystyle \frac{\nu (𝒌𝒒)}{i\omega _{\mathrm{}}ϵ_q^{}}}\left[(𝒌+𝒒)^2{\displaystyle \frac{\{(𝒌+𝒒)(𝒌𝒒)\}^2}{(𝒌𝒒)^2}}\right],`$ (58)
where $`𝒱`$ is the volume of the system, $`\beta =1/T`$ ($`T`$ is the temperature of the system) and $`\omega _{\mathrm{}}=\mathrm{\hspace{0.25em}2}\mathrm{}\pi /\beta `$. $`ϵ_k^{}`$ is considered as
$$ϵ_k^{}=\frac{1}{2\stackrel{~}{m}}𝒌^2+\mathrm{\Sigma }(𝒌)\mu ,$$
(59)
where $`\mu `$ denotes the chemical potential for the scalar field. In the lowest order in interactions, $`\nu (𝒒)=\stackrel{~}{e}^2/𝒒^2`$. Recalling that we have assumed that the variation of the background field, or, that of $`V`$ is small.
We then transform the sum over $`𝒒`$ into an integral representation:
$$\frac{1}{𝒱}\underset{𝒒}{}\frac{d^3𝒒}{(2\pi )^3},$$
(60)
and the sum over $`\mathrm{}`$ in this case yields
$$\underset{\mathrm{}}{}\frac{1}{i\omega _{\mathrm{}}x}=\frac{\beta }{e^{\beta x}1}.$$
(61)
Using these, we can rewrite the self-energy of the scalar field (58) as
$`\mathrm{\Sigma }(𝒌)`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{e}^2}{4\stackrel{~}{m}^2}}{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle d^3𝒒f_q\frac{1}{(𝒌𝒒)^2}\left[(𝒌+𝒒)^2\frac{\{(𝒌+𝒒)(𝒌𝒒)\}^2}{(𝒌𝒒)^2}\right]}`$ (62)
$`=`$ $`{\displaystyle \frac{\stackrel{~}{e}^2}{\stackrel{~}{m}^2}}{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle _0^{\mathrm{}}}𝑑qq^2f_q\left[1{\displaystyle \frac{k^2+q^2}{2kq}}\mathrm{ln}\left|{\displaystyle \frac{kq}{k+q}}\right|\right],`$ (63)
where $`f_q`$ is a distribution function of the scalar field:
$$f_q\frac{1}{e^{\beta ϵ_q^{}}1}.$$
(64)
For a small $`k`$, corresponding to the low energy, we find:
$$1\frac{k^2+q^2}{2kq}\mathrm{ln}\left|\frac{kq}{k+q}\right|\frac{4k^2}{3q^2}.$$
(65)
Since there is no constant term in terms of $`k`$, we can write the self-energy approximately as
$$\mathrm{\Sigma }(𝒌)\frac{1}{2}Bk^2,$$
(66)
where $`B`$ is a constant.
From now on we consider the high-temperature approximation. The particles are still non-relativistic and obey the Maxwell-Boltzmann distribution. Then the distribution function of the scalar field is reduced to
$$f_qe^{\beta ϵ_q^{}}.$$
(67)
Therefore the self-consistent equation is expressed as
$`\mathrm{\Sigma }(𝒌)`$ $``$ $`{\displaystyle \frac{1}{2}}Bk^2`$ (68)
$``$ $`{\displaystyle \frac{\stackrel{~}{e}^2}{4\pi ^2\stackrel{~}{m}^2}}{\displaystyle _0^{\mathrm{}}}𝑑qq^2f_q{\displaystyle \frac{4k^2}{3q^2}}`$ (69)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{4\stackrel{~}{e}^2e^{\beta \mu }}{3(2\pi )^{3/2}\stackrel{~}{m}^2}}\left[\beta \left({\displaystyle \frac{1}{\stackrel{~}{m}}}+B\right)\right]^{1/2}k^2.`$ (70)
Equivalently, we get an equation for the coefficient $`B`$:
$$B=\frac{4\stackrel{~}{e}^2e^{\beta \mu }}{3(2\pi )^{3/2}\stackrel{~}{m}^2}\left[\beta \left(\frac{1}{\stackrel{~}{m}}+B\right)\right]^{1/2}.$$
(71)
The solution of Eq. (71) is
$$B=\frac{4\stackrel{~}{e}^2e^{\beta \mu }}{3(2\pi )^{3/2}\stackrel{~}{m}^2}\sqrt{\frac{\stackrel{~}{m}}{\beta }}\left(\mathrm{cos}\frac{\theta }{3}+\frac{1}{\sqrt{3}}\mathrm{sin}\frac{\theta }{3}\right)^1,$$
(72)
where
$$\theta =\mathrm{arcsin}\frac{2\sqrt{3}\stackrel{~}{e}^2e^{\beta \mu }}{(2\pi )^{3/2}\sqrt{\beta \stackrel{~}{m}}}.$$
(73)
On the other hand, the particle density $`\rho `$ is expressed as
$`\rho `$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle d^3𝒒f_q}`$ (74)
$`=`$ $`{\displaystyle \frac{e^{\beta \mu }}{(2\pi )^{3/2}}}\left[\beta \left({\displaystyle \frac{1}{\stackrel{~}{m}}}+B\right)\right]^{3/2}.`$ (75)
In our model the particle density $`\rho `$ can also be written by $`|\psi |^2`$, therefore
$$|\psi |^2=\rho =e^{\beta \mu }\left(\frac{\stackrel{~}{m}}{2\pi \beta }\right)^{3/2}\left[\mathrm{cos}\frac{\theta }{3}+\frac{1}{\sqrt{3}}\mathrm{sin}\frac{\theta }{3}\right]^3.$$
(76)
Now we have to solve the following equation with $`|\psi |^2`$ given by Eq. (76)
$$^2V+4\pi (1+a^2)m|\psi |^2=0.$$
(77)
To explicitly solve Eq. (77), let us assume spherical symmetry and define the following parameters and the radial coordinate:
$`\rho _0`$ $``$ $`e^{\beta \mu }\left({\displaystyle \frac{m}{2\pi \beta }}\right)^{3/2},`$ (78)
$`\delta `$ $``$ $`2\sqrt{3}G\rho _0\beta ,`$ (79)
$`\stackrel{~}{r}`$ $``$ $`\sqrt{Gm\rho _0}r,`$ (80)
where $`G`$ is the Newton constant, being set to unity until now. Incidentally, $`n_0`$ is the particle density of the ideal gas of particles with mass $`m`$. Then Eq. (77) for $`V`$ is
$`{\displaystyle \frac{1}{\stackrel{~}{r}^2}}\left(\stackrel{~}{r}^2V(\stackrel{~}{r})^{^{}}\right)^{}`$ $`=`$ $`4\pi (1+a^2)V(\stackrel{~}{r})^{\frac{3(3a^2)}{2(1+a^2)}}\left(\mathrm{cos}{\displaystyle \frac{\theta }{3}}+{\displaystyle \frac{1}{\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{\theta }{3}}\right)^3,`$ (81)
where the prime denotes derivative with respect to $`\stackrel{~}{r}`$, and
$$\theta (\stackrel{~}{r})=\mathrm{arcsin}\left[(3a^2)V(\stackrel{~}{r})^{\frac{(13a^2)}{2(1+a^2)}}\delta \right].$$
(82)
We can numerically solve Eq. (81) with Eq. (82) to obtain $`V(\stackrel{~}{r})`$ and the particle density $`\rho =|\psi |^2`$ of the isothermal sphere of “extreme black holes”. Fig. 1 shows the profile of the particle density of the isothermal sphere. The central density is normalized to unity there. As seen from Fig. 1, as $`\delta `$ is the smaller, that is, the temperature is the higher, the particles are bound the more tightly.
We find also that as the coupling constant of the dilaton field is the larger, the particles are bound the more tightly. In particular, there is a critical value for the asymptotic behavior of $`n`$ at a large distance: for $`a^2>1`$, the isothermal sphere at high temperature has a clear “edge”.
For the case $`a^2=3`$, the effective lagrangian represents the free lagrangian, therefore the particles correspond to the free particles.
The physical interpretation on the behavior is given by the fact that the “magnetic” force acts attractively between the particles which move in the opposite direction for $`a^2<3`$ (see the relative sign of the coefficient in Eq. (50)). The thermal average of $`v^2`$ governs the strength of the force; this feature has been suggested in .
We have also seen that there is the critical point at $`a^21`$ whether an extent of the “edge” exists or not. This behavior implies the existence of an effective repulsion in a small range, because similar behavior can be found in the model of boson stars with a repulsive force. This critical nature of the interaction seems to correspond to the fact that the moduli space of two “extreme black holes” has a deficit angle $`\pi `$ for $`a^2=1`$. Further study on this subject is expected.
## IV Eliminating the potentials
We consider the total energy of the system of “extreme black holes”. We can eliminate the field stength and the potential apparently in the expression of the energy using the solution of the field equations expressed by the scalar field. Assuming localized scalar distributions in the resulting expression, we can also obtain the energy for the classical system. For a while, we assume $`a^2/=3`$.
The total energy of the system is written by
$$H=d^3𝒙,$$
(83)
with
$$=\frac{1}{2mV^{(3a^2)/(1+a^2)}}\psi ^{}\left(𝑷+e_0\widehat{𝑨}\right)^2\psi \frac{1}{16\pi }\frac{1+a^2}{3a^2}\frac{1}{V^{2(1a^2)/(1+a^2)}}\widehat{F}^2,$$
(84)
where $`e_0=\sqrt{1+a^2}m`$.
First, we treat the second term of the Hamiltonian density $``$. We use the following definition:
$`𝑱(𝒙)={\displaystyle \frac{}{𝑷}}`$ $`=`$ $`{\displaystyle \frac{1}{mV^{(3a^2)/(1+a^2)}}}\psi ^{}\left(𝑷+e_0\widehat{𝑨}\right)\psi `$ (85)
$``$ $`\psi ^{}(𝒙)𝒗(𝒙)\psi (𝒙).`$ (86)
Then the field equation (50) can be written as
$$\frac{1+a^2}{3a^2}_{\mathrm{}}\left[V^{\frac{2(a^21)}{1+a^2}}\widehat{F}_\mathrm{}i\right]=4\pi e_0𝑱.$$
(87)
Note that
$$\mathbf{}𝑱=_iJ^i=0.$$
(88)
As long as $`𝑱`$ vanishes rapidly at the spatial infinity, the solution is given by
$`{\displaystyle \frac{1+a^2}{3a^2}}V^{\frac{2(a^21)}{1+a^2}}\widehat{F}_\mathrm{}i`$ $`=`$ $`e_0{\displaystyle d^3𝒙^{}\frac{r_{}^{}{}_{}{}^{\mathrm{}}J^i(𝒙^{})r_{}^{}{}_{}{}^{i}J^{\mathrm{}}(𝒙^{})}{|𝒓^{}|^3}},`$ (89)
where $`𝒓^{}=𝒙𝒙^{}`$. Substituting this solution, we obtain
$`{\displaystyle d^3𝒙\frac{1}{16\pi }\frac{1+a^2}{3a^2}\frac{1}{V^{2(1a^2)/(1+a^2)}}\widehat{F}^2}`$ (90)
$`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{3a^2}{1+a^2}}e_0^2{\displaystyle d^3𝒙V^{\frac{2(1a^2)}{1+a^2}}(𝒙)}`$ (92)
$`\times {\displaystyle }d^3𝒙^{}{\displaystyle }d^3𝒙^{\prime \prime }{\displaystyle \frac{𝒓^{}𝒓^{\prime \prime }𝑱(𝒙^{})𝑱(𝒙^{\prime \prime })𝒓^{}𝑱(𝒙^{\prime \prime })𝒓^{\prime \prime }𝑱(𝒙^{})}{|𝒓^{}|^3|𝒓^{\prime \prime }|^3}},`$
where $`𝒓^{\prime \prime }=𝒙𝒙^{\prime \prime }`$.
To garantee the absence of the influence of the spatial infinity explicitly, we use the following identity which holds if $`𝑱`$ vanishes rapidly at the spatial infinity:
$$d^3𝒙^{}\frac{𝒓^{}𝑱(𝒙^{})}{|𝒓^{}|^3}=0,$$
(93)
and modify the equation (92) into
$`{\displaystyle d^3𝒙\frac{1}{16\pi }\frac{1+a^2}{3a^2}\frac{1}{V^{2(1a^2)/(1+a^2)}}\widehat{F}^2}`$ (94)
$`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{3a^2}{1+a^2}}e_0^2{\displaystyle d^3𝒙V^{\frac{2(1a^2)}{1+a^2}}(𝒙)}`$ (96)
$`\times {\displaystyle }d^3𝒙^{}{\displaystyle }d^3𝒙^{\prime \prime }{\displaystyle \frac{𝒓^{}𝒓^{\prime \prime }𝑱(𝒙^{})𝑱(𝒙^{\prime \prime })𝒓^{}𝑱(𝒙^{\prime \prime })𝒓^{\prime \prime }𝑱(𝒙^{})+𝒓^{}𝑱(𝒙^{})𝒓^{\prime \prime }𝑱(𝒙^{\prime \prime })}{|𝒓^{}|^3|𝒓^{\prime \prime }|^3}}`$
$`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{3a^2}{1+a^2}}e_0^2{\displaystyle d^3𝒙V^{\frac{2(1a^2)}{1+a^2}}(𝒙)}`$ (98)
$`\times {\displaystyle }d^3𝒙^{}{\displaystyle }d^3𝒙^{\prime \prime }{\displaystyle \frac{r_{}^{}{}_{}{}^{i}r_{}^{\prime \prime }{}_{}{}^{j}}{|𝒓^{}|^3|𝒓^{\prime \prime }|^3}}(\delta _k\mathrm{}\delta _{ij}\delta _i\mathrm{}\delta _{jk}+\delta _{ik}\delta _j\mathrm{})J^k(𝒙^{})J^{\mathrm{}}(𝒙^{\prime \prime }).`$
Next, we consider the first term of the Hamiltonian density $``$. If we define
$$\rho (𝒙)=\psi ^{}(𝒙)\psi (𝒙),$$
(99)
then the equation (48), which is now written as
$$^2V+4\pi (1+a^2)m\rho =0,$$
(100)
has the solution:
$$V(𝒙)=1+(1+a^2)md^3𝒙^{}\frac{1}{|𝒓^{}|}\rho (𝒙^{}).$$
(101)
We rearrange the first term of $``$ as
$`{\displaystyle \frac{1}{2mV^{(3a^2)/(1+a^2)}}}\psi ^{}\left(𝑷+e_0\widehat{𝑨}\right)^2\psi `$ $`=`$ $`{\displaystyle \frac{1}{2}}mV^{(3a^2)/(1+a^2)}\psi ^{}𝒗^2\psi `$ (102)
$`=`$ $`{\displaystyle \frac{1}{2}}m\psi ^{}𝒗^2\psi +{\displaystyle \frac{1}{2}}m\left[V^{\frac{3a^2}{1+a^2}}1\right]\psi ^{}𝒗^2\psi .`$ (103)
Using Eq (101), we find
$`{\displaystyle \frac{1}{2}}m{\displaystyle d^3𝒙\left[V^{\frac{3a^2}{1+a^2}}(𝒙)1\right]\psi ^{}(𝒙)𝒗^2(𝒙)\psi (𝒙)}`$ (104)
$`=`$ $`{\displaystyle \frac{1}{8\pi }}m{\displaystyle d^3𝒙\left[V^{\frac{3a^2}{1+a^2}}(𝒙)1\right]d^3𝒙^{}\mathbf{}\frac{𝒓^{}}{|𝒓^{}|^3}\psi ^{}(𝒙^{})𝒗^2(𝒙^{})\psi (𝒙^{})}`$ (105)
$`=`$ $`{\displaystyle \frac{3a^2}{8\pi }}m^2{\displaystyle d^3𝒙V^{\frac{2(1a^2)}{1+a^2}}(𝒙)}`$ (107)
$`\times {\displaystyle }d^3𝒙^{}{\displaystyle }d^3𝒙^{\prime \prime }{\displaystyle \frac{𝒓^{}𝒓^{\prime \prime }}{|𝒓^{}|^3|𝒓^{\prime \prime }|^3}}\psi ^{}(𝒙^{\prime \prime })\psi (𝒙^{\prime \prime })\psi ^{}(𝒙^{})𝒗^2(𝒙^{})\psi (𝒙^{}).`$
Consequently, the interaction energy can be written as
$`{\displaystyle d^3𝒙\left\{\frac{1}{2}m\left[V^{\frac{3a^2}{1+a^2}}(𝒙)1\right]\psi ^{}(𝒙)𝒗^2(𝒙)\psi (𝒙)\frac{1}{16\pi }\frac{1+a^2}{3a^2}\frac{1}{V^{2(1a^2)/(1+a^2)}}\widehat{F}^2\right\}}`$ (108)
$`=`$ $`{\displaystyle \frac{3a^2}{8\pi }}m^2{\displaystyle d^3𝒙V^{\frac{2(1a^2)}{1+a^2}}(𝒙)d^3𝒙^{}d^3𝒙^{\prime \prime }\frac{1}{|𝒓^{}|^3|𝒓^{\prime \prime }|^3}}`$ (110)
$`\times \psi ^{}(𝒙^{})\psi ^{}(𝒙^{\prime \prime })\left\{{\displaystyle \frac{1}{2}}𝒓^{}𝒓^{\prime \prime }|𝒗(𝒙^{})𝒗(𝒙^{\prime \prime })|^2(𝒓^{}\times 𝒓^{\prime \prime })(𝒗(𝒙^{})\times 𝒗(𝒙^{\prime \prime }))\right\}\psi (𝒙^{})\psi (𝒙^{\prime \prime }).`$
Because we know that there is no interaction for the $`a^2=3`$ case, we see that this expression holds for any value of the dilaton coupling $`a^2`$.
For $`a^2=1`$, the expression for the total enegy can be simplified into
$`H`$ $`=`$ $`{\displaystyle d^3𝒙\frac{1}{2}m\psi ^{}(𝒙)𝒗^2(𝒙)\psi (𝒙)}`$ (112)
$`+{\displaystyle \frac{1}{8\pi }}m^2{\displaystyle d^3𝒙d^3𝒙^{}d^3𝒙^{\prime \prime }\frac{𝒓^{}𝒓^{\prime \prime }}{|𝒓^{}|^3|𝒓^{\prime \prime }|^3}\psi ^{}(𝒙^{})\psi ^{}(𝒙^{\prime \prime })|𝒗(𝒙^{})𝒗(𝒙^{\prime \prime })|^2\psi (𝒙^{})\psi (𝒙^{\prime \prime })}`$
$`=`$ $`{\displaystyle d^3𝒙\frac{1}{2}m\psi ^{}(𝒙)𝒗^2(𝒙)\psi (𝒙)}`$ (114)
$`+{\displaystyle \frac{1}{2}}m^2{\displaystyle d^3𝒙d^3𝒙^{}\psi ^{}(𝒙)\psi ^{}(𝒙^{})\frac{|𝒗(𝒙)𝒗(𝒙^{})|^2}{|𝒙𝒙^{}|}\psi (𝒙)\psi (𝒙^{})}.`$
Now we consider the classical point-particle limit,
$`𝑱(𝒙)`$ $`=`$ $`\psi ^{}(𝒙)𝒗(𝒙)\psi (𝒙)={\displaystyle \underset{a}{}}\mu _a𝒗_a\delta ^3(𝒓_a),`$ (115)
$`\rho (𝒙)`$ $`=`$ $`\psi ^{}(𝒙)\psi (𝒙)={\displaystyle \underset{a}{}}\mu _a\delta ^3(𝒓_a),`$ (116)
and
$$\psi ^{}(𝒙)𝒗^2(𝒙)\psi (𝒙)=\underset{a}{}\mu _a𝒗_a^2\delta ^3(𝒓_a),$$
(117)
where $`\mu _a`$ and $`𝒗_a`$ are the constant which represents the ratio of the mass and the velocity of the $`a`$-th “extreme black hole” located at $`𝒙_a`$, respectively. We use the notation $`𝒓_a=𝒙𝒙_a`$. Then one can find that the energy of the classical system takes the form:
$`H`$ $`=`$ $`{\displaystyle \underset{a}{}}{\displaystyle \frac{1}{2}}m_a𝒗_a^2`$ (120)
$`+{\displaystyle \frac{3a^2}{8\pi }}{\displaystyle d^3𝒙V^{\frac{2(1a^2)}{1+a^2}}(𝒙)\underset{ab}{}\frac{m_am_b}{|𝒓_a|^3|𝒓_b|^3}}`$
$`\times \left\{{\displaystyle \frac{1}{2}}𝒓_a𝒓_b|𝒗_a𝒗_b|^2(𝒓_a\times 𝒓_b)(𝒗_a\times 𝒗_b)\right\},`$
with
$$V(𝒙)=1+(1+a^2)\underset{c}{}\frac{m_c}{|𝒓_c|},$$
(121)
where the individual mass is $`m_am\mu _a`$. For $`a^2=1`$, the energy has a simple form:
$$H=\underset{a}{}\frac{1}{2}m_a𝒗_a^2+\frac{1}{2}\underset{ab}{}m_am_b\frac{|𝒗_{ab}|^2}{|𝒓_{ab}|},$$
(122)
where $`𝒓_{ab}𝒙_a𝒙_b`$ and $`𝒗_{ab}𝒗_a𝒗_b`$.
Furthermore, we restrict ourselves on the two-body system. We assume that the velocity of the center of mass $`𝐕`$ vanishes:
$$𝐕\frac{m_1𝒗_1+m_2𝒗_2}{M}=\mathrm{𝟎},$$
(123)
where $`M=m_1+m_2`$. The velocity of the relative motion is defined as
$$𝐯𝒗_1𝒗_2.$$
(124)
Then Eq. (120) becomes
$$H=\frac{1}{2}\mu 𝐯^2+\frac{3a^2}{8\pi }\mu M𝐯^2d^3𝒙V^{\frac{2(1a^2)}{1+a^2}}(𝒙)\frac{𝒓_1𝒓_2}{|𝒓_1|^3|𝒓_2|^3},$$
(125)
with
$$V(𝒙)=1+(1+a^2)\left[\frac{m_1}{|𝒓_1|}+\frac{m_2}{|𝒓_2|}\right],$$
(126)
where the reduced mass $`\mu `$ is given by $`\mu =m_1m_2/M`$.
In general, a naive integration in Eq. (125) diverges. Therefore we realize that divergent terms proportional to $`d^3𝒙\delta ^3(𝒙)/|𝒙|^p(p>0)`$ which appear when the integrand is expanded must be regularized. We set them to zero. The prescription is equivalent to carrying out the following replacement in Eq. (125):
$$V^{\frac{2(1a^2)}{1+a^2}}(𝒙)\left[1+(1+a^2)\frac{m_1}{|𝒓_1|}\right]^{\frac{2(1a^2)}{1+a^2}}+\left[1+(1+a^2)\frac{m_2}{|𝒓_2|}\right]^{\frac{2(1a^2)}{1+a^2}}1.$$
(127)
Then we get
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mu 𝐯^2\{1{\displaystyle \frac{M}{\mu }}{\displaystyle \frac{(3a^2)M}{r}}`$ (129)
$`+{\displaystyle \frac{M}{m_1}}[1+(1+a^2){\displaystyle \frac{m_1}{r}}]^{\frac{3a^2}{1+a^2}}+{\displaystyle \frac{M}{m_2}}[1+(1+a^2){\displaystyle \frac{m_2}{r}}]^{\frac{3a^2}{1+a^2}}\},`$
where $`r=|𝒙_1𝒙_2|`$.
Before closing this section, we show the results for $`(N+1)`$ dimensional case. We assume $`a^2/=N`$ here. The basic action is given in Appendix. We find that the effective action is
$$=\frac{1}{2mV^{\frac{Na^2}{N2+a^2}}}\psi ^{}\left(𝑷+e_0\widehat{𝑨}\right)^2\psi \frac{1}{16\pi }\frac{N2+a^2}{Na^2}\frac{1}{V^{\frac{2(1a^2)}{N2+a^2}}}\widehat{F}^2,$$
(130)
where
$$e_0^2=\frac{2(N2+a^2)}{N1}m^2.$$
(131)
Here $`V`$ satisfies
$$^2V+8\pi \frac{N2+a^2}{N1}m|\psi |^2=0.$$
(132)
The derivation of the effective action is shown in Appendix.
Similarly to the previous analysis, we obtain the interaction Hamiltonian:
$`H_{int}`$ $`=`$ $`{\displaystyle }d^N𝒙\{{\displaystyle \frac{1}{2}}m[V^{\frac{Na^2}{N2+a^2}}(𝒙)1]\psi ^{}(𝒙)𝒗^2(𝒙)\psi (𝒙)`$ (134)
$`{\displaystyle \frac{1}{16\pi }}{\displaystyle \frac{N2+a^2}{Na^2}}{\displaystyle \frac{1}{V^{2(1a^2)/(N2+a^2)}}}\widehat{F}^2\}`$
$`=`$ $`{\displaystyle \frac{Na^2}{4(N1)\pi }}\left({\displaystyle \frac{4\pi }{A_{N1}}}\right)^2m^2{\displaystyle d^N𝒙V^{\frac{2(1a^2)}{N2+a^2}}(𝒙)d^N𝒙^{}d^N𝒙^{\prime \prime }\frac{1}{|𝒓^{}|^N|𝒓^{\prime \prime }|^N}}`$ (137)
$`\times \psi ^{}(𝒙^{})\psi ^{}(𝒙^{\prime \prime })\{{\displaystyle \frac{1}{2}}𝒓^{}𝒓^{\prime \prime }|𝒗(𝒙^{})𝒗(𝒙^{\prime \prime })|^2(𝒓^{}𝒗(𝒙^{}))(𝒓^{\prime \prime }𝒗(𝒙^{\prime \prime }))`$
$`+(𝒓^{}𝒗(𝒙^{\prime \prime }))(𝒓^{\prime \prime }𝒗(𝒙^{}))\}\psi (𝒙^{})\psi (𝒙^{\prime \prime }),`$
where
$$𝒗(𝒙)\frac{1}{mV^{(Na^2)/(N2+a^2)}}\left(𝑷+e_0\widehat{𝑨}\right),$$
(138)
and $`A_{N1}=2\pi ^N/\mathrm{\Gamma }(N/2)`$ is the volume of a unit $`N1`$ sphere.
Especially for $`a^2=1`$, we find
$$H_{int}=\frac{1}{2}\frac{4\pi }{A_{N1}}m^2d^N𝒙d^N𝒙^{}\psi ^{}(𝒙)\psi ^{}(𝒙^{})\frac{|𝒗(𝒙)𝒗(𝒙^{})|^2}{(N2)|𝒙𝒙^{}|^{N2}}\psi (𝒙)\psi (𝒙^{}).$$
(139)
The classical limit of total energy of the system is found to be
$`H`$ $`=`$ $`{\displaystyle \underset{a}{}}{\displaystyle \frac{1}{2}}m_a𝒗_a^2`$ (142)
$`+{\displaystyle \frac{Na^2}{4(N1)\pi }}\left({\displaystyle \frac{4\pi }{A_{N1}}}\right)^2{\displaystyle d^N𝒙V^{\frac{2(1a^2)}{N2+a^2}}(𝒙)\underset{ab}{}\frac{m_am_b}{|𝒓_a|^N|𝒓_b|^N}}`$
$`\times \left\{{\displaystyle \frac{1}{2}}𝒓_a𝒓_b|𝒗_a𝒗_b|^2(𝒓_a𝒗_a)(𝒓_b𝒗_b)+(𝒓_a𝒗_b)(𝒓_b𝒗_a)\right\},`$
with
$$V(𝒙)=1+\frac{2(N2+a^2)}{(N1)(N2)}\frac{4\pi }{A_{N1}}\underset{c}{}\frac{m_c}{|𝒓_c|^{N2}}.$$
(143)
For $`a^2=1`$, the total energy of the system has a simple form:
$$H=\underset{a}{}\frac{1}{2}m_a𝒗_a^2+\frac{1}{2}\frac{4\pi }{A_{N1}}\underset{ab}{}m_am_b\frac{|𝒗_{ab}|^2}{(N2)|𝒓_{ab}|^{N2}}.$$
(144)
Finally, we show another expression for the total energy. That is
$$H=\underset{ab}{}v^{ak}v^b\mathrm{}(\delta _k^i\delta _{\mathrm{}}^j+\delta _k\mathrm{}\delta ^{ij}\delta _k^j\delta _{\mathrm{}}^i)_{ai}_{bj}L,$$
(145)
where
$$L=\frac{1}{32\pi }d^N𝒙V^{\frac{2(N1)}{N2+a^2}}(𝒙),$$
(146)
with $`V`$ given by Eq. (143).
## V Conclusion
In this paper, we have derived the effective lagrangian of “extreme black holes” in the low energy limit. At finite temperature, we have obtained a self-consistent equation and then we have seen the structure of the isothermal sphere distribution of “extreme black holes”. At high temperature, the gas of “extreme black holes” have been lumped by the velocity-dependent force.
In future work, we will consider the low-temperature case. It is interesting to investigate whether the condensation of “extreme black holes” may take place. We will also take another posibility for the statistics of “extreme black holes” into consideration. We would like to study the effective theory by means of the lattice calculation, in which the strongly coupled system is appropriately treated.
We have also studied the classical point-particle limit for the energy of the system. Moduli space structure and supersymmetric extensions of the multiple black hole system have recently been studied by many authors. We are interested in such a direction of study and expect some symmetric structure to be found in the effective field theory of multi-black holes.
## Acknowledgements
KS would like to thank Satoru Hirenzaki for reading this manuscript. He also thank Nahomi Kan and Yoshinori Cho for useful comments.
##
In this Appendix, we show the derivation of the effective lagrangian for “extreme black holes” in the $`(N+1)`$ dimensional spacetime in detail.
We start with the action for the charged scalar field:
$`S_m`$ $`=`$ $`{\displaystyle }d^{N+1}x\sqrt{g}[\phi ^{}e^{\frac{2a}{N1}\varphi }g^{\mu \nu }(P_\mu +e_0A_\mu )(P_\nu +e_0A_\nu )\phi m^2e^{\frac{2a}{N1}\varphi }\phi ^{}\phi ].`$ (147)
The total action reads
$`S`$ $`=`$ $`{\displaystyle }d^{N+1}x{\displaystyle \frac{\sqrt{g}}{16\pi }}[R{\displaystyle \frac{4}{N1}}_\mu \varphi ^\mu \varphi e^{\frac{4a}{N1}\varphi }F_{\mu \nu }F^{\mu \nu }]+S_m,`$ (148)
and leads to the field equations:
$`^2\varphi +{\displaystyle \frac{a}{2}}e^{\frac{4a}{N1}\varphi }F^2+\mathrm{\hspace{0.17em}4}\pi a\left[e^{\frac{2a}{N1}\varphi }\phi ^{}(P+e_0A)^2\phi e^{\frac{2a}{N1}\varphi }m^2\phi ^{}\phi \right]=0,`$ (149)
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }R={\displaystyle \frac{4}{N1}}\left[_\mu \varphi _\nu \varphi {\displaystyle \frac{1}{2}}g_{\mu \nu }(\varphi )^2\right]+e^{\frac{4a}{N1}\varphi }\left[\mathrm{\hspace{0.17em}2}F_{\mu \nu }^2{\displaystyle \frac{1}{2}}g_{\mu \nu }F^2\right]`$ (150)
$`+16\pi \{e^{\frac{2a}{N1}\varphi }Re[\phi ^{}(P_\mu +e_0A_\mu )(P_\nu +e_0A_\nu )\phi `$ (151)
$`{\displaystyle \frac{1}{2}}g_{\mu \nu }\phi ^{}(P+e_0A)^2\phi ]{\displaystyle \frac{1}{2}}g_{\mu \nu }e^{\frac{2a}{N1}\varphi }m^2\phi ^{}\phi \},`$ (152)
$`_\mu \left[e^{\frac{4a}{N1}\varphi }F^{\mu \nu }\right]=8\pi e_0e^{\frac{2a}{N1}\varphi }\phi ^{}g^{\nu \lambda }(P_\lambda +e_0A_\lambda )\phi .`$ (153)
The ansätze in the $`(N+1)`$ dimensional case are now:
$`ds^2`$ $`=`$ $`U^2\left(dt+B_idx^i\right)^2+U^{\frac{2}{N2}}d𝒙^2,`$ (154)
$`U(𝒙)`$ $`=`$ $`V(𝒙)^{\frac{N2}{N2+a^2}},`$ (155)
$`e^{\frac{4a}{N1}\varphi }`$ $`=`$ $`V^{\frac{2a^2}{N2+a^2}},`$ (156)
$`A_0(𝒙)`$ $`=`$ $`\sqrt{{\displaystyle \frac{N1}{2(N2+a^2)}}}\left(1{\displaystyle \frac{1}{V}}\right),`$ (157)
$`A_i(𝒙)B_i(𝒙)`$ $`=`$ $`O(v).`$ (158)
In addition, the following charge-mass ratio is assumed:
$$\frac{e_0}{m}=\sqrt{\frac{2(N2+a^2)}{N1}}.$$
(159)
This corresponds to that of the “extreme black holes”.
Now we consider the low energy limit, $`P_0m=Emm`$. Then
$`P_0+e_0A_0=P_0+e_0\sqrt{{\displaystyle \frac{N1}{2(N2+a^2)}}}\left(1{\displaystyle \frac{1}{V}}\right)=P_0+m\left(1{\displaystyle \frac{1}{V}}\right)m{\displaystyle \frac{1}{V}},`$ (160)
$`P_i+e_0A_iB_i(P_0+e_0A_0)P_i+e_0\left(A_i+\sqrt{{\displaystyle \frac{N1}{2(N2+a^2)}}}{\displaystyle \frac{1}{V}}B_i\right)P_i+e_0\widehat{A}_i,`$ (161)
where
$$\widehat{A}_iA_i+\sqrt{\frac{N1}{2(N2+a^2)}}\frac{1}{V}B_i,$$
(162)
and we define
$$\widehat{F}_{ij}_i\widehat{A}_j_j\widehat{A}_i=\overline{F}_{ij}+\sqrt{\frac{N1}{2(N2+a^2)}}\frac{1}{V}G_{ij},$$
(163)
where
$`\overline{F}_{ij}`$ $``$ $`F_{ij}+B_iF_{j0}B_jF_{i0},`$ (164)
$`G_{ij}`$ $``$ $`_iB_j_jB_i.`$ (165)
Using the ansätze and taking the low energy or non-relativistic limit $`P_0m=Emm`$, $`|P_i+e\widehat{A}_i|^2m^2v^2m^2`$, we simplify the field equations.
We reduce the dilaton field equation (149), using Eqs. (154), (161), (162), and (164), to
$`U^{\frac{2}{N2}}^2\varphi +{\displaystyle \frac{a}{2}}e^{\frac{4a}{N1}\varphi }\left[\mathrm{\hspace{0.17em}2}U^2U^{\frac{2}{N2}}(F_{0i})^2+U^{\frac{4}{N2}}\overline{F}^2\right]`$ (166)
$`+\mathrm{\hspace{0.17em}4}\pi a\{e^{\frac{2a}{N1}\varphi }[U^2\phi ^{}(P_0+e_0A_0)^2\phi +U^{\frac{2}{N2}}\phi ^{}(P_i+e_0\widehat{A}_i)^2\phi ]`$ (167)
$`e^{\frac{2a}{N1}\varphi }m^2\phi ^{}\phi \}=0,`$ (168)
where $`\overline{F}^2=\overline{F}_{ij}\overline{F}_{ij}`$. Further using Eqs. (155), (156), and (157), we get
$`{\displaystyle \frac{a(N1)}{2(N2+a^2)}}{\displaystyle \frac{1}{V}}^2V+\mathrm{\hspace{0.17em}4}\pi a\left[e^{\frac{2a}{N1}\varphi }U^2\phi ^{}(P_0+e_0A_0)^2\phi +e^{\frac{2a}{N1}\varphi }m^2\phi ^{}\phi \right]U^{\frac{2}{N2}}`$ (169)
$`=`$ $`{\displaystyle \frac{a}{2}}e^{\frac{4a}{N1}\varphi }U^{\frac{2}{N2}}\overline{F}^2+4\pi aU^{\frac{2}{N2}}\phi ^{}(P_i+e_0\widehat{A}_i)^2\phi .`$ (170)
Finally we use (160) and rearrange the equation, and because the right hand side of Eq. (170) is $`O(v^2)`$, we find that the dilaton equation in the lowest order can be reduced to
$$^2V+16\pi \frac{N2+a^2}{N1}m^2U^{\frac{N}{N2}}|\phi |^2=0.$$
(171)
The time-time component of the gravitational field equation (152) can be treated in the same manner. For the first step, we use the metric ansatz (154) and then get
$`{\displaystyle \frac{N1}{N2}}U^{\frac{2}{N2}}_{\mathrm{}}\left({\displaystyle \frac{_{\mathrm{}}U}{U}}\right)+{\displaystyle \frac{1}{2}}{\displaystyle \frac{N1}{N2}}U^{\frac{2}{N2}}\left({\displaystyle \frac{_{\mathrm{}}U}{U}}\right)^2{\displaystyle \frac{3}{8}}U^2U^{\frac{4}{N2}}G^2`$ (172)
$`=`$ $`{\displaystyle \frac{4}{N1}}\left[{\displaystyle \frac{1}{2}}U^{\frac{2}{N2}}\left(_k\varphi \right)^2\right]+e^{\frac{4a}{N1}\varphi }\left[U^2U^{\frac{2}{N2}}\left(F_{0k}\right)^2{\displaystyle \frac{1}{2}}U^{\frac{4}{N2}}\overline{F}^2\right]`$ (174)
$`+\mathrm{\hspace{0.17em}8}\pi \left\{e^{\frac{2a}{N1}\varphi }\left[U^2\phi ^{}(P_0+e_0A_0)^2\phi U^{\frac{2}{N2}}\phi ^{}(P_i+e_0\widehat{A}_i)^2\phi \right]e^{\frac{2a}{N1}\varphi }m^2\phi ^{}\phi \right\},`$
where $`G^2=G_{ij}G_{ij}`$. Next, we use $`V`$ and obtain
$`{\displaystyle \frac{N1}{N2+a^2}}{\displaystyle \frac{1}{V}}^2V+\mathrm{\hspace{0.17em}8}\pi \left[e^{\frac{2a}{N1}\varphi }U^2\phi ^{}(P_0+e_0A_0)^2\phi +e^{\frac{2a}{N1}\varphi }m^2\phi ^{}\phi \right]U^{\frac{2}{N2}}`$ (175)
$`=`$ $`{\displaystyle \frac{3}{8}}U^2U^{\frac{2}{N2}}G^2{\displaystyle \frac{1}{2}}e^{\frac{4a}{N1}}\overline{F}^28\pi e^{\frac{2a}{N1}\varphi }\phi ^{}(P_i+e_0\widehat{A}_i)^2\phi .`$ (176)
Finally we pick up a part of the lowest order. The reduced equation is the same as Eq. (171).
The temporal component of the electromagnetic field equation (153), in the lowest order, can be read as
$`U^{\frac{2}{N2}}_k\left[e^{\frac{4a}{N1}\varphi }\left(U^2F_{0k}+U^{\frac{2}{N2}}B_{\mathrm{}}\overline{F}_\mathrm{}k\right)\right]`$ (177)
$`=`$ $`8\pi e_0e^{\frac{2a}{N1}\varphi }\left[U^2\phi ^{}(P_0+e_0A_0)\phi +U^{\frac{2}{N2}}B_k\phi ^{}(P_k+e_0\widehat{A}_k)\phi \right].`$ (178)
This equation is equivalent to
$`\sqrt{{\displaystyle \frac{N1}{2(N2+a^2)}}}^2V8\pi e_0e^{\frac{2a}{N1}\varphi }U^{\frac{N}{N2}}U\phi ^{}(P_0+e_0A_0)\phi `$ (179)
$`=`$ $`_k\left[e^{\frac{4a}{N1}\varphi }U^{\frac{2}{N2}}B_{\mathrm{}}\overline{F}_\mathrm{}k\right]+8\pi e_0e^{\frac{2a}{N1}\varphi }B_k\phi ^{}(P_k+e_0\widehat{A}_k)\phi .`$ (180)
The right hand side of this equation is $`O(v^2)`$. Together with Eq. (159), we find that Eq. (180) reduces to Eq. (171) in the lowest order in $`v`$.
From the time-space component of the gravitational field equation (152), we obtain
$`{\displaystyle \frac{1}{2}}UU^{\frac{1}{N2}}_{\mathrm{}}\left(U^2U^{\frac{2}{N2}}G_\mathrm{}i\right)`$ (181)
$`=`$ $`\mathrm{\hspace{0.17em}2}e^{\frac{4a}{N1}\varphi }UU^{\frac{3}{N2}}F_{0k}\overline{F}_{ik}16\pi e^{\frac{2a}{N1}\varphi }UU^{\frac{1}{N2}}\left[\phi ^{}\left(P_0+e_0A_0\right)\left(P_i+e_0\widehat{A}_i\right)\phi \right].`$ (182)
This can be reduced to
$`_{\mathrm{}}\left[V^{\frac{2(a^21)}{N2+a^2}}{\displaystyle \frac{1}{V^2}}G_\mathrm{}i\right]`$ (183)
$`=`$ $`\mathrm{\hspace{0.17em}4}\sqrt{{\displaystyle \frac{N1}{2(N2+a^2)}}}_k\left[V^{\frac{2(a^21)}{N2+a^2}}{\displaystyle \frac{1}{V}}\overline{F}_{ki}\right]+4\sqrt{{\displaystyle \frac{N1}{2(N2+a^2)}}}{\displaystyle \frac{1}{V}}_k\left[V^{\frac{2(a^21)}{N2+a^2}}\overline{F}_{ki}\right]`$ (185)
$`+32\pi e^{\frac{2a}{N1}\varphi }\left[\phi ^{}\left(P_0+e_0A_0\right)\left(P_i+e_0\widehat{A}_i\right)\phi \right]`$
$``$ $`\mathrm{\hspace{0.17em}4}\sqrt{{\displaystyle \frac{N1}{2(N2+a^2)}}}_k\left[V^{\frac{2(a^21)}{N2+a^2}}{\displaystyle \frac{1}{V}}\overline{F}_{ki}\right]+4\sqrt{{\displaystyle \frac{N1}{2(N2+a^2)}}}{\displaystyle \frac{1}{V}}_k\left[V^{\frac{2(a^21)}{N2+a^2}}\overline{F}_{ki}\right]`$ (187)
$`32\pi m{\displaystyle \frac{1}{V}}e^{\frac{2a}{N1}\varphi }\phi ^{}\left(P_i+e_0\widehat{A}_i\right)\phi .`$
On the other hand, the spatial component of the electromagnetic field equation (153) reads
$$_k\left[e^{\frac{4a}{N1}\varphi }U^{\frac{2}{N2}}\overline{F}_{ki}\right]=8\pi e_0e^{\frac{2a}{N1}\varphi }\phi ^{}(P_k+e_0\widehat{A}_k)\phi ,$$
(188)
or equivalently,
$`_k\left[V^{\frac{2(a^21)}{N2+a^2}}\overline{F}_{ki}\right]=8\pi e_0e^{\frac{2a}{N1}\varphi }\phi ^{}(P_k+e_0\widehat{A}_k)\phi .`$ (189)
Finally using the mass-charge relation (159), we have
$`\sqrt{{\displaystyle \frac{N1}{2(N2+a^2)}}}_k\left[V^{\frac{2(a^21)}{N2+a^2}}\overline{F}_{ki}\right]=8\pi me^{\frac{2a}{N1}\varphi }\phi ^{}(P_k+e_0\widehat{A}_k)\phi .`$ (190)
By taking the same low-energy approximation into the total action (148), we obtain the effective lagrangian density $``$, where
$$S=d^{N+1}x.$$
(191)
Note that:
$`R`$ $`=`$ $`{\displaystyle \frac{2}{N2}}U^{\frac{2}{N2}}_{\mathrm{}}\left({\displaystyle \frac{_{\mathrm{}}U}{U}}\right){\displaystyle \frac{N1}{N2}}U^{\frac{2}{N2}}\left({\displaystyle \frac{_{\mathrm{}}U}{U}}\right)^2+{\displaystyle \frac{1}{4}}U^2U^{\frac{4}{N2}}G^2`$ (192)
$`=`$ $`{\displaystyle \frac{2}{N2}}U^{\frac{2}{N2}}_{\mathrm{}}\left({\displaystyle \frac{_{\mathrm{}}U}{U}}\right){\displaystyle \frac{(N1)(N2)}{(N2+a^2)^2}}U^{\frac{2}{N2}}\left({\displaystyle \frac{_{\mathrm{}}V}{V}}\right)^2`$ (194)
$`+{\displaystyle \frac{1}{4}}V^{\frac{2(a^21)}{N2+a^2}}U^{\frac{2}{N2}}{\displaystyle \frac{1}{V^2}}G^2,`$
$`{\displaystyle \frac{4}{N1}}_\mu \varphi ^\mu \varphi `$ $`=`$ $`{\displaystyle \frac{4}{N1}}U^{\frac{2}{N2}}(_{\mathrm{}}\varphi )^2`$ (195)
$`=`$ $`{\displaystyle \frac{(N1)a^2}{(N2+a^2)^2}}U^{\frac{2}{N2}}\left({\displaystyle \frac{_{\mathrm{}}V}{V}}\right)^2,`$ (196)
$`e^{\frac{4a}{N1}\varphi }F^2`$ $`=`$ $`V^{\frac{2a^2}{N2+a^2}}\left[2U^2U^{\frac{2}{N2}}(F_{0i})^2U^{\frac{4}{N2}}\overline{F}^2\right]`$ (197)
$`=`$ $`{\displaystyle \frac{N1}{N2+a^2}}U^{\frac{2}{N2}}\left({\displaystyle \frac{_{\mathrm{}}V}{V}}\right)^2V^{\frac{2(a^21)}{N2+a^2}}U^{\frac{2}{N2}}\overline{F}^2.`$ (198)
Now we find:
$``$ $`=`$ $`\sqrt{g}\{{\displaystyle \frac{1}{16\pi }}[R{\displaystyle \frac{4}{N1}}(\varphi )^2e^{\frac{4a}{N1}\varphi }F^2]`$ (200)
$`+[\phi ^{}e^{\frac{2a}{N1}\varphi }g^{\mu \nu }(P_\mu +e_0A_\mu )(P_\nu +e_0A_\nu )\phi m^2e^{\frac{2a}{N1}\varphi }\phi ^{}\phi \left]\right\}`$
$``$ $`{\displaystyle \frac{1}{16\pi }}V^{\frac{2(a^21)}{N2+a^2}}\left[{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{V^2}}G^2\overline{F}^2\right]`$ (202)
$`+U^{\frac{N}{N2}}\left[V\phi ^{}\left(P_0+e_0A_0\right)^2\phi {\displaystyle \frac{1}{V}}m^2\phi ^{}\phi {\displaystyle \frac{1}{V}}V^{\frac{2(a^21)}{N2+a^2}}\phi ^{}\left(P_i+e_0\widehat{A}_i\right)^2\phi \right]`$
$`=`$ $`{\displaystyle \frac{1}{16\pi }}V^{\frac{2(a^21)}{N2+a^2}}\left[{\displaystyle \frac{N2+a^2}{Na^2}}\left(\widehat{F}^2{\displaystyle \frac{1}{4}}H^2\right)\right]`$ (204)
$`+U^{\frac{N}{N2}}\left[V\phi ^{}\left(P_0+e_0A_0\right)^2\phi {\displaystyle \frac{1}{V}}m^2\phi ^{}\phi {\displaystyle \frac{1}{V}}V^{\frac{2(a^21)}{N2+a^2}}\phi ^{}\left(P_i+e_0\widehat{A}_i\right)^2\phi \right].`$
Here we have defined an antisymmetric tensor field $`H_{ij}`$ as
$$H_{ij}4\sqrt{\frac{N1}{2(N2+a^2)}}\overline{F}_{ij}+\frac{1}{V}G_{ij}.$$
(205)
$`H_{ij}`$ does not couple to the scalar field $`\phi `$, thus we set $`H_{ij}0`$. Then both Eqs. (187) and (190) can be read as
$$\frac{N2+a^2}{Na^2}_{\mathrm{}}\left[V^{\frac{2(a^21)}{N2+a^2}}\widehat{F}_\mathrm{}i\right]=8\pi e_0e^{\frac{2a}{N1}\varphi }\phi ^{}(P_i+e_0\widehat{A}_i)\phi .$$
(206)
To proceed further, we introduce a non-relativistic field $`\psi `$:
$$\psi \sqrt{2m}U^{\frac{N}{2(N2)}}\phi ,$$
(207)
where since the spatial volume measure $`(g^{(N)})^{1/4}=U^{\frac{N}{2(N2)}}`$, we obtain a correct measure for a usual spatial volume.
Finally we get the effective lagrangian density in the low energy limit:
$``$ $`=`$ $`\psi ^{}\left(P_0m\right)\psi {\displaystyle \frac{1}{2mV^{(Na^2)/(N2+a^2)}}}\psi ^{}\left(𝑷+e_0\widehat{𝑨}\right)^2\psi `$ (209)
$`+{\displaystyle \frac{1}{16\pi }}{\displaystyle \frac{N2+a^2}{Na^2}}{\displaystyle \frac{1}{V^{2(1a^2)/(N2+a^2)}}}\widehat{F}^2(a^2/=N),`$
where $`V`$ satisfies the following equation:
$$^2V+8\pi \frac{N2+a^2}{N1}m|\psi |^2=0.$$
(210)
Varying this effective lagrangian (209) with respect to $`\widehat{𝑨}`$, we can derive again the field equation equivalent to Eq. (206) in the low-energy approximation.
For $`a^2=N`$, since the scalar field does not couple to the vector field, the effective lagrangian density at the lowest order is
$$=\psi ^{}\left(P_0m\right)\psi \frac{1}{2m}\psi ^{}𝑷^2\psi (a^2=N).$$
(211)
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# On the transport of magnetic fields by solar-like stratified convection
## 1 Introduction
Magnetic fields play an important rôle for the formation of the spectra of active late type G, K and M dwarf stars (e.g. Schrijver & Rutten Schrijver+Rutten87 (1987), Rutten et al. Rutten+ea89 (1989) and Johns-Krull & Valenti Johns-Krull+Valenti96 (1996)). The most studied star in that respect is the Sun. It is generally assumed that solar active regions are manifestations of a strong toroidal flux system, that is generated and anchored deep below the surface of the Sun, possibly in the undershoot layer below the convection zone proper. Toroidal magnetic strands ascend through the convection zone because of buoyancy and the on the average super-adiabatic stratification. The much weaker poloidal field is assumed to be generated in the convection zone from this toroidal field by a cyclonic effect.
In the mean field dynamo context, the mechanism that takes care of communicating the poloidal field back to the region where the toroidal field is generated, is assumed to be a diffusive coupling of the regions where the $`\omega `$-effect and the cyclonic $`\alpha `$-effect operate (e.g. Parker Parker1993 (1993)). Choudhuri & Dikpati (Choudhuri+ea95 (1995)) have shown that meridional circulation may also couple the two regions, if the time scale of the circulation is shorter than the diffusive time scale, and if the circulation is such that the flow is equator-ward at the bottom of the convection zone. The modeling by diffusive coupling of the regions may be considered to be somewhat unsatisfactory because it relies on a rather ad hoc approach and the meridional circulation approach is equally unsafe since it is not well observed — at best the amplitudes and directions of the flow are indicated (van Ballegooijen Ballegooijen1998 (1998)).
The question of how the poloidal field may return to the region where the generation of the toroidal field supposedly takes place is related to the well-known “buoyancy dilemma” and the generally assumed solution to this problem: Any magnetic field in the convection zone will escape because of its buoyancy but a magnetic field may be stored in the stably stratified region below the convection zone; i.e., in the undershoot layer. In that scenario the convection zone is considered to be a passive one-dimensional medium while the magnetic fields are treated more or less as solid objects that move in it. It is conceivable, however, that some kind of balance may occur between the drag of descending plasma and the buoyancy of the magnetic field embedded in the plasma.
Along this line of thinking Drobyshevski & Yuferev (Drobyshevski+ea74 (1974)) proposed that a downward “topological pumping” of the magnetic field could be occurring, because of the asymmetric nature of the topology of 3D convective flows, i.e. that they consist of networks of descending material embedding regions of ascending material. They investigated the kinematic case at low magnetic Reynolds number by assuming an incompressible simple geometrical velocity flow pattern. Criticism by Parker (Parker1975a (1975)) made Drobyshevski et al. (Drobyshevski+ea80 (1980)) redo the experiments with a more suitable upper boundary condition. Arter et al. (Arter+ea82 (1982)), Arter (Arter1983 (1983)) and Galloway & Proctor (Galloway+ea83 (1983)) extended the work to higher magnetic Reynolds numbers and several different compressible flows. They found that the magnetic energy did indeed increase at the bottom of the domain.
During the last decade several groups have performed more detailed magneto-convective numerical simulations (e.g. Hurlburt & Toomre Hurlburt+ea88 (1988); Brandenburg et al. Brandenburg+ea90 (1990); Jennings et al. Jennings+ea92 (1992); Nordlund et al. Nordlund+92 (1992); Hurlburt et al. Hurlburt+ea94 (1994); Nordlund et al. Nordlund+94a (1994); Brandenburg et al. Brandenburg+96jfm (1996)). Moreover, numerical simulations of stratified convection have shown that trace particles initially placed in a horizontal layer of a highly stratified model on the average are transported downwards, as a result of the asymmetric topology of stratified convection (Stein & Nordlund Stein+Nordlund89b (1989)). In the context of magnetic fields such a tendency for downwards transport of magnetic fields has also been seen in convective dynamo simulations (Nordlund et al. Nordlund+92 (1992); Brandenburg et al. Brandenburg+96jfm (1996)). Most recently Tobias et al. (Tobias+ea98 (1998)) have investigated this effect using their standard method for studying stellar convection and Mcleod (Mcleod98 (1998)) presented speculations along the same line of thought that shall be followed here.
## 2 Model
The objective of the numerical experiments presented here is to study the interaction of magnetic fields and solar-like stratified over-turning convection and differential rotation. The model of the Sun is a “local box” model of a convectively unstable layer (henceforth referred to as the “convection zone”) sandwiched between two stable layers. In order to circumvent problems associated with the very disparate thermal and dynamical time scales, the model has a much higher luminosity than the Sun, and all variables are scaled accordingly. To compare with solar values, the results must be re-scaled as follows. With a flux scale of f$`{}_{\mathrm{scl}}{}^{}=3\times 10^5`$ times the solar flux, the velocity scale factor becomes u$`{}_{\mathrm{scl}}{}^{}=\mathrm{f}_{\mathrm{scl}}^{1/3}67`$, the temperature fluctuation scale factor t$`{}_{\mathrm{scl}}{}^{}=\mathrm{u}_{\mathrm{scl}}^{}{}_{}{}^{2}4.5\times 10^3`$, and the magnetic field strength scale factor b$`{}_{\mathrm{scl}}{}^{}=`$ u<sub>scl</sub>. For convenience, the variables in the model are given in units of time $`u_{\mathrm{time}}=10^3`$ s, length $`u_{\mathrm{length}}=1`$ Mm, and density $`u_\rho =1`$ gcm<sup>-3</sup>. A Kramers’ opacity scaled inversely with the scale factor for the total luminosity is adopted. This ensures that the boundary between the stable layer and the convection zone is at about the same depth in the model as it is in the Sun.
### 2.1 Equations
The full resistive and compressible MHD-equations are solved using the staggered mesh method by Galsgaard and others (e.g. Galsgaard & Nordlund Galsgaard+ea97 (1997), Nordlund & Stein Stein+Nordlund89b (1989), Nordlund+92 (1992) and Nordlund+94a (1994)):
$`{\displaystyle \frac{\rho }{t}}`$ $`=`$ $`\rho 𝐮,`$ (1)
$`{\displaystyle \frac{\rho 𝐮}{t}}`$ $`=`$ $`(\rho \mathrm{𝐮𝐮}\tau )P+𝐅_g+𝐅_{\mathrm{Lorentz}}+𝐅_{\mathrm{ext}},`$ (2)
$`{\displaystyle \frac{𝐁}{t}}`$ $`=`$ $`\times 𝐄,`$ (3)
$`\mu _0𝐣`$ $`=`$ $`\times 𝐁,`$ (4)
$`𝐄`$ $`=`$ $`\eta 𝐣𝐮\times 𝐁,`$ (5)
$`{\displaystyle \frac{e}{t}}`$ $`=`$ $`(e𝐮)+P(𝐮)+Q_{\mathrm{rad}}+Q_{\mathrm{visc}}+Q_{\mathrm{Joule}},`$ (6)
where $`\rho `$ is the mass density, $`𝐮`$ the velocity, $`\tau `$ the viscous stress tensor, $`P`$ the gas pressure, $`𝐁`$ the magnetic field density, $`𝐣`$ the current density, $`𝐄`$ the electric field, $`\eta `$ the magnetic diffusivity, and $`\mu _0`$ is the magnetic vacuum permeability. In the momentum equation (2) $`𝐅_g=\rho 𝐠`$ is the gravitational force, $`𝐅_{\mathrm{Lorentz}}=𝐣\times 𝐁`$ is the Lorentz force and $`𝐅_{\mathrm{ext}}`$ is the sum of other forces associated with the rotation to be discussed below (Sect. 2.4). The gravitational acceleration $`𝐠`$ is along the x-direction (equivalent to the radial direction). Furthermore $`e`$ is the internal thermal energy, $`Q_{\mathrm{visc}}`$ is the viscous heating, $`Q_{\mathrm{Joule}}`$ is the Joule heating, and $`Q_{\mathrm{rad}}`$ is the radiative heating (cooling).
The code uses a finite difference staggered mesh with 6th order derivative operators, 5th order centering operators and a 3rd order time-stepping routine (Hyman Hyman1979 (1979)).
### 2.2 Boundary conditions
In most stars envelope convection is essentially driven by surface cooling. The entropy contrast at the surface is far larger than that at the bottom of the convection zone, if the convection zone extends over several or many pressure scale heights. To model this situation, without having to actually include all layers up to the solar surface, a simple expression for an isothermal cooling layer at the upper boundary of the model was used:
$$Q_{\mathrm{rad}}=\frac{(TT_{\mathrm{top}})}{\tau _{\mathrm{cool}}}f(x),$$
(7)
where $`\tau _{\mathrm{cool}}`$ is the characteristic cooling time, $`T_{\mathrm{top}}`$ is the temperature of the cooling layer, and $`f(x)`$ is a profile function that restricts the effect to a thin surface layer.
Both experiments with closed and open upper boundaries were performed. In order to implement a stable open upper boundary a buffer zone was allocated, where a fiducial electric field is gradually turned on. The sense of the electric field is such that it drags the magnetic field out of the buffer zone, and the magnitude is increased from zero to of the order of $`u_{\mathrm{max}}B`$, where $`u_{\mathrm{max}}`$ is the maximum velocity in the buffer zone. A layer that is about 10 grid zones from the numerical upper boundary may thus be considered as the physical open upper boundary. This layer is far below the real boundary of the solar convection zone (even the numerical upper boundary of the model is far below the photosphere).
### 2.3 Initial conditions
The hydrodynamic part of the initial condition is a snapshot from a well developed stage of a numerical model of the solar convection zone. The initial condition for the magnetic field is given by a unidirectional (“poloidal”) sheet that is placed in the middle of the convection zone. The sheet is initially in a state of isentropic pressure equilibrium with the surroundings.
The model has a high degree of stratification with a density contrast of roughly $`\mathrm{5\hspace{0.17em}10}^3`$ in the convection zone alone. The thickness of the undershoot layer in the model is approximately equal to 0.8 $`\mathrm{H}_\mathrm{P}`$, where $`\mathrm{H}_\mathrm{P}`$ is the pressure scale height at the bottom of the convection zone. This is much larger than the helioseismological upper limits of 0.1 – 0.2 $`\mathrm{H}_\mathrm{P}`$ (see e.g. Christensen-Dalsgaard JCD+ea95 (1995)). Scaling the extension of the undershoot layer naively in proportion to the velocity (Hooke’s law — the retarding force increases approximately linearly with distance), the resulting undershoot thickness, of the order of 0.01 $`\mathrm{H}_\mathrm{P}`$, falls well within the upper limits from observations.
### 2.4 Latitudinal shear
Both experiments with and without differential rotation in the convection zone were performed, with the purpose of illustrating the effects of radial and latitudinal shear and of the Coriolis force on the magnetic field. In the rotational cases background differential rotation and the Coriolis force are included by using the force $`𝐅_{\mathrm{ext}}`$ in the equation of motion (Eq. 2)
$$𝐅_{\mathrm{ext}}=𝐅_{\mathrm{rot}}+𝐅_{\mathrm{Coriolis}}$$
(8)
where $`𝐅_{\mathrm{Coriolis}}=2\rho 𝛀\times 𝐮`$ is the Coriolis force and $`𝐅_{\mathrm{rot}}`$ is a force designed to induce background rotation (see below).
The differential rotation that is implemented in the model may not be identified as rotation around e.g. one of the horizontal axes. Rather the mapping of the rotation is such that the rotation axis is given by a vector in the “meridional” $`(x,y)`$ plane pointing towards the north pole, i.e. a rotation vector $`𝛀=(\mathrm{\Omega }_x,\mathrm{\Omega }_y,0)`$ given by
$`\mathrm{\Omega }_x`$ $`=`$ $`\mathrm{\Omega }(x,\theta )\mathrm{cos}\theta ,`$ (9)
$`\mathrm{\Omega }_y`$ $`=`$ $`\mathrm{\Omega }(x,\theta )\mathrm{sin}\theta ,`$ (10)
where $`\mathrm{\Omega }=\mathrm{\Omega }(x,\theta )`$ is a fit to the observed solar angular frequency (from Dziembowski et al. Dziembowski+ea89 (1989)) and $`\theta [0:2\pi ],\theta (y)=2\pi y/L_y`$ where $`L_y`$ is half the size of the box in the latitudinal direction. This choice means that $`𝛀`$ is parallel to the x-axis (the radial direction) at the poles and parallel to the y-axis (the latitudinal – or poloidal direction) at the equator.
An azimuthal background velocity corresponding to the observed rotation is implemented by adding a Newtonian term in the equation of motion:
$$𝐅_{\mathrm{rot}}=\frac{\rho }{\tau _R}(u_zu_z^0)\widehat{𝐳},$$
(11)
where $`\tau _\mathrm{R}`$ is the time scale of forcing of the background rotation (one turn-over time has proved to be a good choice), $`\widehat{𝐳}`$ is the the unit vector in the azimuthal direction, and $`u_z^0=u_z^0(x,y)`$ is the background rotation speed given by
$$u_z^0=(\mathrm{\Omega }(x,\theta )\mathrm{\Omega }_0)R\mathrm{sin}\theta .$$
(12)
where $`R`$ is a characteristic radius in the convection zone and $`\mathrm{\Omega }_0`$ is the angular rotation frequency evaluated at $`\theta =\frac{\pi }{4}`$.
The differential rotation is then measured relative to the latitude $`\pm 45^\mathrm{o}`$. In the model that latitude corresponds to the position at fractions of 0.25, 0.75, 1.25 and 1.75 of the latitudinal extension of the computational box counting the total length as $`2L_y`$. (see Fig. 2). At that latitude there is no radial shear since we set the internal solid body rotation equal to the in situ azimuthal speed, so that
$$\mathrm{\Omega }=\{\begin{array}{cc}\mathrm{\Omega }(x,\theta ),\hfill & xR\hfill \\ \mathrm{\Omega }_0,\hfill & x<R\hfill \end{array}$$
At other latitudes the peak of the radial shear is located in the undershoot layer.
The strength of the background rotation may also be changed by varying the radius $`R`$ in Eq. 12. Values from 500 to 1400 Mm were used, the higher values in combination with larger values of $`\tau _R`$, in order to still produce the desired amplitude of the differential rotation.
## 3 Results
Several experiments with varying initial magnetic field strengths, numerical resolutions, and upper boundary conditions were performed. First, results from simulations without rotational effects are discussed (see Table 1), and then results from a simulation including latitudinal shear are reported.
Although some of the experiments initially had a maximum field strength of the order of the formal equipartition value, they have been denoted ‘sub-equipartition’ in Table 1 because the distribution of the kinetic energy density $`e_K`$ was much broader than the distribution of the magnetic energy density $`e_M`$, i.e. the peak magnetic energy density was smaller than the peak kinetic energy density even though the most likely value of $`e_M`$ was similar to that of $`e_K`$.
A moderate number of grid points were used in these experiments (a few experiments were also ran at higher resolutions, up to $`145\times 128^2`$, to check for resolution effects). The advantage of performing relatively small numerical experiments, is that it is possible to perform a larger number of experiments, with different setups and with a variety of parameter values. Since it is impossible, with the limits of currently available computer power to accurately reproduce solar conditions, it is necessary to experiment with trade-offs between various constraints.
### 3.1 The rotation-less case
Fig. 1 shows two sets of magnetic field isosurfaces at two different instants in time for an experiment with a sub-equipartition magnetic field (experiment \[i\], in Table 1): The poloidal sheet of magnetic field initially placed in the middle of the convection zone quickly starts to interact with the convection (Fig. 1, left panel). At a subsequent time (Fig. 1, right panel) the magnetic field more or less fills the whole volume of the convection zone and penetrates into the stable layer below.
Fig. 3 shows the horizontally averaged field at 7 equidistant instants of time for experiment \[i\]: The poloidal sheet is spread out, and the distribution of poloidal magnetic flux settles to a characteristic distribution. The highest (horizontally averaged) poloidal flux density occurs in the overshoot layer, and in these particular experiments a significant fraction of the total flux also resides there. In the real Sun, this fraction may be expected to be much smaller, because the undershoot layer is much thinner there. Given the shape of the distribution, with no particular enhancement in the undershoot layer, it is likely the Sun has a correspondingly smooth distribution, with the majority of the poloidal flux residing inside the convection zone.
The ‘pumping effect’ described in the above takes place because of the topology of the over-turning stratified convection: Of all the fluid parcels threaded by magnetic field lines in the initial state about half are initially ascending. However, because of the stratification, most of the ascending fluid parcels have to over-turn and descend and most of these keep descending down to the bottom of the convection zone (Stein & Nordlund Stein+Nordlund89b (1989)). The fluid parcels drag the threading field lines along and hence an appreciable fraction of the field is transported downwards. Fragments of the field that are caught in ascending flows are advected upwards. In the experiments with an open upper boundary some of these fragments escape through the top of the computational domain and flux is systematically lost (see the discussion below).
Even in the cases of super-equipartition fields (e.g. experiments \[iv\] and \[v\]) the magnetic field is pushed downwards in the initial phases until it reaches the undershoot layer where the kinetic energy flux decreases. Because the overall magnetic flux decreases through flux loss at the surface the magnetic field eventually enters a state in which it is below equipartition.
In the case of an initial field below equipartition the pumping effect is — not surprisingly — even more pronounced than for the super-equipartition fields: Fig. 4 shows the results for three cases where the magnetic field was below equipartition (experiments \[i\], \[ii\] and \[iii\]) and the dynamics thus were dominated by the convective motions. The average magnetic field is distributed over the entire convection zone, with maximum (horizontally averaged) flux density in the undershoot layer.
The radial distribution of the magnetic flux present in the convection zone is not very different for the cases with a closed upper boundary (left panel in Fig. 4) and the cases with an open upper boundary (right panel in Fig. 4), but the total amount of magnetic flux is rapidly reduced during the first few turn-over times, when there is a significant loss of poloidal flux through the upper boundary. This does not influence the distribution of flux over depth much, however, it only influences the amount of flux that is available for distribution.
The flux loss in the models with an open upper boundary is strongly exaggerated in comparison to the Sun. The real flux loss may be expected to be significantly smaller than in the models with open upper boundaries discussed here: Much of the weak ascending flux must over-turn rather than reach the solar surface since it is embedded in a fluid of which only a tiny fraction ($``$ 0.1% estimated from mass flux amplitudes) reaches the solar surface. Note that the magnetic field considered is weak and incoherent and does not have sufficient buoyancy to overcome the drag of the fluid motions.
A supplementary numerical experiment confirms the exaggeration of the flux loss: Moving the open magnetic boundary four mesh point upwards, corresponding to about 0.6 density scale heights, reduces the flux loss by about 20%. The new boundary (at $`\rho =2.5\times 10^2`$ gcm<sup>-3</sup>) is still over ten density scale heights (sic) away from the real solar boundary (at $`\rho 3\times 10^7`$ gcm<sup>-3</sup>). Almost all of the flux that is still lost in the modified experiment may thus in the case of the real Sun be expected to turn over before reaching the solar surface.
Fig. 6 shows the average poloidal field strength for the experiments with an open upper boundary (experiments \[iii\], \[iv\] and \[v\]). The relative reduction in the poloidal flux by the escape of magnetic structures through the upper boundary does not depend much on the field strength, for strengths up to several times formal equipartition values.
The rate of flux loss does depend on the field strength but in a rather counterintuitive way: The stronger the initial field the smaller the initial reduction rate. Naively one might expect that if the field strength is large the higher buoyancy would make the field escape faster. The explanation is that the over-turning of the fluid is slowed down in the cases with a stronger initial poloidal flux sheet. The fragmentation of the sheet is also slower in the case with a stronger field.
The magnetic field in the convection zone rapidly becomes very fragmented: While the field in the beginning is uniformly distributed in the sheet, quickly a picture develops where the field in the convection zone becomes very intermittent, while the field that is pumped into the undershoot layer is much more uniform.
The reason for this difference in topology between the magnetic field in the convection zone and in the undershoot layer is that once the magnetic field has been pumped down into the undershoot layer it is less susceptible to fragmentation, since the motions in the stable layer have a much smaller amplitude and are not as systematic as the motions in the convection zone.
The degree of intermittency and fragmentation of the field in the convection zone depends on the strength of the initial poloidal magnetic sheet. The poloidal magnetic sheet is more stable towards the initial fragmentation if its field is stronger since in that case a larger force is needed to over-power the tension in the poloidal field lines.
### 3.2 A case with latitudinal shear
In what follows, results are presented for a particular simulation that includes shear and the Coriolis force (through the force $`𝐅_{\mathrm{ext}}`$ in the equation of motion Eq. 2). The physical size of this experiment was $`381\times 800\times 800`$ Mm and the numerical resolution $`69\times 135^2`$ grid points.
The initial condition of a uniform poloidal sheet placed at a certain depth may be said to be rather arbitrary. However, an arbitrary state is both appropriate and useful in this case, since the system rather quickly enters a state that is independent of the initial condition i.e. there is no internal long-time “memory” in the system.
Initially the pumping effect ensures that an on the average poloidal field is injected into the bottom undershoot layer, and after about four turn-over times the poloidal flux density has a maximum below the convection zone of the model.
After an initial transient process where the differential rotation establishes itself and the over-turning convection distorts the poloidal sheet, the system enters a well-developed state, where the magnetic field displays the structure of the differential rotation. At this point, the magnetic energy and the toroidal field strength have increased to a significantly higher level than their initial values (see the two top panels in Fig. 5), and the rms toroidal field is larger than the rms poloidal field. The average poloidal flux has dropped to about half the initial value due to flux loss at the surface (see the bottom two panels in Fig. 5) but this loss of magnetic flux through the upper boundary is somewhat halted after the formation of a large-scale toroidal field.
At first the growth of the toroidal field is fast while the background rotation increases toward the profile determined by Eq. 12 on the time scale $`\tau _R`$ (cf. Eq. 11). When the background rotation profile is fully attained by the fluid, the toroidal magnetic field begins to increase linearly, with the rate of increase given by the latitudinal shear and the original poloidal field strength. The structure of the magnetic field directly reflects the latitudinal dependence of the background azimuthal velocity field as a result of the latitudinal shear.
In the well-developed state after the main flux loss has taken place, the dominant magnetic field is a strong toroidal flux system located near the bottom of the convection zone, i.e. it is not pumped into the undershoot layer but it does not escape from the convection zone either: The “center of gravity” of the magnetic field is above the bottom of the convection zone (see Fig. 7).
Fig. 8 shows a view of the toroidal flux system: Four toroidal flux streets are formed on each side of the equators, in the regions where there is a maximum shear.
The magnetic field just below the convection zone rotates more slowly than the field just above. This is a mechanism that may twist the magnetic field lines that connect across the equator regions.
As the initial poloidal sheet fragments while being wound up, toroidal flux structures leave the convection zone through the open upper boundary. In a rotating frame of reference vertical motions lead to horizontal motions through the action of the Coriolis force. One special case is the rise of flux loops that are rotated so that they emerge tilted with respect to the latitudinal circle, and another is the excitation of meridional circulation, as a result of the transport of angular momentum: In the simulation both of these effects are found.
The surface of the Sun is an extremely important place: Not only may flux be lost there, but the surface constitutes a “reconnection region” where vertical field lines are effectively “cut over”, and where the remaining “stubs” may be advected passively by the horizontal surface motions. Field lines of course actually continue into the corona and either connect back into the photosphere or connect out into interplanetary space. However, reconnection is observed to proceed so rapidly in the corona and the mass density is so low there, that from the point of view of the sub-surface dynamics the connections above the surface are of little importance.
Fig. 9 shows two panels with a small section of the box, with snapshots of the magnetic field, looking down along the two toroidal flux streets. A magnetic flux structure is sticking out through the surface, and is moving towards the equator between the two toroidal flux streets, but at the same time it is connected through a subsurface structure to one of them. As a result of the drift of the surface-structure, the magnetic field lines that make up the subsurface structure have a tendency to become inclined with respect to the toroidal street to which it connect, and hence a poloidal field component is formed, with opposite polarity relative to the original one.
## 4 Discussion and conclusions
The main conclusions that may be drawn from the study of the interaction of a magnetic sheet with stratified over-turning convection in the absence of rotation are the following:
* Stratified convection induces a strong tendency to transport magnetic flux downwards.
* The distribution of horizontally averaged magnetic flux peaks in the undershoot layer, but the bulk of the flux is in the convection zone.
* The magnetic field that resides in the undershoot layer is considerably less fragmented than the magnetic field in the bulk of the convection zone.
* Unless an open upper boundary is placed sufficiently close to the actual solar surface, i.e. at sufficiently low density, there may be a substantial (and exaggerated) loss of flux through the open boundary as flux is carried around by the over-turning convection.
* The transport properties (both the downwards transport and the surface flux loss) are quite robust and field strengths well in excess of formal equipartition are needed to change the distribution and rates significantly.
Additionally, the results of the simulations including latitudinal shear offer the following conclusions:
* The system rapidly forgets the initial condition and distributes the flux in a generic vertical distribution.
* The latitudinal shear of the differential rotation shapes the magnetic field and creates strong toroidal flux streets located at mid latitudes.
* While magnetic flux indeed is pumped into the undershoot layer, the center of gravity of the magnetic field is above the bottom of the convection zone.
* Magnetic field lines and flux structures penetrating the upper boundary move passively according to the surface motion.
Where do these results leave the “storage problem”, i.e. the problem of explaining how the magnetic field can remain stored while being amplified by differential rotation?
The toroidal magnetic field in these simulations reach peak field strengths of a several tens of kG (scaled to the Sun). The peak field strengths occur near the bottom, but still inside, the convection zone. It is conceivable that emerging flux regions form when buoyancy finally becomes dominant, and that this occurs at field strengths of the order of 100 kG, as has been deduced from emergence and tilt patterns by several investigations. Because of numerical limitations we were not study that process with the current series of experiments—future experiments with higher resolution and even larger density contrasts are needed here.
Besides offering a clue to the operation of the solar dynamo the results presented here may also contribute to the understanding of magnetic field generation in other late type stars. For example; the fact that no undershoot storage is available in the case of magnetically active fully convective M dwarf stars (e.g. Chabrier & Baraffe Chabrier+Baraffe97 (1999)) is generally considered to be a problem (see e.g. Allard et al. Allard+ea97 (1997) and Küker & Rüdiger Kueker+Ruediger99 (1999)) — this ceases to be a problem in the scenario presented here.
We note that, even though the loss of magnetic flux at the upper boundary is exaggerated in the models presented here, such a loss is certainly a real effect that is important to include, since the Sun is known to loose a considerable amount of toroidal magnetic flux during an activity cycle (e.g. van Ballegooijen Ballegooijen1998 (1998)).
Lastly we find it important to emphasize that the visualization of “emerging” flux structures and field lines (Fig. 9) illustrates a mechanism first pointed out by van Ballegooijen (vanBalleprivate95 (1995)) that may be crucial for the reversal of the poloidal field: Flux structures that rise and penetrate the surface effectively results in the field lines being “cut” at the surface, with the leading polarity tending to drift towards the equator, and the following polarity tending to drift towards the pole. The result is that subsurface connections between the following polarity of one emerging structure and the leading polarity of another trailing structure may become tilted in the sense opposite to the tilt associated with the normal winding of the field. When such reversed tilts are caught by the differential rotation, they will effectively lead to “unwinding” and reversal of the poloidal field component. This again is an important topic for future studies.
###### Acknowledgements.
This work was supported in part by the Danish Research Foundation, through its establishment of the Theoretical Astrophysics Center. Computing time at the UNI$``$C computing center was provided by the Danish Natural Science Research Council. SBFD acknowledges support through an EC-TMR grant to the European Solar Magnetometry Network.
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# Semiclassical Dynamics of the Jaynes-Cummings Model
## I Introduction
Since the sixties of the last century the Jaynes-Cummings model is frequently considered as a simple model to describe a two-level atomic system interacting with an electromagnetic field in a cavity; for recent reviews see . Apart from its relevance to quantum optics, in particular laser theory, this integrable quantum model also allows to test approximative methods by comparing them with the exact result. In particular, the “semiclassical” theory has attracted considerable attention where the bosonic field mode is represented by classical c-numbers while the two-level atomic system is represented as a quantum spin-$`\frac{1}{2}`$ -. In this approximation the Heisenberg equations of motion are replaced by linear operator equations for the spin variables and an amplitude equation for the electromagnetic field which is driven by the expectation values of the spin operators. Taking the expectation value of the Heisenberg equations for the spin variables, the optical Bloch equations emerge which describe the dynamics of a classical Bloch vector on the two-sphere . It is well known that this “semiclassical” theory provides results that are equivalent to a full quantum mechanical treatment if the mean number of bosons is very large and fluctuations in the boson number can be neglected .
While this conventional semiclassical approach treats the cavity field just classically, we attempt at a semiclassical theory treating both the atomic and electromagnetic subsystems on an equal footing. Starting from the full quantum model we focus on the most probable paths of the system within the path integral representation and relevant fluctuations about them. Within the scope of (spin) coherent state path integrals we obtain a semiclassical approximation going beyond the classical field approximation. For instance, the approach yields an accurate description of spontaneous emission.
The paper is organized as follows. In Sec.II we first solve the Jaynes-Cummings model exactly with spin coherent state path integrals in a subspace with fixed excitation quanta above the ground state. Then, in Sec.III, we examine a semiclassical description which does not rely on these subspaces and can thus be extended to more complicated Hamiltonians. With coherent state path integrals the leading order of the propagator is determined by solving the Euler-Lagrange equations for the classical path. In Sec.IV we consider contributions from fluctuations about the dominant path and show that they lead to a decay of the excited two-level system by spontaneous emission.
## II The Jaynes-Cummings model
The Jaynes-Cummings model is characterized by the Hamiltonian
$$H=a^{}a+(1+\mathrm{\Delta })S_z+\lambda (aS_++a^{}S_{}),$$
(1)
where $`a`$ is the canonical annihilation operator of a bosonic field mode with frequency $`\omega `$ and $`S_\pm =S_x\pm iS_y`$, $`S_z`$ are operators of a spin-$`\frac{1}{2}`$ describing two levels of an atomic system with energy difference $`\mathrm{}\omega _o`$. There are two dimensionless parameters, the detuning $`\mathrm{\Delta }=(\omega _o\omega )/\omega `$ and the coupling strength $`\lambda =g/\omega `$. We use units with $`\omega =1`$ and $`\mathrm{}=1`$. It is well known that the Jaynes-Cummings model allows apart from $`H`$ for another time independent operator
$$N=a^{}a+S_z,$$
(2)
which measures the number of excitation quanta in the system. Hence, the time evolution operator is of the form
$$U(T)=e^{iHT}=e^{iNT}e^{iCT},$$
(3)
where $`C=HN`$. Representing the spin operators in the eigenbasis of $`S_z`$ formed by the eigenvectors $`|`$ and $`|`$, the first factor in Eq.(3) may be written as
$$e^{iNT}=e^{ia^{}aT}(e^{\frac{i}{2}T}||+e^{+\frac{i}{2}T}||).$$
(4)
Introducing further the eigenkets of $`a^{}a`$, invariant subspaces are distinguished. In particular the kets $`|n1||n1`$ and $`|n||n`$ span the subspace with $`N=(n\frac{1}{2})`$. In this subspace the time independent operator $`C`$ generates $`SU(2)`$ dynamics. This can be seen explicitly by introducing the operators
$`J_x`$ $`=`$ $`{\displaystyle \frac{1}{2}}(|n1n|+|nn1|)`$ (5)
$`J_y`$ $`=`$ $`{\displaystyle \frac{i}{2}}(|n1n|+|nn1|)`$ (6)
$`J_z`$ $`=`$ $`{\displaystyle \frac{1}{2}}(|n1n1||nn|),`$ (7)
describing the angular momentum of a spin-$`\frac{1}{2}`$. In terms of these spin operators we have
$$C=2\lambda \sqrt{n}J_x+\mathrm{\Delta }J_z,$$
(8)
and we see that in this subspace $`C`$ gives indeed rise to pure $`SU(2)`$ dynamics. Accordingly, the propagator may be worked out exactly by a semiclassical approach with path integrals in the spin coherent state representation
$$|\vartheta \phi =e^{i\phi J_z}e^{i\vartheta J_y}|n1.$$
(9)
Following the lines of , we write the spin coherent propagator as a regularized path integral
$$\vartheta ^{\prime \prime }\phi ^{\prime \prime }|e^{iCT}|\vartheta ^{}\phi ^{}=\underset{\nu \mathrm{}}{lim}𝑑\mu \mathrm{exp}\left\{iS[\vartheta (t),\phi (t)]\right\},$$
(10)
with the action
$$S[\vartheta (t),\phi (t)]=_0^T𝑑t\left[\frac{1}{2}\mathrm{cos}(\vartheta )\dot{\phi }C(\vartheta ,\phi )\right].$$
(11)
Here the operator $`C`$ is represented as
$`C(\vartheta ,\phi )`$ $`=`$ $`\vartheta \phi |C|\vartheta \phi .`$ (12)
$`=`$ $`\lambda \sqrt{n}\mathrm{sin}(\vartheta )\mathrm{cos}(\phi )+{\displaystyle \frac{\mathrm{\Delta }}{2}}\mathrm{cos}(\vartheta ).`$ (13)
The spherical Wiener measure
$$d\mu =M\underset{t=0}{\overset{T}{}}d\mathrm{cos}(\vartheta (t))d\phi (t)\mathrm{exp}\left\{\frac{1}{4\nu }_0^T𝑑t\left[\dot{\vartheta }^2+\mathrm{sin}^2(\vartheta )\dot{\phi }^2\right]\right\},$$
(14)
enforces continuous Brownian motion paths on the sphere ($`M`$ is a normalization factor). This measure gives rise to a regularization dependent action
$`S_\nu [\vartheta (t),\phi (t)]`$ $`=`$ $`{\displaystyle _0^T}𝑑t\left\{{\displaystyle \frac{i}{4\nu }}\left[\dot{\vartheta }^2+\mathrm{sin}^2(\vartheta )\dot{\phi }^2\right]+{\displaystyle \frac{1}{2}}\mathrm{cos}(\vartheta )\dot{\phi }C(\vartheta ,\phi )\right\}.`$ (15)
Now, in the semiclassical expansion, we separate the paths $`\mathrm{cos}(\vartheta )=\mathrm{cos}(\vartheta _{cl})+x/\sqrt{s}`$ and $`\phi =\phi _{cl}+y/\sqrt{s}`$ in their classical parts and fluctuations around them. The formal limit of large spin $`s\mathrm{}`$ expresses the classical limit. To lowest order, in the Dominant Path Approximation (DOPA), the semiclassical expansion gives
$$e^{iS_{cl}}=\mathrm{exp}\{i_0^TdtC(\overline{\vartheta }^{\prime \prime }(t),\overline{\phi }^{\prime \prime }(t))\}\vartheta ^{\prime \prime }\phi ^{\prime \prime }|\vartheta ^{}\phi ^{},$$
(16)
where $`(\overline{\vartheta }^{},\overline{\phi }^{})`$ and $`(\overline{\vartheta }^{\prime \prime },\overline{\phi }^{\prime \prime })`$, respectively, describe the starting point and endpoint of the classical trajectory $`(\overline{\vartheta }(t),\overline{\phi }(t))`$, $`0tT`$. For convenience let us introduce the complex variables
$`\zeta `$ $`=`$ $`\mathrm{tan}\left({\displaystyle \frac{\overline{\vartheta }}{2}}\right)e^{i\overline{\phi }}`$ (17)
$`\eta `$ $`=`$ $`\mathrm{tan}\left({\displaystyle \frac{\overline{\vartheta }}{2}}\right)e^{i\overline{\phi }}.`$ (18)
Then, the dominant path is determined by
$`\dot{\zeta }`$ $`=`$ $`i\lambda \sqrt{n}(1\zeta ^2)+i\mathrm{\Delta }\zeta `$ (19)
$`\dot{\eta }`$ $`=`$ $`i\lambda \sqrt{n}(1\eta ^2)i\mathrm{\Delta }\eta ,`$ (20)
with boundary conditions $`\zeta (0)=\zeta ^{}`$ and $`\eta (T)=\eta ^{\prime \prime }`$. Hence, the endpoint of the classical trajectory obeys
$`\zeta (T)`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Omega }_n\zeta ^{}\mathrm{cos}(\mathrm{\Omega }_nT)+i\left[\mathrm{\Delta }\zeta ^{}\lambda \sqrt{n}\right]\mathrm{sin}(\mathrm{\Omega }_nT)}{2\mathrm{\Omega }_n\zeta ^{}\mathrm{cos}(\mathrm{\Omega }_nT)i\left[\lambda \sqrt{n}\zeta ^{}+\mathrm{\Delta }\right]\mathrm{sin}(\mathrm{\Omega }_nT)}}`$ (21)
$`\eta (T)`$ $`=`$ $`\eta ^{\prime \prime },`$ (22)
with the Rabi frequency
$$\mathrm{\Omega }_n=\sqrt{\lambda ^2n+\frac{\mathrm{\Delta }^2}{4}}.$$
(23)
In terms of the complex variables (18) we get
$`C(\zeta (T),\eta ^{\prime \prime })`$ $`=`$ $`\lambda \sqrt{n}{\displaystyle \frac{\zeta (T)+\eta }{1+\zeta (T)\eta ^{\prime \prime }}}+{\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{1\zeta (T)\eta ^{\prime \prime }}{1+\zeta (T)\eta ^{\prime \prime }}}`$ (24)
$`=`$ $`i{\displaystyle \frac{d}{dT}}\mathrm{log}\{(1+\zeta ^{}\eta ^{\prime \prime })\mathrm{cos}(\mathrm{\Omega }_nT)`$ (26)
$`{\displaystyle \frac{i}{\mathrm{\Omega }_n}}[\lambda \sqrt{n}(\zeta ^{}+\eta ^{\prime \prime })+{\displaystyle \frac{\mathrm{\Delta }}{2}}(1\zeta ^{}\eta ^{\prime \prime })]\mathrm{sin}(\mathrm{\Omega }_nT)\}.`$
Now the integral in Eq.(16) is readily solved and the propagator in the DOPA takes the form
$`e^{iS_{cl}}`$ $`=`$ $`a(T)\mathrm{cos}\left({\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\right)\mathrm{cos}\left({\displaystyle \frac{\vartheta ^{}}{2}}\right)e^{\frac{i}{2}(\phi ^{\prime \prime }\phi ^{})}+b(T)\mathrm{cos}\left({\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{\vartheta ^{}}{2}}\right)e^{\frac{i}{2}(\phi ^{\prime \prime }+\phi ^{})}`$ (28)
$`b^{}(T)\mathrm{sin}\left({\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\right)\mathrm{cos}\left({\displaystyle \frac{\vartheta ^{}}{2}}\right)e^{\frac{i}{2}(\phi ^{\prime \prime }+\phi ^{})}+a^{}(T)\mathrm{sin}\left({\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{\vartheta ^{}}{2}}\right)e^{\frac{i}{2}(\phi ^{\prime \prime }\phi ^{})},`$
where
$`a(T)`$ $`=`$ $`\mathrm{cos}(\mathrm{\Omega }_nT)i{\displaystyle \frac{\mathrm{\Delta }}{2\mathrm{\Omega }_n}}\mathrm{sin}(\mathrm{\Omega }_nT)`$ (29)
$`b(T)`$ $`=`$ $`i{\displaystyle \frac{\lambda \sqrt{n}}{\mathrm{\Omega }_n}}\mathrm{sin}(\mathrm{\Omega }_nT).`$ (30)
As discussed elsewhere for pure $`SU(2)`$ dynamics the DOPA is exact and Eqs.(28),(30) give indeed the exact propagator .
In more general situations, such as for the case without rotating wave approximation , the system cannot be separated into invariant subspaces. Therefore it would be interesting to consider a semiclassical expansion that does not rely on the $`SU(2)`$ generators (7).
## III Semiclassical dynamics with coherent state path integrals
In order to formulate a general semiclassical theory for a coupled spin boson problem we make use of product coherent states
$$|\vartheta \phi pq=e^{i\phi S_z}e^{i\vartheta S_y}e^{i(pQqP)}||0.$$
(31)
These states are generated by momentum and space translations of the normalized vacuum state $`|0`$ and $`SU(2)`$ rotations of the $`S_z`$ eigenstate $`|`$. Again, the semiclassical approximation is based on the coherent state path integral representation. We express the propagator as
$$\vartheta ^{\prime \prime }\phi ^{\prime \prime }p^{\prime \prime }q^{\prime \prime }|U(t)|\vartheta ^{}\phi ^{}p^{}q^{}=\underset{\nu _a,\nu _b\mathrm{}}{lim}𝑑\mu _a𝑑\mu _b\mathrm{exp}\left\{iS[p(t),q(t),\vartheta (t),\phi (t)]\right\},$$
(32)
with the action
$$S[p(t),q(t),\vartheta (t),\phi (t)]=_0^T𝑑t\left[\frac{1}{2}(p\dot{q}\dot{p}q)+\frac{1}{2}\mathrm{cos}(\vartheta )\dot{\phi }H(\vartheta ,\phi ,p,q)\right].$$
(33)
For the Jaynes-Cummings model the Hamiltonian takes the form
$`H(\vartheta ,\phi ,p,q)`$ $`=`$ $`\vartheta \phi pq|H|\vartheta \phi pq`$ (34)
$`=`$ $`{\displaystyle \frac{1}{2}}(p^2+q^2)+{\displaystyle \frac{1+\mathrm{\Delta }}{2}}\mathrm{cos}(\vartheta )+{\displaystyle \frac{\lambda }{2\sqrt{2}}}\left[\mathrm{sin}(\vartheta )e^{i\phi }(q+ip)+\mathrm{sin}(\vartheta )e^{i\phi }(qip)\right].`$ (35)
Here, the canonical coherent state path integral is regularized by the flat Wiener measure
$$d\mu _a=M_a\underset{t=0}{\overset{T}{}}dp(t)dq(t)\mathrm{exp}\left\{\frac{1}{2\nu _a}_0^T𝑑t\left[\dot{q}^2+\dot{p}^2\right]\right\},$$
(36)
while the spin paths are again regularized by the spherical Wiener measure
$$d\mu _b=M_b\underset{t=0}{\overset{T}{}}d\mathrm{cos}(\vartheta (t))d\phi (t)\mathrm{exp}\left\{\frac{1}{4\nu _b}_0^T𝑑t\left[\dot{\vartheta }^2+\mathrm{sin}^2(\vartheta )\dot{\phi }^2\right]\right\}.$$
(37)
These measures give rise to the regularization dependent action
$`S_{\nu _a,\nu _b}[\vartheta (t),\phi (t),p(t),q(t)]`$ $`=`$ $`{\displaystyle _0^T}dt\{{\displaystyle \frac{i}{2\nu _a}}[\dot{q}^2+\dot{p}^2]+{\displaystyle \frac{i}{4\nu _b}}[\dot{\vartheta }^2+\mathrm{sin}^2(\vartheta )\dot{\phi }^2]`$ (39)
$`+{\displaystyle \frac{1}{2}}(p\dot{q}\dot{p}q)+{\displaystyle \frac{1}{2}}\mathrm{cos}(\vartheta )\dot{\phi }H(\vartheta ,\phi ,p,q)\}.`$
In the semiclassical expansion we split the paths $`p=p_{cl}+x_a`$, $`q=q_{cl}+y_a`$, $`\mathrm{cos}(\vartheta )=\mathrm{cos}(\vartheta _{cl})+x_b/\sqrt{s}`$ and $`\phi =\phi _{cl}+y_b/\sqrt{s}`$ in their classical parts and fluctuations around them. Restricting ourselves to the DOPA, we obtain the propagator
$`e^{iS_{cl}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{sin}(\vartheta ^{})\mathrm{sin}(\vartheta ^{\prime \prime })}{\mathrm{sin}(\overline{\vartheta }^{})\mathrm{sin}(\overline{\vartheta }^{\prime \prime })}}}\mathrm{exp}\left\{{\displaystyle \frac{1}{2}}\left[q^{\prime \prime }\overline{p}^{\prime \prime }\overline{q}^{\prime \prime }p^{\prime \prime }+\overline{q}^{}p^{}q^{}\overline{p}^{}\right]\right\}`$ (41)
$`\times \mathrm{exp}\left\{i{\displaystyle _0^T}𝑑t\left[{\displaystyle \frac{1}{2}}\mathrm{cos}(\overline{\vartheta })\dot{\overline{\phi }}+{\displaystyle \frac{1}{2}}(\overline{p}\dot{\overline{q}}\dot{\overline{p}}\overline{q})H(\overline{\vartheta },\overline{\phi },\overline{p},\overline{q})\right]\right\}.`$
While for $`\lambda =0`$ this approximation yields the exact propagator, this property is lost for the interacting system. Introducing the complex variables
$`\alpha `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\overline{q}+i\overline{p})`$ (42)
$`\beta `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\overline{q}i\overline{p})`$ (43)
$`\zeta `$ $`=`$ $`\mathrm{tan}\left({\displaystyle \frac{\overline{\vartheta }}{2}}\right)e^{i\overline{\phi }}`$ (44)
$`\eta `$ $`=`$ $`\mathrm{tan}\left({\displaystyle \frac{\overline{\vartheta }}{2}}\right)e^{i\overline{\phi }},`$ (45)
Eq.(41) may be expressed as
$`e^{iS_{cl}}`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \frac{1}{2}}\left[|\alpha ^{}|^2+|\beta ^{\prime \prime }|^2\alpha (T)\beta ^{\prime \prime }\alpha ^{}\beta (0)\right]\right\}\sqrt{{\displaystyle \frac{(1+\zeta ^{}\eta (0))(1+\zeta (T)\eta ^{\prime \prime })}{(1+\zeta ^{}\eta ^{})(1+\zeta ^{\prime \prime }\eta ^{\prime \prime })}}}`$ (47)
$`\times \left({\displaystyle \frac{\zeta ^{\prime \prime }\eta ^{}}{\zeta ^{}\eta ^{\prime \prime }}}\right)^{\frac{1}{4}}\mathrm{exp}\left\{i{\displaystyle _0^T}𝑑t\left[{\displaystyle \frac{i}{2}}\left(\dot{\alpha }\beta \alpha \dot{\beta }\right)+{\displaystyle \frac{i}{2}}{\displaystyle \frac{\dot{\zeta }\eta \zeta \dot{\eta }}{1+\zeta \eta }}H(\alpha ,\beta ,\zeta ,\eta )\right]\right\},`$
where the Hamiltonian (34) reads
$$H(\alpha ,\beta ,\zeta ,\eta )=\alpha \beta +\frac{1+\mathrm{\Delta }}{2}\frac{1\zeta \eta }{1+\zeta \eta }+\lambda \frac{\alpha \zeta +\beta \eta }{1+\zeta \eta }.$$
(48)
Now, the DOPA propagator is determined by the dominant path obeying the classical equations of motion
$`\dot{\alpha }`$ $`=`$ $`i\left[\alpha +\lambda {\displaystyle \frac{\eta }{1+\zeta \eta }}\right]`$ (49)
$`\dot{\beta }`$ $`=`$ $`i\left[\beta +\lambda {\displaystyle \frac{\zeta }{1+\zeta \eta }}\right]`$ (50)
$`\dot{\zeta }`$ $`=`$ $`i\left[(1+\mathrm{\Delta })\zeta \lambda (\beta \alpha \zeta ^2)\right]`$ (51)
$`\dot{\eta }`$ $`=`$ $`i\left[(1+\mathrm{\Delta })\eta \lambda (\alpha \beta \eta ^2)\right],`$ (52)
with the boundary conditions
$`\alpha (0)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(q^{}+ip^{})`$ (53)
$`\beta (T)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(q^{\prime \prime }ip^{\prime \prime })`$ (54)
$`\zeta (0)`$ $`=`$ $`\mathrm{tan}\left({\displaystyle \frac{\vartheta ^{}}{2}}\right)e^{i\phi ^{}}`$ (55)
$`\eta (T)`$ $`=`$ $`\mathrm{tan}\left({\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\right)e^{i\phi ^{\prime \prime }}.`$ (56)
The system of differential equations (52) gives rise to a Hamiltonian vector field. Extending results in , one sees that this Hamiltonian dynamics is identical to the classical mechanics of a spin on the two-sphere coupled to the phase space degrees of freedom of a one dimensional harmonic oscillator. Since the covariant divergence of the Hamiltonian vector field vanishes, this dynamical system is conservative and no attractor can occur. The coupled differential equations (52) with conditions (56) express a nonlinear boundary value problem. We can find a solution exploiting the invariance of the action (33) under phase transformations
$`\zeta `$ $``$ $`\zeta e^{i\mathrm{\Lambda }}`$ (57)
$`\eta `$ $``$ $`\eta e^{i\mathrm{\Lambda }}`$ (58)
$`\alpha `$ $``$ $`\alpha e^{i\mathrm{\Lambda }}`$ (59)
$`\beta `$ $``$ $`\beta e^{i\mathrm{\Lambda }}.`$ (60)
The corresponding integral of motion is
$$N(\alpha ,\beta ,\zeta ,\eta )=\alpha \beta +\frac{1}{2}\frac{1\zeta \eta }{1+\zeta \eta }.$$
(61)
Therefore the Hamiltonian dynamical system becomes integrable by the theorem of Liouville-Arnold . Particularly, by setting $`u=(1\zeta \eta )/(1+\zeta \eta )`$, we reduce the system to a one-dimensional problem of the form
$$\frac{1}{2}\dot{u}^2+V(u)=0,$$
(62)
with the cubic potential
$$V(u)=\lambda ^2(u^3+a_2u^2+a_1u+a_0).$$
(63)
The coefficients read
$`a_0`$ $`=`$ $`2N+2{\displaystyle \frac{C^2}{\lambda ^2}}`$ (64)
$`a_1`$ $`=`$ $`1+2{\displaystyle \frac{\mathrm{\Delta }C}{\lambda ^2}}`$ (65)
$`a_2`$ $`=`$ $`2N+{\displaystyle \frac{\mathrm{\Delta }^2}{2\lambda ^2}},`$ (66)
where
$$C(\alpha ,\beta ,\zeta ,\eta )=\lambda \frac{\alpha \zeta +\beta \eta }{1+\zeta \eta }+\frac{\mathrm{\Delta }}{2}\frac{1\zeta \eta }{1+\zeta \eta }.$$
(67)
Although the potential $`V(u)`$ is time independent, the boundary values (56) enforce the coefficients in Eq.(66) to depend on the end time $`T`$, and the form of $`V(u)`$ changes with $`T`$. Next we set $`v=u+a_2/3`$ and rewrite Eq.(62) as
$$\frac{2}{\lambda ^2}\dot{v}^2=4v^3g_2vg_3.$$
(68)
This is just the differential equation solved by the Weierstrass elliptic function $`\mathrm{}(\frac{\lambda }{\sqrt{2}}t;g_2;g_3)`$ with the invariants
$`g_2`$ $`=`$ $`4\left(a_1{\displaystyle \frac{1}{3}}a_2^2\right)`$ (69)
$`g_3`$ $`=`$ $`{\displaystyle \frac{4}{3}}\left({\displaystyle \frac{2}{9}}a_2^2a_1\right)a_24a_0.`$ (70)
In the following we suppress these invariants in the list of arguments of the function $`\mathrm{}`$. Now, the solution of Eq.(62) becomes
$$u(t)=\frac{a_2}{3}+\mathrm{}\left(A_1+\frac{\lambda }{\sqrt{2}}t\right),$$
(71)
where
$$A_1=\mathrm{}^1\left(\frac{a_2}{3}+\frac{1\zeta ^{}\eta (0)}{1+\zeta ^{}\eta (0)}\right)$$
(72)
is determined by the inverse Weierstrass function $`\mathrm{}^1`$. Making use of the solution (71), the equations of motion lead to elliptic integrals which can be solved in terms of the Weierstrass elliptic functions $`\mathrm{}`$, $`\zeta _w`$ and $`\sigma _w`$ . After some algebra one finds for the field coordinates
$`\alpha (t)`$ $`=`$ $`\alpha ^{}\left[{\displaystyle \frac{1}{2\alpha ^{}\beta (0)}}{\displaystyle \frac{\sigma _w(A_2+A_1)}{\sigma _w(A_2A_1)}}\right]^{1/2}`$ (75)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{\lambda }{\sqrt{2}}}\zeta _w(A_2)t\right\}\mathrm{exp}\left\{i\left(1+{\displaystyle \frac{\mathrm{\Delta }}{2}}\right)t\right\}`$
$`\times \left[\left(2N+{\displaystyle \frac{a_2}{3}}\mathrm{}(A_1+{\displaystyle \frac{\lambda }{\sqrt{2}}}t)\right){\displaystyle \frac{\sigma _w(A_2A_1+\frac{\lambda }{\sqrt{2}}t)}{\sigma _w(A_2+A_1+\frac{\lambda }{\sqrt{2}}t)}}\right]^{1/2},`$
and
$`\beta (t)`$ $`=`$ $`\beta ^{\prime \prime }\left[{\displaystyle \frac{1}{2\alpha (T)\beta ^{\prime \prime }}}{\displaystyle \frac{\sigma _w(A_2A_1\frac{\lambda }{\sqrt{2}}T)}{\sigma _w(A_2+A1+\frac{\lambda }{\sqrt{2}}T)}}\right]^{1/2}`$ (78)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{\lambda }{\sqrt{2}}}\zeta _w(A_2)(Tt)\right\}\mathrm{exp}\left\{i\left(1+{\displaystyle \frac{\mathrm{\Delta }}{2}}\right)(Tt)\right\}`$
$`\times \left[\left(2N+{\displaystyle \frac{a_2}{3}}\mathrm{}(A_1+{\displaystyle \frac{\lambda }{\sqrt{2}}}t)\right){\displaystyle \frac{\sigma _w(A_2+A_1\frac{\lambda }{\sqrt{2}}t)}{\sigma _w(A_2A_1\frac{\lambda }{\sqrt{2}}t)}}\right]^{1/2},`$
where
$$A_2=\mathrm{}^1\left(\frac{a_2}{3}+2N\right).$$
(79)
The spin variables are found to read
$`\zeta (t)`$ $`=`$ $`\zeta ^{}\left[{\displaystyle \frac{1}{\zeta ^{}\eta (0)}}{\displaystyle \frac{\sigma _w(A_3+A_1)}{\sigma _w(A_3A_1)}}{\displaystyle \frac{\sigma _w(A_4+A_1)}{\sigma _w(A_4A_1)}}\right]^{1/2}`$ (82)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{\lambda }{\sqrt{2}}}[\zeta _w(A_3)+\zeta _w(A_4)]t\right\}\mathrm{exp}\{it\}`$
$`\times \left[{\displaystyle \frac{1+\frac{a_2}{3}\mathrm{}(A_1+\frac{\lambda }{\sqrt{2}}t)}{1\frac{a_2}{3}+\mathrm{}(A_1+\frac{\lambda }{\sqrt{2}}t)}}{\displaystyle \frac{\sigma _w(A_3A_1+\frac{\lambda }{\sqrt{2}}t)}{\sigma _w(A_3+A_1+\frac{\lambda }{\sqrt{2}}t)}}{\displaystyle \frac{\sigma _w(A_4A_1+\frac{\lambda }{\sqrt{2}}t)}{\sigma _w(A_4+A_1+\frac{\lambda }{\sqrt{2}}t)}}\right]^{1/2}`$
and
$`\eta (t)`$ $`=`$ $`\eta ^{\prime \prime }\left[{\displaystyle \frac{1}{\zeta (T)\eta ^{\prime \prime }}}{\displaystyle \frac{\sigma _w(A_3+A_1+\frac{\lambda }{\sqrt{2}}T)}{\sigma _w(A_3A_1+\frac{\lambda }{\sqrt{2}}T)}}{\displaystyle \frac{\sigma _w(A_4+A_1+\frac{\lambda }{\sqrt{2}}T)}{\sigma _w(A_4A_1+\frac{\lambda }{\sqrt{2}}T)}}\right]^{1/2}`$ (85)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{\lambda }{\sqrt{2}}}[\zeta _w(A_3)+\zeta _w(A_4)](Tt)\right\}\mathrm{exp}\{i(Tt)\}`$
$`\times \left[{\displaystyle \frac{1+\frac{a_2}{3}\mathrm{}(A_1+\frac{\lambda t}{\sqrt{2}})}{1\frac{a_2}{3}+\mathrm{}(A_1+\frac{\lambda t}{\sqrt{2}})}}{\displaystyle \frac{\sigma _w(A_3A_1+\frac{\lambda }{\sqrt{2}}t)}{\sigma _w(A_3+A_1+\frac{\lambda }{\sqrt{2}}t)}}{\displaystyle \frac{\sigma _w(A_4A_1+\frac{\lambda }{\sqrt{2}}t)}{\sigma _w(A_4+A_1+\frac{\lambda }{\sqrt{2}}t)}}\right]^{1/2}`$
where
$`A_3`$ $`=`$ $`\mathrm{}^1\left({\displaystyle \frac{a_2}{3}}1\right)`$ (86)
$`A_4`$ $`=`$ $`\mathrm{}^1\left({\displaystyle \frac{a_2}{3}}+1\right).`$ (87)
The solutions (75)-(87) give the dominant path in terms of the known initial values $`\alpha (0)=\alpha ^{}`$, $`\zeta (0)=\zeta ^{}`$ and final values $`\beta (T)=\beta ^{\prime \prime }`$, $`\eta (T)=\eta ^{\prime \prime }`$ and as implicit functions of the unknown initial values $`\beta (0)`$, $`\eta (0)`$ and final values $`\alpha (T)`$, $`\zeta (T)`$. Two of these unknowns have to be determined numerically. For instance, from Eqs.(75) and (82) we obtain two transzendental equations for $`\alpha (T)`$ and $`\zeta (T)`$ that can be solved by a root search procedure. Then, the two other unknowns can be found from the two constants $`C`$ and $`N`$.
Having determined the semiclassical trajectory, we may insert the result into Eq.(47) and determine the DOPA-Propagator. Since this propagator obeys a semiclassical Schrödinger equation \[see Appendix A\], an alternative representation of the propagator reads
$$e^{iS_{cl}}=\mathrm{exp}\left\{i_0^T𝑑tH(\alpha (t),\beta ^{\prime \prime },\zeta (t),\eta ^{\prime \prime })\right\}\vartheta ^{\prime \prime }\phi ^{\prime \prime }|\vartheta ^{}\phi ^{}p^{\prime \prime }q^{\prime \prime }|p^{}q^{}.$$
(88)
With this representation the DOPA propagator is just determined by the the endpoint of the classical path.
Although the dynamical system (52) is conservative, it gives rise to stationary states. These are the fix points $`(\alpha ,\beta ,\zeta _N,\eta _N)=(0,0,0,0)`$ and $`(\alpha ,\beta ,\rho _S,\sigma _S)=(0,0,0,0)`$, where $`\rho =1/\zeta `$ and $`\sigma =1/\eta `$. These points correspond to the states $`|\mathrm{\hspace{0.17em}0}`$ and $`|\mathrm{\hspace{0.17em}0}`$ referred to as north pole and south pole, henceforth. For a linear stability analysis we just have to linearize the spin terms since the equations of motion (52) are already linear in the oscillator variables. Expanding about $`(\zeta _N,\eta _N)`$ we find
$$\frac{d}{dt}\left(\begin{array}{c}\delta \zeta \\ \beta \\ \delta \eta \\ \alpha \end{array}\right)=i\left(\begin{array}{cccc}1+\mathrm{\Delta }& \lambda & 0& 0\\ \lambda & 1& 0& 0\\ 0& 0& 1\mathrm{\Delta }& \lambda \\ 0& 0& \lambda & 1\end{array}\right)\left(\begin{array}{c}\delta \zeta \\ \beta \\ \delta \eta \\ \alpha \end{array}\right),$$
(89)
and two invariant subspaces in the variables $`(\delta \zeta ,\beta )`$ and $`(\delta \eta ,\alpha )`$ appear. The solution satisfying the boundary conditions (56) becomes
$`\alpha (t)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}(\mathrm{\Omega }_NT)}}\left\{\alpha ^{}e^{i\omega _mt}\mathrm{cosh}\left[\mathrm{\Omega }_N(Tt)\right]i\eta ^{\prime \prime }e^{i\omega _m(Tt)}\mathrm{sinh}(\mathrm{\Omega }_Nt)\right\}`$ (90)
$`\beta (t)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}(\mathrm{\Omega }_NT)}}\left\{\beta ^{\prime \prime }e^{i\omega _m(Tt)}\mathrm{cosh}(\mathrm{\Omega }_Nt)i\zeta ^{}e^{i\omega _mt}\mathrm{sinh}\left[\mathrm{\Omega }_N(Tt)\right]\right\}`$ (91)
$`\delta \zeta (t)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}(\mathrm{\Omega }_NT)}}\left\{\zeta ^{}e^{i\omega _mt}\mathrm{cosh}\left[\mathrm{\Omega }_N(Tt)\right]i\beta ^{\prime \prime }e^{i\omega _m(Tt)}\mathrm{sinh}(\mathrm{\Omega }_Nt)\right\}`$ (92)
$`\delta \eta (t)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}(\mathrm{\Omega }_NT)}}\left\{\eta ^{\prime \prime }e^{i\omega _m(Tt)}\mathrm{cosh}(\mathrm{\Omega }_Nt)i\alpha ^{}e^{i\omega _mt}\mathrm{sinh}\left[\mathrm{\Omega }_N(Tt)\right]\right\}`$ (93)
with the frequencies
$`\omega _m`$ $`=`$ $`1+{\displaystyle \frac{\mathrm{\Delta }}{2}}`$ (94)
$`\mathrm{\Omega }_N`$ $`=`$ $`\sqrt{\lambda ^2{\displaystyle \frac{\mathrm{\Delta }^2}{4}}}.`$ (95)
Note that for long times the dominant path converges to the corresponding boundary value and no oscillations around the north pole take place anymore.
In the same way, we linearize the motion around the south pole. Now invariant subspaces appear in the variables ($`\delta \rho ,\alpha `$) and ($`\delta \sigma ,\beta `$)
$$\frac{d}{dt}\left(\begin{array}{c}\delta \rho \\ \alpha \\ \delta \sigma \\ \beta \end{array}\right)=i\left(\begin{array}{cccc}1\mathrm{\Delta }& \lambda & 0& 0\\ \lambda & 1& 0& 0\\ 0& 0& 1+\mathrm{\Delta }& \lambda \\ 0& 0& \lambda & 1\end{array}\right)\left(\begin{array}{c}\delta \rho \\ \alpha \\ \delta \sigma \\ \beta \end{array}\right),$$
(96)
with the solution
$`\alpha (t)`$ $`=`$ $`\alpha ^{}e^{i\omega _mt}\mathrm{cos}(\mathrm{\Omega }_St)i{\displaystyle \frac{1}{\zeta ^{}}}e^{i\omega _mt}\mathrm{sin}(\mathrm{\Omega }_St)`$ (97)
$`\beta (t)`$ $`=`$ $`\beta ^{\prime \prime }e^{i\omega _m(Tt)}\mathrm{cos}\left[\mathrm{\Omega }_S(Tt)\right]i{\displaystyle \frac{1}{\eta ^{\prime \prime }}}e^{i\omega _m(Tt)}\mathrm{sin}\left[\mathrm{\Omega }_S(Tt)\right]`$ (98)
$`\delta \rho (t)`$ $`=`$ $`{\displaystyle \frac{1}{\zeta ^{}}}e^{i\omega _mt}\mathrm{cos}(\mathrm{\Omega }_St)i\alpha ^{}e^{i\omega _mt}\mathrm{sin}(\mathrm{\Omega }_St)`$ (99)
$`\delta \sigma (t)`$ $`=`$ $`{\displaystyle \frac{1}{\eta ^{\prime \prime }}}e^{i\omega _m(Tt)}\mathrm{cos}\left[\mathrm{\Omega }_S(Tt)\right]i\beta ^{\prime \prime }e^{i\omega _m(Tt)}\mathrm{sin}\left[\mathrm{\Omega }_S(Tt)\right],`$ (100)
where
$$\mathrm{\Omega }_S=\sqrt{\lambda ^2+\frac{\mathrm{\Delta }^2}{4}}.$$
(101)
Here, the dominant path does not converge for long times but keeps on oscillating around the south pole. North pole and south pole correspond to the local extrema of the cubic potential (63) generated by the coupling of the spin-$`\frac{1}{2}`$ to a vacuum field. Whenever the field becomes filled with bosons, these fix points bifurcate into limit cycles.
The presence of stationary states leads to strong deviations of the DOPA propagator from the exact result for times large compared to $`\omega _o^1`$. In fact, for long times the semiclassical trajectory approaches the saddle point of the cubic potential and stays there for most of the time. For the full quantum problem the state $`|0`$ is not a steady state, rather it will decay by spontaneous emission. In the semiclassical approximation spontaneous emission arises from fluctuations about the classical path that are neglected in the DOPA. Hence, to obtain useful results also for long times, fluctuations about the north pole need to be taken into account.
## IV Fluctuations
The semiclassical expansion of the path integral (10) leads to second order contributions in terms of Gaussian fluctuation path integrals. Denoting by $`(x_a,y_a)`$ and $`(x_b,y_b)`$ deviations from the dominant path variables $`(p,q)`$ and $`(\mathrm{cos}(\vartheta ),\phi )`$, the semiclassical approximation takes the form
$$\vartheta ^{\prime \prime }\phi ^{\prime \prime }p^{\prime \prime }q^{\prime \prime }|U(T)|\vartheta ^{}\phi ^{}p^{}q^{}_{sc}=e^{iS_{cl}}\underset{\nu _a,\nu _b\mathrm{}}{lim}𝑑\mu _a𝑑\mu _b\mathrm{exp}\left\{i\delta ^2S[x_a(t),y_a(t),x_b(t),y_b(t)]\right\},$$
(102)
with the boundary conditions $`x_a(0)=x_a(T)=0`$, $`y_a(0)=y_a(T)=0`$, $`x_b(0)=x_b(T)=0`$, $`y_b(0)=y_b(T)=0`$. Since the canonical Wiener measure (36) is of quadratic form, the measure of the fluctuation path integral becomes
$$d\mu _a=\underset{t=0}{\overset{T}{}}\frac{1}{2\pi }dx_a(t)dy_a(t)\mathrm{exp}\left\{\frac{1}{2\nu _a}_0^T𝑑t\left[\dot{x}_a^2+\dot{y}_a^2\right]\right\},$$
(103)
which is of the same form as the original coherent state path measure. On the other hand, the spin measure (37) is not quadratic, and the dominant path $`(\vartheta (t),\phi (t))`$ cannot be separated from the fluctuation variables $`x_b`$ and $`y_b`$. We have
$`d\mu _b`$ $`=`$ $`{\displaystyle \underset{t=0}{\overset{T}{}}}{\displaystyle \frac{2s+1}{4\pi s}}dx_b(t)dy_b(t)\mathrm{exp}\{{\displaystyle \frac{1}{2\nu _b}}{\displaystyle _0^T}dt[{\displaystyle \frac{\dot{x}_b^2}{\mathrm{sin}^2(\vartheta )}}+\mathrm{sin}^2(\vartheta )\dot{y}_b^2`$ (105)
$`\mathrm{cos}(\vartheta )\dot{\phi }x_b\dot{y}_b+2{\displaystyle \frac{\mathrm{cos}(\vartheta )\dot{\vartheta }}{\mathrm{sin}^3(\vartheta )}}\dot{x}_bx_b+(\dot{\phi }^2{\displaystyle \frac{2\mathrm{cos}^2(\vartheta )\dot{\vartheta }^2+\mathrm{sin}(2\vartheta )\ddot{\vartheta }}{2\mathrm{sin}^4(\vartheta )}})x_b^2]\},`$
and the regularization of the fluctuation path integral becomes in general time dependent. However, when the dominant spin path is strictly independent of time, $`(\vartheta (t),\phi (t))=(\vartheta _o,\phi _o)`$, the measure (105) simplifies considerably and we get
$`d\mu _b`$ $`=`$ $`{\displaystyle \underset{t=0}{\overset{T}{}}}{\displaystyle \frac{2s+1}{4\pi s}}dx_b(t)dy_b(t)\mathrm{exp}\left\{{\displaystyle \frac{1}{2\nu _b}}{\displaystyle _0^T}𝑑t\left[{\displaystyle \frac{\dot{x}_b^2}{\mathrm{sin}^2(\vartheta _o)}}+\mathrm{sin}^2(\vartheta _o)\dot{y}_b^2\right]\right\}.`$ (106)
Then, after a canonical transformation
$`\stackrel{~}{x}_b`$ $`=`$ $`{\displaystyle \frac{x_b}{\mathrm{sin}(\vartheta _o)}}`$ (107)
$`\stackrel{~}{y}_b`$ $`=`$ $`\mathrm{sin}(\vartheta _o)y_b,`$ (108)
the measure(106) takes for large $`s`$ the form of the canonical measure (103)
$$d\mu _b=\underset{t=0}{\overset{T}{}}\frac{1}{2\pi }d\stackrel{~}{x}_b(t)d\stackrel{~}{y}_b(t)\mathrm{exp}\left\{\frac{1}{2\nu _b}_0^T𝑑t\left[\dot{\stackrel{~}{x}}_b^2+\dot{\stackrel{~}{y}}_b^2\right]\right\}.$$
(109)
Both measures give rise to the regularization dependent second order variational action
$`\delta ^2S_{\nu _a,\nu _b}[x_a(t),y_a(t),\stackrel{~}{x}_b(t),\stackrel{~}{y}_b(t)]={\displaystyle _0^T}dt[{\displaystyle \frac{i}{2\nu _a}}(\dot{x}_a^2+\dot{y}_a^2)+{\displaystyle \frac{i}{2\nu _b}}(\dot{\stackrel{~}{x}}_b^2+\dot{\stackrel{~}{y}}_b^2)`$ (110)
$`+{\displaystyle \frac{1}{2}}(x_a\dot{y}_a\dot{x}_ay_a)+{\displaystyle \frac{1}{2}}(\stackrel{~}{x}_b\dot{\stackrel{~}{y}}_b\dot{\stackrel{~}{x}}_b\stackrel{~}{y}_b)H_o(x_a,y_a,\stackrel{~}{x}_b,\stackrel{~}{y}_b,t)],`$ (111)
where the Hamiltonian $`H_o(t)`$ is determined by the second order contributions of the Hamiltonian $`H`$ expanded around the dominant path
$`H_o(x_a,y_a,\stackrel{~}{x}_b,\stackrel{~}{y}_b,t)`$ $`=`$ $`a_1(t)x_a^2+a_2(t)x_ay_a+a_3(t)y_a^2`$ (114)
$`+b_1(t)\stackrel{~}{x}_b^2+b_2(t)\stackrel{~}{x}_b\stackrel{~}{y}_b+b_3(t)\stackrel{~}{y}_b^2`$
$`+c_1(t)x_a\stackrel{~}{x}_b+c_2(t)x_a\stackrel{~}{y}_b+c_3(t)y_a\stackrel{~}{x}_b+c_4(t)y_a\stackrel{~}{y}_b,`$
with the coefficients
$`a_1(t)={\displaystyle \frac{1}{2}}{\displaystyle \frac{^2H}{\overline{q}^2}},a_2(t)={\displaystyle \frac{^2H}{\overline{q}\overline{p}}},a_3(t)={\displaystyle \frac{1}{2}}{\displaystyle \frac{^2H}{\overline{p}^2}}`$ (115)
$`b_1(t)={\displaystyle \frac{\mathrm{sin}^2(\vartheta _o)}{2s}}{\displaystyle \frac{^2H}{\mathrm{cos}(\overline{\vartheta })^2}},b_2(t)={\displaystyle \frac{1}{s}}{\displaystyle \frac{^2H}{\overline{\phi }\mathrm{cos}(\overline{\vartheta })}},b_3(t)={\displaystyle \frac{1}{2s\mathrm{sin}^2(\vartheta _o)}}{\displaystyle \frac{^2H}{\overline{\phi }^2}}`$ (116)
$`c_1(t)={\displaystyle \frac{\mathrm{sin}(\vartheta _o)}{\sqrt{s}}}{\displaystyle \frac{^2H}{\overline{p}\mathrm{cos}(\overline{\vartheta })}},c_2(t)={\displaystyle \frac{1}{\sqrt{s}\mathrm{sin}(\vartheta _o)}}{\displaystyle \frac{^2H}{\overline{p}\overline{\phi }}},`$ (117)
$`c_3(t)={\displaystyle \frac{\mathrm{sin}(\vartheta _o)}{\sqrt{s}}}{\displaystyle \frac{^2H}{\overline{q}\mathrm{cos}(\overline{\vartheta })}},c_4(t)={\displaystyle \frac{1}{\sqrt{s}\mathrm{sin}(\vartheta _o)}}{\displaystyle \frac{^2H}{\overline{q}\overline{\phi }}}.`$ (118)
For large $`s`$, starting and end points in the coherent state fluctuation path integral (102) parameterize states $`|x_a(0)y_a(0)x_b(0)y_b(0)`$ and $`|x_a(T)y_a(T)x_b(T)y_b(T)`$ which correspond to product vacuum states. Propagators leading to stationary saddle points $`(\vartheta _o,\phi _o)`$ may be represented now as
$$\vartheta _o\phi _op^{\prime \prime }q^{\prime \prime }|U(T)|\vartheta _o\phi _op^{}q^{}_{sc}=e^{iS_{cl}}\mathrm{0\hspace{0.17em}0}|U_o(T)|\mathrm{0\hspace{0.17em}0},$$
(119)
with the unitary time evolution operator
$$U_o(T)=𝒯_t\mathrm{exp}\left\{i_0^T𝑑tH_o(t)\right\},$$
(120)
determined by the quadratic Hamiltonian
$`H_o(t)`$ $`=`$ $`a_1(t)(Q_a^2{\displaystyle \frac{1}{2}})+a_2(t)(P_aQ_a+Q_aP_a)+a_3(t)(P_a^2{\displaystyle \frac{1}{2}})`$ (123)
$`+b_1(t)(Q_b^2{\displaystyle \frac{1}{2}})+b_2(t)(P_bQ_b+Q_bP_b)+b_3(t)(P_b^2{\displaystyle \frac{1}{2}})`$
$`+c_1(t)P_aP_b+c_2(t)P_aQ_b+c_3(t)Q_aP_b+c_4(t)Q_aQ_b,`$
describing two driven coupled oscillators.
As we have seen in the previous section, for the Jaynes-Cummmings model the north pole $`|0`$ becomes a steady state in the DOPA, and it is essential to take fluctuations about this state into account. Unfortunately, the description of the spin degrees of freedom with spherical coordinates leads to coordinate singularities. Particularly, the azimuthal angle $`\phi `$ is undefined at the poles of the two-sphere. To calculate fluctuations about the north pole accurately, we change the coordinate system by a rotation. Since rotations are isometrical canonical transformations, the spin path measure (37) stays invariant but the kinematical term is not preserved. Instead a phase factor appears which only vanishes if starting and endpoint of the spin coordinates are identical.
Within the DOPA, the probability amplitude to remain at the north pole is just a phase factor
$$0|U(T)|0_{DOPA}=e^{iS_{cl}}=\mathrm{exp}\{i\frac{1+\mathrm{\Delta }}{2}T\},$$
(124)
and the north pole becomes a steady state. Taking now Gaussian fluctuations into account we have
$$0|U(T)|0_{sc}=e^{iS_{cl}}\mathrm{0\hspace{0.17em}0}|U_o(T)|\mathrm{0\hspace{0.17em}0},$$
(125)
where the the vacuum amplitude is determined by the time independent Hamiltonian
$$H_o=\frac{1}{2}(P_a^2+Q_a^21)\frac{1+\mathrm{\Delta }}{2}(P_b^2+Q_b^21)+\lambda (P_bQ_a+Q_bP_a).$$
(126)
For convenience we represent the operators $`Q_a,P_a`$ and $`Q_b,P_b`$ by corresponding creation and annihilation operators $`a,a^{}`$ and $`b,b^{}`$
$$H_o=1+(1+\frac{\mathrm{\Delta }}{2})(aa^{}b^{}b)\frac{\mathrm{\Delta }}{2}(aa^{}+b^{}b)i\lambda (aba^{}b^{}).$$
(127)
Since $`aa^{}b^{}b`$ commutes with $`H_o`$, we rewrite the time evolution operator in the form
$$U_o(T)=\mathrm{exp}\left\{i\frac{\mathrm{\Delta }T}{2}\right\}\mathrm{exp}\left\{i(1+\frac{\mathrm{\Delta }}{2})a^{}aT\right\}\mathrm{exp}\left\{i(1+\frac{\mathrm{\Delta }}{2})b^{}bT\right\}U_1(T),$$
(128)
with $`U_1(T)=\mathrm{exp}(iH_1T)`$ and
$$H_1=\left[\frac{\mathrm{\Delta }}{2}(aa^{}+b^{}b)+i\lambda (aba^{}b^{})\right].$$
(129)
The operators $`aa^{}+b^{}b`$, $`ab`$ and $`a^{}b^{}`$ span the three dimensional $`su(1,1)`$ Lie algebra with commutators
$`[aa^{}+b^{}b,ab]`$ $`=`$ $`2ab`$ (130)
$`[aa^{}+b^{}b,a^{}b^{}]`$ $`=`$ $`2a^{}b^{}`$ (131)
$`[ab,a^{}b^{}]`$ $`=`$ $`aa^{}+b^{}b.`$ (132)
For this algebra there is a decomposition into one-dimensional $`SU(1,1)`$ transformations which holds for the whole group, i.e. for all times . We start with the ansatz
$$U_1(T)=\mathrm{exp}\{\mu (T)a^{}b^{}\}\mathrm{exp}\{\nu (T)ab\}\mathrm{exp}\{\xi (T)(aa^{}+b^{}b)\},$$
(133)
which results in the vacuum amplitude
$$\mathrm{0\hspace{0.17em}0}|U_o(T)|\mathrm{0\hspace{0.17em}0}=\mathrm{exp}\left\{i\frac{\mathrm{\Delta }}{2}T+\xi (T)\right\}.$$
(134)
Then, requiring that $`U_1(T)`$ obeys the Schrödinger equation $`d/dTU_1(T)=iH_1U_1(T)`$, we get the relation
$`i{\displaystyle \frac{\mathrm{\Delta }}{2}}(aa^{}+b^{}b)+\lambda (a^{}b^{}ab)`$ $`=`$ $`\dot{\mu }a^{}b^{}+\dot{\nu }e^{\mu a^{}b^{}}abe^{\mu a^{}b^{}}`$ (136)
$`+\dot{\xi }e^{\mu a^{}b^{}}e^{\nu ab}(aa^{}+b^{}b)e^{\nu ab}e^{\mu a^{}b^{}},`$
where we have made use of the Baker-Campbell-Hausdorff formula. Further, the commutation relations (132) imply
$`e^{\nu ab}(aa^{}+b^{}b)e^{\nu ab}`$ $`=`$ $`aa^{}+b^{}b+2\nu ab`$ (137)
$`e^{\mu a^{}b^{}}(aa^{}+b^{}b)e^{\mu a^{}b^{}}`$ $`=`$ $`aa^{}+b^{}b2\mu a^{}b^{}`$ (138)
$`e^{\mu a^{}b^{}}abe^{\mu a^{}b^{}}`$ $`=`$ $`ab\mu (aa^{}+b^{}b)+\mu ^2a^{}b^{}.`$ (139)
Now, Eq.(136) determines the time rate of change of the functions $`\mu `$, $`\nu `$ and $`\xi `$ by the linear equations
$$\left(\begin{array}{c}\lambda \\ \lambda \\ i\mathrm{\Delta }\end{array}\right)=\left(\begin{array}{ccc}1& \mu ^2& 2\mu (1\mu \nu )\\ 0& 1& 2\nu \\ 0& 2\mu & 2(12\mu \nu )\end{array}\right)\left(\begin{array}{c}\dot{\mu }\\ \dot{\nu }\\ \dot{\xi }\end{array}\right),$$
(140)
which are readily solved with the initial conditions $`\mu (0)=0`$, $`\nu (0)=0`$ and $`\xi (0)=0`$. In particular, we get for the function $`\xi (T)`$ in Eq.(133)
$$\xi (T)=i\mathrm{\Delta }T+\mathrm{log}\left[\mathrm{cos}(\mathrm{\Omega }T)i\frac{\mathrm{\Delta }}{2}\mathrm{sin}(\mathrm{\Omega }T)\right],$$
(141)
with the Rabi frequency
$$\mathrm{\Omega }=\sqrt{\lambda ^2+\frac{\mathrm{\Delta }^2}{4}}.$$
(142)
Hence, the vacuum amplitude (134) becomes
$$\mathrm{0\hspace{0.17em}0}|U_o(T)|\mathrm{0\hspace{0.17em}0}=\mathrm{exp}\left\{i\frac{\mathrm{\Delta }T}{2}\right\}\left[\mathrm{cos}(\mathrm{\Omega }T)i\frac{\mathrm{\Delta }}{2}\mathrm{sin}(\mathrm{\Omega }T)\right].$$
(143)
and the semiclassical propagator with fluctuations
$$0|U(T)|0_{sc}=e^{\frac{i}{2}T}[\mathrm{cos}(\mathrm{\Omega }T)i\frac{\mathrm{\Delta }}{2}\mathrm{sin}(\mathrm{\Omega }T)]$$
(144)
includes spontaneous emission leading to an instability of the north pole. Eq.(144) gives the exact matrix element of the propagator sandwiched between north pole states.
When the field is initially and finally not in the vacuum state, the semiclassical propagator (102) is no longer characterized by a fix point path. An evaluation of the fluctuations about the semiclassical path would then require numerical methods beyond the scope of this article.
###### Acknowledgements.
The authors would like to thank Joachim Ankerhold and Jürgen Stockburger for valuable discussions. This work was supported by the Deutsche Forschungsgemeinschaft (Bonn) through the Schwerpunktprogramm “Zeitabhängige Phänomene und Methoden in Quantensystemen der Physik und Chemie”.
## A Semiclassical Schrödinger equation
Here we derive the semiclassical Schrödinger equation for the DOPA propagator given in Eq.(47). The time rate of change is readily evaluated, and after an integration by parts it may be expressed as
$`{\displaystyle \frac{}{T}}e^{iS_{cl}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{{\displaystyle \frac{\alpha (T,T)}{T}}\beta ^{\prime \prime }\alpha ^{}{\displaystyle \frac{\beta (0,T)}{T}}{\displaystyle \frac{\zeta ^{}\frac{\eta (0,T)}{T}}{1+\zeta ^{}\eta (0,T)}}{\displaystyle \frac{\frac{\zeta (T,T)}{T}\eta ^{\prime \prime }}{1+\zeta (T,T)\eta ^{\prime \prime }}}.`$ (A7)
$`{\displaystyle \frac{\alpha (t,T)}{t}}|_T\beta ^{\prime \prime }+\alpha (T,T){\displaystyle \frac{\beta (t,T)}{t}}|_T{\displaystyle \frac{\frac{\zeta (t,T)}{t}|_T\eta ^{\prime \prime }\zeta (T,T)\frac{\eta (t,T)}{t}|_T}{1+\zeta (T,T)\eta ^{\prime \prime }}}`$
$`2iH(\alpha (T,T),\beta ^{\prime \prime },\zeta (T,T),\eta ^{\prime \prime })`$
$`\left[{\displaystyle \frac{\alpha (t,T)}{T}}\beta (t,T)\alpha (t,T){\displaystyle \frac{\beta (t,T)}{T}}+{\displaystyle \frac{\frac{\zeta (t,T)}{T}\eta (t,T)\zeta (t,T)\frac{\eta (t,T)}{T}}{1+\zeta (T,T)\eta ^{\prime \prime }}}\right]_{t=0}^{t=T}`$
$`{\displaystyle _0^T}dt[{\displaystyle \frac{\alpha (t,T)}{T}}({\displaystyle \frac{\beta (t,T)}{t}}i{\displaystyle \frac{H}{\alpha }}){\displaystyle \frac{\beta (t,T)}{T}}({\displaystyle \frac{\alpha (t,T)}{t}}+i{\displaystyle \frac{H}{\beta }}).`$
$`.{\displaystyle \frac{\zeta (t,T)}{T}}({\displaystyle \frac{\frac{\eta (t,T)}{t}}{(1+\zeta (t,T)\eta (t,T))^2}}i{\displaystyle \frac{H}{\zeta }})`$
$`{\displaystyle \frac{\eta (t,T)}{T}}({\displaystyle \frac{\frac{\zeta (t,T)}{t}}{(1+\zeta (t,T)\eta (t,T))^2}}+i{\displaystyle \frac{H}{\eta }})]\}e^{iS_{cl}}.`$
Using the classical equations of motions, the integral is found to vanish. Then, we rewrite the remaining parts in the form
$`{\displaystyle \frac{}{T}}e^{iS_{cl}}`$ $`=`$ $`iH(\alpha (T,T),\beta ^{\prime \prime },\zeta (T,T),\eta ^{\prime \prime })`$ (A14)
$`{\displaystyle \frac{1}{2}}[\beta ^{\prime \prime }({\displaystyle \frac{\alpha (T,T)}{T}}+{\displaystyle \frac{\alpha (t,T)}{t}}|_T+{\displaystyle \frac{\alpha (t,T)}{T}}|_T)`$
$`+\alpha ^{}\left({\displaystyle \frac{\beta (0,T)}{T}}+{\displaystyle \frac{\beta (t,T)}{T}}|_0\right)`$
$`+\alpha (T,T)\left({\displaystyle \frac{\beta (t,T)}{t}}|_T{\displaystyle \frac{\beta (t,T)}{T}}|_T\right)\beta (0,T){\displaystyle \frac{\alpha (t,T)}{T}}|_0`$
$`+{\displaystyle \frac{\eta ^{\prime \prime }\left(\frac{\zeta (T,T)}{T}+\frac{\zeta (t,T)}{t}|_T+\frac{\zeta (t,T)}{T}|_T\right)+\zeta (T,T)\left(\frac{\eta (t,T)}{t}|_T\frac{\eta (t,T)}{T}|_T\right)}{1+\zeta (T,T)\eta ^{\prime \prime }}}`$
$`+{\displaystyle \frac{\zeta ^{}\left(\frac{\eta (0,T)}{T}+\frac{\eta (t,T)}{T}|_0\right)\eta (0,T)\frac{\zeta (0,T)}{T}|_0}{1+\zeta ^{}\eta (0,T)}}]e^{iS_{cl}},`$
where most of the terms on the right hand site vanish. Finally we get
$$\frac{}{T}e^{iS_{cl}}=iH(\alpha (T),\beta ^{\prime \prime },\zeta (T),\eta ^{\prime \prime })e^{iS_{cl}}.$$
(A15)
Note that the matrix element of the Hamiltonian at the endpoint of the dominant path $`(\alpha (T),\beta ^{\prime \prime },\zeta (T),\eta ^{\prime \prime })`$ generates the time rate of change of the DOPA propagator and not the matrix element of the final state $`|\vartheta ^{\prime \prime }\phi ^{\prime \prime }p^{\prime \prime }q^{\prime \prime }`$. For a spin-$`\frac{1}{2}`$ coupled to a classical field this Schrödinger equation generates the exact quantum mechanics .
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# Dislocation theory as a 3-dimensional translation gauge theory*footnote **footnote *Dedicated to Professor Ekkehart Kröner, Stuttgart, on the occasion of his 80th birthday.
## I Introduction
The field theory of defects in crystals has an old history. First, Kondo and Bilby et al. described independently a dislocation theory in the language of differential geometry. They proved the equivalence of dislocation density to Cartan’s torsion. Kröner and Seeger completed this theory to a non-linear theory of elasticity with dislocations and internal stress. Recent developments of this theory are given in .
The first step towards a gauge theory of defects had been taken by Turski . He derived the equilibrium equations by means of a variational principle.
Edelen et al. developed a theory of dislocations and disclinations as a gauge theory of the Euclidean group $`SO(3)\times T(3)`$ and a pure dislocation theory as a translational or $`T(3)`$-gauge theory. Unfortunately, they did not distinguish between disclinations in solid and in liquid crystals. Moreover, they did not refer to the gauge theory of the Poincaré group $`SO(1,3)\times T(4)`$ which was very well developed in gravity . Their theory is claimed to describe the dynamics of defects. However, they have ignored the following effects: Dislocations can move in two different modes, called glide and climb. If dislocations climb, they interchange with point defects such as vacancies and/or interstitials. Therefore, it is necessary to include the point defects into a dynamical theory of dislocations. Furthermore, the motion of dislocation is highly dissipative (friction and radiation damping).
Other proposals for gauge theories of dislocations were put forward by Gairola , Kleinert, Trzȩsowski , and, quite recently, by Malyshev .
The classical theory of dislocations deals with the incompatibility and the equilibrium conditions (see, e.g., ). The drawback of this theory is that a field equation is missing. Therefore, one has to take as ansatz for the dislocation density a singular $`\delta `$-function. The price to be paid is that the elastic energy is singular. Thus the information about the dislocation core is lost.
The aim of this paper is to develop a $`T(3)`$-gauge theory of dislocations similar to gravity as $`T(4)`$-gauge theory. These gauge theories are included in the metric-affine gauge theory (MAG) as given by Hehl et al. , where the metric, the coframe, and the connection are independent field variables. We use the tool of MAG in order to derive the gauge theory of dislocations in an elastoplastic material as field theory of elastoplasticity.
We obtain a system of static field equations which determine both, the elastic and the plastic fields. Additionally, the equations contain the equilibrium condition and the Bianchi identity for the torsion (or dislocation density). As shown later, our field equations differ from the field equations proposed by Edelen et al. . One reason is that we have used a different and more physically motivated gauge Lagrangian. For compact notation, we will use the calculus of exterior differential forms.
## II Geometrical framework
In order to gauge the translation group, it is convenient to start with the affine group. Therefore, we consider the affine group in three dimensions $`A(3)=GL(3)\times T(3)`$, where the symbol $`\times `$ denotes the semidirect product.
The mathematical domain of gauge theories is the theory of fibre bundles ; for rigorous definitions of fibre bundles and connections, see . In the case of the affine group, the corresponding fibre bundle is the bundle of affine frames $`AM`$. The connection 1-form $`\stackrel{~}{\stackrel{~}{\omega }}`$ on $`AM`$ is called generalized affine connection. It is an $`𝔞(3)`$-valued 1-form of type ad, where $`𝔞(3)`$ is the Lie algebra of $`A(3)`$. The corresponding 2-form $`\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}`$ is the curvature 2-form on $`AM`$. It is given by Cartan’s structure equation
$`\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}:=\stackrel{~}{\stackrel{~}{\text{D}}}\stackrel{~}{\stackrel{~}{\omega }}\text{d}\stackrel{~}{\stackrel{~}{\omega }}+\stackrel{~}{\stackrel{~}{\omega }}\stackrel{~}{\stackrel{~}{\omega }}.`$ (1)
The symbol $``$ denotes the exterior product, d is the exterior derivative and D the gauge covariant exterior derivative.
Using the subbundle map $`\gamma :LMAM`$, $`\stackrel{~}{\stackrel{~}{\omega }}`$ can be pulled back to the bundle of linear frames $`LM`$,
$`\gamma ^{}\stackrel{~}{\stackrel{~}{\omega }}=\omega +\varphi =\omega ^{ab}L_{ab}+\varphi ^aP_a,`$ (2)
where $`\omega `$ is a $`𝔤l(3)`$-valued 1-form on $`LM`$, the linear connection, and $`\varphi `$ is an $`^3`$-valued tensorial 1-form on $`LM`$, the translational connection. Here $`P_a`$ are the generators of translations $`^3`$ and the generators $`L_{ab}`$ span the Lie algebra $`𝔤l(3)`$. The pull back of $`\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}`$ is given by
$`\gamma ^{}\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}=\mathrm{\Omega }+\mathrm{\Phi }=\mathrm{\Omega }^{ab}L_{ab}+\mathrm{\Phi }^aP_a,`$ (3)
where
$`\mathrm{\Omega }:=\text{D}\omega =\text{d}\omega +\omega \omega `$ (4)
is the curvature of the linear connection $`\omega `$. The translational curvature $`\mathrm{\Phi }`$ reads
$`\mathrm{\Phi }:=\text{D}\varphi =\text{d}\varphi +\omega \varphi .`$ (5)
In particular, the translational connection can be decomposed into the soldering 1-form $`\vartheta `$ and an additional 1-form $`\chi `$,
$`\varphi =\vartheta +\chi .`$ (6)
In order to ensure the tensorial transformation behaviour of $`\vartheta =\vartheta ^aP_a`$, which is an $`^3`$-valued tensorial 1-form of type id, it is convenient to follow Trautman . We introduce the vector-valued 0-form $`\xi \xi ^aP_a`$ and identify $`\chi =\text{D}\xi `$. With these identifications, the soldering-form turns out to be
$`\vartheta =\stackrel{~}{\stackrel{~}{\text{D}}}\xi \text{D}\xi +\varphi .`$ (7)
It transforms as a vector-valued 1-form. With (5) and (7), the torsion 2-form, in the framework of the affine group, reads
$`TT^aP_a:=\text{D}\vartheta =\left(\mathrm{\Phi }^a+\mathrm{\Omega }_b^a\xi ^b\right)P_a.`$ (8)
Since the translation group is a subgroup of $`A(3)`$, we obtain, with the choice $`\omega 0`$, the translational part of the generalized affine connection as
$`\varphi =\vartheta \text{d}\xi .`$ (9)
Some remarks on the generalized affine connection in relation to dislocations have been made by Mistura .
## III Physical identifications and classical elasticity
The soldering form is given with respect to the translational part of the generalized affine connection $`\varphi ^a`$ by
$`\vartheta ^a=\text{d}\xi ^a+\varphi ^a.`$ (10)
Eq. (10) is the key formula for the $`T(3)`$-gauge theory. It is valid in a Weitzenböck space ($`\mathrm{\Omega }_b^a=0`$) under the gauge condition $`\omega _b^a0`$. The coupling in eq. (10) between the translational gauge potential $`\varphi ^a`$ and the vector field $`\xi ^a`$ is characteristic for the translation gauge theory as a special case of the metric–affine gauge theory. The reason lies in the nature of the affine group itself.
At first, in agreement with Edelen , we identify the translational connection or the translational gauge potential $`\varphi ^a`$ as the gauge potential of the dislocations. This identification is justified by the following ideas. A gauge potential is a 1-form. If we form the corresponding 2-form, there results the gauge field strength. In our case of a Weitzenböck space, in the gauge chosen, the torsion 2-form is equal to the object of anholonomity. Therefore, one obtains
$`T^a=\text{d}\vartheta ^a\text{d}\varphi ^a={\displaystyle \frac{1}{2}}T_{ij}^a\text{d}x^i\text{d}x^j,`$ (11)
which can be identified with the dislocation density. One recovers the conventional dislocation density tensor $`\alpha _i^a`$ from $`T^a`$ by means of $`\alpha _i^a=\frac{1}{2}ϵ_i^{jk}T_{jk}^a`$. Thus, the dislocation density or the torsion represents the gauge field strength of the translation group. In the framework of $`T(3)`$-gauge theory ($`\omega _b^a0`$), the torsion 2-form satisfies the Bianchi identity
$`\text{d}T^a=0.`$ (12)
The indices $`a,b,c,\mathrm{}=1,2,3`$ denote the (anholonomic) material or the final coordinates and $`i,j,k,\mathrm{}=1,2,3`$ the (holonomic) Cartesian coordinates of the reference system (defect-free or ideal reference system)A field of coframes $`\vartheta `$ is holonomic if $`\text{d}\vartheta =0`$ and anholonomic if $`\text{d}\vartheta 0`$.. Now, we identify $`\xi ^a`$ with the diffeomorphisms of the material space into the Euclidean space in which the crystal is embedded. We introduce the displacement field $`u^a`$ and write $`\xi ^a`$ in terms of it as $`\xi ^a=\delta _i^ax^i+u^a`$. With this identification, the soldering form is specified
$`\vartheta ^a=B_i^a\text{d}x^i=\left(\delta _i^a+_iu^a+\varphi _i^a\right)\text{d}x^i.`$ (13)
The following consideration justifies the last identification: In a material with compatible distortion $`\varphi ^a=0`$, the soldering form is $`\vartheta ^a=\text{d}\xi ^a`$. The metric of the final state is the Cauchy-Green strain tensor $`G`$ which, in our case (teleparallelism), is given by
$`G=\delta _{ab}\vartheta ^a\vartheta ^b=B_i^aB_{aj}\text{d}x^i\text{d}x^j=g_{ij}(x)\text{d}x^i\text{d}x^j,\delta _{ab}=\text{diag}(+++).`$ (14)
The engineering strain tensor is
$`2E=G1=(g_{ij}\delta _{ij})\text{d}x^i\text{d}x^j.`$ (15)
It is obvious that the soldering 1-form corresponds to the distortion 1-form. If the distortion is compatible, the distortion is given as a deformation gradient of a vector field. This vector field is $`\xi ^a`$.
In the general case of elastoplasticity, eqs. (10) and (13) describe an incompatible distortion with the distortion 1-form $`𝜷\mathit{\vartheta }`$. The cause of plasticity are defects and the material gives rise to a specific elastic response. The plastic distortion is given by the dislocation gauge potential $`𝜷^{\mathrm{pl}}\mathit{\varphi }`$. Finally, the total distortion $`𝜷^\mathrm{T}`$ contains elastic and plastic contributions according to
$`𝜷^\mathrm{T}=𝜷^{\mathrm{pl}}+𝜷.`$ (16)
Since it is compatible, we have
$`\text{d}𝜷^\mathrm{T}=0,\text{d}𝜷^{\mathrm{pl}}=\text{d}𝜷.`$ (17)
Thus the physical space (crystal) is determined by the two fields $`\text{d}\xi ^a`$ and $`\varphi ^a`$.
Let us now discuss the concepts of internal and external observers as introduced by Kröner . The internal observer lives in the crystal and uses $`\vartheta ^a`$ as coframe. Consequently, he can detect defects due to $`\text{d}\vartheta ^a=T^a`$. But he misses the information about the holonomic coordinate system and is unable to detect compatible deformations. The external observer lives in the external space in which the crystal is embedded. He has more information available than the internal observer because he knows the holonomic coordinate system and is able to detect compatible deformations. The external observer measures the torsion in his coordinate system as $`T^i=B_a^iT^a`$.
In gravity, the physical meaning of $`\xi ^a`$ is not well-understood . It is obvious that we are internal observers in our universe and cannot detect $`\xi ^a`$ or $`\text{d}\xi ^a`$. The internal observer in gravity would use the Cartan or affine connection ($`\vartheta ^a,\omega _b^a`$) instead of the generalized affine connection. Does, perhaps, this idea also lead to an understanding of $`\xi ^a`$ in gravity?
In order to simplify the formulas, we assume a linear constitutive law (stress-strain relation). But our general considerations are not restricted to this assumption. The elastic (anisotropic) behaviour of the material can be described by the elasticity tensor
$`C=C^{ijkl}_i_j_k_l.`$ (18)
$`C`$ is a contravariant tensor of fourth rank with the symmetries
$`C^{ijkl}=C^{jikl}=C^{ijlk}=C^{klij}.`$ (19)
The elastic energy (potential energy) contains the constitutive law and is given by
$`W={\displaystyle \frac{1}{2}}C:(EE)={\displaystyle \frac{1}{2}}C^{ijkl}E_{ij}E_{kl},`$ (20)
where the symbol $`:`$ denotes double contraction. Since we consider a static and not a dynamical theory, the elastic Lagrangian is given by means of the potential energy
$`_{\mathrm{el}}={}_{}{}^{}W=W\text{d}v_\mathrm{E},`$ (21)
where the volume 3-form is defined by
$`\text{d}v_\mathrm{E}:={\displaystyle \frac{1}{3!}}\eta _{abc}\vartheta ^a\vartheta ^b\vartheta ^c,`$ (22)
with $`\eta _{abc}:=\sqrt{|B|}ϵ_{abc}`$, and $`ϵ_{abc}`$ as the Levi-Civita symbol. Furthermore, $`\eta _a:=e_a\text{d}v_\mathrm{E}`$, $`\eta _{ab}:=e_ae_b\text{d}v_\mathrm{E}`$, and $`\eta _{abc}:=e_ae_be_c\text{d}v_\mathrm{E}`$, where $``$ denotes the interior product with
$`e_a\vartheta ^b=B_a^iB^b_i=\delta _a^b,e_a=B_a^i_i.`$ (23)
The symbol denotes the Hodge star operator which, in three dimensions, defines the dual $`(3p)`$-form of a given $`p`$-form.
## IV Gauge theory of dislocations
Let us now derive the gauge Lagrangian of the translation group $`T(3)`$ in analogy to gravity . We make the most general Yang-Mills ansatz which is quadratic in the corresponding field strength $`T^a=\text{d}\vartheta ^a`$,
$`V_{}={\displaystyle \frac{1}{2}}T^aH_a.`$ (24)
The simplest choice is $`H_a=\frac{1}{\mathrm{}}{}_{}{}^{}T_{a}^{}`$ which was used by Edelen for a dislocation gauge theory. It is well-known that, in gravity, this Lagrangian does not yield Einstein’s theory . Recently, Malyshev discussed the gauge Lagrangian used by Edelen . He showed that it does not lead to the correct solutions for edge dislocations within a linear approximation. Accordingly, we will use the most general Lagrangian for an isotropic material with the 1-form
$`H_a={\displaystyle \frac{1}{\mathrm{}}}{\displaystyle {}_{}{}^{}\underset{I=1}{\overset{3}{}}}a_I^{(I)}T_a.`$ (25)
Here $`\mathrm{}`$ is the coupling constant of the theory. It is obvious that eq. (25) is a constitutive law for an isotropic material. This Lagrangian has to be invariant under local $`SO(3)`$-transformations in order to obtain the teleparallel version of the Hilbert-Einstein Lagrangian, since the Einstein theory of gravity can be viewed as the gauge theory of the translation group in four dimension. In this framework, the vierbein fields are the translational gauge potentials and the field strength is given in terms of the anholonomity .
The gauge Lagrangian (25) was used by Katanaev and Volovich to describe dislocations, but without combining it with an elastic Lagrangian. Thus their field equations do not really describe dislocations in an elastic material.
The condition of local $`SO(3)`$-invariance yields the following three parameters (see, e.g., ,),
$`a_1=1,a_2=2,a_3={\displaystyle \frac{1}{2}}.`$ (26)
Consequently
$`H_a={\displaystyle \frac{1}{\mathrm{}}}{}_{}{}^{}(^{(1)}T_a+2^{(2)}T_a+{\displaystyle \frac{1}{2}}^{(3)}T_a).`$ (27)
The three pieces $`{}_{}{}^{(I)}T_{a}^{}`$ are irreducible with respect to $`SO(3)`$. Thus the torsion or dislocation density reads $`T^a=^{(1)}T^a+^{(2)}T^a+^{(3)}T^a`$, with the number of independent components $`9=531`$, where (for notations see )
$`{}_{}{}^{(1)}T_{}^{a}`$ $`:=T^a^{(2)}T^a^{(3)}T^a`$ $`\text{(tentor)},`$ (28)
$`{}_{}{}^{(2)}T_{}^{a}`$ $`:={\displaystyle \frac{1}{2}}\vartheta ^a(e_bT^b)`$ $`\text{(trator)},`$ (29)
$`{}_{}{}^{(3)}T_{a}^{}`$ $`:={\displaystyle \frac{1}{3}}e_a(\vartheta ^bT_b)`$ (axitor). (30)
In order to obtain these three irreducible pieces, one can use the standard method of Young tableaux. In the field theory of dislocations, the axitor describes three vertical bands of screw dislocations .
With the identity ($`R^{ab}`$=Riemann-Cartan curvature 2-form, $`\stackrel{~}{R}^{ab}`$=Riemann-Christoffel curvature 2-form)
$`\stackrel{~}{R}^{ab}\eta _{ab}=R^{ab}\eta _{ab}+T^a{}_{}{}^{}(^{(1)}T_a+2^{(2)}T_a+{\displaystyle \frac{1}{2}}^{(3)}T_a)+2\text{d}(\vartheta ^a{}_{}{}^{}T_{a}^{}),`$ (31)
in the case of teleparallelism $`R^{ab}=0`$ and dropping the surface term, we find the equivalence to the Hilbert-Einstein Lagrangian $`V_{\mathrm{GR}}`$,
$`V_{}=V_{\mathrm{GR}}={\displaystyle \frac{1}{2\mathrm{}}}\stackrel{~}{R}^{ab}\eta _{ab}{\displaystyle \frac{1}{2\mathrm{}}}{}_{}{}^{}\stackrel{~}{R}.`$ (32)
Here $`\stackrel{~}{R}:=e_ae_b\stackrel{~}{R}^{ab}`$. Finally, the total Lagrangian reads
$`=V_{}+_{\mathrm{el}}.`$ (33)
Thus, for the screw dislocations in an isotropic material, one needs only one material constant. Accordingly, in this specific case, the ansatz of Edelen et al. coincides with the more general ansatz (25). Moreover, in the framework of Edelen’s dislocation theory, Osipov showed that the second order corrections to the stress field of a screw dislocation far from the core are in good agreement with the formulas given by Kröner and Seeger .
Let us now turn to the anisotropic case. The linear constitutive law for an anisotropic material reads
$`H_a={\displaystyle \frac{1}{2}}{}_{}{}^{}(\kappa _{aij}^{bkl}T_{bkl}\text{d}x^i\text{d}x^j).`$ (34)
A non-linear constitutive law, in analogy to non-linear (Born-Infeld type) electrodynamics, is given by
$`H_a={\displaystyle \frac{1}{2}}{}_{}{}^{}(\kappa _{aij}^{bkl}T_{bkl}\text{d}x^i\text{d}x^j+\stackrel{~}{\kappa }_{aij}^{bklcmn}T_{bkl}T_{cmn}\text{d}x^i\text{d}x^j+\mathrm{}),`$ (35)
where $`\kappa _{aij}^{bkl}`$ and $`\stackrel{~}{\kappa }_{aij}^{bklcmn}`$ are the constitutive functions.
The elastic stress tensors, or the corresponding stress forms, are the currents of the elastic field. Therefore, the stress 2-form is defined by
$`\mathrm{\Sigma }_a:={\displaystyle \frac{\delta _{\mathrm{el}}}{\delta \vartheta ^a}}.`$ (36)
In local components, this stress 2-form reads,
$`\mathrm{\Sigma }_a={\displaystyle \frac{1}{2}}\left(\sigma ^{kl}B_{ak}\eta _{lmn}{\displaystyle \frac{1}{2}}\sigma ^{kl}E_{kl}\eta _{amn}\right)\text{d}x^m\text{d}x^n,`$ (37)
with $`\sigma ^{kl}=C^{ijkl}E_{ij}`$. The first Piola-Kirchhoff stress tensor is included in the 1-form dual to the stress 2-form
$`{}_{}{}^{}\mathrm{\Sigma }_{a}^{}=t_{al}\text{d}x^l=\left(\sigma _l^kB_{ak}{\displaystyle \frac{1}{2}}\sigma ^{ij}E_{ij}B_{al}\right)\text{d}x^l,`$ (38)
where $`t_{al}`$ is the first Piola-Kirchhoff stress tensor. The second Piola-Kirchhoff stress tensor, derived from the first Piola-Kirchhoff stress tensor, is given by
$`t_{ac}=B_c^lt_{al}=\sigma _{ac}{\displaystyle \frac{1}{2}}\delta _{ac}\sigma ^{ij}E_{ij}.`$ (39)
Here $`t_{ac}`$ corresponds to the Maxwell tensor of elasticity for the incompatible case. The symmetric Cauchy stress tensor arises as
$`t_{ij}=B_i^at_{aj}=\sigma _{ij}{\displaystyle \frac{1}{2}}g_{ij}\sigma ^{kl}E_{kl}.`$ (40)
If dislocations are present, the stress forms are the elastic responses of the body to the dislocations.
The gauge field momentum 1-form (nowadays called excitation) is defined by
$`H_a:={\displaystyle \frac{V_{}}{T^a}},`$ (41)
i.e., it is the specific “response” quantity of the gauge Lagrangian $`V_{}`$ (and not of the elastic Lagrangian $`_{\mathrm{el}}`$) to $`T^a`$. It has the dimension of a moment stress. Thus, we interpret $`H_a`$ as a moment or couple stress originating in the dislocation core. The moment stress 1-form, in local components for an anisotropic linear material, reads
$`H_a=H_{an}\text{d}x^n=\kappa _a^{ijbkl}T_{bkl}\eta _{ijn}\text{d}x^n.`$ (42)
On the other hand, the translation gauge current or the stress 2-form of the gauge field,
$`h_a:={\displaystyle \frac{V_{}}{\vartheta ^a}}=e_aV_{}+(e_aT^b)H_b,`$ (43)
explicitly reads
$`h_a={\displaystyle \frac{1}{2}}\left[(e_aT^b\right)H_bT^b(e_aH_b)].`$ (44)
Whereas $`\mathrm{\Sigma }_a`$ is the stress 2-form of the elastic material, $`h_a`$ is the response stress 2-form of the dislocations or the plastic stress 2-form. Accordingly, $`h_a`$ is the internal stress caused by dislocations.
The field equations are derived by varying the total Lagrangian with respect to the elastic field $`\xi ^a`$ and the gauge potential $`\varphi ^a`$:
$`{\displaystyle \frac{\delta }{\delta \xi ^a}}`$ $`\text{d}\mathrm{\Sigma }_a+\text{d}h_a=0,`$ (45)
$`{\displaystyle \frac{\delta }{\delta \varphi ^a}}`$ $`\text{d}H_ah_a=\mathrm{\Sigma }_a.`$ (46)
Eq. (46) is the Yang-Mills type gauge field equation in the theory of elastoplasticity. Both stress 2-forms, $`\mathrm{\Sigma }_a`$ and $`h_a`$, are sources in eq. (46). Eq. (45) is a consequence of (46) and the Poincaré lemma $`\text{d}\text{d}=0`$. The Euler-Lagrange equations can be interpreted as equilibrium equations. Eq. (45) describes the force and eq. (46) the moment equilibrium. The fields to be determined from the field equations are $`\xi ^a`$ or $`u^a`$ and $`\varphi ^a`$. Due to the nonlinear geometrical character of elastoplasticity, the field equations are represented by a coupled system of nonlinear partial differential equations. In the framework of MAG, eq. (45) is the matter field and eq. (46) the first gauge field equation (see ).
The force density 3-forms are defined as follows: The Peach-Koehler force (elastic force acting on a dislocation), which is analogously defined as the Lorentz force in Maxwell’s theory, reads
$`f_a^{\mathrm{el}}:=\text{d}\mathrm{\Sigma }_a=(e_aT^b)\mathrm{\Sigma }_b.`$ (47)
The response force to dislocations is
$`f_a^\mathrm{g}:=\text{d}h_a=(e_aT^b)h_b.`$ (48)
Eq. (45) describes the force equilibrium between the Peach-Koehler force and the dislocation-response force
$`0=f_a=f_a^\mathrm{g}+f_a^{\mathrm{el}}=\text{d}h_a+\text{d}\mathrm{\Sigma }_a,`$ (49)
or is, equivalently, interpreted as the equilibrium condition of an elastic body containing dislocations. Both these forces are configurational forces. It is to be emphasized that (47) and (48), in the framework of field theory, are the first Noether identities for the elastic and the gauge fields.
For vanishing dislocation density ($`\varphi ^a=0`$ and $`T^a=0`$), the field equations reduce to the equilibrium condition of the classical elasticity theory under zero external forces
$`\text{d}\mathrm{\Sigma }_a=0^it_{ai}=0.`$ (50)
## V Relations to other descriptions of dislocation theory
In an alternative description of teleparallelism one replaces $`T^a\text{d}\vartheta ^a`$ by the Levi-Civita connection $`\stackrel{~}{\omega }_b^a`$ of the metric (Cauchy-Green tensor) $`G=\delta _{ab}\vartheta ^a\vartheta ^b`$. One applies Cartan’s first structure equation
$`\text{d}\vartheta ^a=\stackrel{~}{\omega }_b^a\vartheta ^b,`$ (51)
which yields
$`\stackrel{~}{\omega }_{ab}={\displaystyle \frac{1}{2}}(e_aT_be_bT_a(e_ae_bT_c)\vartheta ^c).`$ (52)
The corresponding Riemannian curvature 2-form reads
$`\stackrel{~}{R}_{ab}=\text{d}\stackrel{~}{\omega }_{ab}+\stackrel{~}{\omega }_{ac}\stackrel{~}{\omega }_b^c.`$ (53)
Eventually, we get the corresponding field equation by substituting eq. (53) into the Hilbert-Einstein Lagrangian $`V_{\mathrm{GR}}`$ of eq. (32). After variation, one recovers the Einstein type field equation
$`{\displaystyle \frac{1}{2}}\eta _{abc}\stackrel{~}{R}^{bc}=\mathrm{}\mathrm{\Sigma }_a`$ (54)
used by Malyshev for dislocation theory. A disadvantage of this description is that the original structures are blurred and that the field equation does not have the Yang-Mills form.
Now, we can completely translate our formulas of the gauge theoretical description of dislocation theory into the coordinate system of the external observer (holonomic coordinates). The linear connection of the external observer is defined as pure gauge,
$`\omega _j^i=B_a^i\text{d}B_j^a=B_a^i_kB_j^a\text{d}x^k.`$ (55)
In holonomic coordinates, the torsion is given by
$`T^i=B_a^i_{[j}B_{k]}^a\text{d}x^j\text{d}x^k=\omega _k^i\text{d}x^k.`$ (56)
Cartan’s torsion tensor $`T_{jk}^i`$ is the antisymmetric part<sup>§</sup><sup>§</sup>§We will be using the notation $`A_{[ij]}\frac{1}{2}(A_{ij}A_{ji})`$. of the components of the connection $`\omega _{jk}^i`$, namely $`\frac{1}{2}T_{jk}^i\omega _{[jk]}^i`$. By means of the torsion (56) and the condition of vanishing nonmetricity,
$`\text{D}g_{ij}\text{d}g_{ij}\omega _i^kg_{jk}\omega _j^kg_{ik}=0.`$ (57)
The connection $`\omega _j^i`$ can be decomposed into the Levi-Civita connection $`\stackrel{~}{\omega }_j^i`$ and the contortion $`\tau _j^i`$, which is a tensorial 1-form of type ad,
$$\omega _j^i=\stackrel{~}{\omega }_j^i\tau _j^i.$$
(58)
The local components of the Levi-Civita connection and of the contortion read, respectively,
$$\stackrel{~}{\omega }_j^i=\frac{1}{2}g^{il}\left(_jg_{lk}+_kg_{lj}_lg_{jk}\right)\text{d}x^k$$
(59)
and
$$\tau _j^i=\frac{1}{2}\left(T_{jk}^i+T_{jk}^iT_{kj}^i\right)\text{d}x^k.$$
(60)
We can resolve (60) with respect to the torsion,
$$\tau _j^i\text{d}x^j=\frac{1}{2}T_{jk}^i\text{d}x^j\text{d}x^k=T^i.$$
(61)
The forces are
$`f_i^{\mathrm{el}}=\text{D}\mathrm{\Sigma }_i\text{d}\mathrm{\Sigma }_i+\omega _i^j\mathrm{\Sigma }_j=T_i^j\mathrm{\Sigma }_j`$ (62)
and
$`f_i^\mathrm{g}=\text{D}h_i=T_i^jh_j.`$ (63)
In holonomic coordinates, the field equations read
$`\text{D}\mathrm{\Sigma }_i+\text{D}h_i=0\text{(force equilibrium)},`$ (64)
$`\text{D}H_i=h_i+\mathrm{\Sigma }_i\text{(moment equilibrium)}.`$ (65)
Using the condition of teleparallelism $`R^{ij}=0`$, in linear approximation we obtain Kröner’s incompatibility equation
$`\mathrm{inc}E\times E\times =\eta .`$ (66)
The symmetric second rank tensor $`\eta `$ is called incompatibility tensor; it encompasses the dislocation density. With $`H_i=0`$, we recover a dislocation theory without moment stress which is in agreement with the dislocation theory given by Kröner and Seeger .
## VI Conclusion
We have proposed a dislocation gauge theory in an elastoplastic material. The basic equations (12), (45), (46) and (49) of dislocation gauge theory can be summarized in axiomatic way in analogy to the Maxwell theory (for an axiomatic formulation of Maxwell’s theory, see, e.g., ). As soon as the Lagrangian is specified, one can find the basic laws. As gauge Lagrangians we use the teleparallel one, which is equivalent to the Hilbert-Einstein Lagrangian, and some Lagrangians for anisotropic constitutive laws. For linear constitutive laws, the total Lagrangian has the following symbolic form
$`(\text{strain})^2+(\text{dislocation density})^2.`$ (67)
The first law expresses the force equilibrium
$`\text{d}\mathrm{\Sigma }_a^\mathrm{T}=0\text{with}\mathrm{\Sigma }_a^\mathrm{T}=\mathrm{\Sigma }_a+h_a.`$ (68)
A consequence of eq. (68), in analogy to the inhomogeneous Maxwell equation, is the inhomogeneous Yang-Mills equation
$`\text{d}H_a=\mathrm{\Sigma }_a^\mathrm{T}\text{(moment equilibrium)}.`$ (69)
The definition of the elastoplastic force is the second law
$`f_a=(e_aT^b)\mathrm{\Sigma }_b^\mathrm{T}.`$ (70)
The conservation law of dislocation density (homogeneous Yang-Mills equation or Bianchi identity) is the third law
$`\text{d}T^a=0.`$ (71)
The constitutive laws $`\sigma E`$ and $`HT`$, which are the physical input from experimental data, are the fourth law.
We have compared our proposal with Kröner’s geometric theory of dislocation. Our gauge theoretical formulation of dislocation theory includes Kröner’s basic equations and is thus a straightforward description of dislocation theory with moment stress as given previously in .
Kröner , Stojanović , and, later, Kleinert have introduced the concept of a double gauge theory of dislocations, that is the stress tensor can be considered as an Einstein tensor of a formal stress space with torsion. From the geometrical and field theoretical point of view, there is no need to interpret the stress as an Einstein tensor. The generalization of the stress tensor is the energy-momentum tensor. The stress tensor is nothing else but the source of the Einstein tensor, and the equilibrium condition
$`\text{div}\sigma =0`$ (72)
is the first Noether identity. Therefore, we did not make use of the concepts of formal stress space and strain space.
Now some remarks are in place with respect to the moment stress caused by disclinations. The generalization of (material) moment stress is the hypermomentum $`\mathrm{\Delta }_{ab}=\delta _{\mathrm{el}}/\delta \omega ^{ab}`$ which contains the spin current . The hypermomentum is the source of the $`GL(n,)`$-gauge field in MAG and it does not appear in a pure translation gauge theory. Therefore, there are no degrees of freedom for describing spin-disclination in a $`T(3)`$-gauge theory, in agreement with Kröner . Spin-disclinations are defects in materials with microstructure such as liquid crystals or magnetic spin systems, where spin-moment stress occurs. A gauge theory of materials with microstructure was given by Lagoudas .
The goal of this paper was the formulation of a dislocation theory as a gauge theory in analogy to gravity. We combined the physical ideas of Kröner’s geometrical theory with the framework of MAG. We found a dislocation theory with moment stress $`H_a`$ and two kinds of force stresses $`\mathrm{\Sigma }_a`$ and $`h_a`$. In this picture, the dislocation is a source of nontrivial torsion in a Weitzenböck space or of Riemannian curvature in a Riemann space, respectively. The elastic material plays the role of a kind of an (an)isotropic “ether” in analogy to the vacuum in gravity theory. The metric (Cauchy-Green strain tensor) $`g_{ij}`$ is an effective quantity determined by $`\varphi ^a`$ and $`\text{d}\xi ^a`$. However, it is not a gauge potential.
###### Acknowledgements.
The author is grateful to Profs. Friedrich W. Hehl, Ekkehart Kröner, Luciano Mistura, Alfred Seeger, and Hans-Rainer Trebin and to Dr. Gerald Wagner for many helpful discussions, furthermore to Prof. Friedrich W. Hehl for useful comments on an earlier version of this paper. M.L. acknowledges the support by the Max-Planck-Institut für Metallforschung and the Institut für Theoretische & Angewandte Physik, University of Stuttgart during his stay in Stuttgart.
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# Wavelets as basis functions in canonical quantization
## I INTRODUCTION
Canonical quantization of electromagnetic field is traditionally done using plane waves. The field is enclosed into a cubic cavity and the field operators are expanded using eigenfunctions of the cavity, i.e., plane waves. After quantization it is possible to take the limit of infinite cavity and eliminate the unphysical finite size cavity. The inconvenience of the use of plane waves is that they are delocalized in real space. This makes it difficult to formulate for example photodetection theories. Photodetectors measure the field locally and are sensitive over a finite bandwidth of frequencies. It is possible to use any complete set of basis functions in the quantization , so a formulation with a more localized basis functions would be desirable.
The theory of multiresolution analysis (MRA) has been under intensive study during the recent years . A complete set of basis functions in MRA are called wavelets. Wavelets are localized in real as well as in Fourier space. They are parameterized by scale and translation parameters which both get integer values. The translation parameter translates the wavelet in real space. For large negative values of the scale parameter the wavelet is wide and for large scale parameter values narrow. Typical wavelets are orthogonal with respect to both indices. There are many wavelets with different characteristics. Differences include how well wavelets are localized, what kind of Fourier transforms they have, whether they are real or complex and how symmetric they are. It is interesting that some wavelets have compact support, i.e., they are zero outside a certain finite length interval. These kind of wavelets do not have analytical expressions. The main theory of multiresolution analysis is the same for all useful wavelets.
In this paper we show how wavelets can be used as basis functions in the canonical quantization. Real and orthonormal wavelets are used. New mode functions and operators are linear transforms of plane waves and the corresponding operators. Different vector valued mode functions for electric and magnetic fields are obtained. New creation and annihilation operators are the same for both fields and satisfy bosonic commutation relations. This means that formalism of the new operators remains the same and makes it easy to use the new basis. New mode functions are localized and have similar properties as wavelets, for examble they are parameterized by scale and translation parameters.
In Sec. II we give a short introduction to the theory of multiresolution analysis, scaling functions and wavelets. In Sec. III we derive equations for field operators in wavelet basis. In Sec. IV the theory developed in earlier sections is applied to some simple quantum mechanical simulations and the wavelet and plane wave bases are compared. Finally in Sec. V we give our conclusions and suggest several generalizations of the theory developed in this paper.
## II Basic properties of scaling functions and wavelets
### A Multiresolution analysis and wavelets
In the following we give a brief introduction to the basic properties of wavelets. The discussion follows books . We start with the scaling function $`\varphi _l(t)=\varphi (tl),lZ,\varphi L^2`$, which spans the function space $`A_0`$
$$f(t)=\underset{l}{}a_l\varphi _l(t),f(t)A_0.$$
(1)
Introducing a new index $`s`$ by the formula
$$\varphi _{s,l}(t)=2^{s/2}\varphi (2^stl),s,lZ$$
(2)
we can define different function spaces $`A_s`$. If $`s`$ is large, $`s0`$, the scaling function is narrow and peaked around some center point. Large negative values $`s0`$ gives a function which is wide. The parameter $`k`$ translates the scaling function in a given scale by a factor $`2^s`$.
In order to the scaling function to satisfy a multiresolution analysis (MRA) the function space $`A_{s+1}`$ must include the function space with a lower index $`A_s`$
$$A_sA_{s+1},sZ$$
(3)
When the parameter $`s`$ approaches infinity we get a space of square integrable functions $`A_{\mathrm{}}=L^2`$ and when $`s`$ approaches minus infinity the result is a null space $`A_{\mathrm{}}=\{0\}`$. Because of multiresolution analysis, it is possible to expand every function in a function space $`A_s`$ using basis functions of function space $`A_{s+1}`$. Specifically we can expand scaling functions $`\varphi _{s,l}(t)`$ in $`A_s`$ using basis functions of $`A_{s+1}`$. This gives
$$\varphi (t)=\underset{n}{}h(n)\sqrt{2}\varphi (2tn),nZ$$
(4)
with some coefficients $`h(n)`$. We have chosen $`s=0`$ and used Eq. (2).
The wavelet space is defined to be a difference space between different function spaces spanned by the scaling functions. We define $`W_s`$ as
$$A_{s+1}=A_sW_s.$$
(5)
Intuitively $`W_s`$ contains functions which must be added to $`A_s`$ in order to get $`A_{s+1}`$. The basis functions in the function space $`W_s`$ are called wavelets $`\psi (t)`$. All functions in $`W_s`$ can be expanded using translations of a fundamental wavelet, which are obtained using the same equation as for scaling functions, Eq. (2). Because $`W_sA_{s+1}`$, the wavelet function in $`W_s`$ can be expanded using scaling functions in $`A_{s+1}`$. We get
$$\psi (t)=\underset{n}{}h_1(n)\sqrt{2}\varphi (2tn),nZ,$$
(6)
which corresponds to the equation (4) for $`\varphi `$. The coefficients $`h(n)`$ and $`h_1(n)`$ are not independent. They satisfy the relation
$$h_1(n)=(1)^nh(1n).$$
(7)
Using Eq. (3) in Eq. (5) we can decompose the right hand side into subspaces with lower indices. After one iteration we get
$$A_{s+2}=A_{s+1}W_{s+1}=A_sW_sW_{s+1}.$$
(8)
We can do the iteration repeatedly and in the infinite limit we get
$$L^2=A_sW_sW_{s+1}\mathrm{}$$
(9)
where $`A_{\mathrm{}}=L^2`$ has been used. This means that every function which belongs to $`L^2`$ can be expanded as
$`f(t)={\displaystyle \underset{l=\mathrm{}}{\overset{\mathrm{}}{}}}c_{s,l}2^{s/2}\varphi (2^stl)`$ (10)
$`+`$ $`{\displaystyle \underset{s^{}=s}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=\mathrm{}}{\overset{\mathrm{}}{}}}d_{s^{},l}2^{s^{}/2}\psi (2^s^{}tl).`$ (11)
The parameter $`s`$ is totally arbitrary. We can take a limit $`s\mathrm{}`$ and because $`A_{\mathrm{}}=\{0\}`$ the scaling function space can be eliminated giving
$`f(t)`$ $`=`$ $`{\displaystyle \underset{s=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=\mathrm{}}{\overset{\mathrm{}}{}}}d_{s,l}\psi _{s,l}(t)`$ (12)
$`=`$ $`{\displaystyle \underset{s=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=\mathrm{}}{\overset{\mathrm{}}{}}}d_{s,l}2^{s/2}\psi (2^stl),`$ (13)
i.e., any $`L^2`$-function can be expanded using translated and scaled wavelets. In this paper we use only real orthonormal wavelets which have the property
$$\underset{\mathrm{}}{\overset{\mathrm{}}{}}\psi _{sl}(t)\psi _{s^{}l^{}}(t)𝑑t=\delta _{ss^{}}\delta _{ll^{}}.$$
(14)
Multiplying Eq. (12) by parts with $`\psi _{s^{}l^{}}(t)`$ and integrating $`\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑t`$ we get the coefficients
$$d_{s,l}=\underset{\mathrm{}}{\overset{\mathrm{}}{}}f(t)\psi _{s,l}(t)𝑑t.$$
(15)
It follows from the multiresolution analysis that the integral over the wavelet gives zero, i.e.
$$\underset{\mathrm{}}{\overset{\mathrm{}}{}}\psi (t)𝑑t=0.$$
(16)
This means that wavelets must have some kind of oscillating structure as the name suggests.
So far we have not defined any wavelets explicitly. There are several methods to construct wavelets for different purposes. One method is to divide the Fourier space in such a way that the resulting wavelets fulfill multiresolution analysis requirements and are orthogonal. The wavelets obtained using this method can be well localized but do not have a compact support. Another method is to derive wavelets based on the filter coefficients. The filter coefficients of typical wavelets satisfy the fundamental condition
$$\underset{n}{}h(2n)=\underset{n}{}h(2n+1)=\frac{1}{\sqrt{2}}$$
(17)
and are orthogonal
$$\underset{n}{}h(n)h(n2k)=\delta _{k0}.$$
(18)
If the length of the filter coefficient sequence is long enough these two conditions do not determine the coefficients uniquely. The remaining degrees of freedom can be used to give additional desirable properties to wavelets. Typical choices are to demand the wavelet or scaling function to be as smooth as possible. These kind of wavelets can have a compact support, i.e., they are zero outside a specific region.
### B A few examples of wavelets
One of the best known scaling function is a sinc-function
$$\varphi (x)=\mathrm{sinc}(x)=\frac{\mathrm{sin}(\pi x)}{\pi x}$$
(19)
with $`\mathrm{sinc}(0)=1`$. The corresponding wavelet is $`\psi (x)=2\varphi (2x)\varphi (x)`$. Figure 1 shows the scaling function and the wavelet which are both localized around $`x=0`$. Both functions are not so well localized as would be desirable for many practical applications. The Fourier transforms of the scaling and wavelet functions are
$`\stackrel{~}{\varphi }(\omega )=\{\begin{array}{cc}\frac{1}{\sqrt{2\pi }},\hfill & \omega <|\pi |\hfill \\ 0,\hfill & \text{otherwise}\hfill \end{array}`$ (22)
$`\stackrel{~}{\psi }(\omega )=\{\begin{array}{cc}\frac{1}{\sqrt{2\pi }},\hfill & \pi <|\omega |<2\pi \hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}`$ (25)
The reason why oscillations do not decay rapidly is because of the abrupt changes of the Fourier transforms at $`|\omega |=\pi `$ and $`|\omega |=2\pi `$. One specialty of the sinc-family of wavelets is that the division of frequency space to different scales is orthogonal, i.e., there is no overlap of the Fourier-transforms of wavelets with different scale parameters $`s`$. Because of poor localization in real space, Shannon wavelets are rarely used.
It is possible to smoothen the change in Fourier transforms in such a way that orthogonality and multiresolution analysis requirements are preserved. The resulting wavelets are called Meyer-wavelets. There are several different Meyer scaling and wavelet functions depending on the smoothening function. The Meyer scaling and wavelet functions used in this paper are shown in Fig. 2. They are clearly more localized than the corresponding Shannon functions. Absolute values of the Fourier transforms of scaling and wavelet functions are shown in Fig. 3. The Fourier transform of a scaling function $`\varphi `$ is flat around $`\omega =0`$. Around frequencies $`\omega =\pm \pi `$ the Fourier transform decays to zero and is zero at larger values. The smoothing compared to Shannon case is clearly seen in both functions. For wavelets two different scales are shown. As parameter $`s`$ increases the transform is wider and shifted further away from the origin.
The length of filter sequences $`h(n)`$ and $`h_1(n)`$ in Eq. (4) and (6) for sinc and Meyer scaling fuctions and wavelets are infinite. As a result of this wavelets do not have a compact support but are nonzero, although possibly very small, over the whole $`x`$-axis. It is possible to construct wavelets which have compact support, i.e., they are exactly zero, not just exponentially small, outside a specific region. The length of the filter sequence for these kind of wavelets is finite and the corresponding wavelets and scaling functions do not have analytical forms. One group of these wavelets are the Daubechies wavelets. In this group of wavelets there are several wavelet families which can be characterized by the lengths of filter sequences. As the length of the sequence increases the wavelets and scaling functions become smoother. Figure 4 shows Daubechies scaling and wavelet functions with a filter coefficient length $`N=6`$ in Eq. (2) and (6). Both functions are nonzero in the interval $`0x5`$. They are continuous and derivable. The functions are not symmetric. The general behavior is that the scaling function is concentrated to the left of the nonzero interval and the wavelet is significantly nonzero in the middle.
### C Wavelets in higher dimensions
The discussion above has been about one dimensional wavelets. We want to use wavelets as basis functions in canonical quantization so three dimensional wavelets must be used. One method to construct multidimensional wavelets is to use products of one dimensional scaling and wavelet functions . The following products are all wavelets in two dimensions
$`\psi ^1(𝐫)`$ $`=`$ $`\varphi (x)\psi (y)`$ (26)
$`\psi ^2(𝐫)`$ $`=`$ $`\psi (x)\varphi (y)`$ (27)
$`\psi ^3(𝐫)`$ $`=`$ $`\psi (x)\psi (y).`$ (28)
The integral over all functions above vanishes as is required for wavelets. The two dimensional scaling function is a product of two one dimensional scaling functions
$$\varphi (𝐫)=\varphi (x)\varphi (y).$$
(29)
The wavelets with different scaling and translation parameters become
$$\psi _{s𝐥}^i(𝐫)=\psi _{s0}^i(𝐫2^s𝐥)=2^s\psi _{00}^i(2^s𝐫𝐥).$$
(30)
The translation parameter $`𝐥`$ is now a vector with two components which give the translation in $`x`$ and $`y`$ directions. The scaling parameter is a scalar, so scaling is the same in both directions. Note that the scaling of the wavelet is different compared to the one dimensional case. Any two-dimensional $`L^2`$-function can be expanded as
$$f(𝐫)=\underset{s𝐥i}{}d_{s𝐥}^i\psi _{s𝐥}^i(𝐫)\underset{s=\mathrm{}}{\overset{\mathrm{}}{}}\underset{l_x=\mathrm{}}{\overset{\mathrm{}}{}}\underset{l_y=\mathrm{}}{\overset{\mathrm{}}{}}\underset{i=1}{\overset{3}{}}d_{s𝐥}^i\psi _{s𝐥}^i(𝐫).$$
(31)
Shannon and Meyer wavelets were constructed by dividing the Fourier space. For these wavelets the three different wavelets in two dimensions, obtained by multiplying one dimensional scaling functions and wavelets, are significantly nonzero in different regions of $`𝐤`$-space. Figure 5 shows nonzero regions for three different wavelets. Wavelets $`\psi ^1(𝐫)`$ and $`\psi ^2(𝐫)`$ are nonzero in regions $`|k_x|k_{min}`$, $`k_{min}|k_y|k_{max}`$ and $`|k_y|k_{min}`$, $`k_{min}|k_x|k_{max}`$ respectively. Wavelet $`\psi ^3(𝐫)`$ is nonzero in regions $`k_{min}|k_x|k_{max}`$, $`k_{min}|k_y|k_{max}`$. With a smaller scaling index $`s`$, the scaling function in the center is split into similar regions with different wavelets. With larger parameter $`s`$, larger $`𝐤`$-values are divided to regions. The division of frequency space is exact only for sinc-wavelets. For all other wavelets the Fourier transform has nonzero values also outside of its main region. How well the Fourier transform is localized into the regions depends on the type of a wavelet used. For wavelets with compact support and small filter sequence length the division described above is not very well.
Three dimensional wavelets can be constructed in a similar way as two dimensional ones. We get one scaling function and seven wavelet functions which are products of one dimensional scaling and wavelet functions
$`\varphi (𝐫)`$ $`=`$ $`\varphi (x)\varphi (y)\varphi (z)\psi ^1(𝐫)=\varphi (x)\varphi (y)\psi (z)`$ (32)
$`\psi ^2(𝐫)`$ $`=`$ $`\varphi (x)\psi (y)\varphi (z)\psi ^3(𝐫)=\varphi (x)\psi (y)\psi (z)`$ (33)
$`\psi ^4(𝐫)`$ $`=`$ $`\psi (x)\varphi (y)\varphi (z)\psi ^5(𝐫)=\psi (x)\varphi (y)\psi (z)`$ (34)
$`\psi ^6(𝐫)`$ $`=`$ $`\psi (x)\psi (y)\varphi (z)\psi ^7(𝐫)=\psi (x)\psi (y)\psi (z).`$ (35)
The division of $`𝐤`$-space is a direct generalization of the two dimensional case. The scaling relation becomes
$$\psi _{s𝐥}^i(𝐫)=\psi _{s0}^i(𝐫2^s𝐥)=2^{3s/2}\psi _{00}^i(2^s𝐫𝐥).$$
(36)
The expansion of $`L^2`$-function is given by Eq. (31) with the difference that the $`𝐥`$-vector is now three dimensional and index $`i`$ gets values from one to seven.
### D Expansion of a plane wave using wavelets
In the following section we need expansions of plane waves using wavelets. The normalized plane wave in a box of side $`L`$ can be expanded as
$$\frac{1}{L^{3/2}}e^{i𝐤𝐫}=\underset{s𝐥i}{}d_{𝐤,s𝐥}^i\psi _{s𝐥}^i(𝐫).$$
(37)
Multiplying by parts with $`\psi _{s^{}𝐥^{}}^i^{}(𝐫)`$ and using orthogonality of wavelets gives the coefficients
$$d_{𝐤,s𝐥}^i=\frac{1}{L^{3/2}}\psi _{s𝐥}^i(𝐫)e^{i𝐤𝐫}d^3𝐫=d_{k_x,sl_x}^{i_x}d_{k_y,sl_y}^{i_y}d_{k_z,sl_z}^{i_z}.$$
(38)
The factorization is possible because multidimensional wavelets are products of one dimensional scaling functions and wavelets. The index $`i_{x,y,z}`$ is either $`0`$ or $`1`$ denoting scaling and wavelet functions respectively. Similar factorization is possible in Eq. (37). Using equation (36) we get after the change of integration variable
$$d_{𝐤,s𝐥}^i=\mathrm{exp}(i2^s𝐤𝐥)d_{𝐤,s0}^i=2^{3s/2}\mathrm{exp}(i2^s𝐤𝐥)d_{2^s𝐤,00}^i.$$
(39)
Thus every coefficient $`d_{𝐤,s𝐥}^i`$ can be calculated using the translated and scaled Fourier transform of a fundamental wavelet. The scaling parameter $`s`$ changes also the absolute value of the coefficients whereas the translation parameter $`𝐥`$ gives only a phase shift. Because of orthogonality of plane waves and wavelets we get the following relations
$`{\displaystyle \underset{s𝐥i}{}}d_{𝐤,s𝐥}^id_{𝐤^{},s𝐥}^i=\delta _{\mathrm{𝐤𝐤}^{}}`$ (40)
$`{\displaystyle \underset{𝐤}{}}d_{𝐤,s𝐥}^id_{𝐤,s^{}𝐥^{}}^i^{}=\delta _{ss^{}}\delta _{\mathrm{𝐥𝐥}^{}}\delta _{ii^{}}.`$ (41)
Next we study the parameters $`d_{𝐤,s𝐥}^i`$ more closely. We restrict the discussion to one dimension, i.e., the parameters $`𝐤`$ and $`𝐥`$ are one dimensional. Figure 6 shows the absolute values of the coefficients with several different $`k`$-values for Meyer wavelets. Note that according to Eq. (38) the translation parameter gives only a phase shift. The coefficients of a plane wave are shown with three different $`k`$-values. In all cases coefficients with only a few scale parameters are nonzero. For small $`k`$-values $`k=0.39`$, two scales with parameters $`s<0`$ have nonzero values. Remember that the Fourier transform of the fundamental ($`s=0`$) Meyer wavelet is mainly nonzero in the interval $`\pi <k<2\pi `$. The $`k`$-value $`k=4.71=1.5\pi `$ is carefully chosen to be in the middle of the scale $`s=0`$ interval. It is seen that coefficients for all other $`s`$-parameters are zero. With large $`k`$-values the peaks are centered at a higher scale as can be seen from the last peak. Peaks with large scale parameter are smaller than peaks at small scale. This is because of the factor $`2^{3s/2}`$ in equation (39). One could think that at large scale the peaks are lower because there are more $`l`$-values which share the contribution from the single plane wave. Remember that at high scale values the shift of a wavelet is smaller when the translation parameter is changed by one unit.
## III Canonical quantization using wavelets
Traditionally canonical quantization of field is done using plane waves as basis functions . The field is enclosed to a cubic of side $`L`$. All functions can be expanded using the eigenfunctions of the quantization volume, i.e., plane waves. The quantization itself is done by introducing bosonic operators for every basis function. The vector potential of field can be expanded
$$\widehat{𝐀}=\frac{1}{L^{3/2}}\underset{𝐤\sigma }{}\left(\frac{\mathrm{}}{2\omega _𝐤ϵ_0}\right)^{1/2}(\widehat{a}_{𝐤\sigma }ϵ_{𝐤\sigma }e^{i𝐤𝐫}+\mathrm{h}.\mathrm{c}.).$$
(42)
The vector $`ϵ_{𝐤\sigma }`$ gives the polarization of the plane wave with wave vector $`𝐤`$ and polarization $`\sigma `$. Operators $`\widehat{a}_{𝐤\sigma }`$ and $`\widehat{a}_{𝐤\sigma }^{}`$ are bosonic annihilation and creation operators. The parameter $`𝐤`$ gets values
$$𝐤=(\frac{2\pi n_1}{L},\frac{2\pi n_2}{L},\frac{2\pi n_3}{L}),n_{1,2,3}=0,\pm 1,\pm 2\mathrm{}$$
(43)
For electric and magnetic fields we get
$`\widehat{𝐄}(𝐫)`$ $`=`$ $`{\displaystyle \frac{\widehat{𝐀}(𝐫)}{t}}`$ (44)
$`=`$ $`{\displaystyle \frac{1}{L^{3/2}}}{\displaystyle \underset{𝐤\sigma }{}}\left({\displaystyle \frac{\mathrm{}\omega _𝐤}{2ϵ_0}}\right)^{1/2}(i\widehat{a}_{𝐤\sigma }ϵ_{𝐤\sigma }e^{i𝐤𝐫}+\mathrm{h}.\mathrm{c}.)`$ (45)
$`\widehat{𝐁}(𝐫)`$ $`=`$ $`\times \widehat{𝐀}(𝐫)`$ (46)
$`=`$ $`{\displaystyle \frac{1}{L^{3/2}}}{\displaystyle \underset{𝐤\sigma }{}}\left({\displaystyle \frac{\mathrm{}}{2ϵ_0\omega _𝐤}}\right)^{1/2}(i\widehat{a}_{𝐤\sigma }(𝐤\times ϵ_{𝐤\sigma })e^{i𝐤𝐫}+\mathrm{h}.\mathrm{c}.).`$ (47)
Field operators above are in plane wave basis, which is parameterized by $`𝐤`$-vectors.
Next we change the basis from plane waves to wavelets. We could just expand the plane waves using wavelets and use this expansion in Eqs. (44) and (46). However this approach would not allow the separation of mode functions and operators. Therefore we proceed by inserting a unity in the form $`\underset{𝐤}{}\delta _{\mathrm{𝐤𝐤}^{}}`$ to expansions (44) and (46). The sum is now divided in such a way that all other $`𝐤`$-values except the operator $`𝐤`$ are changed to $`𝐤^{}`$. After the use of Eq. (40) for the delta function we get
$`\widehat{𝐄}(𝐫)`$ $`=`$ $`{\displaystyle \frac{1}{L^{3/2}}}{\displaystyle \underset{𝐤\sigma }{}}\left({\displaystyle \frac{\mathrm{}\omega _𝐤}{2ϵ_0}}\right)^{1/2}(i\widehat{a}_{𝐤\sigma }ϵ_{𝐤\sigma }e^{i𝐤𝐫}+\mathrm{h}.\mathrm{c}.)`$ (48)
$`=`$ $`{\displaystyle \frac{1}{L^{3/2}}}{\displaystyle \underset{s𝐥i}{}}{\displaystyle \underset{\mathrm{𝐤𝐤}^{}\sigma }{}}\left({\displaystyle \frac{\mathrm{}\omega _𝐤^{}}{2ϵ_0}}\right)^{1/2}(i\widehat{a}_{𝐤\sigma }ϵ_{𝐤^{}\sigma }d_{𝐤^{},s𝐥}^id_{𝐤,s𝐥}^ie^{i𝐤^{}𝐫}+\mathrm{h}.\mathrm{c}.)`$ (49)
$`=`$ $`{\displaystyle \underset{s𝐥i}{}}{\displaystyle \underset{\sigma }{}}\left({\displaystyle \frac{i}{L^{3/2}}}{\displaystyle \underset{𝐤^{}}{}}\left({\displaystyle \frac{\mathrm{}\omega _𝐤^{}}{2ϵ_0}}\right)^{1/2}ϵ_{𝐤^{}\sigma }d_{𝐤^{},s𝐥}^ie^{i𝐤^{}𝐫}\right)\left({\displaystyle \underset{𝐤}{}}d_{𝐤,s𝐥}^i\widehat{a}_{𝐤\sigma }\right)+\mathrm{h}.\mathrm{c}.`$ (50)
$`=`$ $`{\displaystyle \underset{s𝐥i}{}}{\displaystyle \underset{\sigma }{}}(\widehat{b}_{s𝐥,\sigma }^i𝐮_{s𝐥,\sigma }^{iE}(𝐫)+\mathrm{h}.\mathrm{c}.),`$ (51)
where the new basis functions and operators are
$`𝐮_{s𝐥,\sigma }^{iE}(𝐫)`$ $`=`$ $`{\displaystyle \frac{i}{L^{3/2}}}{\displaystyle \underset{𝐤}{}}\left({\displaystyle \frac{\mathrm{}\omega _𝐤}{2ϵ_0}}\right)^{1/2}d_{𝐤,s𝐥}^iϵ_{𝐤\sigma }e^{i𝐤𝐫}`$ (52)
$`=`$ $`𝐮_{s0}^{iE}(𝐫2^s𝐥)=2^{2s}𝐮_{00,\sigma }^{iE}(2^s𝐫𝐥)`$ (53)
$`\widehat{b}_{s𝐥,\sigma }^i`$ $`=`$ $`{\displaystyle \underset{𝐤}{}}d_{𝐤,s𝐥}^i\widehat{a}_{𝐤\sigma }.`$ (54)
Similarly for the magnetic field we get
$$\widehat{𝐁}(𝐫)=\underset{s𝐥i}{}\underset{\sigma }{}(\widehat{b}_{s𝐥,\sigma }^i𝐮_{s𝐥,\sigma }^{iB}(𝐫)+\mathrm{h}.\mathrm{c}.),$$
(55)
where the operator $`\widehat{b}_{s𝐥,\sigma }^i`$ is given by equation (54) and
$`𝐮_{s𝐥,\sigma }^{iB}(𝐫)`$ $`=`$ $`{\displaystyle \frac{i}{L^{3/2}}}{\displaystyle \underset{𝐤}{}}\left({\displaystyle \frac{\mathrm{}}{2ϵ_0\omega _𝐤}}\right)^{1/2}d_{𝐤,s𝐥}^i(𝐤\times ϵ_{𝐤\sigma })e^{i𝐤𝐫}`$ (56)
$`=`$ $`𝐮_{s0}^{iB}(𝐫2^s𝐥)=2^{2s}𝐮_{00,\sigma }^{iB}(2^s𝐫𝐥).`$ (57)
The new mode functions behave in the same way as wavelets when the indices are changed. The translation parameter $`𝐥`$ translates the mode functions in three dimensions and the parameter $`s`$ compresses and stretches them. Even though the quantization in the wavelet basis is in free space the operators and mode functions are expanded using a countable set of basis functions. On the contrary to plane waves the polarization of the new mode functions is not constant. As is the case with wavelets, the integral of the mode functions over quantization volume vanishes
$$𝐮_{s𝐥,\sigma }^{iE}(𝐫)d^3𝐫=𝐮_{s𝐥,\sigma }^{iB}(𝐫)d^3𝐫=0.$$
(58)
This means that also the mode functions have similar ’waviness’ as wavelets.
It is easy to show that the new operators $`\widehat{b}_{s𝐥,\sigma }^i`$ and $`\widehat{b}_{s𝐥,\sigma }^i`$ obey bosonic commutation relations
$`[b_{s𝐥,\sigma }^i,b_{s^{}𝐥^{},\sigma ^{}}^i^{}]=[b_{s𝐥,\sigma }^i,b_{s^{}𝐥^{},\sigma ^{}}^i^{}]=0,`$ (59)
$`[b_{s𝐥,\sigma }^i,b_{s^{}𝐥^{},\sigma ^{}}^i^{}]=\delta _{ii^{}}\delta _{ss^{}}\delta _{\mathrm{𝐥𝐥}^{}}\delta _{\sigma \sigma ^{}}.`$ (60)
Because of this the operators act like annihilation and creation operators for wavelet modes. This means that the formalism in wavelet basis is the same as with plane wave operators. Multiplying both sides of Eq. (54) by $`d_{𝐤,s𝐥}^i`$ and performing the sum $`\underset{s𝐥i}{}`$ we get with the help of the orthogonality integral
$$\widehat{a}_{𝐤\sigma }=\underset{s𝐥i}{}d_{𝐤,s𝐥}^i\widehat{b}_{s𝐥,\sigma }^i.$$
(61)
Similarly it is possible to expand plane wave basis functions using wavelet mode functions.
The field Hamiltonian in wavelet basis is obtained as an integral of the energy density over the quantization volume
$$\widehat{H}_F=\left(\frac{1}{2}ϵ_0\widehat{E}^2(𝐫)+\frac{1}{2\mu _0}\widehat{B}^2(𝐫)\right)d^3𝐫,$$
(62)
where expansions (52) and (56) are used for electric and magnetic field operators. In the calculation of (62) the following integrals are needed (last one is listed only for completeness)
$`{\displaystyle 𝐮_{s𝐥,\sigma }^{iE}(𝐫)𝐮_{s^{}𝐥^{},\sigma ^{}}^{i^{}E}(𝐫)d^3𝐫}=\delta _{\sigma \sigma ^{}}{\displaystyle \frac{\mathrm{}}{2ϵ_0}}w_{s^{}𝐥^{},s𝐥}^{i^{},i}`$ (63)
$`{\displaystyle 𝐮_{s𝐥,\sigma }^{iB}(𝐫)𝐮_{s^{}𝐥^{},\sigma ^{}}^{i^{}B}(𝐫)d^3𝐫}=\delta _{\sigma \sigma ^{}}{\displaystyle \frac{\mathrm{}}{2ϵ_0c^2}}w_{s^{}𝐥^{},s𝐥}^{i^{},i}`$ (64)
$`{\displaystyle 𝐮_{s𝐥,\sigma }^{iE}(𝐫)𝐮_{s^{}𝐥^{},\sigma ^{}}^{i^{}B}(𝐫)d^3𝐫}=0,`$ (65)
where the real coupling constants of the modes are
$`w_{s𝐥,s^{}𝐥^{}}^{ii^{}}=w_{s^{}𝐥^{},s𝐥}^{i^{}i}`$ $`=`$ $`{\displaystyle \underset{𝐤}{}}\omega _𝐤d_{𝐤,s𝐥}^id_{𝐤,s^{}𝐥^{}}^i^{}`$ (66)
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle \underset{𝐤}{}}{\displaystyle d^3𝐫d^3𝐫^{}\omega _𝐤\mathrm{\Psi }_{s𝐥}^i(𝐫)\mathrm{\Psi }_{s^{}𝐥^{}}^i^{}(𝐫^{})e^{i𝐤(𝐫^{}𝐫)}}.`$ (67)
Using (39) we get
$`w_{s𝐥,s^{}𝐥^{}}^{ii^{}}`$ $`=`$ $`{\displaystyle \underset{𝐤}{}}\omega _𝐤d_{𝐤,s0}^id_{𝐤,s^{}0}^i^{}\mathrm{exp}(i𝐤(2^s^{}𝐥^{}2^s𝐥))`$ (68)
$`=`$ $`F_{ss^{}}^{ii^{}}(2^s^{}𝐥^{}2^s𝐥),`$ (69)
where
$$F_{ss^{}}^{ii^{}}(𝐱)=\underset{𝐤}{}\omega _𝐤d_{𝐤,s0}^id_{𝐤,s^{}0}^i^{}e^{i𝐤𝐱}.$$
(70)
The identity (68) shows that the coupling constant is a function of the difference of the scaled translation parameters only. Using equations (63) and (64) and the identity
$$𝐮_{s𝐥,\sigma }^{iE}(𝐫)𝐮_{s^{}𝐥^{},\sigma ^{}}^{i^{}E}(𝐫)d^3𝐫=c^2𝐮_{s𝐥,\sigma }^{iB}(𝐫)𝐮_{s^{}𝐥^{},\sigma ^{}}^{i^{}B}(𝐫)d^3𝐫$$
(71)
we get after a straightforward calculation for the free field Hamiltonian
$$\widehat{H}_F=\underset{\sigma }{}\underset{s𝐥i}{}\underset{s^{}𝐥^{}i^{}}{}\mathrm{}w_{s𝐥,s^{}𝐥^{}}^{i,i^{}}\widehat{b}_{s𝐥}^i\widehat{b}_{s^{}𝐥^{}}^i^{}+\frac{\mathrm{}}{2}\underset{\sigma }{}\underset{s𝐥i}{}w_{s𝐥,s𝐥}^{i,i}.$$
(72)
The modes are coupled with the coupling constant $`w_{s𝐥,s^{}𝐥^{}}^{i,i^{}}`$. The diagonal elements $`w_{s𝐥,s𝐥}^{i,i}`$ give the energy of the field state which has only one mode with parameters $`s`$, $`𝐥`$ and $`i`$ excited, i.e., energy of one wavelet quantum. From Eq. (66) and (67) it is seen that the coupling is nonzero only if the corresponding wavelets and Fourier transforms have a nonzero overlap. Because the wavelets are localized both in real and Fourier space this means that for majority of mode pairs the coupling is zero. The detailed structure of the coupling constants $`w_{s𝐥,s^{}𝐥^{}}^{i,i^{}}`$ depends on the wavelet used. The second term in Eq. (72) is the result of the self-coupling of the modes and gives the zero point energy of the field. As is the case for plane waves the energy per mode is half of the unit excitation field.
## IV Numerical simulations using wavelet modes
### A Conventions and wavelets used
In this section we give examples how to apply the theory developed in the last section. We do the calculations in two rather than in three dimensions because visualization is easier and all essential features are included. Thus the $`𝐤`$-vector is restricted to xy plane $`𝐤=k_x𝐞_1+k_y𝐞_2`$ and the polarization vector is $`ϵ_{𝐤1}=𝐞_3`$. Here we exclude the other polarization $`ϵ_{𝐤2}`$. The cross product in the expansion of the magnetic field (46) becomes $`𝐤\times ϵ_{𝐤1}=k_y𝐞_1+k_x𝐞_2`$. The choice of the field configuration considered is the same as in our earlier paper . The mode functions for electric and magnetic fields become
$`u_{s𝐥,1}^{iE}(𝐫)`$ $`=`$ $`{\displaystyle \frac{i}{L}}{\displaystyle \underset{𝐤}{}}\left({\displaystyle \frac{\mathrm{}\omega _𝐤}{2ϵ_0}}\right)^{1/2}d_{𝐤,s𝐥}^i𝐞_3e^{i𝐤𝐫}`$ (73)
$`u_{s𝐥,1}^{iB}(𝐫)`$ $`=`$ (75)
$`{\displaystyle \frac{i}{L}}{\displaystyle \underset{𝐤}{}}\left({\displaystyle \frac{\mathrm{}}{2ϵ_0\omega _𝐤}}\right)^{1/2}d_{𝐤,s𝐥}^i(k_y𝐞_1+k_x𝐞_2)e^{i𝐤𝐫},`$
where the summation over $`𝐤`$ is now over the xy plane. In the following we choose units such that $`\mathrm{}=ϵ_0=c=1`$. In all examples in this section we use the Meyer wavelets shown in Fig. 2. Meyer wavelets were designed using smoothened Fourier transform of Shannon wavelets, so the division of the Fourier-plane shown in Fig. 5 is rather valid. The Fourier transforms of the mode functions are quite similar to the Fourier transforms of the corresponding wavelets. In all simulations the parameters are chosen in such a way that it is necessary to consider only two scales $`s=0`$ and $`s=1`$. This means that the frequencies $`|k_x|<\pi `$, $`|k_y|<\pi `$ and $`|k_x|>4\pi `$, $`|k_y|>4\pi `$ are excluded. One has to note that in two dimensions the scaling of the wavelet mode functions become
$$𝐮_{s𝐥,\sigma }^{iE,B}(𝐫)=𝐮_{s0,\sigma }^{iE,B}(𝐫2^s𝐥)=2^{3s/2}𝐮_{00,\sigma }^{iE,B}(2^s𝐫𝐥),$$
(76)
instead of the scaling in three dimensions given by Eqs. (53) and (57).
The absolute values of the electric and magnetic mode functions with indices $`s=0`$ and $`l_x=l_y=0`$ are shown in Figs. 7 and 8. On the left in both figures there is the mode function with index $`i=1`$ and on the right $`i=3`$. The $`i=2`$ mode function is the same as $`i=1`$ case with axes $`x`$ and $`y`$ changed. All mode functions are clearly localized at the origin and have wavelet type of oscillating structure. Because the Fourier transform of $`i=3`$ two dimensional wavelet is nonzero at higher frequencies than Fourier transforms of $`i=1`$ and $`i=2`$ wavelets, the oscillations of $`i=3`$ mode function have smaller details. As was explained earlier the parameters $`s`$ and $`𝐥`$ scale and translate mode functions.
### B A Gaussian photon
In this section we study a one photon field which is a superposition of different single excitation plane wave modes. The distribution of the absolute values of the mode coefficients is a Gaussian centered at some point $`𝐤_0`$ with a variance $`\mathrm{\Delta }_k^2`$ in $`𝐤`$-space. Thus the field state is
$$|\mathrm{\Psi }=\underset{𝐤}{}c_𝐤|1_𝐤,\{0\}=\underset{𝐤}{}c_𝐤\widehat{a}_𝐤^{}|\{0\}$$
(77)
where
$$c_𝐤=(2\pi \mathrm{\Delta }_k^2)^{1/2}e^{i𝐤𝐫_0}\mathrm{exp}\left(\frac{(𝐤𝐤_0)^2}{4\mathrm{\Delta }_k^2}\right).$$
(78)
The transformation to wavelet basis can be done using the expansion of operators $`\widehat{a}_𝐤`$ in a wavelet basis (54). The state vector in wavelet basis becomes
$$|\mathrm{\Psi }=\underset{s𝐥i}{}c_{s𝐥}\widehat{b}_{s𝐥}^i|\{0\},$$
(79)
where
$$c_{s𝐥}^i=\underset{𝐤}{}d_{𝐤,s𝐥}^ic_𝐤.$$
(80)
Because coefficients $`d_{𝐤,s𝐥}^i`$ and $`c_𝐤`$ factorize, it is possible to factorize the sum (80).
First we briefly study a one dimensional Gaussian distribution. We have a Gaussian photon with parameters $`k_0=18.0`$, $`x_0=10.0`$ and $`\mathrm{\Delta }_k^2=10.0`$. The width of the distribution is relatively large so one would expect several scales to have nonzero coefficients. Figure 9 shows coefficients of three different scales as a function of a translation parameter. The peak on the left has the scale parameter $`s=1`$. When the translation parameter is changed by one, the wavelet and mode functions are translated by $`2^s`$ in real space. At scale $`s=1`$ the translation unit is $`2^1=0.5`$. The distribution is centered at $`x_0=10.0`$ so coefficients at $`s=1`$ are centered at $`l=20`$. When $`s=2`$ and $`s=3`$ the coefficients around $`l=40`$ and $`l=80`$ have nonzero values. It is seen that a Gaussian state with the parameters used has the biggest contribution at a scale $`s=2`$.
Next we study the time evolution of a two dimensional Gaussian photon. Figure 10 shows the time evolution of the energy density distribution of a photon at two different times. The parameters in Eq. (78) for the initial state at $`t=0.0`$ are $`r_x=4.0`$, $`r_y=0.0`$, $`k_x=7.0`$, $`k_y=0.0`$ and $`\mathrm{\Delta }_k^2=0.25`$. It is seen that the energy density profile of the single excitation field which has a Gaussian distribution in Fourier space is also Gaussian. At a later time $`t=6.0`$ the photon has propagated to the right and the intensity profile has spread a little. The time evolution is qualitatively the same as obtained in the paper using plane wave quantization.
On the contrary to plane wave modes the wavelet modes are coupled also for the free field. This makes the operation of the free field Hamiltonian to the statevector, which is needed in the numerical integration, slower. However, most of the coupling constants $`w_{s𝐥,s^{}𝐥^{}}^{i,i^{}}`$ are zero. Use of this fact makes the operation much faster. The function $`F_{ss^{}}^{ii^{}}(𝐱)`$ in the calculation of the coupling coefficients can be used to save memory. In general fewer wavelet mode functions are needed to represent a localized field state than plane wave mode functions.
### C Wavelet mode function initial state
Next we study the time evolution of the field state which at time $`t=0.0`$ has only one wavelet mode excited. The mode which is initially excited has parameters $`s=1`$, $`𝐥=0`$ and $`i=3`$. The intensity of the field state is shown in Fig. 11. It is clearly localized around $`𝐫=0`$ which is expected based on the parameters chosen. The Fourier transform for $`i=3`$ mode function is divided into four parts at the corners of the frequency interval with a specific scaling index $`s`$, as is explained earlier and shown in Fig. 5. The intensity profile at time $`t=7.5`$ is shown in Fig. 11. The intensity is divided to four main parts which all propagate away from the origin. This is understandable based on the Fourier transform of the mode function.
### D Decay of a two level atom
Next we couple a two level atom to the two dimensional field. The interaction Hamiltonian in dipole approximation can be written as
$$\widehat{H}_I=\widehat{𝐃}\widehat{𝐄}(𝐫_0),$$
(81)
where $`\widehat{𝐃}`$ is the electric dipole moment operator of the two level atom
$$\widehat{𝐃}=(D\widehat{\sigma }_++D^{}\widehat{\sigma }_{})𝐞_3$$
(82)
and $`𝐫_0`$ the position of the atom. We take the direction of the dipole operator to be in the $`z`$ direction. Using the wavelet expansion of the field (51) we get
$$\widehat{H}_I=\underset{s𝐥i}{}\underset{\sigma }{}(D𝐮_{s𝐥,\sigma }^{iE}(𝐫_0)\widehat{\sigma }_+\widehat{b}_{s𝐥,\sigma }^i+D^{}𝐮_{s𝐥,\sigma }^{iE}(𝐫_0)\widehat{\sigma }_{}\widehat{b}_{s𝐥,\sigma }^i).$$
(83)
In the Hamiltonian the rotating wave approximation (RWA) has been used. The approximation can be done if the mode functions used are well localized in Fourier space. This is the case with Meyer wavelets. The scale parameter $`s`$ determines the frequency interval where the Fourier transform of the mode function is nonzero. Typically the atom interacts with mode functions of only a few, maybe one or two, scales. The mode functions are also spatially localized. Because the interaction is proportional to the electric mode function evaluated at a position of the atom, only mode functions which are centered close to the atom interact with it. In this respect the situation is different compared to the plane wave mode functions which are delocalized. In general the atom is coupled to fewer mode functions than in the plane wave quantization.
Fig. 12 shows the logarithm of the excitation probability of the atom as function of time. The resonance frequency of the atom is $`\omega =10.0`$ and the dipole coupling D=0.06. The atom decays energy to the wavelet modes. The decay is clearly exponential with a decay constant $`\mathrm{\Gamma }=0.18`$. This corresponds to the theoretical value $`\mathrm{\Gamma }=\frac{1}{2}D^2\omega ^2`$.
## V CONCLUSION
In this paper we have shown how wavelets can be used as basis functions in canonical quantization. Different mode functions for electric and magnetic fields, which are localized both in real and Fourier space, are obtained. Mode functions as well as new operators in wavelet basis are linear transforms of plane wave mode functions and operators. The new annihilation and creation operators satisfy bosonic commutation relations. Because the formalism remains the same, it is easy to change the basis from plane waves to wavelets. We have applied the theory to a few example simulations and showed that the new basis gives well known results in all cases. In this paper we have used wavelet basis for bosonic operators in canonical quantization. The same methods can be used to change basis also for fermionic operators. A localized basis is beneficial for many solid state and semiconductor physics problems.
There are several generalizations and improvements of the theory described in this paper. Complex and biorthogonal wavelets have some benefits compared to real and orthonormal wavelets, for example they can be symmetric. One generalization is to use wavelet packets or multiwavelets instead of wavelets. Finally it would be interesting to compare characteristics of mode functions of different wavelets. It is also possible to construct new wavelets, which have desirable properties, for different problems.
## VI ACKNOWLEDGEMENTS
We thank the Academy of Finland (project 43336) for financial support. Computers of the Center for Scientific Computing (CSC) were used in the simulations. The C++ class library ’blitz’ developed by Todd Veldhuizen was used (http://oonumerics.org/blitz/). We thank A. R. Baghai-Wadji, N. Lütkenhaus and K.-A. Suominen for discussions and comments. More figures of different wavelet mode functions can be found from the page http://tftsg6.hip.helsinki.fi/~mhavukai/wavelets/ .
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# PROTON SPIN STRUCTURE AND THE LOW-ENERGY 𝑝𝑝→𝑝𝑝𝜂' REACTION
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# 1 Introduction
## 1 Introduction
Let $`V`$ be a simple vertex operator algebra (see \[B\], \[FLM\]), and $`G`$ a finite automorphism group. A major problem in orbifold conformal field theory is to understand the module category for the vertex operator algebra $`V^G`$ of $`G`$-invariants. In the case $`V`$ is holomorphic, this is related to the quantum double (see \[DPR\] and \[DM5\]). The main feature in the study of orbifold theory is the appearance of dual pairs. It is proved in \[DLM1\] and \[DM2\] that $`G`$ and $`V^G`$ form a dual pair in the sense of Howe \[H1\]-\[H2\]. More precisely, it is shown in \[DLM1\] that all the irreducible $`G`$-modules occur in $`V`$ and the space of multiplicity of each irreducible $`G`$-module in $`V`$ is an irreducible $`V^G`$-module. Moreover, inequivalent irreducible $`G`$-modules produce inequivalent $`V^G`$-modules in this way.
In this paper we extend the duality result in \[DLM1\] to any irreducible $`V`$-module and obtain again several duality theorems of Schur-Weyl type. In the process we realize that it is better to consider a finite set of inequivalent irreducible modules which is $`G`$-stable instead of one single module. The general setting is more natural and the results are more beautiful.
More explicitly, for an irreducible $`V`$-module $`M=(M,Y_M)`$ (see Section 4 for the details of the definition of a module) and $`gG`$ we define a new irreducible $`V`$-module $`Mg=(Mg,Y_{Mg}).`$ Here $`Mg`$ is equal to $`M`$ as a vector space and $`Y_{Mg}(v,z)=Y_M(gv,z)`$ for $`vV`$ following \[DM1\]. A set $`𝒮`$ of irreducible $`V`$-modules is called $`G`$-stable if for any $`MS`$ and $`gG`$ there exists $`N𝒮`$ such that $`MgN.`$ Then every irreducible $`V`$-module $`M`$ produces naturally such a set by collecting all $`Mg`$ for $`gG.`$ Now we take a finite $`G`$-stable set $`𝒮`$ consisting of inequivalent irreducible $`V`$-modules. We construct a finite dimensional semisimple associative algebra $`A_\alpha (G,𝒮)`$ for a suitable 2-cocycle on $`G`$ (which is uniquely determined up to isomorphism by $`V,`$ $`G`$ and $`𝒮`$) such that $`(A_\alpha (G,𝒮),V^G)`$ forms a dual pair on $`=_{M𝒮}M`$ in the precise sense of Howe \[H1\]-\[H2\]. That is, each simple $`A_\alpha (G,𝒮)`$ occurs in $``$ and its multiplicity space is an irreducible $`V^G`$-module. Moreover, the different multiplicity spaces are inequivalent $`V^G`$-modules.
These duality results not only tell us the complete reducibility of every irreducible $`V`$-module as a $`V^G`$-module, but also provide an equivalence between the $`A_\alpha (G,𝒮)`$-module category and a subcategory of $`V^G`$-modules generated by the irreducible submodules of $``$ by sending each simple $`A_\alpha (G,𝒮)`$-module to its multiplicity space. This kind of idea has appeared in \[DPR\] and \[DM5\] in the study of holomorphic orbifold conformal field theory.
Although the algebra $`A_\alpha (G,𝒮)`$ mentioned above appears naturally in the theory of vertex operator algebras, the construction itself is totally canonical and abstract. In fact, one can define $`A_\alpha (G,𝒮)`$ for any finite group $`G,`$ any finite right $`G`$-set $`𝒮`$ and a suitable 2-cocycle on $`G`$ (see Section 3). This algebra is essentially the crossed product in the theory of Hopf algebra (see \[S\], \[BCM\] and \[DT\]). It turns out that $`A_\alpha (G,𝒮)`$ is the right algebra in the study of general orbifold theory. In the case that $`𝒮`$ is the dual basis of $`C[G]^{},`$ $`A_\alpha (G,𝒮)`$ is exactly the twisted double $`D_\alpha (G)`$ introduced in \[M\] and \[DM5\] in the study of holomorphic orbifold theory. The twisted double $`D_\alpha (G)`$ is conjecturally isomorphic to the twisted quantum double introduced in \[DPR\] and \[D\].
A main tool in the proof of the main theorems is a series of associative algebras $`A_n(V)`$ constructed in \[DLM4\] for nonnegative integers $`n.`$ The original motivation for introducing and studying the $`A_n(V)`$ comes from the representation theory of vertex operator algebras. Let $`M=_{n0}M(n)`$ be an admissible $`V`$-module with $`M(0)0`$ (see Section 4 for the definition of admissible module). The algebra $`A_n(V)`$ is a suitable quotient of $`V`$ and “takes care” of the first $`n+1`$ pieces of $`M`$: each $`M(m)`$ is a module for $`A_n(V).`$ In the case $`n=0,`$ $`A_0(V)`$ reduces to the associative algebra $`A(V)`$ introduced previously in \[Z\]. The main results concerning $`A_n(V)`$ are summarized in Section 4. The result that we often use in this paper is the following: $`M`$ is irreducible if and only if each $`M(n)`$ is a simple $`A_n(V)`$-module for all $`n0`$ \[DM4\]. This key result allows us to reduce an infinite dimensional problem ($`M`$ is infinite dimensional) to a finite dimensional problem (each $`M(n)`$ is finite dimensional in our main theorems).
This paper is organized as follows. In Section 2 we review twisted group algebras. Section 3 is about the algebra $`A_\alpha (G,𝒮)`$ based on a finite group $`G,`$ a right $`G`$-set $`𝒮`$ and a suitable 2-cocycle $`\alpha `$ on $`G.`$ We construct all simple modules for $`A_\alpha (G,𝒮)`$ explicitly by using twisted group algebras. A basis for the center is also given for the future study. In Section 4 we recall the various notions of modules for a vertex operator algebra and review the theory on associative algebras $`A_n(V).`$ Section 5 is the first part on duality. We show that if $`M`$ is an irreducible $`V`$-module which is $`G`$-stable in the sense that $`MgM`$ for all $`g,`$ then there exists a 2-cocycle $`\alpha _MZ^2(G,C^{})`$ such that $`(C^{\alpha _M}[G],V^G)`$ forms a dual pair on $`M`$ where $`C^{\alpha _M}[G]`$ is the twisted group algebra. Section 6 is a continuation of Section 5 without assuming that $`M`$ is $`G`$-stable. The algebra $`A_\alpha (G,S)`$ based on a $`G`$-stable set $`𝒮`$ of $`V`$-modules enters the picture naturally. The two main theorems of this paper are proved in this section.
We would like to thank A. Wassermann for pointing out a gap in the proof of Theorem 2.4 (i) of \[DLM1\]. This gap has been filled in this paper (see Remark 5.2).
## 2 Twisted group algebras
In this elementary section we review some results on twisted group algebras which will be used in Section 3.
Throughout this paper, $`F^{}`$ denotes the multiplicative group of a field $`F`$ and $`G`$ a finite group which acts trivially on $`F^{}`$. Let $`\alpha `$ be any element in $`Z^2(G,F^{})`$ . Thus $`\alpha `$ is a map
$$\alpha :G\times GF^{}$$
which satisfies the following properties:
$`\alpha (x,1)=\alpha (1,x)=1\text{f}orallxG`$
$`\alpha (x,y)\alpha (xy,z)=\alpha (y,z)\alpha (x,yz)\text{f}orallx,y,zG.`$
###### Definition 2.1
The twisted group algebra $`F^\alpha [G]`$ of $`G`$ over $`F`$ is defined to be the vector space over $`F`$ with basis $`\{\overline{g}|gG\}`$. Multiplication in $`F^\alpha [G]`$ is defined distributively by using
$$(a\overline{x})(b\overline{y})=ab\alpha (x,y)\overline{xy}(a,bF,x,yG).$$
It is easy to verify that $`F^\alpha [G]`$ is an $`F`$-algebra with $`\overline{1}`$ as the identity element and with
$$\overline{g}^1=\alpha (g^1,g)^1\overline{g^1}=\alpha (g,g^1)^1\overline{g^1}\text{f}orallgG.$$
###### Lemma 2.2
\[K\] If char $`F`$ does not divide $`|G|`$, then $`F^\alpha [G]`$ is semisimple.
Our aim for this section is to find an explicit basis for the center $`Z(F^\alpha [G])`$.
###### Definition 2.3
An element $`gG`$ is said to be $`\alpha \mathrm{r}\mathrm{e}\mathrm{g}\mathrm{u}\mathrm{l}\mathrm{a}\mathrm{r}`$ if
$$\alpha (g,x)=\alpha (x,g)\text{f}orallxC_G(g).$$
Thus $`g`$ is $`\alpha `$-regular if and only if $`\overline{g}\overline{x}=\overline{x}\overline{g}`$ for all $`xC_G(g).`$
It is clear that identity element of $`G`$ is $`\alpha `$-regular for any cocycle $`\alpha `$.
###### Lemma 2.4
\[K\] For any $`\alpha Z^2(G,F^{})`$, the following properties hold:
i) An element $`gG`$ is $`\alpha `$-regular if and only if it is $`\beta `$-regular for any cocycle $`\beta `$ cohomologous to $`\alpha `$.
ii) If $`gG`$ is $`\alpha `$-regular, then so is any conjugate of $`g`$.
###### Definition 2.5
Let $`C`$ be a conjugacy class of $`G`$ and let $`gC`$. We say that $`C`$ is $`\alpha `$-regular if $`g`$ is $`\alpha `$-regular.
###### Remark 2.6
If we replace $`\overline{g}`$ by $`\stackrel{~}{g}=\lambda (g)\overline{g},\lambda (g)F^{},\lambda (1)=1`$, yields
$$\stackrel{~}{x}\stackrel{~}{y}=\alpha (x,y)\lambda (x)\lambda (y)\lambda (xy)^1\stackrel{~}{xy}\text{f}orallx,yG.$$
Hence, performing a diagonal change of basis $`\{\overline{g}|gG\}`$ results in replacing $`\alpha `$ by a cohomologous cocycle. The advantage of this observation is that in choosing a distinguished cocycle cohomologous to $`\alpha `$, we can work exclusively within the algebra $`F^\alpha [G]`$.
###### Definition 2.7
We say that $`\alpha Z^2(G,F^{})`$ is a normal cocycle if $`\alpha (x,g)=\alpha (xgx^1,x)`$ for all $`xG`$ and all $`\alpha `$-regular $`gG`$.
Therefore $`\alpha `$ is normal if and only if $`\overline{x}\overline{g}\overline{x}^1=\overline{xgx^1}`$ for all $`xG`$ and all $`\alpha `$-regular $`gG`$.
###### Lemma 2.8
\[K\] Let $`\{g_1,\mathrm{},g_r\}`$ be a full set of representatives for the $`\alpha `$-regular conjugacy classes of $`G`$ and, for each $`i\{1,\mathrm{},r\}`$, let $`T_i`$ be a left transversal for $`C_G(g_i)`$ in $`G`$. Put $`\stackrel{~}{tg_it^1}=\overline{t}\overline{g_i}\overline{t}^1,1ir`$, and $`\stackrel{~}{g}=\overline{g}`$ if $`g`$ is not $`\alpha `$-regular . Then $`\stackrel{~}{x}\stackrel{~}{g}\stackrel{~}{x}^1=\stackrel{~}{xgx^1}`$ for all $`xG`$ and all $`\alpha `$-regular $`gG`$. Thus any cocycle $`\alpha Z^2(G,F^{})`$ is cohomologous to a normal cocycle.
###### Theorem 2.9
\[K\] Let $`F`$ be an arbitrary field, let $`\alpha Z^2(G,F^{})`$ and let $`\{g_1,\mathrm{},g_r\}`$ be a full set of representatives for the $`\alpha `$-regular conjugacy classes of $`G`$. Denote by $`T_i`$ a left transversal for $`C_G(g_i)`$ in $`G`$ and by $`C_i`$ the $`\alpha `$-regular conjugacy class of $`G`$ containing $`g_i,1ir`$. Then the elements
$$z_i=\underset{tT_i}{}\overline{t}\overline{g_i}\overline{t}^1(1ir)$$
constitute an $`F`$-basis $`Z(F^\alpha [G])`$. In particular, if $`\alpha `$ is a normal cocycle, then the elements
$$z_i=\underset{gC_i}{}\overline{g}(1ir)$$
constitute an $`F`$-basis $`Z(F^\alpha [G])`$.
## 3 The Generalized twisted double
In this section we construct a finite dimensional semisimple associative algebra $`A_\alpha (G,S)`$ over $`C`$ associated to a finite group $`G,`$ a finite right $`G`$-set $`𝒮`$ and a suitable 2-cocycle $`\alpha `$. Although this construction in the present form is influenced by the work in \[DPR\], \[M\] and \[DM5\] and Section 6 of this paper, its origin goes back to the earlier work of \[S\], \[BCM\] and \[DT\] in Hopf algebra, where it is called crossed product. In the case the $`G`$-set $`𝒮`$ consists of a dual basis of $`C[G]^{},`$ $`A_\alpha (G,𝒮)`$ is exactly the twisted double studied in \[M\] and \[DM5\]. This should explain why we call $`A_\alpha (G,𝒮)`$ a generalized twisted double. The exposition of this section follows \[M\] closely and ideas are also similar.
Let $`G`$ be any finite group with identity $`1`$ and $`𝒮`$ be a finite right $`G`$-set. Set
$$C𝒮=\underset{s𝒮}{}Ce(s)$$
which is an associative algebra under product $`e(s)e(t)=\delta _{s,t}e(t)`$ and isomorphic to $`C^{|S|}.`$ The right action of $`G`$ on $`𝒮`$ induces a right $`G`$-module structure on $`C𝒮.`$ Let
$$𝒰(C𝒮)=\{\underset{s𝒮}{}\lambda _se(s)|\lambda _sC^{}\}$$
be the set of units of $`C𝒮.`$ It is easy to check that $`𝒰(𝒞𝒮)`$ is a multiplicative right $`G`$-module.
From now on we will fix an element $`\alpha Z^2(G,𝒰(𝒞𝒮))`$ such that $`\alpha (h,1)=\alpha (1,h)=_{s𝒮}e(s)`$ for all $`h`$ in $`G`$. Then $`\alpha `$ defines functions $`\alpha _s:G\times GC^{}`$ for $`s𝒮`$ such that
$$\alpha (h,k)=\underset{sS}{}\alpha _s(h,k)e(s)$$
for $`h,kG.`$ The cocycle condition $`\alpha (hk,l)\alpha (h,k)^l=\alpha (h,kl)\alpha (k,l)`$ implies immediately that
$$\alpha _s(hk,l)\alpha _{(s)l^1}(h,k)=\alpha _s(h,kl)\alpha _s(k,l)$$
for $`h,k,lG`$ and $`s𝒮.`$
The main object that we study in this section is the vector space $`𝒜_\alpha (G,𝒮)=C[G]C𝒮`$ with a basis $`ge(s)`$ for $`gG`$, and $`s𝒮`$.
###### Proposition 3.1
i) The $`𝒜_\alpha (G,𝒮)`$ is an associative algebra under multiplication defined by
$$ge(s)he(t)=\alpha _t(g,h)ghe(sh)e(t).$$
ii) If $`\alpha `$ is cohomologous to $`\beta `$ then $`𝒜_\alpha (G,𝒮)𝒜_\beta (G,𝒮)`$ as algebras.
Proof: For i), we first prove that the product is associative. Let $`g,h,kG`$ and $`s,t,u𝒮.`$ Then
$`(ge(s)he(t))ke(u)`$
$`=(\alpha _t(g,h)ghe(sh)e(t))ke(u)`$
$`=\delta _{sh,t}\alpha _t(g,h)ghe(t)ke(u)`$
$`=\delta _{sh,t}\alpha _t(g,h)\alpha _u(gh,k)ghke(tk)e(u)`$
$`=\delta _{sh,t}\delta _{tk,u}\alpha _t(g,h)\alpha _u(gh,k)ghke(u)`$
$`=\delta _{sh,t}\delta _{tk,u}\alpha _{uk^1}(g,h)\alpha _u(gh,k)ghke(u)`$
$`=\delta _{sh,t}\delta _{tk,u}\alpha _u(g,hk)\alpha _u(h,k)ghke(u)`$
and
$`ge(s)(he(t)ke(u))`$
$`=ge(s)(\alpha _u(h,k)hke(tk)e(u))`$
$`=ge(s)(\delta _{tk,u}\alpha _u(h,k)hke(u))`$
$`=\delta _{tk,u}\alpha _u(g,hk)\alpha _u(h,k)ghke(s(hk))e(u)`$
$`=\delta _{tk,u}\delta _{s(hk),u}\alpha _u(g,hk)\alpha _u(h,k)ghke(u).`$
Thus $`(ge(s)he(t))ke(u)=ge(s)(he(t)ke(u))`$ and associativity holds.
A straightforward verification
$`ge(s){\displaystyle \underset{t𝒮}{}}1e(t)`$
$`={\displaystyle \underset{t𝒮}{}}\alpha _t(g,1)ge(s)e(t)`$
$`=ge(s)`$
and
$`({\displaystyle \underset{t𝒮}{}}1e(t))ge(s)`$
$`={\displaystyle \underset{t𝒮}{}}\alpha _s(1,g)ge(tg)e(s)`$
$`=ge(s)`$
shows that $`_{t𝒮}1e(t)`$ is an identity on $`𝒜_\alpha (G,𝒮)`$.
For ii) since $`\beta `$ is cohomologous to $`\alpha `$, so there exists a map $`\lambda :G𝒰(C𝒮)`$ such that $`\beta (x,y)=\alpha (x,y)\lambda (x)^y\lambda (xy)^1\lambda (y)`$ for all $`x,yG`$. For each $`xG`$ , we can rewrite $`\lambda (x)`$ as $`_{s𝒮}\lambda _s(x)e(s)`$. Therefore, we have $`\beta _s(x,y)=\alpha _s(x,y)\lambda _{sy^1}(x)\lambda _s(xy)^1\lambda _s(y)`$, for any $`s𝒮,x,yG.`$ Let $`f:𝒜_\alpha (G,𝒮)𝒜_\beta (G,𝒮)`$ defined by $`f(ge(s))=\lambda _s(g)^1(ge(s)).`$ Since
$`f(ge(s)he(t))`$
$`=f(\alpha _t(g,h)ghe(sh)e(t))`$
$`=\delta _{sh,t}\lambda _t(gh)^1\alpha _t(g,h)ghe(t)`$
and
$`f(ge(s))f(he(t))`$
$`=\lambda _s(g)^1ge(s)\lambda _t(h)^1he(t)`$
$`=\lambda _s(g)^1\lambda _t(h)^1\beta _t(g,h)ghe(sh)e(t),`$
$`𝒜_\alpha (G,𝒮)`$ and $`𝒜_\beta (G,𝒮)`$ are isomorphic as algebras. $`\mathrm{}`$
###### Remark 3.2
The algebra $`A_\alpha (G,𝒮)`$ is essentially the crossed product $`C𝒮\mathrm{}C[G]`$ except that $`C𝒮`$ in our setting is a right $`C[G]`$-module instead of left $`C[G]`$-module in the setting of crossed product. So one could have used the results from \[S\], \[BCM\] and \[DT\] to conclude that $`A_\alpha (G,𝒮)`$ is an associative algebra after a careful identifying two different settings. Since the proof is not very long we chose to give a complete proof.
###### Remark 3.3
If we take $`𝒮=\{e(g)|gG\}`$ to be the dual basis of $`C[G]^{}`$ which is a right $`G`$-set via $`e(g)h=e(h^1gh)`$ for $`g,hG`$ we get the associative algebra $`D_\alpha (G)`$ in \[DM5\] which is a deformation of the Drinfeld’s double $`D(G).`$ It turns out that there exists $`\alpha `$ determined naturally by the “twisted representations” such that $`D_\alpha (G)`$ is the right algebra in the study of holomorphic orbifold conformal field theory (see \[DVVV\] and \[DM5\]).
For each $`s𝒮`$, let $`G_s=\{hG|sh=s\}`$ be the stabilizer of $`s.`$ Note that $`Res_{G_s}^G\alpha _s`$ is in $`Z^2(G_s,C^{}).`$ Let $`𝒪_s`$ be the orbit of $`s`$ under $`G`$ and $`G=_{i=1}^kG_sg_i`$ be a right coset decomposition with $`g_1=1.`$ Then $`𝒪_s=\{sg_i|i=1,\mathrm{},k\}`$ and $`G_{sg_i}=g_i^1G_sg_i.`$ We define several subspaces of $`𝒜_\alpha (G,𝒮):`$
$$S(s)=ae(s)|aG_s,N(s)=ae(s)|aGG_s,$$
$$D(s)=ae(s)|aG,D(𝒪_s)=ae(sg_i)|i=1,\mathrm{},k,aG.$$
Then $`D(s)=S(s)N(s)`$.
Decompose $`𝒮`$ into a disjoint union of orbits $`𝒮=_{jJ}𝒪_j`$. Let $`s_j`$ be a representative element of $`𝒪_j.`$ Then $`𝒪_j=\{s_jh|hG\}`$ and $`𝒜_\alpha (G,𝒮)=_{jJ}D(𝒪_{s_j}).`$
###### Lemma 3.4
Let $`s𝒮`$ and $`G=_{i=1}^kG_sg_i.`$ Then
1) $`S(s)`$ is a subalgebra of $`𝒜_\alpha (G,𝒮)`$ isomorphic to $`C^{\alpha _s}[G_s]`$ where $`C^{\alpha _s}[G_s]`$ is the twisted group algebra.
2) $`N(s)`$ is a 2-sided nilpotent ideal of $`D(s)`$ and $`D(s)N(s)=0.`$
3) $`D(𝒪_s)=_{i=1}^kD(sg_i)`$ is a direct sum of left ideals.
4) Each $`D(𝒪_s)`$ is a 2-sided ideal of $`𝒜_\alpha (G,𝒮)`$ and $`𝒜_\alpha (G,𝒮)=_{jJ}D(𝒪_{s_j}).`$ Moreover, $`D(𝒪_s)`$ has identity element $`_{t𝒪_s}1e(t).`$
Proof: 1) Clearly, $`S(s)`$ is a subalgebra of $`𝒜_\alpha (G,𝒮).`$ Let $`\rho :S(s)C^{\alpha _s}[G_s]`$ be a linear isomorphism determined by $`\rho (ae(s))=a`$. Then
$`\rho (ae(s)be(s))`$
$`=\rho (\alpha _s(a,b)abe(sb)e(s))`$
$`=\rho (\alpha _s(a,b)abe(s))`$
$`=\alpha _s(a,b)ab`$
$`=\rho (ae(s))\rho (be(s)).`$
So, $`\rho `$ is an algebra isomorphism.
2) Since for $`bGG_s,`$ $`sbs`$. We immediately have $`ae(s)be(s)=\alpha _s(a,b)abe((s)b)e(s)=0`$ for $`ae(s)D(s)`$ and $`be(s)N(s)`$.
It remains to show that $`N(s)`$ is right ideal. Take $`ae(s)N(s)`$ and $`be(s)S(s).`$ Then $`ae(s)be(s)=\alpha _s(a,b)abe((s)b)e(s)=\alpha _s(a,b)abe(s).`$ If $`ab`$ lies in $`G_s`$ so does $`a`$. Thus $`abGG_s`$ and $`ae(s)be(s)N(s)`$.
3)-4) are clear. $`\mathrm{}`$
Let $`A`$ be an algebra (eg., associative algebra, Lie algebra or vertex operator algebra). We denote the module category of $`A`$ by $`A`$-Mod.
For convenience we set $`d_{g,s}=ge(s)`$ for $`gG`$ and $`sS.`$
###### Theorem 3.5
The functors
$`M`$ $`\stackrel{f}{}D(s)`$ $`{\displaystyle \underset{S(s)}{}}M`$
$`N`$ $`\stackrel{g}{}`$ $`d_{1,s}N`$
($`MC^{\alpha _s}[G_s]`$-Mod, $`ND(𝒪_s)`$-Mod) define an equivalence between categories $`C^{\alpha _s}[G_s]`$-Mod and $`D(𝒪_s)`$-Mod. In particular, the simple $`C^{\alpha _s}[G_s]`$-modules are mapped to simple $`D(𝒪_s)`$-modules and conversely.
Proof: Recall that $`G=_{i=1}^kG_sg_i`$ be a coset decomposition with $`g_1=1.`$ Then $`G=_{i=1}^kg_i^1G_s.`$ Note that $`g_i^1ae(s)`$ can be rewritten as $`\alpha _s(g_i^1,a)^1g_i^1e(s)ae(s)`$ for any $`i`$ and $`aG_s.`$ Let $`MC^{\alpha _s}[G_s]`$-Mod. Then an arbitrary vector in $`D(s)_{S(s)}M`$ has expression $`_{i=1}^kd_{g_i^1,s}_{S(s)}m_i`$ for some $`m_iM.`$ From
$$d_{1,s}(\underset{i=1}{\overset{k}{}}d_{g_i^1,s}_{S(s)}m_i)=d_{1,s}_{S(s)}m_1$$
we see that $`d_{1,s}D(s)_{S(s)}M`$ is a subset of $`d_{1,s}_{S(s)}M.`$ On the other hand, $`d_{1,s}_{S(s)}m=d_{1,s}(d_{1,s}_{S(s)}m)`$ belongs to $`d_{1,s}(D(s)_{S(s)}M)`$. Thus $`d_{1,s}_{S(s)}M`$ and $`d_{1,s}(D(s)_{S(s)}M)`$ are equal. Since $`d_{1,s}_{S(s)}M`$ is isomorphic to $`M`$ as $`C^{\alpha _s}[G_s]`$-modules . This implies that $`gf(M)`$ and $`M`$ are isomorphic as $`C^{\alpha _s}[G_s]`$-modules.
Now take $`ND(𝒪_s)`$-Mod. We are going to show that $`D(s)_{S(s)}d_{1,s}N`$ and $`N`$ are isomorphic as $`D(𝒪_s)`$-modules. For this purpose we define a linear map
$`\rho :D(s){\displaystyle \underset{S(s)}{}}d_{1,s}N`$ $``$ $`N`$
$`d_{b,s}_{S(s)}d_{1,s}n`$ $``$ $`d_{b,s}n.`$
Noting that $`\rho (d_{b,s}_{S(s)}d_{1,s}n)=d_{b,s}d_{1,s}n`$ we immediately see that $`\rho `$ is well defined. The fact that $`\rho `$ is a $`D(𝒪_s)`$-module homomorphism follows from
$`\rho (d_{a,sg_i}d_{b,s}_{S(s)}d_{1,s}n)`$
$`=\rho (\delta _{sg_ib,s}\alpha _s(a,b)d_{ab,s}_{S(s)}d_{1,s}n)`$
$`=\delta _{sg_ib,s}\alpha _s(a,b)d_{ab,s}n`$
and
$`d_{a,sg_i}\rho (d_{b,s}_{S(s)}d_{1,s}n)`$
$`=d_{a,sg_i}d_{b,s}n`$
$`=\delta _{sg_ib,s}\alpha _s(a,b)d_{ab,s}n`$
for $`a,bG,`$ $`1ik`$ and $`nN.`$
It remains to show that $`\rho `$ is a bijection. As $`_{i=1}^kd_{1,sg_i}`$ is the identity of $`D(𝒪_s),`$ we have $`n=_{i=1}^kd_{1,sg_i}n`$. For fixed $`i`$ one can easily see that
$$\rho (d_{g_i^1,s}_{S(s)}\alpha _{sg_i}(g_i^1,g_i)^1d_{g_i,sg_i}n)=d_{1,sg_i}n.$$
Thus $`\rho `$ is onto.
Since $`d_{1,sg_i}d_{1,sg_j}=\delta _{i,j}d_{1,sg_i}`$ we have
$$N=\underset{i=1}{\overset{k}{}}d_{1,sg_i}N.$$
Let $`b=g_i^1ag_i^1G_s.`$ Then $`\rho (d_{b,s}_{S(s)}d_{1,s}n)=d_{1,sg_i}d_{b,s}nd_{1,sg_i}N`$ for $`nN.`$ In order to prove that $`\rho `$ is one to one, it is enough to show that if $`\rho (_pd_{b_p,s}_{S(s)}d_{1,s}n_p)=0`$ for some $`b_pg_i^1G_s`$ and $`n_pN,`$ then $`_pd_{b_p,s}_{S(s)}d_{1,s}n_p=0.`$ Using the relation $`d_{g_i^1a,s}=\alpha _s(g_i^1,a)^1d_{g_i^1,s}d_{a,s}`$ for $`aG_s`$ we can rewrite $`_pd_{b_p,s}_{S(s)}d_{1,s}n_p`$ as $`d_{g_i^1,s}_{S(s)}d_{1,s}n`$ for some $`nN.`$ Thus $`\rho (d_{g_i^1,s}_{S(s)}d_{1,s}n)=d_{g_i^1,s}n=0.`$ Applying $`\alpha _s(g_i,g_i^1)^1d_{g_i,sg_i}`$ to $`d_{g_i^1,s}n`$ yields $`d_{1,s}n=0.`$ This shows that $`\rho `$ is one to one. This completes the proof that $`\rho `$ is an isomorphism. $`\mathrm{}`$
###### Theorem 3.6
We have
i) Algebra $`D(𝒪_s)`$ is semisimple for $`s𝒮`$ and simple $`D(𝒪_s)`$-modules are precisely $`\mathrm{Ind}_{S(s)}^{D(s)}M`$ where $`M`$ ranges over the simple $`C^{\alpha _s}[G_s]`$-modules.
ii) $`A_\alpha (G,𝒮)`$ is semisimple and simple $`A_\alpha (G,𝒮)`$-modules are precisely $`\mathrm{Ind}_{S(s_j)}^{D(s_j)}M`$ where $`M`$ ranges over the simple $`C^{\alpha _{s_j}}[G_{s_j}]`$-modules and $`jJ.`$
Proof: i) follows from Theorem 3.5 and that fact that $`C^{\alpha _s}[G_s]`$-Mod is a semisimple category. ii) follows from Lemma 3.4 and i). $`\mathrm{}`$
###### Remark 3.7
If we only wanted to know the semisimplicity of $`A_\alpha (G,𝒮),`$ we could find it in the literature of Hopf algebra when we regard $`A_\alpha (G,𝒮)`$ as a crossed product (see Remark 3.2). But in the later sections we need to know the explicit structure of simple $`A_\alpha (G,𝒮)`$-modules given in Theorem 3.5. The semisimplicity of $`A_\alpha (G,𝒮)`$ is a trivial corollary of Theorem 3.5.
Next we determine the center $`Z(D(𝒪_s))`$ of $`D(𝒪_s)`$. However, we are not going to use it in this paper but still include it for general interest. We certainly expect that the result on center will be used in our future study on general orbifold conformal field theory.
Since any cocycle in $`Z^2(G_s,C^{})`$ is cohomologous to a normal cocycle (cf. Lemma 2.8). Then we may assume $`\alpha _s`$ is a normal cocycle for all $`sS.`$ We also assume that $`\alpha _{sg_i}(g_j^1hg_i,g_i^1ag_i)=\alpha _{sg_i}(g_j^1hah^1g_j,g_j^1hg_i)`$ for all $`1i,jk,`$ $`hG_s`$ and $`\alpha _s`$-regular $`aG_s`$. In the orbifold theory, these conditions are satisfied by chosing $`\alpha _s`$ carefully.
Let $`\{l_1,\mathrm{},l_r\}`$ be a full set of representatives for the $`\alpha _s`$-regular conjugacy classes of $`G_s`$ and for each $`t`$ in $`\{1,\mathrm{},r\}`$, let $`L_t`$ be $`\alpha _s`$-regular conjugacy class of $`G_s`$ containing $`l_t,1tr.`$ Set
$$Z(L_t)=\underset{aL_t}{}\underset{i=1}{\overset{k}{}}g_i^1ag_ie(sg_i).$$
###### Lemma 3.8
$`Z(L_t)`$ is an center element of $`D(𝒪_s).`$
Proof: Let $`be(sg_j)D(𝒪_s)`$. Then there exists $`1i^{}k`$ and $`hG_s`$ such that $`b=g_i^{}^1hg_j.`$ We have
$`be(sg_j)Z(L_t)`$
$`=be(sg_j){\displaystyle \underset{aL_t}{}}{\displaystyle \underset{i=1}{\overset{k}{}}}g_i^1ag_ie(sg_i)`$
$`={\displaystyle \underset{aL_t}{}}{\displaystyle \underset{i=1}{\overset{k}{}}}\alpha _{sg_i}(b,g_i^1ag_i)bg_i^1ag_ie(sg_jg_i^1ag_i)e(sg_i)`$
$`={\displaystyle \underset{aL_t}{}}\alpha _{sg_j}(b,g_j^1ag_j)bg_j^1ag_je(sg_j)`$
$`={\displaystyle \underset{aL_t}{}}\alpha _{sg_j}(g_i^{}^1hg_j,g_j^1ag_j)g_i^{}^1hag_je(sg_j)`$
and
$`Z(L_t)be(sg_j)`$
$`={\displaystyle \underset{aL_t}{}}{\displaystyle \underset{i=1}{\overset{k}{}}}g_i^1ag_ie(sg_i)be(sg_j)`$
$`={\displaystyle \underset{aL_t}{}}{\displaystyle \underset{i=1}{\overset{k}{}}}\alpha _{sg_j}(g_i^1ag_i,b)g_i^1ag_ibe(sg_ib)e(sg_j)`$
$`={\displaystyle \underset{aL_t}{}}\alpha _{sg_j}(g_i^{}^1ag_i^{},g_i^{}^1hg_j)g_i^{}^1ahg_je(sg_j).`$
Now the assertion follows from
$`{\displaystyle \underset{aL_t}{}}\alpha _{sg_j}(g_i^1hg_j,g_j^1ag_j)hah^1`$
$`={\displaystyle \underset{aL_t}{}}\alpha _{sg_j}(g_i^1hah^1g_i,g_i^1hg_j)hah^1`$
$`={\displaystyle \underset{aL_t}{}}\alpha _{sg_j}(g_i^1ag_i,g_i^1hg_j)a.`$
$`\mathrm{}`$
###### Theorem 3.9
Let $`\{l_1,\mathrm{},l_r\}`$ be a full set of representatives for the $`\alpha _s`$-regular conjugacy classes of $`G_s`$ and for each $`t\{1,\mathrm{},r\}`$, let $`L_t`$ be a $`\alpha _s`$-regular conjugacy class of $`G_s`$ containing $`l_t,1tr`$. Then the elements $`Z(L_t)`$ constitute a $`C`$-basis of $`Z(D(𝒪_s))`$.
Proof: By Theorem 3.5 the dimension of $`Z(D(𝒪_s))`$ equals to the number of inequivalent irreducible $`C^{\alpha _s}[G_s]`$-modules which is $`r`$ (see Theorem 2.9). $`\mathrm{}`$
The following corollary is immediate.
###### Corollary 3.10
Let $`\{l_1^j,\mathrm{},l_{r_j}^j\}`$ be a full set of representatives for the $`\alpha _{s_j}`$-regular conjugacy classes of $`G_{s_j}`$ and for each $`t_j\{1,\mathrm{},r_j\}`$, let $`L_{t_j}`$ be $`\alpha _{s_j}`$-regular conjugacy class of $`G_{s_j}`$ containing $`l_{t_j},1t_jr_j.`$ Set
$$Z(L_{t_j})=\underset{aL_{t_j}}{}\underset{i_j=1}{\overset{k_j}{}}g_{i_j}^1ag_{i_j}e(sg_{i_j}).$$
Then for all $`jJ`$, for all $`t_j\{1,\mathrm{},r_j\}`$, $`Z(L_{t_j})`$ constitutes a $`C`$-basis for $`Z(A_\alpha (G,𝒮)).`$
## 4 Modules for vertex operator algebras and related results
In this section we turn our attention to the theory of vertex operator algebras. In particular we shall defines various notion of modules for a vertex operator algebra $`V`$ following \[FLM\], \[DLM2\] and \[Z\]. We also recall from \[DLM4\] the associative algebras $`A_n(V)`$ for any nonnegative integer $`n`$ and relevant results. These results will be used extensively in Sections 5 and 6 to study dual pairs arising from an action of a finite group $`G`$ on $`V.`$
Let $`V=(V,Y,\mathrm{𝟏},\omega )`$ be a vertex operator algebra (see \[B\] and \[FLM\]). For a vector space $`W`$, let $`W\{z\}`$ be the space of $`W`$-valued formal series in arbitrary complex powers of $`z`$. We present three different notion of modules (cf. \[FLM\], \[DLM2\] and \[Z\]).
###### Definition 4.1
A weak $`V`$-module $`M`$ is a vector space equipped with a linear map
$$\begin{array}{ccc}V\hfill & \hfill & (\mathrm{End}M)\{z\}\hfill \\ v\hfill & \hfill & Y_M(v,z)=\underset{nQ}{}v_nz^{n1}(v_n\mathrm{End}M)\hfill \end{array}$$
which satisfies the following properties for all $`uV`$, $`vV,`$ $`wM`$,
$`u_lw=0\text{for}l>>0`$ (4.1)
$`Y_M(\mathrm{𝟏},z)=1;`$ (4.2)
$$\begin{array}{c}z_0^1\delta \left(\frac{z_1z_2}{z_0}\right)Y_M(u,z_1)Y_M(v,z_2)z_0^1\delta \left(\frac{z_2z_1}{z_0}\right)Y_M(v,z_2)Y_M(u,z_1)\\ =z_2^1\delta \left(\frac{z_1z_0}{z_2}\right)Y_M(Y(u,z_0)v,z_2),\end{array}$$
(4.3)
where $`\delta (z)=_{nZ}z^n`$ and all binomial expressions (here and below) are to be expanded in nonnegative integral powers of the second variable. Elementary properties of the $`\delta `$-function can be found in \[FLM\] and \[FHL\]
(4.3) is called the Jacobi identity. One can prove that the Jacobi identity is equivalent to the following associativity formula (see \[FLM\] and \[FHL\])
$`(z_0+z_2)^kY_M(u,z_0+z_2)Y_M(v,z_2)w=(z_2+z_0)^kY_M(Y(u,z_0)v,z_2)w`$ (4.4)
and commutator relation
$`[Y_M(u,z_1),Y_M(v,z_2)]`$
$`=\mathrm{Res}_{z_0}z_2^1\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)Y_M(Y(u,z_0)v,z_2)`$ (4.5)
where $`wM`$ and $`k`$ is a nonnegative integer such that $`z^kY_M(u,z)w`$ involves only nonnegative integral powers of $`z.`$
We may also deduce from (4.1)-(4.3) the usual Virasoro algebra axioms (see \[DLM2\]). Namely, if $`Y_M(\omega ,z)=_{nZ}L(n)z^{n2}`$ then
$$[L(m),L(n)]=(mn)L(m+n)+\frac{1}{12}(m^3m)\delta _{m+n,0}(\text{rank}V)$$
(4.6)
and
$$\frac{d}{dz}Y_M(v,z)=Y_M(L(1)v,z).$$
(4.7)
Suppose that $`(M_i,Y_i)`$ are two weak $`V`$-modules, $`i`$=1,2. A homomorphism from $`M_1`$ to $`M_2`$ is a linear map $`f`$:$`M_1M_2`$ which satisfied $`fY_{M_1}(v,z)=Y_{M_2}(v,z)f`$ for all $`vV`$. We call $`f`$ an isomorphism if $`f`$ is also a linear isomorphism.
###### Definition 4.2
A (ordinary) $`V`$-module is a weak $`V`$-module $`M`$ which carries a $`C`$-grading induced by the spectrum of $`L(0)`$. Then $`M=_{\lambda C}M_\lambda `$ where $`M_\lambda =\{wM|L(0)w=\lambda w\},`$ dim$`M_\lambda <\mathrm{}.`$ Moreover, for fixed $`\lambda ,M_{n+\lambda }=0`$ for all small enough integers $`n.`$
The notion of module here is essentially the notion of module given in \[FLM\].
###### Definition 4.3
An admissible $`V`$-module is a weak $`V`$-module $`M`$ which carries a $`Z_+`$ grading $`M=_{nZ_+}M(n)`$ satisfying the following condition:
$$v_mM(n)M(n+\mathrm{wt}vm1)$$
for homogeneous $`vV,`$ where $`Z_+`$ is the set of the nonnegative integers.
The notion of admissible module here is the notion of module in \[Z\]. Using a grading shift we can always arrange the grading on $`M`$ so that $`M(0)0.`$ This shift is important in the study of algebra $`A_n(V)`$ below.
It is not too hard to see that any $`V`$-module is an admissible $`V`$-module. So there is a natural identification of the category of $`V`$-modules with a subcategory of the category of admissible $`V`$-modules.
###### Definition 4.4
$`V`$ is called rational if every admissible $`V`$-module is a direct sum of irreducible admissible $`V`$-modules.
It is proved in \[DLM3\] that if $`V`$ is rational then there are only finitely many inequivalent irreducible admissible $`V`$-modules and each irreducible admissible $`V`$-module is an ordinary module.
The following proposition can be found in \[L\] and \[DM2\].
###### Proposition 4.5
If $`M`$ is a simple weak $`V`$module then $`M`$ is spanned by $`\{u_nm|uV,nQ\}`$ where $`mM`$ is a fixed nonzero vector.
We now recall the associative algebra $`A_n(V)`$ as constructed in \[DLM4\]. Let $`O_n(V)`$ be the linear span of all $`u_nv`$ and $`L(1)u+L(0)u`$ where for homogeneous $`uV`$ and $`vV,`$
$$u_nv=\mathrm{Res}_zY(u,z)v\frac{(1+z)^{\mathrm{wt}u+n}}{z^{2n+2}}.$$
(4.8)
Define the linear space $`A_n(V)`$ to be the quotient $`V/O_n(V).`$ We also define a second product $`_n`$ on $`V`$ for $`u`$ and $`v`$ as above:
$$u_nv=\underset{m=0}{\overset{n}{}}(1)^m\left(\genfrac{}{}{0pt}{}{m+n}{n}\right)\mathrm{Res}_zY(u,z)\frac{(1+z)^{\mathrm{wt}u+n}}{z^{n+m+1}}v.$$
(4.9)
Extend linearly to obtain a bilinear product on $`V`$.
Let $`M=_{nZ_+}M(n)`$ be an admissible $`V`$-module. Following \[Z\] we define weight zero operator $`o_M(v)=v_{\mathrm{wt}v1}`$ on $`M`$ for homogeneous $`v`$ and extend $`o_M(v)`$ to all $`v`$ by linearity. It is clear from the definition that $`o_M(v)M(n)M(n)`$ for all $`n.`$
###### Theorem 4.6
Let $`V`$ be a vertex operator algebra and $`M`$ an admissible $`V`$-module. Then
1) The product $`_n`$ induces the structure of an associative algebra on $`A_n(V)`$ with identity $`\mathrm{𝟏}+O_n(V).`$ Moreover $`\omega +O_n(V)`$ is a central element of $`A_n(V).`$
2) For $`0mn,`$ the map $`\psi _n:v+O_n(V)o_M(v)`$ from $`A_n(V)`$ to $`\mathrm{End}M(m)`$ makes $`M(m)`$ an $`A_n(V)`$-module.
3) $`M`$ is irreducible if and only if $`M(n)`$ is a simple $`A_n(V)`$-module for all $`n.`$
4) The identity map on $`V`$ induces an onto algebra homomorphism from $`A_n(V)`$ to $`A_m(V)`$ for $`0mn.`$
5) Two irreducible admissible $`V`$-modules $`M^1`$ and $`M^2`$ with $`M^1(0)0`$ and $`M^2(0)0`$ are isomorphic if and only if $`M^1(n)`$ and $`M^2(n)`$ are isomorphic $`A_n(V)`$-modules for any $`n0`$.
Parts 1)-2) and 4)-5) are proved in \[DLM4\], and 3) is given in \[DM4\].
## 5 Dual pair I
In the next two sections we assume that $`V`$ is a simple vertex operator algebra. Our main goal is to generalize the duality result obtained in \[DLM1\] on $`V`$ to an arbitrary irreducible $`V`$-module. More precisely, let $`G`$ be a finite automorphism group of $`V`$ and $`M`$ an irreducible $`V`$-module. Then $`M`$ is a module for the vertex operator subalgebra $`V^G`$ of $`G`$-invariants. There are three questions one can ask: (1) Is $`M`$ is a completely reducible $`V^G`$-module? (2) What are irreducible modules for $`V^G`$ which occur as submodules of $`M\mathrm{?}`$ (3) What are the relation between irreducible $`V^G`$-submodules of $`M`$ and $`V^G`$-submodules of another irreducible $`V`$-module $`N\mathrm{?}`$ These questions will be answered completely in this section and the next section.
In this section we deal with the case that $`M`$ is $`G`$-stable (see the definition below). The general case will be treated in the next section. The basis ideas in the proof of main theorems come from \[DLM1\].
Let $`(M,Y_M)`$ be an irreducible $`V`$-module and $`gG`$. Following \[DM1\] we set $`Mg=M`$ as vector spaces, and $`Y_{Mg}(v,z)=Y_M(gv,z)`$. Note that $`Mg`$ is also an irreducible $`V`$-module. We assume in this section that $`M`$ is $`G`$-stable in the sense that for any $`gG,`$ $`Mg`$ and $`M`$ are isomorphic. If $`M=V`$ this assumption is always true.
For $`gG`$ there is a linear isomorphism $`\varphi (g):MM`$ satisfying
$$\varphi (g)Y_M(v,z)\varphi (g)^1=Y_{Mg}(v,z)=Y_M(gv,z)$$
(5.10)
for $`vV.`$ The simplicity of $`M`$ together with Schur’s lemma shows that $`g\varphi (g)`$ is a projective representation of $`G`$ on $`M.`$ Let $`\alpha _M`$ be the corresponding 2-cocycle in $`C^2(G,C^\times ).`$ Then $`M`$ is a module for $`C^{\alpha _M}[G]`$ where $`C^{\alpha _M}[G]`$ is the twisted group algebra. It is worth pointing out that if $`M=V`$ we can take $`\varphi (g)=g`$ and $`C^{\alpha _M}[G]=C[G].`$
Let $`\mathrm{\Lambda }_{G,\alpha _M}`$ be the set of all irreducible characters $`\lambda `$ of $`C^{\alpha _M}[G]`$. We denote the corresponding simple module by $`W_\lambda .`$ Again if $`M=V`$, $`\mathrm{\Lambda }_{V,\alpha _V}=\widehat{G}`$ is the set of all irreducible characters of $`G.`$ Note that $`M`$ is a semisimple $`C^{\alpha _M}[G]`$-module. Let $`M^\lambda `$ be the sum of simple $`C^{\alpha _M}[G]`$-submodules of $`M`$ isomorphic to $`W_\lambda .`$ Then $`M=_{\lambda \mathrm{\Lambda }_{G,\alpha _M}}M^\lambda `$. Moreover, $`M^\lambda =W_\lambda M_\lambda `$ where $`M_\lambda =\mathrm{hom}_{C^{\alpha _M}[G]}(W_\lambda ,M)`$ is the multiplicity of $`W_\lambda `$ in $`M.`$ As in \[DLM1\], we can, in fact, realize $`M_\lambda `$ as a subspace of $`M`$ in the following way. Let $`wW_\lambda `$ be a fixed nonzero vector. Then we can identify $`\mathrm{hom}_{C^{\alpha _M}[G]}(W_\lambda ,M)`$ with the subspace
$$\{f(w)|f\mathrm{hom}_{C^{\alpha _M}[G]}(W_\lambda ,M)\}$$
of $`M^\lambda .`$
We also set $`V^G=\{vV|gv=v,gG\}.`$ Then $`V^G`$ is a simple vertex operator subalgebra of $`V`$ (see \[DM2\]) and $`\varphi (g)Y_M(v,z)\varphi (g)^1=Y_M(gv,z)=Y_M(v,z)`$ for all $`vV^G.`$ Thus the actions of $`C^{\alpha _M}[G]`$ and $`V^G`$ on $`M`$ are commutative. This shows that both $`M^\lambda `$ and $`M_\lambda `$ are ordinary $`V^G`$-modules.
###### Lemma 5.1
If $`M_\lambda 0`$ then $`M_\lambda `$ is an irreducible $`V^G`$-module.
Proof: By Theorem 4.6 3) it is enough to prove that $`M_\lambda (n)=M_\lambda M(n)`$ is a simple $`A_n(V^G)`$-module for all $`n0.`$ It is equivalent to show that for any $`nZ,n0`$ $`M_\lambda (n)`$ is generated by any nonzero vector in $`M_\lambda (n)`$ as an $`A_n(V^G)`$-module.
Note from the definition of $`A_n(V)`$ that $`A_n(V)`$ is a $`G`$-module via $`g(v+O_n(V))=gv+O_n(V)`$ for $`gG`$ and $`vV.`$ Also $`\mathrm{End}M(n)`$ is an $`G`$-module via $`gf=\varphi (g)f\varphi (g)^1`$ for $`gG`$ and $`f\mathrm{End}M(n).`$ Then the algebra homomorphism $`v+O_n(V)o_M(v)`$ from $`A_n(V)`$ to $`\mathrm{End}M(n)`$ is a $`G`$-homomorphism as $`\varphi (g)o_M(v)\varphi (g)^1=o_M(gv).`$ From Proposition 4.5 and Theorem 4.6 3) we see that $`\mathrm{End}M(n)=\{o_M(v)|vV\}.`$ Since $`G`$ is a finite group we immediately have $`\psi (A_n(V)^G)=(\mathrm{End}M(n))^G=\{o_M(v)|vV^G\}.`$
Let $`x,yM_\lambda (n)`$ be linearly independent. Since $`M(n)`$ is a $`C^{\alpha _M}[G]`$-module, we can write $`M(n)`$ as a direct sum $`W_\lambda xW_\lambda yW`$ where $`W`$ is a $`C^{\alpha _M}[G]`$-submodule of $`M(n).`$ Define a map $`\beta \mathrm{End}M(n)`$ such that $`\beta (ux+vy+w)=vx+uy+w`$ for $`u,vW_\lambda `$ and $`wW.`$ Then $`\beta (\mathrm{End}M(n))^G.`$ Thus there exists vectors $`vV^G`$ such that $`o_M(v)x=y.`$ This implies that $`M_\lambda (n)`$ is an simple $`A_n(V)^G`$-module. Since $`M_\lambda (n)`$ is also an $`A_n(V^G)`$-module with the same action, we see that $`M_\lambda (n)`$ is a simple $`A_n(V^G)`$-module. This can be also explained by noting that the identity on $`V^G`$ induces an onto algebra homomorphism from $`A_n(V^G)`$ to $`A_n(V)^G.`$ $`\mathrm{}`$
###### Remark 5.2
If we take $`M=V`$ in Lemma 5.1, then each $`V_\lambda `$ is an irreducible $`V^G`$-module for any irreducible character $`\lambda `$ of $`G.`$ This is the part (1) of Theorem 2.4 in \[DLM1\], where Lemma 2.2 of \[DLM1\] is used. It is pointed out to us by Wassermann that there is a gap in Lemma 2.2 of \[DLM1\]. The special case of Lemma 5.1 now fixes the problem in \[DLM1\]. (The group $`G`$ in \[DLM1\] can be a compact group. But the correction present here works for such $`G`$.) The proofs for parts (2) and (3) of Theorem 2.4 in \[DLM1\] are correct.
For the purpose of continuous discussion we now recall Theorem 2.4 from \[DLM1\] (also see Remark 5.2).
###### Theorem 5.3
Let $`V`$ be a simple vertex operator algebra and $`G`$ a finite automorphism group. Then
1) $`V^\lambda `$ is nonzero for any $`\lambda \widehat{G}.`$
2) Each $`V_\lambda `$ is an irreducible $`V^G`$-module.
3) $`V_\lambda `$ and $`V_\mu `$ are equivalent $`V^G`$-module if and only if $`\lambda =\mu .`$
The main result in this section is the following.
###### Theorem 5.4
With the same notation as above we have:
1) $`M^\lambda `$ is nonzero for any $`\lambda \mathrm{\Lambda }_{G,\alpha _M}.`$
2) Each $`M_\lambda `$ is an irreducible $`V^G`$-module.
3) $`M_\lambda `$ and $`M_\gamma `$ are equivalent $`V^G`$-module if and only if $`\lambda =\gamma .`$
That is, $`V^G`$ and $`C^{\alpha _M}[G]`$ form a dual pair on $`M`$ in the sense of Howe (see \[H1\] and \[H2\]).
Let $`\gamma \widehat{G}`$ and $`\mu \mathrm{\Lambda }_{G,\alpha _M}`$ such that $`M^\mu 0.`$ By Theorem 5.3 we can regard $`W_\gamma `$ and $`W_\mu `$ as a $`C[G]`$-submodule of $`V`$ and a $`C^{\alpha _M}[G]`$-submodule of $`M`$, respectively. In fact, we can assume that $`W_\gamma `$ and $`W_\mu `$ are homogeneous subspaces of $`V`$ and $`M,`$ respectively. Following \[DM3\] we define a subspace $`Z_s`$ of $`M`$ for any $`sZ`$ by
$$Z_s=\{\underset{ms}{}u_mw|uW_\gamma ,wW_\mu \}.$$
Then it is easy to see that $`Z_s`$ is a $`C^{\alpha _M}[G]`$-submodule of $`M.`$ We also define a map
$$\psi _s:W_\gamma W_\mu Z_s$$
such that $`\psi _s(uw)=_{ms}u_mw.`$ Recall twisted group algebra from Section 2. In particular, $`\{\overline{g}|gG\}`$ is a basis of $`C^{\alpha _M}[G].`$ Then $`W_\gamma W_\mu `$ is a $`C^{\alpha _M}[G]`$-module with $`\overline{g}`$ acting as $`g\varphi (g)`$ and $`\psi _s`$ is a $`C^{\alpha _M}[G]`$-module homomorphism.
###### Lemma 5.5
For all sufficiently small $`s,`$ $`\psi _s`$ is a $`C^{\alpha _M}[G]`$-module isomorphism.
This result in the case $`M=V`$ is proved in \[DM2\] and \[DM3\]. The proof is similar in the general case. Here we give an outline of the proof and we refer the reader to \[DM2\] and \[DM3\] for the details. From the associativity formula (4.4) and commutator formula (4.5) we can prove the following fact. Let $`u^iV`$ for $`i=1,\mathrm{},n`$ and $`w^1,\mathrm{},w^nM`$ be linearly independent. Then $`_{i=1}^nY_M(u^i,z)w^i0`$ (see the proof of Lemma 3.1 of \[DM2\]; also see Proposition 11.9 of \[DL\]). Then following the proof of Lemma 2.2 of \[DM3\] for the case $`M=V`$ gives the result.
We now prove Theorem 5.4. Part 2) is Lemma 5.1. In order to prove 1) we first note that there exists $`\mu \mathrm{\Lambda }_{M,\alpha _M}`$ such that $`M^\mu 0.`$ Let $`\mu ^{}`$ be the character of $`C^{\alpha _M^1}[G]`$ dual to $`\mu .`$ That is, the corresponding $`C^{\alpha _M^1}[G]`$-module is exactly $`W_\lambda ^{}=\mathrm{Hom}_C(W_\lambda ,C).`$ Then $`\mu ^{}\lambda `$ is a character of $`G`$ for any $`\lambda \mathrm{\Lambda }_{G,\alpha _M}.`$ Let $`\gamma `$ be an irreducible character of $`G`$ such that $`\gamma \mu ^{}\lambda .`$ Thus
$$\mathrm{Hom}_{C^{\alpha _M}[G]}(\lambda ,\mu \gamma )=\mathrm{Hom}_{C[G]}(\mu ^{}\lambda ,\gamma )0.$$
By Lemma 5.5, for small enough $`s,`$ the submodule $`Z_s`$ of $`M`$ contains a submodule isomorphic to $`W_\lambda .`$ That is, $`M^\lambda 0`$ for all $`\lambda \mathrm{\Lambda }_{G,\alpha _M}.`$
Finally we prove 3). Let $`\lambda ,\gamma \mathrm{\Lambda }_{G,\alpha _M}`$ are different. We can take $`nZ`$ such that $`M^\lambda (n)=M^\lambda M(n)0.`$ Then $`M(n)`$ is a direct sum of $`C^{\alpha _M}[G]`$-modules
$$M(n)=M^\lambda (n)W$$
for some suitable $`C^{\alpha _M}[G]`$ submodule $`W`$ of $`M(n).`$
Define $`\beta \mathrm{End}M(n)`$ such that it is the identity on $`M^\lambda (n)`$ and zero on $`W.`$ Then $`\beta (\mathrm{End}M(n))^G`$. From the proof of lemma 5.1 there exists $`vV^G`$ such that $`o_M(v)=\beta `$ is the identity on $`M^\lambda (n)`$ and zero on $`W.`$ Thus, there is no $`A_n(V^G)`$-module homomorphism between $`M_\lambda (n)`$ and $`M_\gamma (n).`$ That is, $`M_\lambda (n)`$ and $`M_\gamma (n)`$ are inequivalent $`A_n(V^G)`$-modules.
Since $`M`$ is irreducible $`V`$-module there exists $`cC`$ such that $`M=_{n0}M_{c+n}`$ with $`M_c0`$ where $`M_{c+n}`$ is the eigenspace for $`L(0)`$ with eigenvalue $`c+n`$ (see Definition 4.2). Thus we can take $`M(n)=M_{c+n}.`$ If $`dim(M_\lambda )_{c+m}dim(M_\gamma )_{c+m}`$ it is clear that $`M_\lambda `$ and $`M_\gamma `$ are nonisomorphic $`V^G`$-module. Otherwise by Theorem 4.6, $`M_\lambda `$ is not isomorphic to $`M_\gamma `$ as $`V^G`$-modules. $`\mathrm{}`$
## 6 Dual pair II
Let $`V`$ be a simple vertex operator algebra as in the last section and $`G`$ finite automorphism group of $`V.`$ Let $`M`$ be an irreducible $`V`$-module. But we do not assume that $`M`$ is $`G`$-stable. We set
$$G_M=\{gG|MhM\}$$
which is a subgroup of $`G.`$ Recall Theorem 5.4. One of the main results in this section is the following result of Schur-Weyl type:
###### Theorem 6.1
With the same notation as above we have:
1) $`M^\lambda `$ is nonzero for any $`\lambda \mathrm{\Lambda }_{G_M,\alpha _M}.`$
2) Each $`M_\lambda `$ is an irreducible $`V^G`$-module.
3) $`M_\lambda `$ and $`M_\gamma `$ are equivalent $`V^G`$-module if and only if $`\lambda =\gamma .`$
That is, $`V^G`$ and $`C^{\alpha _M}[G_M]`$ form a dual pair on $`M.`$
This result generalizes Theorem 2.4 of \[DLM1\] (a case $`M=V`$) and sharpens Theorem 5.4. In particular, the result shows that $`M`$ is a completely reducible $`V^G`$-module.
In order to prove Theorem 6.1 we need a general setting and we will prove a stronger result.
###### Definition 6.2
A set $`𝒮`$ of irreducible $`V`$-modules is called stable if for any $`M𝒮`$ and $`gG`$ there exists $`N𝒮`$ such that $`MgN.`$
Assume that $`𝒮`$ is a finite set of inequivalent irreducible $`V`$-module which is $`G`$-stable. Since $`M(g_1g_2)`$ and $`(Mg_1)g_2`$ are isomorphic $`V`$-module for $`g_1,g_2G`$ and $`M𝒮`$ we can define an right action of $`G`$ on $`𝒮.`$ Let $`M𝒮`$ and $`gG`$ we define $`Mg=N`$ if $`MgN.`$ It is clear that this action makes $`𝒮`$ a right $`G`$-set.
###### Remark 6.3
1) For any irreducible $`V`$-module $`M`$ consider the coset decomposition $`G=_{i=1}^kG_Mg_i.`$ Then $`𝒮=\{Mg_i|i=1,\mathrm{},k\}`$ is a such right $`G`$-set which will be used to prove Theorem 6.1.
2) If $`V`$ is rational, a complete list of inequivalent irreducible $`V`$-modules forms a such right $`G`$-set.
Let $`M𝒮`$ and $`xG`$. Then there exists $`N𝒮`$ such that $`NMx.`$ That is, there is a linear map $`\varphi _N(x):NM`$ satisfying the condition: $`\varphi _N(x)Y_N(v,z)\varphi _N(x)^1=Y_M(xv,z).`$ By simplicity of $`N`$, there exists $`\alpha _N(y,x)C^{}`$ such that $`\varphi _M(y)\varphi _N(x)=\alpha _N(y,x)\varphi _N(yx).`$ Moreover, for $`x,y,zG`$ we have
$$\alpha _N(z,yx)\alpha _N(y,x)=\alpha _M(z,y)\alpha _N(zy,x).$$
As in Section 3 we set $`C𝒮=_{M𝒮}Ce(M)`$ and $`e(M)e(N)=\delta _{M,N}e(M)`$. Let $`U(C𝒮)=\{_{M𝒮}\lambda _Me(M)|\lambda _MC^{}\}`$ and $`\alpha (h,k)=_{M𝒮}\alpha _M(h,k)e(M).`$ It is easy to check that $`\alpha (hk,l)\alpha (h,k)^l=\alpha (h,kl)\alpha (k,l).`$ So $`\alpha Z^2(G,U(C𝒮))`$ is a 2-cocycle. Following Section 3, we construct an associative algebra $`A_\alpha (G,𝒮)`$ with multiplication defined by
$$ae(M)be(N)=\alpha _N(a,b)abe(Mb)e(N)$$
for $`a,bG`$ and $`M,N𝒮.`$
We define an action of $`A_\alpha (G,𝒮)`$ on $`_{N𝒮}N`$ in the following way: for $`M,N𝒮`$ and $`wN`$ we set
$$ae(M)w=\delta _{M,N}\varphi _M(a)w$$
where $`\varphi _N(a):NNa^1.`$
For $`N𝒮`$ we let $`𝒪_N=\{Ng|gG\}`$ be the orbit of $`N`$ under $`G.`$
###### Lemma 6.4
With the action defined above, $`_{N𝒮}N`$ becomes a module for $`A_\alpha (G,𝒮).`$ Moreover $`_{M𝒪_N}M`$ is a submodule for any $`N.`$
Proof: The proof is a straightforward computation: for $`a,bG`$ and $`M,N𝒮`$ and $`mN,`$
$`(be(L)ae(M))m`$
$`=(\delta _{La,M}\alpha _M(b,a)bae(M))m`$
$`=\delta _{La,M}\delta _{M,N}\alpha _M(b,a)\varphi _M(ba)m`$
$`=\delta _{La,M}\delta _{M,N}\varphi _{Ma^1}(b)\varphi _M(a)m`$
and
$`be(L)(ae(M)m)`$
$`=be(L)(\delta _{M,N}\varphi _M(a)m)`$
$`=\delta _{L,Ma^1}\delta _{M,N}\varphi _L(b)\varphi _M(a)m`$
$`=\delta _{L,Ma^1}\delta _{M,N}\varphi _{Ma^1}(b)\varphi _M(a)m.`$
Thus $`_{M𝒮}M`$ is $`A_\alpha (G,𝒮)`$-module. It is clear that $`_{M𝒪_N}M`$ is a submodule. $`\mathrm{}`$
Let $`M𝒮.`$ Recall from Section 3 that $`D(M)=ae(M)|aG`$ and $`S(M)=ae(M)|aG_M`$.
###### Proposition 6.5
Let $`N𝒮.`$ Then $`N`$ is a $`S(N)`$-module and there is an $`A_\alpha (G,𝒮)`$-modules isomorphism between $`D(N)_{S(N)}N`$ and $`_{M𝒪_N}M`$ determined by $`\mathrm{\Psi }:ae(N)m\varphi _N(a)m.`$
Proof: To prove $`N`$ is a $`S(N)`$-module, it is enough to prove that $`N`$ is $`S(N)`$-stable. But this clear as $`ae(N)w=\varphi _N(a)wN`$ for $`aG_N`$ and $`wN.`$
Next we prove that $`\mathrm{\Psi }`$ is well defined. We must show that $`\mathrm{\Psi }(d_{a,N}d_{b,N}w)=\mathrm{\Psi }(d_{a,N}d_{b,N}w)`$ for $`wN,`$ $`aG,bG_N`$ where $`d_{a,N}=ae(N).`$ This is clear as
$`\mathrm{\Psi }(d_{a,N}d_{b,N}w)`$
$`=\alpha _N(a,b)\mathrm{\Psi }(d_{ab,N}w)`$
$`=\alpha _N(a,b)\varphi _N(ab)w`$
$`=\varphi _N(a)\varphi _N(b)w`$
$`=\varphi _N(a)d_{b,N}w`$
$`=\mathrm{\Psi }(d_{a,N}d_{b,N}w).`$
The fact that $`\mathrm{\Psi }`$ is a module homomorphism follows from
$`\mathrm{\Psi }(d_{b,M}d_{a,N}w)`$
$`=\mathrm{\Psi }(\delta _{Ma,N}\alpha _N(b,a)d_{ba,N}w)`$
$`=\delta _{Ma,N}\alpha _N(b,a)\varphi _N(ba)w)`$
$`=\delta _{Ma,N}\varphi _{Na^1}(b)\varphi _N(a)w`$
and
$`d_{b,M}\mathrm{\Psi }(d_{a,N}w)=d_{b,M}\varphi _N(a)w`$
$`=\delta _{M,Na^1}\varphi _{Na^1}(b)\varphi _N(a)w`$
where $`a,bG.`$
In order to show that $`\mathrm{\Psi }`$ is a bijection we construct an inverse of $`\mathrm{\Psi }.`$
Suppose that $`G=_{i_N=1}^{k_N}G_Ng_{i_N}`$ is a coset decomposition. Then
$$\underset{M𝒪_N}{}M=\underset{i_N=1}{\overset{k_N}{}}\varphi _N(g_{i_N}^1)N.$$
Define
$`{\displaystyle \underset{M𝒪_N}{}}M`$ $`\stackrel{\chi }{}`$ $`D(N){\displaystyle \underset{S(N)}{}}N`$
$`{\displaystyle \underset{i_N=1}{\overset{k_N}{}}}\varphi _N(g_{i_N}^1)m_{i_N}`$ $``$ $`{\displaystyle \underset{i_N=1}{\overset{k_N}{}}}d_{g_{i_N}^1,N}m_{i_N}.`$
It is straightforward to show that $`\mathrm{\Psi }\chi =Id_{_{M𝒪_N}M}`$ and $`\chi \mathrm{\Psi }=Id_{D(N)_{S(N)}N}`$. So, $`D(N)_{S(N)}N`$ and $`_{M𝒪_N}M`$ are isomorphic as $`A_\alpha (G,)`$-modules. $`\mathrm{}`$
Set
$$=\underset{M𝒮}{}M.$$
###### Proposition 6.6
1) Every simple $`A_\alpha (G,𝒮)`$-module occurs as a submodule of $``$.
2) The actions of $`A_\alpha (G,)`$ and $`V^G`$ on $``$ commute.
Proof: Part 1) follows from Theorems 3.6, 5.4 and a natural identification of $`S(N)`$ with $`C^{\alpha _N}[G_N]`$ for $`N𝒮.`$ Part 2) follows from the relation
$$ae(M)\overline{Y}(v,z)=\overline{Y}(av,z)ae(M)$$
on $``$ for $`aG,`$ $`M𝒮`$ and $`vV`$ where we have used $`\overline{Y}(v,z)`$ denote the operator on $``$ which acts on $`M`$ by $`Y_M(v,z).`$ $`\mathrm{}`$
###### Lemma 6.7
Set $`B=G\times U(C𝒮)`$. Then $`B`$ is a group under multiplication defined by
$$(y,v)(x,u)=(yx,\alpha (y,x)v^xu)$$
This result is standard as $`\alpha `$ is a 2-cocycle.
###### Lemma 6.8
$``$ is a $`B`$-module via $`(x,_{M𝒮}\lambda _Me(M))=_{M𝒮}\varphi _M(x)\lambda _M.`$
This result is immediate by noting that the subset $`B^{}=\{au|aG,uU(C𝒮)\}`$ of $`A_\alpha (G,𝒮)`$ is a multiplicative group isomorphic to $`B`$ and the action of $`B`$ on $``$ is exactly the action of $`B^{}`$ on $`.`$
###### Lemma 6.9
For any $`(x,u)B`$, $`vV,`$ we have
$$(x,u)\overline{Y}(v,z)(x,u)^1=\overline{Y}(xv,z)$$
on $`.`$
Proof: Let $`b=(x,_{M𝒮}\lambda _Me(M)))B`$. Then
$`b\overline{Y}(v,z)=(x,{\displaystyle \underset{M𝒮}{}}\lambda _Me(M)))\overline{Y}(v,z)`$
$`=\overline{Y}(xv,z)(x,{\displaystyle \underset{M𝒮}{}}\lambda _Me(M)))`$
(see the proof of Proposition 6.6 2)) as required. $`\mathrm{}`$
###### Lemma 6.10
$`_{M𝒮}\mathrm{End}(M)`$ is a $`B`$-module via
$$(x,\underset{M𝒮}{}\lambda _Me(M))\underset{M𝒮}{}f_M=\underset{M𝒮}{}\varphi _M(x)f_M\varphi _M(x)^1.$$
Furthermore, for each orbit $`𝒪`$ and $`n0,`$ $`_{M𝒪}\mathrm{End}(M(n))`$ is a $`B`$-submodule. In particular, $`_{M𝒪}\mathrm{End}(M(n))`$ is a $`G`$-module.
Proof: Let $`(x,_{M𝒮}\lambda _Me(M))`$, $`(y,_{M𝒮}\beta _Me(M))B.`$ Since
$`[(y,{\displaystyle \underset{M𝒮}{}}\beta _Me(M))(x,{\displaystyle \underset{M𝒮}{}}\lambda _Me(M))]{\displaystyle \underset{M𝒮}{}}f_M`$
$`=(yx,{\displaystyle \underset{M𝒮}{}}\beta _{Mx^1}\lambda _M\alpha _M(y,x)e(M)){\displaystyle \underset{M𝒮}{}}f_M`$
$`={\displaystyle \underset{M𝒮}{}}\varphi _M(yx)f_M\varphi _M(yx)^1`$
$`={\displaystyle \underset{M𝒮}{}}\varphi _{Mx^1}(y)\varphi _M(x)f_M\varphi _M(x)^1\varphi _{Mx^1}(y)^1`$
and
$`(y,{\displaystyle \underset{M𝒮}{}}\beta _Me(M))[(x,{\displaystyle \underset{M𝒮}{}}\lambda _Me(M)){\displaystyle \underset{M𝒮}{}}f_M]`$
$`(y,{\displaystyle \underset{M𝒮}{}}\beta _Me(M)){\displaystyle \underset{M𝒮}{}}\varphi _M(x)f_M\varphi _M(x)^1`$
$`={\displaystyle \underset{M𝒮}{}}\varphi _{Mx^1}(y)\varphi _M(x)f_M\varphi _M(x)^1\varphi _{Mx^1}(y)^1,`$
$`_{M𝒮}\mathrm{End}(M)`$ is a $`B`$-module. The other assertions in the lemma are clear. $`\mathrm{}`$
###### Remark 6.11
It is worth to point out that $`\varphi _M(x)f_M\varphi _M(x)^1`$ does not lie in $`\mathrm{End}M`$ for $`f_M\mathrm{End}M`$ and $`xG`$ unless $`xG_M.`$ In general, $`\varphi _M(x)f_M\varphi _M(x)^1`$ is an element of $`\mathrm{End}(Mx^1).`$
For any nonnegative $`nZ,`$ let $`\sigma _n`$ be a map from $`A_n(V)`$ to $`_{M𝒮}_{0mn}\mathrm{End}M(m)`$ defined by $`\sigma _n(v+O_n(V))=_{M𝒮}o_M(v).`$
###### Lemma 6.12
The map $`\sigma _n`$ is a $`G`$-module epimorphism. In particular, $`\sigma (A_n(V)^G)=(_{M𝒮,0mn}\mathrm{End}M(m))^G`$
Proof: Let $`gG`$ and $`vV.`$ Then
$`\sigma _n(g(v+O_n(v)))=\sigma _n(gv+O_n(V))`$
$`={\displaystyle \underset{M𝒮}{}}o_M(gv)`$
$`={\displaystyle \underset{M𝒮}{}}\varphi _M(g)o_M(v)\varphi _M(g)^1.`$
That is, $`\sigma _n`$ is a $`G`$-homomorphism.
In order to see that $`\sigma _n`$ is onto we note that all $`M𝒮`$ are inequivalent. We assume that $`M(0)0`$ for all $`M𝒮.`$ By Theorem 4.6, $`\{M(m)|M𝒮,0mn\}`$ is a set of finite dimensional inequivalent $`A_n(V)`$-modules. Let $`K_{M(m)}`$ be the kernel of $`A_n(V)`$ on $`M(m).`$ Then $`A_n(V)/K_{M(m)}`$ is isomorphic to $`\mathrm{End}M(m)`$ and $`A_n(V)/K_n`$ is isomorphic to the direct sum $`_{M𝒮}_{0mn}\mathrm{End}M(m)`$ where $`K_n=_{M,m}K_{M(m)}.`$ Thus $`\sigma _n`$ is onto. $`\mathrm{}`$
Now we are in the position to state and to prove the main result in this paper. Let $`𝒮=_{jJ}𝒪_j`$ be an orbit decomposition and fix $`M^j𝒪_j.`$ For convenience, we set $`G_j=G_{M^j}`$ and $`\mathrm{\Lambda }_j=\mathrm{\Lambda }_{M^j,\alpha _{M^j}}.`$ Then by Theorem 5.4 we have a decomposition
$$M^j=\underset{\lambda \mathrm{\Lambda }_j}{}W_\lambda M_\lambda ^j$$
where $`M_\lambda ^j`$ is an irreducible $`V^{G_j}`$-module. Thanks to Propositions 6.5 and 6.6 we have
$$=\underset{jJ}{}\underset{\lambda \mathrm{\Lambda }_j}{}(\mathrm{Ind}_{S(M^j)}^{D(M^j)}W_\lambda )M_\lambda ^j$$
as a $`A_\alpha (G,𝒮)V^G`$-module. For $`jJ`$ and $`\lambda \mathrm{\Lambda }_j`$ we set $`W_{j,\lambda }=\mathrm{Ind}_{S(M^j)}^{D(M^j)}W_\lambda .`$ Then by Theorems 3.6, $`W_{j,\lambda }`$ forms a complete list of simple $`A_\alpha (G,𝒮)`$-modules.
###### Theorem 6.13
As a $`A_\alpha (G,𝒮)V^G`$-module,
$$=\underset{jJ,\lambda \mathrm{\Lambda }_j}{}W_{j,\lambda }M_\lambda ^j.$$
Moreover,
1) Each $`M_\lambda ^j`$ is a nonzero irreducible $`V^G`$-module.
2) $`M_{\lambda _1}^{j_1}`$ and $`M_{\lambda _2}^{j_2}`$ are isomorphic $`V^G`$-module if and only if $`j_1=j_2`$ and $`\lambda _1=\lambda _2.`$
That is, $`(A_\alpha (G,𝒮),V^G)`$ forms a dual pair on $`.`$
Proof: 1) We have already seen from Theorem 5.4 that each $`M_\lambda ^j`$ is nonzero. Again by Theorem 4.6 3) we only need to show that each $`M_\lambda ^j(n)=M^j(n)M_\lambda ^j`$ is a simple $`A_n(V^G)`$-module.
We now fix $`j.`$ Suppose $`G=_{i=1}^kG_jg_i`$ be the coset decomposition such that $`g_1=1.`$ For each $`vV^{G_j}`$ we set $`f^v_{M𝒮}\mathrm{End}M(n)`$ such that $`f^v`$ acts on $`M^jg_i^1`$ as
$$\varphi _{M^j}(g_i)o_{M^j}(v)\varphi _{M^j}(g_i)^1=o_{M^jg_i^1}(g_iv)$$
for $`i=1,\mathrm{},k,`$ acts on any other $`M(n)`$ as zero. Then it is clear from Lemma 6.10 that $`f^v(_{M𝒮}\mathrm{End}M(n))^G.`$ By Lemma 6.12 there exists $`uV^G`$ such that $`f^u=f^v.`$ This shows that for any $`vV^{G_j}`$ there exists $`uV^G`$ such that $`o_{M^j}(v)=o_{M^j}(u)`$ on $`M^j(n).`$ Since $`M_\lambda ^j(n)`$ is a simple module for $`A_n(V^{G_j})`$ by Theorems 4.6 and 5.4 we see immediately that $`M_\lambda ^j(n)`$ is a simple $`A_n(V^G)`$-module.
2) We take $`(j_1,\lambda _1)(j_2,\lambda _2).`$ We must to prove that $`M_{\lambda _1}^{j_1}`$ and $`M_{\lambda _2}^{j_2}`$ are inequivalent. Let $`n0`$ such that both $`M_{\lambda _1}^{j_1}(n)`$ and $`M_{\lambda _2}^{j_2}(n)`$ are nonzero. Then
$$\underset{0mn}{}(m)=\underset{0mn}{}W_{j_1,\lambda _1}M_{\lambda _1}^{j_1}(m)W$$
for a suitable $`A_\alpha (G,𝒮)`$-module $`W.`$ Define $`f`$ in $`_{M𝒮,0mn}\mathrm{End}M(m)`$ such that $`f=1`$ on $`_{0mn}W_{j_1,\lambda _1}M_{\lambda _1}^{j_1}(m)`$ and $`f=0`$ on $`W.`$ Again it is obvious that
$$f(\underset{M𝒮,0mn}{}\mathrm{End}M(m))^G.$$
Using Lemma 6.12 we find $`vV^G`$ such that $`\sigma _n(v)=f.`$ That is $`o_{M_{\lambda _1}^{j_1}}(v)=1`$ on $`_{0mn}M_{\lambda _1}^{j_1}(m)`$ and $`o_{M_{\lambda _2}^{j_2}}(v)=0`$ on $`_{0mn}M_{\lambda _2}^{j_2}(m).`$ That is, $`M_{\lambda _1}^{j_1}(s)`$ and $`M_{\lambda _2}^{j_2}(t)`$ are nonisomorphic $`A_n(V^G)`$-modules for $`0s,tn`$ if either $`M_{\lambda _1}^{j_1}(s)`$ or $`M_{\lambda _2}^{j_2}(t)`$ is nonzero. In particular, $`M_{\lambda _1}^{j_1}(s_0)`$ or $`M_{\lambda _2}^{j_2}(t_0)`$ are nonisomorphic $`A_n(V^G)`$-modules where $`s_0,t_00`$ such that $`M_{\lambda _1}^{j_1}(s)=0`$ and $`M_{\lambda _2}^{j_2}(t)=0`$ for all $`ss_0`$ and $`tt_0.`$ Our choices of $`s_0`$ and $`t_0`$ then assert that $`M_{\lambda _1}^{j_1}(s_0)`$ or $`M_{\lambda _2}^{j_2}(t_0)`$ are, in fact, inequivalent $`A_0(V^G)`$-modules. Thus by Theorem 4.6 $`M_{\lambda _1}^{j_1}`$ and $`M_{\lambda _2}^{j_2}`$ are inequivalent $`V^G`$-modules. $`\mathrm{}`$
Now Theorem 6.1 follows from Theorem 6.13 immediately by taking the right $`G`$-set $`𝒮`$ to be the $`G`$-orbit $`\{Mg_i|i=1,\mathrm{},l\}`$ where $`G=_{i=1}^lG_Mg_i`$ is the coset decomposition of $`G`$ with respect to the stabilizer $`G_M=\{gG|MgM\}.`$
We end this paper with the following general discussion: If $`V`$ has only finitely many inequivalent irreducible modules (this is the case when $`V`$ is rational; see the discussion after Definition 4.4), then a complete list of irreducible $`V`$-modules is a right $`G`$-set. Theorem 6.13 then tells us not only $`=_{M𝒮}M`$ is completely reducible but also gives an equivalence between $`A_\alpha (G,𝒮)`$-module category and a subcategory of admissible $`V^G`$-modules generated by the irreducible submodules occurring in $``$ by sending $`W_{j,\lambda }`$ to $`M_\lambda ^j.`$ We expect to use this result to determine the module category for $`V^G`$ when $`V`$ is rational in the future.
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# 1 Introduction
## 1 Introduction
One of the oldest and deepest ideas in gauge theories is the conjectured duality between a loop description, based on the Wilson loops, or holonomies, of a gauge theory and a string description, in which the position of a string of quantized electric flux are taken as the fundamental coordinates. Taken into the gravitational context, this suggests that there should be a duality between string theory and loop quantum gravity, as the latter is based on the quantum dynamics of the Wilson loops of the spacetime connection-. In this paper we propose a specific form for such a string/loop duality, by finding evidence that a particular extension of loop quantum gravity is dual to the standard dWHN-BFSS and IKKT matrix models. We do this by arguing that both arise from different compactifications of a single matrix model, called the cubic matrix model.
In a recent paper the cubic matrix models were proposed, and we presented some evidence that one of them has compactifications which reproduce the IKKT and dWHN-BFSS matrix models. In this paper we study a new class of compactifications of these models, which lead directly to a background independent description of the theory. This turns out to be an extension of loop quantum gravity in which the usual $`SU(2)`$ algebra has been replaced by a supersymmetric and quantum extension of $`SU(16)`$. The $`SU(16)`$ means that the model extends the symmetries of the $`9`$ dimensional Clifford algebra, while the quantum group extension means that it is also in a class of theories previously proposed as background independent membrane field theories. In we argued that when based on an $`SU(16)`$ symmetry this may provide a background independent formulation of $``$ theory. Here we provide independent evidence for this claim by showing how it can be obtained by compactification of a model that we argue has other compactifications that lead to the standard $``$ theory matrix models. This suggests that the full duality group of $``$ theory relates the standard background dependent descriptions of $``$ theory to this new class of background independent theories.
The model we defined in is based on two simple ideas. First the degrees of freedom are an $`N\times N`$ matrix, $`M`$, whose elements themselves are matrices, valued in $`Osp(1|32)`$. Thus, the matrix elements refer to elements of an algebra rather than to positions in a flat background geometry. The different background geometries must then arise from expansions around classical solutions that break the algebra to one that generates the symmetry group of a background spacetime. The algebra $`Osp(1|32)`$ is chosen because there is evidence that it may be the full symmetry group of $``$ theory. Second, the action has the simplest non-trivial form possible, which is the trace of the cube of $`M`$.
In we argued that this theory has a compactification which break the symmetries of $`M`$ down to the Super-Poincare group in $`9+1`$ dimensions. The one loop effective action describing fluctuations around the classical solution representing the compactification was argued to reproduce the $`IKKT`$ matrix model. Another compactification was described, based on a classical solution that breaks the symmetry to the Super-Euclidean group of the light cone gauge of $`10+1`$ dimensional Minkowski spacetime. We argued in that the one loop effective action includes the dWHN-BFFS matrix model<sup>1</sup><sup>1</sup>1These claims are supported by one loop calculations that will be described in detail in ..
The model has, however, many classical solutions in which the $`Osp(1|32)`$ is broken to a group which is not the symmetry group of any spacetime. The small fluctuations around one of these solutions describes a quantum mechanical system which cannot be interpreted in terms of a background spacetime. We may call these non-geometrical compactifications. As these incorporate more of the $`Osp(1|32)`$ group which has been conjectured to be the full symmetry group of $``$ theory, it is necessary to understand these solutions to understand the full physical content of the theory and the full range of duality transformations which operate on its solutions.
Among the nongeometrical compactifications of the cubic matrix model are a special set whose continuum limits are Chern-Simons field theories. As these are topological field theories, they are independent of any metric. Further, quantum Chern-Simons theory is exactly solvable and its state space is understood in terms of the representation theory of quantum groups. This gives rise to a completely algebraic description of the physics of these non-geometrical compactifications. This can be understood in several different ways, one of which involves spaces of conformal blocks, or intertwiners of quantum groups, on 2-surfaces. As was pointed out first by Crane, and Kauffman, the $`k\mathrm{}`$ limit of this structure (where $`k`$ is the level, or the coupling constant of the Chern-Simons theory) reproduces the spin networks which label a normalizable bases of diffeomorphism invariant states in loop quantum gravity descriptions of quantum gravity and supergravity. We will see here that this can be used to derive an extension of loop quantum gravity from the Chern-Simons compactifications of the cubic matrix model.
To make these claims precise we must distinguish two meanings of background independence. We can have a quantum field theory that depends on the topology and differential structure of a manifold of some given dimension, but is independent of any classical fields. Such theories must then have active diffeomorphisms as part of their gauge group. As a result all fields are represented by quantum operators which are subject to evolution by dynamical laws. We may call these field-background independent theories. Loop quantum gravity has provided examples of such theories, in several different dimensions, with and without supersymmetry.
However, one can require that the theory not depend on even the dimension, topology or differential structure of a manifold. We may call such a theory manifold-background independent. A theory with both properties may be called totally-background independent.
If string theory has a background independent formulation it must be totally background independent, because its different solutions are defined on manifolds of different dimension and topology. Loop quantum gravity, in the form originally given in , is not such a theory. This is because its derivation through a rigorous quantization of general relativity and supergravity required that the manifold, dimension and differential structure be fixed. However, it is possible to extend the structure of loop quantum gravity to make it totally background independent. As explained in this requires thickening the spin networks to 2-surfaces labeled by conformal blocks, or intertwiners of a quantum group. This extension is also motivated by the fact that it solves a key problem in loop quantum gravity, because it introduces certain terms in the Hamiltonian or action which are required both for the existence of a classical limit and the recovery of spacetime diffeomorphism invariance.
It was noted in that the resulting theories may be understood as background independent membrane field theories, and that the $`SL(2,Z)`$ duality group can be easily represented in a background independent fashion. This recalls the old arguments that the fundamental objects in $``$ theory should be membranes. To realize this idea in the present context we must invent a theory in which the embeddings of the membrane in different background can be extracted from the intertwiners of a quantum group, by symmetry breaking. A detailed conjecture for how to do this was proposed in . It was found there that to match the dWHN-BFSS matrix model the quantum group must be a quantum deformation of a super-Lie algebra that extends $`SU(16)`$. Here we will derive a closely related theory from a non-geometrical compactification of the cubic matrix model, but with a simpler dynamics, which is also very reminiscent of the dynamics generated by the Hamiltonian constraint in quantum gravity and supergravity.
The version of the cubic matrix model we study here was introduced in , and is characterized by the fact that it uses a complexification of the $`Osp(1|32)`$ degrees of freedom which is represented in terms of $`SU(16,16|1)`$ matrices. We posit here the simplest possible action, given by
$$I=\frac{k}{4\pi }TrM^3$$
(1)
where $`M_{iA}^{jB}`$ is a double matrix with $`i,j=1,\mathrm{},N`$ and $`A,B=1,\mathrm{},33`$. For fixed $`i,j`$, $`M_A^B`$ is an element in the $`33`$ dimensional adjoint representation of $`SU(16,16|1)`$. This was called the gauged action in .
It is not difficult to carry out the same steps as described in to show that this model has compactifications that reduce at the one loop level to the IKKT and BFSS models. This will be discussed very briefly in section 8 below and described in full detail elsewhere. We may note that these compactifications involve breaking the $`SU(16,16|1)`$ symmetry to $`SO(9,1)`$ and $`SO(9)`$, respectively. The non-geometric compactifications we will study break the symmetry in a different way.
$$SU(16,16)Sp(2)SU(16)$$
(2)
In the next section we introduce the model and some of its properties. In section 3 we describe a Chern-Simons compactification leading to an $`SU(16)`$ Chern-Simons theory. In section 4 we show how related compactifications give rise to a set of interacting Chern-Simons theories. In sections 5 to 7 we describe the quantization of these compactifications and show how they reproduce the extension of loop quantum gravity described in . An argument for how the dWHN-BFFS theory may be recovered from the effective action in a different limit of the theory is sketched in section 8. This allows us, in section 9 to describe the explicit duality that holds between the standard background dependent matrix models of $``$ theory and the background independent description derived here. This allows us to answer questions such as what corresponds to $`D0`$ branes in the background independent language. The paper closes with a brief mention of some of the important open problems raised by the results reported here.
## 2 The $`SU(16,16|1)`$ matrix theory
We begin by describing the model and its basic properties. We recall first the definition of $`Osp(1|32)`$, which consists of supermatrices of $`32`$ even dimensions and one odd dimension, $`M_A^B`$ that preserve the graded antisymmetric metric
$$G_A^B=\left[\begin{array}{ccc}0& I& 0\\ I& 0& 0\\ 0& 0& 1\end{array}\right]$$
(3)
(where the first two rows and columns are $`16\times 16`$ bosonic blocks, while the third row and column is in the one fermionic coordinate) so that
$$MG=GM^T$$
(4)
where $`T`$ stands for transpose. We complexify this by considering complex matrices of the same form, which satisfy
$$MG=GM^{}$$
(5)
where $``$ means hermitian conjugate. These generate a supergroup by $`R=expM`$, which satisfies $`RGR^{}=G`$. This group is $`SU(16,16|1)`$, as can be seen from the fact that $`ıG_A^B`$ is an hermitian metric of signature $`(16,16)`$. We may decompose $`M_A^B`$ as follows,
$$M_A^B=\left(\begin{array}{ccc}A_2+Y& A_{}& \mathrm{\Psi }\\ A_+& A_2+Y& \mathrm{\Phi }\\ \mathrm{\Phi }^{}& \mathrm{\Psi }^{}& 0\end{array}\right)$$
(6)
where the $`A_2,A_\pm `$ are three $`16\times 16`$ hermitian matrices, $`Y`$ is a tracefree $`16\times 16`$ antihermitian matrix and $`\mathrm{\Psi }^A`$ and $`\mathrm{\Phi }^A^{}`$ are $`16`$ component spinors. We will also find it convenient to use $`A_0=A_+A_{}`$ and $`A_1=A_++A_{}`$. The components of $`Y`$ and $`A_a`$ for $`a=0,1,2`$ are even Grassman variables, while the components of $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are odd Grassman. We will decompose the $`33`$ dimensional indices $`A,B,\mathrm{}`$ as $`A=(P,P^{},)`$ where $`P=1,\mathrm{},16,`$, $`P^{}=1^{},\mathrm{},16^{},`$ and $``$ is the lone fermionic index. We will find it very useful to consider the decomposition
$$SU(16,16)Sp(2)SU(16)$$
(7)
where the $`SU(16)`$ is generated by
$$W_A^B=\left(\begin{array}{ccc}Y& 0& 0\\ 0& Y& 0\\ 0& 0& 0\end{array}\right)$$
(8)
and the generators of $`Sp(2)`$ are given by
$$\tau ^0=\left(\begin{array}{ccc}0& I& 0\\ I& 0& 0\\ 0& 0& 0\end{array}\right),\tau ^1=\left(\begin{array}{ccc}0& I& 0\\ I& 0& 0\\ 0& 0& 0\end{array}\right),\tau ^2=\left(\begin{array}{ccc}I& 0& 0\\ 0& I& 0\\ 0& 0& 0\end{array}\right)$$
(9)
These satisfy
$$\frac{1}{32}Tr\tau ^a\tau ^b=\eta ^{ab}$$
(10)
where $`\eta ^{ab}`$ is the $`2+1`$ dimensional metric
$$\eta ^{ab}=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)$$
(11)
thus showing $`Sp(2)=SO(2,1)`$.
The $`SU(16)`$ transformations can be extended to the superalgebra $`SU(16|1)`$, leading to a reduction of superalgebras,
$$SU(16,16|1)Sp(2)SU(16|1).$$
(12)
We now extend the matrix by considering each entry to be an $`N\times N`$ matrix parameterized by $`i,j=1,\mathrm{},N`$. We thus have a double matrix $`M_{Ai}^{Bj}`$. We define the action
$$I^{gauged}=\frac{k}{4\pi }TrM^3$$
(13)
where the multiplication and trace is on both sets of indices. $`k`$ is the coupling constant of the theory. This action is somewhat simpler than that studied in and has more gauge symmetry, as we can see by writing it more explicitly
$$I^{gauged}=\frac{k}{4\pi }\underset{ijk}{}Tr_{SU(16,16|1)}M_i^jM_j^kM_k^i$$
(14)
where for each $`i,j,k`$ the multiplications and trace are over the $`SU(16,16|1)`$ indices. The action is then invariant under transformations
$$M_i^jU(i)M_i^jU(j)^{}$$
(15)
where, for each $`i`$, $`U(i)SU(16,16|1)`$. Similarly, for each $`A`$ and $`B`$ we have the gauge symmetry, suppressing now the $`SU(N)`$ indices
$$M_A^BV(A)M_A^BV(B)^{}$$
(16)
where $`V(A)SU(N)`$. For this reason we call it a gauged matrix action.
It is useful to decompose the action in terms of the variables defined in (6). We have
$`I^{gauged}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}Tr_{ij}Tr_{PQ}\left\{ϵ^{abc}(A_aA_bA_c)+YA_aA_b\eta ^{ab}+Y^3\right\}`$ (17)
$`+3Tr_{ij}\left\{\mathrm{\Phi }^PA_P^Q\mathrm{\Phi }_Q\mathrm{\Psi }^PA_{+P}^Q\mathrm{\Psi }_QA_{2PQ}\{\mathrm{\Psi }^Q,\mathrm{\Phi }^P\}+Y_{PQ}[\mathrm{\Psi }^Q,\mathrm{\Phi }^P]\right\}`$
We see explicitly here the decomposition into $`Sp(2)SU(16|1)`$.
It is also interesting to note that when the matrix elements are restricted to be real, so that the symmetry is reduced from $`SU(16,16|1)`$ to $`Osp(1|32)`$ the action has a translation invariance given by
$$A_{aPi}^{Qj}A_{aPi}^{Qj}=A_{aPi}^{Qj}+\delta _i^jV_{aP}^Q$$
(18)
with $`V_{aPQ}=V_{aPQ}^{}`$. This is reduced from the symmetry of the model studied in but the remaining translation symmetry includes that of the dWHN-BFSS and IKKT models, when we identify the fields of those models with components of $`M_A^B`$ as described in .
It is also of interest to consider the theory without the $`Y`$ degrees of freedom. This theory is invariant under the $`Sp(2)SU(16|1)`$ subalgebra of $`SU(16,16|1)`$. Its action is given by
$$I^{MCS}=I_{Y=0}^{gauged}=\frac{k}{4\pi }Tr_{ij}\left\{ϵ^{abc}Tr_{PQ}(A_aA_bA_c)+\chi ^P\tau ^aA_{aP}^Q\chi _Q\right\}$$
(19)
We will call this matrix Chern-Simons theory in the following.
## 3 A simple Chern-Simons compactification
The Chern-Simons compactification comes about because for each $`a`$, $`A_a=(A_{aP}^Q)_i^j`$ is an $`N\times N`$ matrix of $`16`$ dimensional hermitian matrices. There are then compactifications of the model in which the $`ϵ^{abc}(A_aA_bA_c)`$ term in the action becomes an $`SU(16)`$ Chern-Simons theory. To see this we compactify the theory on a three-torus, making use of a modification of the route studied in . We consider the classical solutions given by
$$A_{aiP}^{jQ}=\delta _P^Q(_a)_i^j$$
(20)
$$Y=0$$
(21)
with the fermionic fields vanishing. We divide the indices so as to make three derivative operators. We choose $`i=i_0,i_1,i_2`$ where $`i_a=M_a,\mathrm{},0,1,\mathrm{},M_a`$ such that $`N=_{a=0,1,2}(2M_a+1)`$. We then choose
$$(_0)_i^j=(_0)_{i_0i_1i_2}^{j_0j_1j_2}=i_0\delta _{i_0}^{j_0}\delta _{i_1}^{j_1}\delta _{i_2}^{j_2}$$
(22)
and similarly for the other two derivative operators. Clearly we have $`[_a,_b]=0`$. We then expand around this classical solution, using the usual matrix compactification trick, defining variables
$$A_{aiP}^{jQ}=\delta _P^Q(_a)_i^j+a_{aiP}^{jQ}$$
(23)
Following the usual translation into continuum fields, we find in the limit $`M_a\mathrm{}`$
$`I^{gauged}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{T^3}}Tr_{SU(16)}\left\{ada+{\displaystyle \frac{2}{3}}a^3+\chi ^{}\tau ^a𝒟_a\chi \right\}`$ (24)
$`+(Y_P^Q)^3+Y_{PQ}\eta ^{ab}(𝒟_a\stackrel{~}{a}_b)^{QP}+Y_{PQ}[\mathrm{\Psi }^P,\mathrm{\Phi }^Q]`$
where $`\stackrel{~}{a}_b^{PQ}`$ is tracefree.
If we neglect the coupling to the $`Y`$ field this is a Chern-Simons theory on $`T^3`$.
Note that had we begun with the bosonic part matrix Chern-Simons theory (19), whose action is just $`ϵ^{abc}Tr(A_aA_bA_c)`$, the result would just be the first two terms of (24). In this case something remarkable has happened, which is that the dependence on a background metric on the torus, which is implicit in the definition of the compactification, has gone away in the limit $`M_a\mathrm{}`$, leaving a topological field theory. We may note that with the ultraviolet cutoff $`l_{Planck}`$ held fixed, this is equivalent to the limit in which all three compactification radii are taken to infinity. The fact that a topological field theory emerges from the limit is a consequence of the large amount of gauge symmetry in the original action, together with the taking of a limit in which the length scales set by the compactification radii are removed.
The coupling to the spinor variables $`\chi `$ depends on the background $`\tau ^a`$ which define the compactification directions. These result in a supersymmetrization of the state space of Chern-Simons theory as we will discuss in .
The theory with the $`Y`$ terms present is more subtle, but it is possible that they will not change the topological character of the theory in the limit of infinite compactification radii, up to renormalizations of the coupling constant of Chern-Simons theory. The reason is that we can use the gauge symmetry (15) to set $`Y_i^j=0`$ locally in $`i`$ and $`j`$. Given any sequence of values for $`i`$, with no repeats, given by $`i_\mu =i_1,\mathrm{},i_L`$, we can set $`Y_{i_\mu }^{i_{\mu +1}}=0`$. To see this note that by (15)
$$\delta Y_{iP}^{jQ}=\delta M_{i\alpha P}^{j\alpha Q}=U(i)_{\alpha P}^{\gamma R}M_{i\gamma R}^{j\alpha Q}M_{i\alpha P}^{j\gamma R}U(j)_{\gamma R}^{\alpha Q}$$
(25)
Thus, $`Y_i^j`$ transforms like a gauge potential, it can be set to zero along any open curve in the space labeled by $`i,j`$ but its effect on closed curves in those indices must be taken into account. In a matrix compactification the indices $`i`$ and $`j`$ refer to Fourier modes, in particular the local field is constructed from sums such as $`Y(x)=_kY_{i_0}^{i_o+k}e^{ıkx}`$. This means we have a kind of gauge invariance in momentum space. As a result, we can eliminate the $`Y`$’s from any scattering matrix involving other degrees of freedom. But we cannot eliminate closed loops of $`Y`$’s in any Feynman diagram, as they will correspond to a closed loop of matrix entries. The effect on the classical equations of motion will only be through the $`Y`$ equation itself. The bosonic equations of motion which follow from (24) have the form
$$ϵ^{abc}F_{bc}\chi ^{}\tau ^c\chi \eta ^{cd}𝒟_dY=0$$
(26)
$$Y_{PQ}^2+\left(𝒟_aa_b\right)_{PQ}\eta ^{ab}+[\mathrm{\Psi }_P,\mathrm{\Phi }_Q]=0$$
(27)
If we set $`Y=0`$ at a point we have the equations of Chern-Simons theory plus the condition
$$\left(𝒟_aa_b\right)_{PQ}\eta ^{ab}+[\mathrm{\Psi }_P,\mathrm{\Phi }_Q]=0$$
(28)
But this differs only by the fermion term from the standard gauge fixing term which is used to define Chern-Simons theory perturbatively. This means that at least up to the effects of closed loops in the $`Y`$ variables, the theory given by (24) is a supersymmetric extension of the $`SU(16)`$ Chern-Simons theory.
This does not completely resolve the question of the influence of the $`Y`$ degrees of freedom. However, it suggests that the effect of the $`Y`$’s on the Chern-Simons compactification can only be to renormalize the coupling of the Chern-Simons theory. For this reason we will neglect the effects of the $`Y`$’s in the next sections where we study the physics of the Chern-Simons compactifications.
## 4 The $`(\text{Chern-Simons})^P`$ compactification
We now introduce a different set of compactifications that reduce the theory to a coupled set of $`P`$ Chern-Simons theories, for any $`P`$. The idea is to blow up each of the $`i_aj_a`$ entries of the previous compactification into $`P\times P`$ blocks. We then have a symmetry, for each $`i_a,j_a`$, which is $`Sp(2)SU(16)U(P)`$. When we make a compactification on the three torus, as just described, we find a $`3`$ dimensional quantum field theory with symmetry $`SU(16)U(P)`$. There are however, several different ways of taking the limit that defines the field theory, which result in a different set of fluctuating fields with different symmetry. One way, which we will sketch in section 8, preserves the $`U(P)`$ symmetry and breaks the $`SU(16)`$ down to the $`SO(9)`$ symmetry of the lightcone gauge of $`10+1`$ dimensional spacetime.
Here we will consider a way to preserve the full $`SU(16)`$, but break the $`U(P)`$ symmetry completely. This leads not to one $`3d`$ quantum field theory with $`U(P)`$ symmetry, but to $`P`$ $`3d`$ field theories, each of which defines a Chern-Simons theory. These interact via bi-local fields that create and annihilate pairs of punctures that join the tori on which these Chern-Simons theories are defined. As we will show in the next several sections, this gives a background independent theory. This theory is a version of the background independent membrane field theory which was proposed in as a background independent version of $``$ theory.
We begin the demonstration by choosing integer factors $`M_a`$ and $`P`$ such that $`_a(2M_a+1)P=N`$. We will write
$$i=i_0,i_1,i_2,I$$
(29)
with $`I=1,\mathrm{},P`$ and $`i_a=M_a,\mathrm{},M_a`$.
We then decompose $`M_I^J`$ according to (with all the other indices suppressed)
$$M_I^J=\left[\begin{array}{cccc}A^1& B_1^2& B_1^3& \mathrm{}\\ B_2^1& A^2& B_2^3& \mathrm{}\\ \mathrm{}& \mathrm{}& A^3& \mathrm{}\end{array}\right]$$
(30)
That is, we define
$$A_{aIi_0i_1i_2P}^{J=Ij_0j_1j_2Q}=A_{ai_0i_1i_2P}^{Ij_0j_1j_2Q}$$
(31)
and for the off diagonal terms,
$$A_{aIi_0i_1i_2P}^{JIj_0j_1j_2Q}\tau _\alpha ^{a\beta }=B_{\alpha Ii_0i_1i_2P}^{\beta Jj_0j_1j_2Q}$$
(32)
We continue to neglect the role of the $`Y`$ field.. The dynamics is given by (with the $`i_a`$ indices suppressed)
$`I^{gauged}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}Tr_{i_a}\{{\displaystyle \underset{I}{}}ϵ^{abc}(A_{aP}^{IQ}A_{bQ}^{IR}A_{cR}^{IP})+\chi _{\alpha I}^P\tau _\beta ^{a\alpha }A_{aP}^{IQ}\chi _Q^{\beta I}`$ (33)
$`+{\displaystyle \underset{JI}{}}[3B_{IP\alpha }^{JQ\beta }\tau _\beta ^{b\gamma }A_{bQ}^{JR}B_{JR\gamma }^{IP\alpha }`$
$`+\chi _{\alpha I}^{PJ}\tau _\beta ^{a\alpha }A_{aP}^{JQ}\chi _{QJ}^{\beta I}+\chi _\alpha ^{PJ}B_{PJ}^{IQ}\chi _{QI}^{\beta J}+\chi _{\alpha I}^{PJ}B_{PJ}^{IQ}\chi _Q^{\beta I}]`$
$`+{\displaystyle \underset{JIK}{}}[B_{IP\alpha }^{JQ\beta }B_{JQ\beta }^{KR\gamma }B_{KR\gamma }^{IP\alpha }+\chi _{\alpha I}^{PJ}B_{PJ}^{KQ}\chi _{QK}^{\beta I}]\}`$
We see we get $`P`$ theories described by an $`A^I`$ coupled by terms involving the $`B_I^J`$ and $`\chi _{QK}^{\beta I}`$ variables.
The theory has a $`U(P)`$ symmetry mixing up the $`I,J`$ indices. We now define a class of compactifications which can lead to a quantum theory that breaks this symmetry. To do this we will define a Chern-Simons compactification to an $`SU(16)`$ Chern-Simons theory separately for each $`I`$. For each of the $`P`$ diagonal $`A^I`$’s we define the compactification to an $`SU(16)`$ gauge theory described in the last section. This leads, for each $`I`$, to a Chern-Simons theory on a three torus based on the fusion algebra of $`SU(16)`$.
This multiple-Chern-Simons compactification also induces a transformation on the $`B`$ variables. We find (for $`JI`$) that when $`M_a\mathrm{}`$ the $`B_{I\alpha _Ii}^{J\beta _Jj}`$ become bilocal fields whose domain are pairs of the $`T^3`$ on which the Chern-Simons theories are defined.
$$B_{I\alpha _Ii_0i_1i_2P}^{J\beta _Jj_0j_1j_2Q}B_{I\alpha _IP}^{J\beta _JQ}(x_I,x_J)$$
(34)
We also have bilocal fermionic fields given by
$$\chi _{QIi_0i_1i_2}^{\alpha Jj_0j_1j_2}\chi _{QI}^{\alpha J}(x_I,x_J)$$
(35)
Thus after the $`P`$ simultaneous three-torus compactifications and the $`M_a\mathrm{}`$ limits the action of the theory becomes
$`I^{cgauged}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle \underset{I}{}}S_{CS}^I(𝒜^I,\chi ^I)+{\displaystyle \underset{JI}{}}{\displaystyle \frac{k}{4\pi }}{\displaystyle _{T_I^3}}d^3x_I{\displaystyle _{T_J^3}}d^3x_J\{Tr_G\left(B_{I\alpha }^{J\beta }(x_I,x_J)\tau _\beta ^{\overline{a}\gamma }𝒟_{\overline{a}}^JB_{J\gamma }^{I\alpha }(x_J,x_I)\right)`$ (36)
$`+\chi _{\alpha I}^{PJ}(x_I,x_J)\tau _\beta ^{a\alpha }𝒟_{aP}^{JQ}\chi _{QJ}^{\beta I}(x_J,x_I)+\chi _\alpha ^{PJ}(x_{J)}B_{PJ}^{IQ}(x_J,x_I)\chi _{QI}^{\beta J}(x_I,x_J)`$
$`+\chi _{\alpha I}^{PJ}(x_I,x_J)B_{PJ}^{IQ}(x_J,x_I)\chi _Q^{\beta I}(x_I)\}`$
$`+{\displaystyle \underset{IJK}{}}{\displaystyle \frac{k}{4\pi }}{\displaystyle _{T_I^3}}d^3x_I{\displaystyle _{T_J^3}}d^3x_J{\displaystyle _{T_K^3}}d^3x_KTr_G(B_{I\alpha }^{J\beta }(x_I,x_J)B_{J\beta }^{K\gamma }(x_J,x_K)B_{K\gamma }^{I\alpha }(x_K,x_I)`$
$`+\chi _{\alpha I}^{PJ}(x_I,x_J)B_{PJ}^{KQ}(x_J,x_K)\chi _{QK}^{\beta I}(x_K,x_I))`$
where
$$S_{CS}^I(𝒜^I,\chi ^I)=_{T^3}Tr_{SU(16)}\left\{ada+\frac{2}{3}a^3+\chi ^{}\tau ^a𝒟_a\chi \right\}$$
(37)
Thus, we have $`P`$ Chern-Simons theories, each defined on a distinct 3-torus, which interact with each other via the bi-local fields $`B_{I\alpha }^{J\beta }(x_I,x_J)`$.
## 5 Hamiltonian dynamics of the $`(\text{Chern-Simons})^P`$ compactification
We now study the dynamics of the theory defined by the multiple Chern-Simons compactification. To simplify the discussion we will ignore the fermion fields and consider only the bosonic parts of (36). The fermion terms give a supersymmetric completion of the structure we will describe here; details of this will be described elsewhere.
To uncover the dynamics of the theory in this compactification we make a $`2+1`$ splitting in each of the $`P`$ 3-tori. This gives coordinates $`x_I^a=(t_I,x_I^m)`$, where we define spatial indices $`l,m,n=1,2`$. We also let the compactification radius for each time coordinate go to infinity, so that the Hamiltonian theory becomes defined on a domain which is $`P`$ copies of $`R\times T^2`$. The action then takes the form,
$`I^{cgauged}`$ $`=`$ $`{\displaystyle \underset{I}{}}{\displaystyle \frac{k}{4\pi }}{\displaystyle 𝑑t_Id^2x_ITr_{SU(16)}\left[ϵ^{mn}A_m^I\dot{A}_n^IA_0^I𝒢^I\right]}`$
$`+`$ $`{\displaystyle \underset{J>I}{}}{\displaystyle \frac{k}{4\pi }}{\displaystyle }dt_Id^2x_Idt_Jd^2x_JTr\{(B_J^I(x_J,t_J;x_I,t_I))_\alpha ^\beta \tau _\beta ^{0\gamma }({\displaystyle \frac{}{t_I}}{\displaystyle \frac{}{t_J}})B_{I\gamma }^{J\alpha }(x_I,t_I,x_J,t_J)`$
$`+`$ $`(B_J^I(x_J,t_J;x_I,t_I))_\alpha ^\beta \tau _\beta ^{m\gamma }(𝒟_m^I𝒟_m^J)B_{I\gamma }^{J\alpha }(x_I,t,x_J,s)]\}`$
$`+`$ $`{\displaystyle \underset{IJK}{}}{\displaystyle \frac{k}{4\pi }}{\displaystyle _{T_I^3}}d^3x_I{\displaystyle _{T_J^3}}d^3x_J{\displaystyle _{T_K^3}}d^3x_KTr\left(B_{I\alpha }^{J\beta }(x_I,x_J)B_{J\beta }^{K\gamma }(x_J,x_K)B_{K\gamma }^{I\alpha }(x_K,x_I)\right)`$
where the Gauss’s law constraint is
$`𝒢^I(x_I,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}ϵ^{mn}F_{mn}^I(x_I,t)+{\displaystyle \underset{J>I}{}}{\displaystyle 𝑑t_Jd^2x_J(B_J^I(x_J,t_J;x_I,t_I))_\alpha ^\beta \tau _\beta ^{0\gamma }B_{I\gamma }^{J\alpha }(x_I,t_I,x_J,t_J)}`$ (39)
$``$ $`{\displaystyle \underset{K<I}{}}{\displaystyle 𝑑t_Kd^2x_J(B_I^K(x_I,t_I;x_K,t_K))_\alpha ^\beta \tau _\beta ^{0\gamma }B_{K\gamma }^{I\alpha }(x_K,t_K,x_I,t_I)}`$
We see that the theory appears to be non-local in time. We will shortly see that this is an expression of a many fingered time gauge invariance, in that the time can be evolved independently on each torus.
The $`A_m^I`$ have conventional momenta,
$$\pi ^{Im}(x_i,t)\frac{\delta I^{cubic}}{\dot{𝒜}_m^I(x_I,t)}=\frac{k}{2\pi }ϵ^{mn}A_n^I(x_I,t)$$
(40)
The $`B_I^J`$’s depend on two time coordinates, one in each of the two tori they live in. The momenta then similarly depend on two time and two space variables. As we see from (5) the momenta of the $`B_I^J`$’s depends on the difference of its two time coordinates,
$$\mathrm{\Pi }_J^I(x_I,t_I;x_J,t_J)_{\alpha P}^{\gamma Q}\frac{\delta I^{cubic}}{(t_It_J)B_{I\gamma Q}^{J\alpha P}(x_I,t_I,x_J,t_J)}=\frac{k}{4\pi }(B_J^I)_{\alpha P}^{\beta Q}(x_J,t_J,x_I,t_I)\tau _\beta ^{0\gamma }$$
(41)
To define a polarization we pick an arbitrary ordering of the $`P`$ tori. This breaks the $`U(P)`$ symmetry. We then take the variables $`B_{I\gamma }^{J\alpha }(x_I,t,x_J,s)`$ for $`J>I`$ to be configuration variables, while the momenta are coded in the $`B_I^J`$ for $`J<I`$.
Note also that there is some freedom of the density transformation properties of the fields defined by the compactification. This must be defined so that all the integrands are densities. A sensible choice seems to be to define the compactification in such a way that the $`B_{I\gamma }^{J\alpha }(x_I^a,x_J^b)`$ are densities on the second spacetime variables $`x_J^b`$ and ordinary functions on the first. The momenta will then have the opposite density weights, i.e. weight one in the first slot and zero in the second.
Finally, we write the action in Hamiltonian form
$`I^{gauged}`$ $`=`$ $`{\displaystyle \underset{I}{}}{\displaystyle 𝑑t_Id^2x_ITr_{SU(16)}\left[\pi ^{Im}\dot{A}_m^IA_0^I𝒢^I\right]}`$ (43)
$`+{\displaystyle \underset{J>I}{}}{\displaystyle 𝑑t_Id^2x_I𝑑t_Jd^2x_JTr_{SU(16)}\left[\mathrm{\Pi }_I^J(x_I,t_I;x_J,t_J)_\alpha ^\beta \left(\frac{}{t_I}\frac{}{t_J}\right)B_{I\gamma }^{J\alpha }(x_I,t_I,x_J,t_J)\right]}`$
$`+{\displaystyle \underset{J>I}{}}{\displaystyle 𝑑t_I𝑑t_J_2^{IJ}(t_I,t_J)}+{\displaystyle \underset{I<J<K}{}}{\displaystyle 𝑑t_I𝑑t_J𝑑t_K_3^{IJK}(t_I,t_J,t_K)}`$
where the Gauss’s law constraint is now
$$𝒢^I(x_I,t)=\frac{1}{2}ϵ^{mn}F_{mn}^I(x_I,t)j^I(x_I,t)=0$$
(44)
where the $`SU(16)`$ current for each Chern-Simons theory is given by
$$j_Q^{KP}(x_K,t_K)=\underset{I<K}{}d^3x_ITr_{Sp(2)}\mathrm{\Pi }_{IR}^{KP}B_{IQ}^{KR}\underset{K<J}{}d^3x_JTr_{Sp(2)}\mathrm{\Pi }_{KR}^{JP}B_{KQ}^{JR}$$
(45)
and the two and three time Hamiltonians are given by
$$_2^{IJ}(t_I,t_J)=d^2x_Id^2x_JTr_{SU(16)}\overline{\mathrm{\Pi }}_I^J(x_I,t;x_J,s)_\alpha ^\beta \tau _\beta ^{m\gamma }\left(𝒟_m^I𝒟_m^J\right)B_{I\gamma }^{J\alpha }(x_I,t,x_J,s)$$
(46)
$$_3^{IJK}(t_I,t_J,t_K)=3d^2x_Id^2x_Jd^2x_KTr\left\{B_{I\gamma }^{K\alpha }\overline{\mathrm{\Pi }}_{J\beta }^{K\gamma }\overline{\mathrm{\Pi }}_{I\gamma }^{J\alpha }+B_{I\alpha }^{J\beta }B_{J\beta }^{K\gamma }\overline{\mathrm{\Pi }}_{I\gamma }^{K\alpha }\right\}$$
(47)
Here $`\overline{\mathrm{\Pi }}_{I\alpha }^{J\beta }=\mathrm{\Pi }_{I\alpha }^{J\gamma }\tau _\gamma ^{0\beta }`$.
## 6 Quantization of the $`(CS)^P`$ compactification
What we have is a multi-time theory in which a set of $`P`$ Chern-Simons theories, each with a time coordinate, are coupled at all pairs and triples of times. This can be treated to some extent like a many fingered time, as we will now see.
The canonical commutation relations for the $`B`$’s have the structure, for $`I<J`$ and $`K<L`$,
$$[B_{I\alpha }^{J\beta }(x_I,t,x_J,s)_{AB},\mathrm{\Pi }_K^L(x_K,t;x_L,u)_\gamma ^{\delta CD}]=\delta _I^K\delta _L^J\delta (s,u)\delta ^2(x_K,x_I)\delta ^2(x_J,x_L)\delta _\alpha ^\delta \delta _\gamma ^\beta \delta _{AB}^{CD}$$
(48)
To realize these quantum mechanically we are going to have to have a Hilbert space for every set of $`P`$ times $`t_I`$. We will call this space $`H(t_I)`$. For each set of $`t_I`$ we will have operators that act on this space that satisfy (48). This kind of structure is common to certain histories formulations of quantum theory, such as those studied in . In those papers it is shown that this kind of multitime quantum theory is natural for histories formulations of systems with spacetime diffeomorphism invariance.
The form of the Poisson brackets suggests that we quantize in a generalization of the connection representation in which the states are of the form
$$\mathrm{\Psi }[\{t_I\}]=\mathrm{\Psi }[A_I(t_I),B_I^J(t_I,t_J)]$$
(49)
for $`J>I`$.
The first step in the quantization will be to find solutions to the Gauss’s law constraint (44). To define this we note that the states are defined to be functions only of the configuration variables at a set of $`P`$ fixed times $`t_I`$. Thus, when acting on a state $`\mathrm{\Psi }(t_I)`$ the integral $`𝑑s`$ in 45 will for each $`J`$ only pick up times in that list. The $`\delta (s,t)`$ in (48) absorbs the integral over $`s`$; thus, the action of $`\widehat{j}^I(x_I)`$ on states has the form
$`\widehat{j}^I(x_I,t_I)`$ $`\mathrm{\Psi }[A^I(t_I),B_I^J(t_I,t_J)]`$
$`\left\{{\displaystyle \underset{J>I}{}}{\displaystyle d^2x_JB_I^J(x_I,t_I,x_J,t_J)\frac{\delta }{\delta B_I^J(x_I,t_I,x_J,t_J)}}\right\}\mathrm{\Psi }[A^I(t_I),B_I^J(t_I,t_J)]`$
A natural set of solutions to these equations may be obtained by making an ansatz
$$\mathrm{\Psi }[A_I(t_I),B_I^J(t_I,t_J)]=\chi [B_I^J]\underset{K}{}\mathrm{\Phi }_K[A^K]$$
(51)
To solve this we want to find the action of the operator representing the current (45), which acts as
$`\widehat{j}_Q^{KP}(x_K)\chi [B_I^J]`$ (52)
$`=\left\{{\displaystyle \underset{I<K}{}}{\displaystyle d^3X_ITr_{Sp(2)}B_{IQ}^{KR}\frac{\delta }{\delta B_{IR}^{KP}}}{\displaystyle \underset{K<J}{}}{\displaystyle d^3X_JTr_{Sp(2)}B_{KQ}^{JR}\frac{\delta }{\delta B_{KR}^{JP}}}\right\}\chi [B_I^J].`$
It is clear that there are solutions to the Gauss law which involve finite products of $`B_I^J`$’s. In such a state, for each pair of tori $`T_I^2`$ and $`T_J^2`$, there will be a finite number of points $`(p_I^a,p_J^a)T_I^2\times T_J^2`$ for $`a=1,\mathrm{},n_{IJ}`$ on which there are $`B_I^J`$’s in the product (51). We will call these punctures. The current will then have the form
$$𝒥_Q^{KP}(x_K)=\underset{I<K}{}\underset{a=1}{\overset{n_{IK}}{}}\delta ^2(p_K^a,x_K)𝒥^a[B_I^K]+\underset{K<J}{}\underset{b=1}{\overset{n_{KJ}}{}}\delta ^2(p_K^b,x_K)𝒥^b[B_K^J]$$
(53)
where the currents $`𝒥^a[B_I^K]`$ depend on the $`B`$’s as indicated. As a consequence $`\mathrm{\Phi }_K(A^K)`$ satisfies the condition
$$\frac{k}{2\pi }\widehat{F}_{12}^K(x_K)\mathrm{\Phi }_K(A^K)=\left(\underset{I<K}{}\underset{a=1}{\overset{n_{IK}}{}}\delta ^2(p_K^a,x_K)𝒥^a[B_I^K]\underset{K<J}{}\underset{b=1}{\overset{n_{KJ}}{}}\delta ^2(p_K^b,x_K)𝒥^b[B_K^J]\right)\mathrm{\Phi }_K(A^K)$$
(54)
But this is a familiar equation from quantum Chern-Simons theory . The connection on each tori is flat except for a finite set of punctures at which there is a delta function contribution, which depends on the $`B`$’s. We can solve this in the standard way, expressing the solutions in terms of conformal blocks or intertwiners of $`G`$ on the punctured two torus. The dependence of the states on the $`B`$’s will be expressed in terms of the representations labeling the punctures. As a consequence the punctures will satisfy braid statistics. This means care must be taken if we create two punctures on top of each other. A simple solution to this problem which we will employ is to construct the states with all punctures distinct and use the recoupling relations of the quantum group associated to $`SU(16)`$ to extend their values to the cases of multiple punctures at the same point.
But once the punctures are distinct we can use the gauge symmetry in (15) to reduce the number of independent components of the $`B_I^J`$’s. Taking $`U(i)`$ in (15) to be valued in $`Sp(2)`$ we have transformations,
$$A_a^I(x_I)\tau ^aU_I^1(x_I)A_a^I(x_I)\tau ^aU_I(x_I)$$
(55)
$$B_I^J(x_I,x_J)U_I^1(x_I)B_I^J(x_I,x_J)U_J(x_J)$$
(56)
Since the punctures are distinct we can use this gauge freedom to diagonalize the $`B_I^J`$. Given the definition of the momenta, this gives us, for $`I<J`$
$$B_I^J=\left[\begin{array}{cc}b_{+I}^J& 0\\ 0& b_I^J\end{array}\right]$$
(57)
while for $`J>I`$ we define momenta<sup>2</sup><sup>2</sup>2It is convenient here to take for the time variable the coefficient of $`\tau ^2`$ in the parameterization we used previously.,
$$B_J^I=\left[\begin{array}{cc}p_{+I}^J& 0\\ 0& p_I^J\end{array}\right]$$
(58)
The commutation relations are then.
$$[b_{\pm I}^J(x_I,t,x_J,s)_{PQ},p_{\pm ^{}K}^L(x_K,t;x_L,u)^{RS}]=\pm ı\delta _I^K\delta _L^J\delta _{\pm ,\pm ^{}}\delta (s,u)\delta ^2(x_K,x_I)\delta ^2(x_J,x_L)\delta _{PQ}^{RS}$$
(59)
The two degrees of freedom of the $`B_I^J`$ correspond to the two cases in which the current flows from $`I`$ to $`J`$ (with $`I<J`$) in a positive or a negative sense. The current is now,
$$j_Q^{KP}(x_K,t_K)=\underset{\pm }{}(\pm 1)\left\{\underset{I<K}{}d^3x_Ip_{\pm IR}^{KP}b_{\pm IQ}^{KR}\underset{K<J}{}d^3x_Jp_{\pm KR}^{JP}b_{\pm KQ}^{JR}\right\}$$
(60)
We then work with the reduced Hilbert space
$$\mathrm{\Psi }[A_1^I(x_I,t_I),b_{\pm I}^J(x_I,t_I;x_J,t_J)]$$
(61)
for all ordered pairs $`J>I`$. On this space we will have the kinematical operators,
$$\widehat{A}_1^I(x_I,t_I)\mathrm{\Psi }=A_1^I(x_I,t_I)\mathrm{\Psi }$$
(62)
$$\widehat{A}_2^I(x_I,t_I)\mathrm{\Psi }=ı\mathrm{}\frac{\delta }{\delta A_1^I(x_I,t_I)}\mathrm{\Psi }$$
(63)
$$\widehat{b}_{\pm I}^J(x_I,t_I;x_J,t_J)\mathrm{\Psi }=b_{\pm I}^J(x_I,t_I;x_J,t_J)\mathrm{\Psi }$$
(64)
$$\widehat{p}_{\pm I}^J(x_I,t_I;x_J,t_J)\mathrm{\Psi }=\pm ı\mathrm{}\frac{\delta }{\delta b_{\pm I}^J(x_I,t_I;x_J,t_J)}\mathrm{\Psi }$$
(65)
We can then describe the solutions to the Gauss’s law constraint as follows.
Let us then pick $`R`$ pairs of punctures, each member of a pair is a point on a distinct 2-torus. The pairs are labeled by $`w=1,\mathrm{},R`$ and each pair is given by a point on the 2-torus’s $`I_w`$ and $`J_w`$ labeled by,
$$p_w=(p_w^+,p_w^{})T_{I_w}^2\times T_{J_w}^2$$
(66)
Each puncture has a polarity $`ϵ(w)=\pm `$. We then consider an ansatz for the wavefunction,
$$\mathrm{\Psi }[A^I,b_{\pm I}^J]=\underset{w}{}b_{ϵ(w)I_w}^{J_w}(p_w^+,p_w^{})\underset{I}{}\mathrm{\Phi }_I(A^I)]$$
(67)
These satisfy
$$\widehat{j}^I(x_I)\mathrm{\Psi }[A^I,b_{\pm I}^J]=\underset{w}{}ϵ(w)\left(𝒥_w\delta _{II_w}\delta ^2(x_I,p_w^+)𝒥_w\delta _{IJ_w}\delta ^2(x_I,p_w^{})\right)\mathrm{\Psi }[A^I,b_{\pm I}^J]$$
(68)
where $`𝒥_wSU(16)`$. This means that $`\mathrm{\Phi }_I(A^I)`$ satisfies the condition
$$\frac{k}{2\pi }\widehat{F}_{12}^I(x_I)\mathrm{\Phi }_I(A^I)=\underset{w}{}ϵ(w)\left(𝒥_w\delta _{II_w}\delta ^2(x_I,p_w^+)𝒥_w\delta _{IJ_w}\delta ^2(x_I,p_w^{})\right)\mathrm{\Phi }_I(A^I)$$
(69)
We can solve this by using the usual methods from Chern-Simons theory . We may picture a typical state in terms of punctured tori whose punctures are connected in pairs as in Figure (1).
Finally, we may note that once we have fixed the $`B_I^J`$ to the diagonal gauge the quadratic terms in the Hamiltonian, $`H_{IJ}^2`$, all vanish. This suggests that the role of these terms is in implementing spatial diffeomorphism invariance on each of the 2-tori. Since the choice of representation of the $`Sp(2)`$ matrices defined which components of the matrix variables corresponded to the spatial coordinates of the tori it is not surprising that there is in this way a coupling between the $`Sp(2)`$ gauge freedom and spatial diffeomorphisms.
We may then finally describe the space of gauge invariant and diffeomorphism invariant states as follows. We can first describe an auxilliary linear space $`H_{aux}^P`$ constructed from all ways to join $`P`$ tori along pairs of punctures. Given a set of $`R`$ punctures labeled by $`\sigma =(I_w,p_w^+,J_w,p_w^{})`$ we join our $`P`$ tori into a set of compact 2-surfaces by joining them along each pair of punctures in the set. For each $`w`$ we draw a circle around the puncture $`p_w^+`$ on the $`I_w`$’th torus and another circle around the puncture $`p_w^{}`$ on the $`J_w`$’th torus. We then remove the interiors of each circle and join the two boundaries together. The resulting single circle $`c_w`$ is oriented and labeled with an $`s`$ , corresponding to a projection operator which restricts the current flowing accross $`c_w`$ to be in the spinor representation.
Once all these operations are completed we have a (generally disjoint) 2-surface we may call $`𝒮_\sigma `$. The result of our construction is a Hilbert space $`𝒱_\sigma ^P`$ which is the space of intertwiners on $`𝒮_\sigma `$ subject to the condition that there are projection operators for the spinor representation on the $`R`$ oriented circles $`c_w`$.
For each $`P`$ and $`R`$ the space of states lives in the auxillary Hilbert space given by
$$_{aux}^{P,R}=\underset{\sigma }{}𝒱_\sigma ^P$$
(70)
where the sum is over all sets $`\sigma `$ of $`R`$ pairs of punctures. For each $`P`$ we may let $`R\mathrm{}`$ to define
$$_{aux}^P=\underset{R\mathrm{}}{lim}_{aux}^{P,R}.$$
(71)
The auxilliary space (71) is larger than the physical state space in that a given conformal block may be created many different ways by joining states on punctured tori in this way. As a result there is a non-trivial action of the group, $``$ of modular transformations, or recoupling relations of the conformal blocks on $`_{aux}^P`$. (There are also linear relations between states in different $`_{aux}^P`$ for different $`P`$.) We then define the physical state space to be,
$$_{phys}^P=_{aux}^P/.$$
(72)
It is not difficult to argue that any conformal block on any compact surface of genus higher than $`1`$ can be built up in this way. As a result we have
$$_{phys}^P=\underset{g1}{}𝒱_g$$
(73)
where $`𝒱_g`$ is the space of intertiners on the compact surface of genus $`g`$. For each $`P`$ this is the physical Hilbert space of the theory. We note that becaus there is no limit on the numbers of punctures all the Hilbert spaces are identical for $`P2`$.
The translation from states on punctured tori to states on compact surfaces is illustrated in Figure (2).
We note that each $`𝒮_\sigma `$ can be decomposed non-uniquely into a set of $`P`$ 2-tori joined on pairs of circles. Given a basis for each of the spaces of intertwiners on the punctured two tori then each such decomposition gives a basis for $`_{aux}^P`$. We arbitrarily fix one such basis, which we will call the reference basis. This is necessary because the construction of the Hilbert space from a set of solutions to the quantum Gauss law constraint produces such a basis. As we will see Hamiltonian comes expressed in that basis.
Finally we note that the state space ( described in 72) is almost the same as that discussed in The main difference is that the zero genus state is excluded from the present theory, and in the proposal of we considered only states arising from connected two-surfaces. This means that the state space we have arrived at is a natural background independent extension of the spin network states of quantum general relativity and supergravity, with $`SU(16)`$ in place of $`SU(2)`$ or $`SU(2|1)`$.
## 7 Action of the Hamiltonian and evolution rules
We now construct the action of the remaining cubic term (47) in the Hamiltonian on the Hilbert space (72). The hamiltonian is constructed in terms of the operators $`b_{\pm I}^J`$ and $`p_{\pm I}^J`$ which means their action is defined on the auxillary space (71). The action on the physical space (72) is defined by imposing equivalence under modular transformations. We have
$`{\displaystyle \underset{I<J<K}{}}`$ $`{\displaystyle }`$ $`dt_Idt_Jdt_KH^{IJK}(t_I,t_J,t_K)\mathrm{\Psi }(\{t_I\})=`$ (74)
$`3{\displaystyle \underset{\pm }{}}{\displaystyle d^2x_Id^2x_Jd^2x_KTr_G\left\{\underset{I<J<K}{}b_{\pm I}^Jb_{\pm J}^Kp_{\pm I}^K+\frac{4\pi }{k}\underset{I<K<J}{}b_{\pm I}^Jp_{\pm K}^Jp_{\pm I}^K\right\}}`$
The action of $`H_3`$ on states is given in Figure (3).
$`H_3`$ has two kinds of actions. Two of the terms eliminate two links labeled by the $`b`$’s and creates a new one as shown in Fig. (3). In this case the action of $`H`$ eliminates a positive and negative puncture on a single torus. It then induces a map from $`𝒱_{s\overline{s},r_2,\mathrm{}}𝒱_{r_2,\mathrm{}}`$, where $`s`$ is the fundamental representation, which is defined by taking the intertwiner $`s\overline{s}Id`$. The other terms act in the opposite way to eliminate a $`b_I^K`$ joining tori $`I`$ and $`K`$ and create, for every distinct torus $`IJK`$ a link from $`I`$ to $`J`$ and a link from $`J`$ to $`K`$. The action on the space of intertwiners on $`J`$ then acts in the opposite way, through the channel $`Ids\overline{s}`$.
Translated to an action on states on compact two surfaces rather than punctured tori the action of $`H_3`$ is illustrated in Figure (4).
We may note that the action has exactly the form of the Hamiltonian constraint for quantum general relativity, found in <sup>3</sup><sup>3</sup>3different variants are discussed in . with the replacement of $`SU(2)`$ by the quantum deformation of $`SU(16)`$. This can be seen in Figure (5) where we indicate schematically the action of the Hamiltonian constraint of general relativity found in .
This means that if we restrict $`SU(16)`$ to an $`SU(2)`$ subgroup, the dynamics predicted by the cubic matrix model will be of the same form as that found from canonical quantization of quantum general relativity. In terms of the dual picture of this gives the $`13`$ and $`31`$ Pachner moves. Both the $`13`$ and $`31`$ moves arise in the theory with the same amplitude, showing that the Hamiltonian is hermitian.
The fact that the present theory is at finite $`k`$ means, moreover, that the difficulty found for quantum general relativity, in which there is no long range propagation of information, and hence no chance of a classical limit need not be present. The reason, as pointed out in is that at finite $`k`$ the missing moves which are required to have propagation of information are the $`22`$ Pachner moves, shown in Figure (6). However, these moves are present for all theories with finite $`k`$ as they are just changes of basis in the space of intertwiners as illustrated in Figure (6). Thus, if we consider a sequence of moves for finite $`k`$ induced by $`H_3`$ we may insert at any point the change of basis shown in Fig. 5. For finite $`k`$ this is just a change of basis and the history in terms of finite $`k`$ states is unchanged whether we do it or not. But this operation does not commute with the limit $`k\mathrm{}`$. This means that the histories that must be included in the sum over histories at $`k\mathrm{}`$ on the spin network states includes the insertion of all possible finite $`k`$ change of basis moves. Since these are no longer changes of bases for $`k=\mathrm{}`$ this is equivalent to the insertion in the $`k=\mathrm{}`$ hamiltonian constraint of the change of basis move, with the amplitudes given by the $`k\mathrm{}`$ limit of the $`6j`$ symbol. This is equivalent to saying that the theory has these $`22`$ moves with the amplitudes given by the $`BF`$ or Crane-Yetter-Ooguri theory, while the $`13`$ and $`31`$ moves have the non-topological form given by $`H_3`$.
However, with the inclusion of these $`22`$ moves the problem with lack of long range correlations found in is solved, as pointed out in .
Finally, we note that the fact that the $`22`$ moves have a different amplitude than the $`13`$ move means that there is a dependence of the amplitudes on the causal structure of the quantum history. The theory does not have the crossing symmetry required by a Euclidean quantum gravity theory, it is then intrinsically a causal theory in the sense of .
## 8 Compactification to the dWHN-BFSS model
In this section we describe very briefly a different compactifcations which appears to lead to the dWHN-BFSS model at the one loop level. This will be discussed in detail in . The required compactification involves a simultaneous compactification on three directions, but it is different than the one described in section 4. There we considered a compactification that broke the $`U(P)`$ symmetry involving the indices $`I,J,\mathrm{}`$ while preserving the whole $`SU(16)`$ symmetry. Now we consider a compactification that does the reverse: we preserve the $`U(P)`$ symmetry, but break the $`SU(16)`$ symmetry.
To do this we break the fields in terms of $`SU(16)`$ traces and trace-free parts:
$$A_{aIi_0i_1i_2P}^{Jj_0j_1j_2Q}=\delta _P^Q𝒜_{aIi_0i_1i_2}^{Jj_0j_1j_2}+\stackrel{~}{A}_{aIi_0i_1i_2P}^{Jj_0j_1j_2Q}$$
(75)
where $`\stackrel{~}{A}_{aIi_0i_1i_2P}^{Jj_0j_1j_2P}=0`$.
We begin as before by expanding around the classical solution
$$𝒜_{aIi_0i_1i_2}^{Jj_0j_1j_2}=\delta _I^J(_a)_{i_0i_1i_2}^{j_0j_1j_2};\stackrel{~}{A}_{aIi_0i_1i_2P}^{Jj_0j_1j_2Q}=0$$
(76)
However, what is different in this case is that we define the compactification in such a way that all of the $`SU(16)`$ components of the fields becoming $`SU(P)`$ matrices. To do this we employ a parameterization that preserves the $`U(P)`$ gauge symmetry which is
$$A_{aIi_0i_1i_2P}^{Jj_0j_1j_2Q}=\delta _P^Q\left(\delta _I^J_{ai_0i_1i_2}^{j_0j_1j_2}+a_{aIi_0i_1i_2}^{Jj_0j_1j_2}\right)+\stackrel{~}{A}_{aIi_0i_1i_2P}^{Jj_0j_1j_2Q}$$
(77)
We will also restrict the fields to real values, so that the symmetry is reduced from $`SU(16,16|1)`$ to $`Osp(1|32)`$.
We now compactify on the three directions associated with $`x^+`$, $`x^{}`$ and $`x^2`$. This is similar to the compactification we used up till this point, but we will treat it in a rather different manner. Rather than considering the limit in which all three compactification radii go to infinity, which leads us to a Chern-Simons type theory, we will keep the compactification radii $`L_{2,}=(2M_{2,}+1)l_{Pl}`$ small, while taking only the third radius $`L_+=(2M_++1)l_{Pl}`$ large. This breaks the $`SO(2,1)=Sp(2)`$ symmetry. The gauged action can be written as
$`I^{gauged}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle _{T^3}}dx_+dx_{}dx_2Tr_{U(P)}\{Y^{CS}(a)`$
$`+ϵ^{abc}\left(3\stackrel{~}{A}_{aP}^Q𝒟_b\stackrel{~}{A}_{cQ}^R+\stackrel{~}{A}_{aP}^Q\stackrel{~}{A}_{bQ}^R\stackrel{~}{A}_{cR}^P\right)`$
$`+Y^3+\eta ^{ab}\stackrel{~}{Y}_P^Q(\{𝒟_a,\stackrel{~}{A}_{bQ}^P\}+\stackrel{~}{A}_{aQ}^R\stackrel{~}{A}_{bR}^P)`$
$`+\chi ^P\tau _{PQ}^a𝒟_a\chi ^Q\}`$
although we must remember that as we have not taken the compactification radii to infinity what is really meant is the cutoff version of this theory with $`2M_a+1`$ modes in each direction. An important consequence of not taking the limit of infinite compactification radii is that the theory knows about the $`2+1`$ dimensional metric $`\eta ^{ab}`$.
The next step is to compute the effective action. This will be described in , here we only discuss how the symmetries of the theory constrain its form.
The form of the effective potential will be governed by the symmetries of the action which are left unbroken by the compactification. The local bosonic symmetries includes $`U(P)`$. In the case of compactifications with independent compactification radii $`L_+L_{}L_2`$, the $`SU(16,16)`$ symmetry has been broken down into a subgroup, which is $`SO(9)`$. This can be seen from the fact that in a standard parameterization of the components of the matrix $`M`$ in terms of $`32`$ component $`\mathrm{\Gamma }`$ matrices for $`Spin(10,1)`$, the $`+,`$ and $`2`$ components that defined the three $`\tau ^a`$ matrices that generate the $`Sp(2)`$ symmetry are in the $`0`$, $`10`$ and $`010`$ boost directions .
There are in addition the remaining translation symmetries (18) involving the tracefree parts, $`\stackrel{~}{A}_{aP}^Q`$. These are symmetries of the action in the case that we restrict all fields to the real slice. This is one reason we have restricted the fields to the real case. Finally, one can check that $`16`$ of the $`32`$ supersymmetries have been broken, so that there remain $`16`$ unbroken supersymmetry generators.
The lowest dimensional terms that can appear in the effective action, consistent with these symmetries are,
$`I^{1loop}`$ $`=`$ $`{\displaystyle _{T^3}}dx_+dx_{}dx_2Tr_{U(P)}\{f_{ab}f_{cd}\eta ^{ac}\eta ^{bd}+[𝒟_a,\stackrel{~}{A}_b][𝒟_c,\stackrel{~}{A}_d]\eta ^{ac}\eta ^{bd}`$
$`+(𝒟_a\stackrel{~}{A}_b\eta ^{ab})^2+fermions+interactions\}`$
We now consider the limit
$$L_{},L_2l_{Pl},L_+\mathrm{}.$$
(78)
At the same time we boost to the infinite momentum frame in the $`+`$ direction. All fields decouple except those that have the maximal number of $`_+`$ derivatives. These turn out to be the $`\stackrel{~}{A}_{}`$ and $`\mathrm{\Psi }`$ fields.
In this limit the only terms that survive in the kinetic energy are
$$I^{1loop}=𝑑x_+Tr_{U(P)}\left\{(_+\stackrel{~}{A}_{})^2+\mathrm{\Psi }^P_+\mathrm{\Psi }_P\right\}$$
(79)
We see that the fields which remain consist of one $`16`$ component fermion together with the symmetric tensor $`A_{}^{PQ}`$. This decomposes as
$$A_{PQ}=\mathrm{\Gamma }_{PQ}^\mu X_\mu +\mathrm{\Gamma }_{PQ}^{\mu \nu \rho \sigma }X_{\mu \nu \rho \sigma }$$
(80)
The $`X_\mu `$ are the nine transverse matrices of the dWHN-BFSS theory, which correspond to the $`D0`$ brane coordinates as well as to the light cone gauge coordinates of the embedding of a membrane. The $`X_{\mu \nu \rho \sigma }`$ are additional degrees of freedom that do not appear in that model.
The leading terms in a derivative expansion of the effective action will be completed by interaction terms amongst these degrees of freedom. These will be determined by the unbroken symmetries, which are precisely the symmetries of the dWHN-BFSS theory. The action involving $`\mathrm{\Psi }^P`$ and $`X^\mu `$ with these symmetries is exactly the dWHN-BFSS theory. One can check that this is extended by terms involving the four form field $`X_{\mu \nu \rho \sigma }`$. The full action invariant under the translations (18), $`SO(9)`$ rotations and $`16`$ supersymmetry generators is
$`I^{1loop}`$ $`=`$ $`{\displaystyle }dx_+Tr_{U(P)}\{(_+\stackrel{~}{A}_{})^2+\mathrm{\Psi }^P_+\mathrm{\Psi }_P`$ (82)
$`+\mathrm{\Psi }^P[\stackrel{~}{A}_{PQ},\mathrm{\Psi }^Q]+\stackrel{~}{A}_{[P}^Q\stackrel{~}{A}_{R]Q}\stackrel{~}{A}_{}^{PS}\stackrel{~}{A}_S^R\}`$
$`=`$ $`{\displaystyle }dx_+Tr_{U(P)}\{(_+X_\mu )^2+(_+X_{\mu \nu \rho \sigma })^2+\mathrm{\Psi }^P_+\mathrm{\Psi }_P`$
$`+\mathrm{\Psi }^P\mathrm{\Gamma }^{\mu PQ}[X_\mu ,\mathrm{\Psi }^Q]+\mathrm{\Psi }^P\mathrm{\Gamma }^{\mu \nu \rho \sigma PQ}[X_{\mu \nu \rho \sigma },\mathrm{\Psi }^Q]`$
$`+[X^\mu ,X^\nu ][X_\mu ,X_\nu ]`$
$`+[X^{\mu \nu \rho \sigma },X^{\alpha \beta \gamma \delta }][X_{\mu \nu \rho \sigma },X_{\alpha \beta \gamma \delta }]+[X^\mu ,X^{\alpha \beta \gamma \delta }][X_\mu ,X_{\alpha \beta \gamma \delta }]\}`$
One can check that the $`X_{\alpha \beta \gamma \delta }`$ terms do not break the supersymmetry of the dWHN-BFSS theory. Instead,the symmetry algebra is extended by central charges that exist for any $`P`$. In the standard dWHN-BFSS model these central charges appear only in the limit $`P\mathrm{}`$. In that case the central charges are interpreted to describe certain components the $`5`$-brane. This suggests that the new degrees of freedom $`X_{\alpha \beta \gamma \delta }`$ provide an additional description of the $`5`$-brane, wrapped on the $`X^+`$ and $`X^{}`$ directions. This will be discussed in more detail elsewhere.
## 9 The duality between strings and loops
From the results of the last several sections it is possible to describe precisely how the BFSS matrix model and the $`SU(16)_q`$ extension of loop quantum gravity give equivalent descriptions of the same degrees of freedom. We can do this by tracing how each description is derived from a compactification of the cubic matrix model.
In both cases we make use of the decomposition $`SU(16,16|1)Sp(2)SU(16|1)`$ and decompose the fields in terms of the basic degrees of freedom,
$$A_{aIi_aP}^{Jj_bQ}$$
(83)
The $`a`$ is the $`Sp(2)`$ index and the $`P`$ and $`Q`$ are $`SU(16)`$ indices. We have then decomposed the $`N\times N`$ components of the matrices in terms of two sets of indices so that
$$i=(i_a,I)$$
(84)
where $`i_a=1,\mathrm{},M_a`$ are the indices that will give rise to compactification on a three torus, and $`I=1,\mathrm{},P`$ defines a remaining $`U(P)`$ symmetry.
In both cases we begin by defining a compactification on a $`3`$-torus, leaving aside a $`U(P)`$ symmetry. Each compactification is then defined by an expansion around the classical solution
$$(A_a^0)_{Ii_aP}^{Jj_bQ}=\delta _I^J\delta _P^Q(_a)_{i_a}^{j_a}$$
(85)
Each then defines a fluctuating field
$$A_{aIi_aP}^{Jj_bQ}=\delta _I^J\delta _P^Q(_a)_{i_a}^{j_a}+𝒜_{aIi_aP}^{Jj_bQ}$$
(86)
The stringy and loopy descriptions then part company. In each case one of the two remaining symmetries, which are $`U(P)`$ and $`SU(16)`$, are broken and the other is treated as a gauge symmetry in the background created by the compactification. The two compactifications differ as to which components of the fluctuating field $`𝒜_{aIi_aP}^{Jj_bQ}`$ are treated as a gauge field and which are taken as a matter-like field.
There is freedom to chose which symmetry appears as an ordinary gauge symmetry because before compactification the fields have three index types: $`A_{ai_aIP}^{j_aJQ}`$. The $`i_a`$ describe the fourier modes of the spatial dependence on the $`3`$ torus, while the $`I`$ and $`P`$ parameterize respectively $`SU(P)`$ and $`SU(16)`$. The gauge symmetry as it appears as a local symmetry in the three dimensional field theory on the $`T^3`$’s comes from a symmetry of the form of either (16) or (15). But there are now three sets of indices, which must be grouped into two sets. The first set determines which indices the gauge parameters depend on, while the second determines the gauge group. For example if we group the $`i_a`$ with the $`I`$ then we have gauge transformations of the form
$$A_{i_aI}^{j_aJ}U(i_aI)A_{i_aI}^{j_aJ}U(j_aJ)^{}$$
(87)
where $`U(i_aI)SU(16)`$. This gives an $`SU(16)`$ gauge symmetry that acts locally in each of the tori labeled by $`I`$. This is the symmetry of the background independent multi-tori compactifications. In this representation of the theory the $`U(P)`$ symmetry is hidden. It is possible that it is expressed by the recoupling relations of the intertwiners of $`SU(16)_q`$ (equivalently the modular group of the 2-surfaces acting on the conformal blocks).
In this case we defined the components of the fluctuating fields as
$$𝒜_{aIi_aP}^{Jj_bQ}=\delta _I^JA_{aPi_a}^{IQj_b}+B_{aIi_aP}^{Jj_bQ}$$
(88)
We treated the diagonal components $`A_{aPi_a}^{IQj_b}`$ as gauge fields for the $`SU(16)`$ symmetry. There are $`P`$ of them, corresponding to the diagonal components of $`U(P)`$ matrices, so we have a separate gauge invariance on each of $`P`$ tori. The off-diagonal (in $`U(P)`$) components $`B_{aIi_aP}^{Jj_bQ}`$ defined only for $`IJ`$ define interactions between the $`P`$ Chern-Simons theories defined by the diagonal components $`a_{aPi_a}^{IQj_b}`$. The $`U(P)`$ symmetry is broken by the choice of polarization, which requires an arbitrary ordering of the $`P`$ Chern-Simons theories.
The other choice is to make manifest the gauge symmetry that comes from grouping the $`i_a`$ indices with the $`SU(16)`$ indices. In this case we write,
$$A_{i_aP}^{j_aQ}W(i_aP)A_{i_aP}^{j_aQ}W(j_aQ)^{}$$
(89)
where $`W(i_aP)SU(P)`$. In this case we get a local $`SU(P)`$ gauge symmetry on a single torus. This leads to the case of the stringy compactification, which is expressed in terms of a local $`SU(P)`$ gauge theory, which we argued leads to the dWHN-BFSS matrix model.
In this case we we defined the components of the fluctuating field as
$$𝒜_{aIi_aP}^{Jj_bQ}=\delta _P^Qa_{aIi_a}^{Jj_b}+\stackrel{~}{𝒜}_{aIi_aP}^{Jj_bQ}$$
(90)
where $`\stackrel{~}{𝒜}_{aIi_aP}^{Jj_bP}=0`$. The scalar (in $`SU(16)`$ terms) components, $`a_{aIi_a}^{Jj_b}`$ define the fluctuations in the compactification radii, and then must become the string theory moduli. The theory is treated from that point on as an $`U(P)`$ gauge theory, where the scalars $`a_{aIi_a}^{Jj_b}`$ play the role of gauge fields in the compactified directions and the $`SU(16|1)`$ symmetry is broken down to the superpoincare algebra in $`10+1`$ dimensions or the super-euclidean algebra in $`9`$ dimensions.
The different theories are treated differently in other ways, particularly in the fact that to get the standard matrix models two or three of the compactified radii must be taken to the Planck scale, leading to a low energy theory defined in either $`0+1`$ or $`0+0`$ spacetime dimensions (in the dWHN-BFSS and IKKT cases, respectively.) But the essential difference between them is defined by these two different coordinatizations of the fluctuating fields around the classical solution (85).
As a consequence of the identifications defined here we have a genuine translation between the loopy and stringy descriptions of the kinematics and dynamics of the cubic matrix model. This may be used to translate problems from one description to the other. For example, the $`P\mathrm{}`$ limit which is problematic for the conventional matrix models is clearly related to the limit of loop quantum gravity in which the universe grows infinitely large. On the other side we can say that the continuum limit of the loopy description may involve a restoration of the $`U(P)`$ symmetry which is broken by the quantization described here. The correspondence may also be used to translate the $`D`$-brane description of black hole horizons into the loop quantum gravity description, which is based on describing the state space on the horizon in terms of conformal blocks. This may make possible a description of black holes which is not restricted to the near extremal case of positive specific heat. Similarly, we may try to use the correspondence to extend the description of boundaries with non-zero cosmological constant given in to arrive at a detailed description of the $`AdS/CFT`$ correspondence in $`3+1`$ dimensions.
The duality can also be expressed by considering how the fundamental excitations are described in the two pictures. In the multi-Chern-Simons compactifications, the fundamental degrees of freedom are the pairs of punctures which are created or destroyed by the dynamical terms in the Hamiltonian. It is then tempting to see them as the background independent analogue of $`D_0`$ branes. There is, in fact, some more direct evidence that these excitations are related to $`D_0`$ branes, which is described in . We then have two different representations of these degrees of freedom, either in terms of the light cone gauge components of the matrices $`A_I^J`$, which lead to the standard matrix description of $`D0`$ branes, or as the operators which create and annihilate punctures which join the $`P`$ tori. By tracing through the correspondence we have just discussed we can see explicitly how the two descriptions of the fundamental degrees of freedom may be translated into each other. It is then perhaps fitting to call this the string/loop duality.
## 10 Conclusions
In this paper we have studied a new kind of compactification for $``$ theory, which is defined, not in terms of a background geometry, but in terms of an algebra, that is the fusion algebra of conformal blocks for the quantum deformation of $`SU(16)`$ on arbitrary 2-surfaces. We may then call this an algebraic compactification. It is clear from the construction that many other algebraic compactifications can be defined corresponding to the reduction of the representation theory of quantum deformed $`SU(16)`$ to the representation theory of its subalgebras.
There are many issues raised by this formulation. Some are technical. The most pressing of these include the need to study in detail the quantum group extension of $`SU(16)`$ defined by a Chern-Simons theory with a connection valued in that algebra and the need to develop the details of its representation theory. Another set of issues to be discussed in involve the details of the supersymmetric extension of the results described here. We want to understand also the way in which the modular group acting on compact two surfaces is related to the $`U(P)`$ symmetry that was broken in the multi-Chern-Simons compactification. The computation of the effective action which we sketched above will be discussed in detail in . The relationship between the real and complex forms of the theory needs more investigation. Finally, there remain subtle issues associated with the role of the $`Y`$ degrees of freedom.
On the conceptual side, the role of the many fingered time of the form found here needs further study. The action appears to be non-local in time, but the multi-time quantization discussed here seems to lead to a sensible notion of a history. One may begin with an initial state at a time $`t_I=0`$ for all $`t_I`$ and then act repeatedly by the cubic hamiltonian to evolve, on each action, triples of Hilbert spaces forward in time. In this way one can generate a history, which seems to be of a form which is closely related to that described in . In that case acting repeatedly with the moves generates a causal history, of the general form studied in . The causal structure results from the ordering dependence of the action of the different terms in the Hamiltonian, the time labels themselves seem irrevelevant apart from ordering. This structure is very suggestive but deserves further study. The relationship to the histories projection operator formulation developed in is also very suggestive.
Another interesting issue to discuss is the relationship of this formulation with the holographic principle and the related problem of the interpretation of quantum theories of cosmology. The construction in was motivated in part by the holographic principle, which we showed in appears naturally in loop quantum gravity when boundaries are considered. It is interesting to consider each one of the Hilbert spaces $`H(t_I)`$ to be a screen, in a background independent formulation of the holographic principle of the kind described in . We may then to try to follow the argument there and construct the quantum geometry by defining the area of each screen to be the log of the dimension of its space of intertwiners, as in .
When we take the limit $`k\mathrm{}`$ we reproduce, as argued in exactly the spin network states of quantum general relativity or supergravity. Moreover, the cubic action gives rise to evolution rules which are exactly of the same form as follow from a first principles, canonical quantization of general relativity. The main difference is that because the theory is defined at finite $`k`$, new evolution rules must appear in the $`k\mathrm{}`$ limit, exactly of the form required to cure a major problem of loop quantum gravity, which is the absence of long ranged correlations at zero cosmological constant.
It is very interesting to note that there are other derivations of path integral formulations for loop quantum gravity from matrix models. It is possible that there is a direct derivation from the cubic matrix model to matrix models of the form used in those derivations, which is induced by quantum corrections to the present model. This is presently under investigation. Also of interest is the question of whether the topological field theory parameterization of the degrees of freedom of $`11`$ dimensional supergravity introduced in can be derived directly from the cubic matrix model studied here.
Finally, any approach to a background independent formulation of $``$ theory must answer the question of how the particular structures which seem required by string theory for perturbative consistency of a quantum theory of gravity are picked out at the more fundamental, background independent level. It is quite possible that the theory presented here has a more fundamental formulation in which the choice of algebra is not arbitrary. The possibility of a reformulation in terms of Jordan algebras and octonions comes naturally to mind. The $`SU(16,16)`$ structure reminds one of an octonionic extension of twistor theory. There is also an intriguing similarity to Chern-Simons inspired formulations of string field theory that deserves further investigation.
## ACKNOWLEDGEMENTS
I am grateful to Louis Crane, Bartomeu Fiol, Willy Fischler, Mike Green, Chris Hull, Chris Isham, Yi Ling, Fotini Markopoulou, George Minic, Mike Reisenberger, Kelle Stelle and Dennis Sullivan for discussions, suggestions and encouragement during the course of this work. In addition, Richard Levine contributed a number of extremely helpful observations and corrections. I am grateful also for hospitality at the Institute in Theoretical Physics in Santa Barbara, where this work was begun, and to the theory group at Rutgers University, where it was completed. This work was supported by the NSF through grant PHY95-14240 and by an SPG grant. I would also like to thank the Jesse Phillips Foundation for support and encouragement.
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# Quantum Coherence and Decoherence in Magnetic Nanostructures
## I Introduction
Is quantum computing physics or engineering? Some time ago I asked this question to Murray Gell-Mann who replied that he had a formal proof that it is physics: articles on quantum computing are published by The Physical Review Letters. In the spirit of his answer I will concentrate on the physics aspects of the problem, leaving out the engineering aspects. The purpose of this article is to discuss the prospect of using small magnetic clusters as qubits.
Qubits based upon quantum superposition of the $`|>`$ and $`|>`$ states of individual electrons and nuclei have been discussed during all years of quantum computermania . In parallel, but with little overlap, there has been intensive theoretical and experimental research on spin tunneling between the $`|>`$ and $`|>`$ states of molecular magnets and in ferromagnetic and antiferromagnetic nanoparticles . These are composite objects with total spin $`S`$ ranging from $`S=10`$ for Fe-8 and Mn-12 molecular clusters to $`S`$ of a few thousand in nanoparticles. Experimental study of the magnetization reversal in individual nanoparticles gives evidence that very small particles are uniformly magnetized . At low enough temperature they, like molecular magnets, possess a large fixed-length spin formed by the strong exchange interaction.
Quantum tunneling between the $`|>`$ and $`|>`$ states of spin-10 molecular nanomagnets has been unambiguously established in experiment . Early measurements of low temperature magnetic relaxation and more recent measurements of individual nanoparticles also provided evidence of quantum tunneling of spin. There is experimental evidence of quantum coherent oscillations between the $`|>`$ and $`|>`$ states in nanoparticles and molecular magnets .
On one hand, individual nanoparticles and high-spin molecules, because of their large magnetic moments, must be easier to operate as qubits than individual electron and nuclear spins. In fact, the techniques of measuring individual particles of $`S10^210^3`$ already exist . On the other hand, as the size of the system increases, its interaction with the dissipative environment also increases and it is not obvious whether the decoherence in large spin systems can be made low enough to allow their application as qubits. To achieve this goal the decoherence rate $`\gamma `$ must be made small compared to the frequency $`\mathrm{\Delta }/\mathrm{}`$ of coherent oscillations between the $`|>`$ and $`|>`$ states. For a large spin system it requires effort for two reasons. Firstly, the interaction with the environment is proportional to the size of the system. Secondly, $`\mathrm{log}(\mathrm{\Delta })`$ scales linearly with $`S`$ so that, except for some special cases, $`\mathrm{\Delta }`$ becomes immeasurably small for $`S30`$.
In this article we will make suggestions on how to obtain large $`\mathrm{\Delta }`$ for a relatively large spin and how to achieve the condition $`\gamma <<\mathrm{\Delta }/\mathrm{}`$. It is our believe that this is doable in spin systems, in part due to their special properties with respect to time reversal. Spin Hamiltonians and tunneling rates will be considered in Section 2. Mechanisms of decoherence will be discussed in Section 3.
## II Coherence
### A Ferromagnetic clusters
The basic Hamiltonian for a ferromagnetic cluster of spin $`S`$ that exhibits quantum coherence is
$$H=AS_z^2+V,$$
(1)
where $`A`$ is a positive constant and $`V`$ is an operator that does not commute with $`S_z`$ but is invariant under the transformation $`S_zS_z`$. The first term in Eq.(1) is typically produced by a crystal field or by the shape anisotropy of the cluster. At $`V=0`$ the Hamiltonian (1) has a double degenerate ground state corresponding to two opposite orientations, $`|>`$ and $`|>`$, of $`𝐒`$ along the Z-axis. In terms of the magnetic quantum number, $`S_z|m>=m|m>`$, these ground states are $`|S>`$ and $`|S>`$. Note that this degeneracy is independent of any geometry of the problem including the shape of the magnetic cluster and is solely due to the odd symmetry of $`𝐒`$ with respect to time reversal.
Our case of interest will be an integer large $`S`$. In that case the problem is almost classical so that even a large perturbation $`V0`$ generates only a small probability of tunneling between the $`|>`$ and $`|>`$ states. The degeneracy of the ground state is then removed and the new ground state can be approximated by
$$|0>=\frac{1}{\sqrt{2}}(|>+|>).$$
(2)
It is separated from the first excited state,
$$|1>=\frac{1}{\sqrt{2}}(|>|>),$$
(3)
by the energy gap $`\mathrm{\Delta }`$ which is determined by the strength and the nature of $`V`$ and is small compared to the scale $`A`$ of the energy levels of $`AS_z^2`$. The probability of finding $`𝐒`$ looking up (down) then oscillates in time according $`\mathrm{cos}(\mathrm{\Delta }t/\mathrm{})`$, which is quantum coherence.
One of the possible forms of $`V`$ is $`V=g\mu _BH_xS_x`$ induced by the external field applied along the X-axis. If the field is small, the tunneling splitting can be obtained by the perturbation theory
$$\mathrm{\Delta }=\frac{4AS}{(2S1)!}\left(\frac{g\mu _B|H_x|}{2A}\right)^{2S}.$$
(4)
¿From practical point of view this case is not very promising since the inevitably present weak misorientation of the field, resulting in $`H_z0`$, will destroy the coherence.
A more promising case, which to a good approximation corresponds to the Fe-8 spin-10 molecular nanomagnet, is $`V=BS_x^2`$. At $`|B|<<A`$ the perturbation theory gives
$$\mathrm{\Delta }=8A\frac{(2S)!}{[(S1)!]^2}\left(\frac{|B|}{16A}\right)^S.$$
(5)
In Fe-8 $`\mathrm{\Delta }/\mathrm{}`$ is of order $`10^4`$s<sup>-1</sup>.
For large $`S`$ and arbitrary $`B`$ the tunneling splitting has been computed by the instanton method and by mapping the spin problem onto a particle problem
$$\mathrm{\Delta }=16\pi ^{1/2}S^{3/2}|B|\left[\frac{A}{|B|}\left(1+\frac{A}{|B|}\right)\right]^{3/4}\left[\left(1+\frac{A}{|B|}\right)^{1/2}+\left(\frac{A}{|B|}\right)^{1/2}\right]^{(2S+1)}.$$
(6)
¿From practical point of view, in ferromagnetic nanoparticles with $`S>>1`$ the case of interest is $`|B|>>A`$. In that case Eq.(6) gives
$$\mathrm{\Delta }=16\pi ^{1/2}S^{3/2}A^{3/4}|B|^{1/4}\mathrm{exp}\left[2S\left(\frac{A}{|B|}\right)^{1/2}\right],$$
(7)
and one can see that the effect of large $`S`$ in the exponent is suppressed by a a small factor $`(A/|B|)^{1/2}`$. There are two ways to achieve this suppression and to increase $`\mathrm{\Delta }`$. The first is to use magnetic clusters with very strong easy plane anisotropy and relatively weak easy axis anisotropy in that plane. Particles of Tb and Dy may satisfy this condition. The second way is to place the particle above the surface of a superconductor . In that case the magnetic dipole interaction of $`𝐒`$ with its mirror image inside the superconductor effectively reduces the uniaxial anisotropy $`A`$. The particle and the superconductor should be selected such that the magnetic field induced by the particle, $`H\mathrm{\hspace{0.17em}4}\pi \mu _BS`$, does not exceed the first critical field of the superconductor, $`H_{c1}`$. Manipulating the distance between the particle and the superconductor, one can achieve the condition $`A<<|B|`$. This can be done by, e.g., controlling the distance with a ferroelectric buffer in the electric field. Note that in the absence of the external magnetic field the odd symmetry of $`𝐒`$ with respect to time reversal preserves the coherence of such a setup independently of the shape of the superconducting surface and electric fields in the problem.
The tunneling of $`𝐒`$ also can be induced by its hyperfine interaction with nuclear spins, $`V=B𝐒𝐈`$, where $`𝐈=𝐈_i`$ is the total nuclear spin of the cluster obtained by summing over spins of individual nuclei. This problem is rather involved. It has been studied in Ref. and is relevant to tunneling in Mn-12. The total Hamiltonian conserves the magnitude of $`I`$ and the Z-projection of $`𝐒+𝐈`$. In the millikelvin range nuclear spins must order, developing $`I_{max}=|I_i|`$. It is easy to see that the problem is the one of quantum coherence only if $`I_{max}=S`$. In that case the classical ground states correspond to $`𝐒`$ and $`𝐈`$ of equal length looking opposite to each other along the Z-axis, $`|S>|I_{max}>`$ and $`|S>|I_{max}>`$. Tunneling removes the degeneracy of the ground state. The corresponding splitting can be obtained by the perturbation theory for $`B<<A`$:
$$\mathrm{\Delta }=8(A+B)S^2\left[\frac{B}{2(A+B)}\right]^{2S}.$$
(8)
In Mn-12 $`S=10`$ while $`I_{max}=30`$, so that the coherence of the above type is impossible. It is not out of the question, however, that chemists will produce a molecular cluster with $`S=I`$ in the future.
### B Antiferromagnetic clusters
Tunneling in antiferromagnetic clusters turns out to be much stronger than in ferromagnetic clusters, making them promising candidates for quantum coherence. Consider an anisotropic antiferromagnetic cluster with two compensated sublattices of spin $`𝐒_1`$ and $`𝐒_2`$, described by the Hamiltonian
$$H=A(S_{1z}^2+S_{2z}^2)+B𝐒_1𝐒_2$$
(9)
with positive $`A`$ and $`B`$ satisfying $`A<<B`$.
Let us show that this model can be mapped onto the model with strong transverse anisotropy. The Lagrangian corresponding to Eq.(9) is
$$L=S(\dot{\varphi }_1\mathrm{cos}\theta _1+\dot{\varphi }_2\mathrm{cos}\theta _2)S(\dot{\varphi }_1+\dot{\varphi }_2)+AS^2(\mathrm{cos}^2\theta _1+\mathrm{cos}^2\theta _2)BS^2\mathrm{cos}\theta _1\mathrm{cos}\theta _2BS^2\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{cos}(\varphi _1\varphi _2),$$
(10)
where $`\varphi _1`$, $`\theta _1`$, $`\varphi _2`$, $`\theta _2`$ are spherical coordinates of vectors $`𝐒_1`$ and $`𝐒_2`$ of fixed length $`S`$. The path integral over these angles is dominated by $`\mathrm{cos}\theta _1=\mathrm{cos}\theta _2\mathrm{cos}\theta `$, $`\mathrm{sin}\theta _1=\mathrm{sin}\theta _2\mathrm{sin}\theta `$. Introducing $`\varphi =\varphi _1\varphi _2`$, one obtains, up to a phase term, the effective Lagrangian
$$L_{eff}=S\dot{\varphi }\mathrm{cos}\theta +(2A+B)S^2\mathrm{cos}^2\theta BS^2\mathrm{sin}^2\theta \mathrm{cos}\varphi .$$
(11)
With the notations $`\varphi /2=\mathrm{\Phi }`$ and $`2S=\sigma `$, it can be transformed into
$$L_{eff}=\sigma \dot{\mathrm{\Phi }}\mathrm{cos}\theta +\frac{A}{2}\sigma ^2\mathrm{cos}^2\theta \frac{B}{2}\sigma ^2\mathrm{sin}^2\theta \mathrm{cos}^2\mathrm{\Phi },$$
(12)
which is equivalent to the Hamiltonian
$$H=\frac{A}{2}\sigma _z^2+\frac{B}{2}\sigma _x^2.$$
(13)
One can then use the known result, Eq.(7), to obtain the tunneling splitting,
$$\mathrm{\Delta }=32(2\pi )^{1/2}S^{3/2}A^{3/4}B^{1/4}\mathrm{exp}\left[4S\left(\frac{A}{B}\right)^{1/2}\right].$$
(14)
Here $`B`$ is the exchange constant which is typically $`10^410^6`$ the anisotropy constant $`A`$. Consequently, antiferromagnetic particles consisting of a few thousand magnetic atoms can exhibit a significant tunneling rate between the $`|>`$ and $`|>`$ states.
One should notice, that, in order to manipulate these states by the magnetic field, some magnetic non-compensation of the sublattices is needed. In this case tunneling is possible only due to the presence in the Hamiltonian of the transverse field or the transverse anisotropy, e.g., $`b(S_{1x}^2+S_{2x}^2)`$. Let the non-compensated spin be $`s`$. It has been demonstrated that the antiferromagnetic tunneling with $`\mathrm{\Delta }\mathrm{exp}[4S(A/B)^{1/2}]`$ holds upto $`s(b/B)^{1/2}S<<S`$. At greater $`s`$ it switches to the ferromagnetic tunneling with $`\mathrm{\Delta }\mathrm{exp}[2s(A/b)^{1/2}]`$. Strongly non-compensated ferrimagnetic clusters, like, e.g., Mn-12, are always in the ferromagnetic tunneling regime, while weakly non-compensated ferritin particles can be in the antiferromagnetic tunneling regime.
## III Decoherence
A magnetic cluster of the type described above will always be imbedded in a non-magnetic solid dissipative environment. The potential for low decoherence arises from a number of reasons. The first of them is that strong electrostatic interactions are involved, through exchange couplings, only in the formation of the single spin $`𝐒`$ of the cluster. All other interactions of $`𝐒`$ have relativistic smallness of order $`(v/c)^2`$ to a some power. Due to this fact the ferromagnetic resonance in some materials has a quality factor of one million. The second reason is selection rules for spins due to the time reversal symmetry discussed below.
In quantum computation one is interested in creating an arbitrary superposition of the $`|>`$ and $`|>`$ states,
$$|\mathrm{\Psi }>=C_1|>+C_2|>.$$
(15)
Using equations (2) and (3) this state can be re-written in terms of $`|0>`$ and $`|1>`$:
$$|\mathrm{\Psi }>=C_1^{}|0>+C_2^{}|1>,$$
(16)
where $`C_1^{}=(C_1+C_2)/\sqrt{2}`$ and $`C_2^{}=(C_1C_2)/\sqrt{2}`$. It is clear, therefore, that the spontaneous decay of the excited state $`|1>`$ into the ground state $`|0>`$, accompanied by the emission of the energy quantum $`\mathrm{\Delta }`$, should be a major concern for preserving quantum coherence.
Note that, in principle, there may be decohering processes involving other excited levels of the spin Hamiltonian. For large spin, due to a small tunneling rate, such levels are separated from $`|0>`$ and $`|1>`$ by the energy gap (say $`A`$) that is large compared to $`\mathrm{\Delta }`$. An example of such a process would be an Orbach two-phonon process corresponding to the transition $`|1>|A>`$ caused by the absorption of a phonon, followed by the spontaneous decay $`|A>|0>`$ with the emission of a phonon. Also, there may be processes involving the excited states of the environment. An example of such a process would be a two-phonon Raman process that corresponds to the emission and absorption of two real phonons satisfying $`\mathrm{}(\omega _1\omega _2)=\mathrm{\Delta }`$. All such processes are strongly temperature-dependent. Their strength is measured by $`\mathrm{exp}(A/k_BT)`$ or by some high power of $`T/\mathrm{\Theta }`$ where $`\mathrm{\Theta }10^2K`$ is the Debye temperature . In the millikelvin range the rate of such processes is negligible. Consequently, they are of little concern for the decoherence. On the contrary, processes of the spontaneous decay $`|1>|0>`$ can exist even at T=0. One should therefore concentrate on such processes.
As follows from the previous section, cases of interest for quantum coherence are described by Hamiltonians which contain even powers of spin operators. Due to the time reversal symmetry, this will always be the case in the absence of the external magnetic field. Consequently, the states $`|0>`$ and $`|1>`$ most generally can be written as
$`|0>`$ $`=`$ $`{\displaystyle \underset{m=S}{\overset{S}{}}}\alpha _m|m>`$ (17)
$`|1>`$ $`=`$ $`{\displaystyle \underset{m=S}{\overset{S}{}}}\beta _m|m>,`$ (18)
where $`\alpha _m=\alpha _m`$ and $`\beta _m=\beta _m`$. These equations are exact, as compared to the approximate equations (2) and (3). They simply reflect the fact that $`|m>`$ is a complete set of vectors in the Gilbert space of the spin Hamiltonian and that $`|0>`$ and $`|1>`$ have different symmetry with respect to time reversal.
Let $`K`$ be the antilinear antiunitary operator of time reversal. The spin operator is odd with respect to time reversal, $`K𝐒K^{}=𝐒`$. On the contrary, the spin Hamiltonian that only contains even powers of components of $`𝐒`$ is even with respect to time reversal, $`KHK^{}=H`$. Consider now a decohering operator $`D`$. $`D`$ can be due to the interaction of $`𝐒`$ with phonons, electromagnetic fields, nuclear spins, etc. These interactions have different symmetry with respect to time reversal. Let $`D_o`$ and $`D_e`$ be time-odd and time-even operators respectively. Since the state $`|O>`$ is even with respect to time reversal and $`|1>`$ is odd, the following general statement is true
$$<0|D_e|1>=0.$$
(19)
The spin-phonon interaction is of the form
$$H_{sp}=a_{iklm}S_iS_k\frac{u_l}{r_m}+h.c.,$$
(20)
where $`𝐮`$ is the lattice displacement and $`a_{iklm}`$ is a tensor reflecting the symmetry of the lattice. In terms of the operators of creation and annihilation of phonons
$$𝐮=\frac{i}{(2MN)^{1/2}}\underset{𝐤,\lambda }{}\frac{𝐞_{𝐤,\lambda }e^{i𝐤𝐫}}{(\omega _{k\lambda })^{1/2}}(a_{k\lambda }a_{k\lambda }^{}),$$
(21)
where $`M`$ is the unit cell mass, $`N`$ is the number of cells in the lattice, $`𝐞_{𝐤,\lambda }`$ is the phonon polarization vector, $`\lambda =t,t,l`$, and $`\omega _{k\lambda }=v_\lambda k`$ is the phonon frequency ($`v_\lambda `$ being the speed of sound). The spin-phonon interaction given by Eq.(19) describes transitions between different spin states accompanied by the emission and absorption of phonons. For instance, in the absence of tunneling the excited state $`|S1>`$ of the Hamiltonian $`H=AS_z^2`$, corresponding to the ferromagnetic resonance, can relax to the ground state $`|S>`$ at $`T=0`$ by spontaneously emitting a phonon at a rate
$$\gamma =C\frac{A^2S^2\mathrm{\Delta }^3}{\mathrm{}^4\rho v^5},$$
(22)
where $`\rho `$ is the mass density of the lattice and $`C`$ is a constant of order unity. For, e.g., $`A\mathrm{\hspace{0.17em}1}K`$, $`S\mathrm{\hspace{0.17em}10}^3`$, $`(\mathrm{\Delta }/\mathrm{})\mathrm{\hspace{0.17em}10}^9s^1`$, $`\rho \mathrm{\hspace{0.17em}1}g/cm^3`$, and $`v\mathrm{\hspace{0.17em}10}^5cm/s`$, Eq.(21) gives $`\gamma \mathrm{\hspace{0.17em}10}^3s^1`$.
We should now notice that $`KH_{sp}K^{}=H_{sp}`$, that is, the spin-phonon operator is even with respect to time reversal. Consequently, $`<0|H_{sp}|1>=0`$. This also can be checked by the direct calculation of matrix elements of $`H_{sp}`$ using expressions (17). Thus, $`|1>`$ cannot spontaneously decay into $`|0>`$ with an emission of a phonon and we conclude that in a millikelvin range the spontaneous emission of phonons cannot decohere the coherent superposition of the $`|>`$ and $`|>`$ spin states. This is a strong statement which requires some clarification. Indeed, the spin-phonon interaction originates from the spin-orbit coupling of the form $`H_{so}𝐋𝐒`$ where $`𝐋`$ is the orbital momentum. Its density can be presented as $`\rho ϵ_{ikl}r_k\dot{u}_l`$. This operator seems to have non-zero matrix elements between $`|0>`$ and $`|1>`$ given by equations (17). However, this is only because the formulation of the spin tunneling problem presented above does not account for the conservation of the total angular momentum. In fact, the coherence is possible only if $`𝐋+𝐒=0`$ . Consequently, as $`𝐒`$ tunnels between $`|>`$ and $`|>`$, so does the mechanical angular momentum of the body. This is an analogue of the Mössbauer effect for spin tunneling. The energy associated with the mechanical rotation is $`\mathrm{}^2S^2/2I_{in}`$ where $`I_{in}`$ is the moment of inertia of the solid matrix containing the magnetic cluster. This energy should not exceed $`\mathrm{\Delta }`$, otherwise it would be energetically favorable for $`𝐒`$ to localize in the $`|>`$ or $`|>`$ state. Since $`I_{in}`$ scales as the fifth power of the size of the matrix, it is easy to see that such localization may occur in a free particle of size less then 10 nm, while for bigger systems the conservation of the total angular momentum is rather formal than practical question. With the account of the momentum conservation, the ground state and the first excited state are $`\frac{1}{\sqrt{2}}(|S>|L>\pm |S>|L>)`$, where the absolute values of $`𝐒`$ and $`𝐋`$ are equal. Now, again, the $`|0>`$ state is even with respect to time reversal while the $`|1>`$ state is odd. The spin-orbit operator is even, $`KH_{so}K^{}=H_{so}`$, because $`𝐋`$ is odd, $`K𝐋K^{}=𝐋`$. Thus, as before, $`<0|H_{so}|1>=0`$.
As has been shown in the previous section (see also Refs. 21,22), nuclear spins always destroy the coherence at $`H=0`$ unless tunneling is induced by the hyperfine interaction with $`I=S`$. Except for that exotic possibility, nuclear spins must be always eliminated from magnetic qubits by the isotopic purification. Similarly, the presence of free non-superconducting electrons in the sample will decohere tunneling through the spin scattering of electrons, $`D_e𝐬𝐒`$. Although, this operator is time-even, free electrons incidentally passing through the magnetic cluster, perturb $`|0>`$ and $`|1>`$, breaking their properties with respect to time reversal. Thus, strongly insulating materials should be chosen for magnetic qubits. The effect of incidental phonons due to, e.g., relaxation of elastic stresses in the matrix (the $`1/f`$ noise), should be similar to the effect of incidental electrons in perturbing $`|\mathrm{\Psi }>`$. Thus, the perfection of the lattice should be given a serious thought when manufacturing magnetic or any other qubits. One should then worry about the decohering effect of spin interactions which are odd with respect to time reversal. These are Zeeman terms due to magnetic fields. For example, $`D_o=g\mu _BS_zH_z(t)`$ has a non-zero matrix element between $`|0>`$ and $`|1>`$. The effort should be made, therefore, to shield the magnetic qubit from the magnetic fields during the process of quantum computation. This can be done by placing the magnetic cluster inside a nanoscopic superconducting ring. Such a ring may be used to control and measure the states of the qubit. Connecting such rings by superconducting lines may be the way to make a miltiqubit system.
## Acknowledgements
I am grateful to Joe Birman, Lev Bulaevski, and Dima Garanin for discussions. This work has been supported by the U.S. National Science Foundation Grant No. DMR-9978882.
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# Independent freezing of charge and spin dynamics in La1.5Sr0.5CoO4.
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## Abstract
We present elastic and quasielastic neutron scattering measurements characterizing peculiar short-range charge-orbital and spin order in the layered perovskite material $`\mathrm{La}_{1.5}\mathrm{Sr}_{0.5}\mathrm{CoO}_4`$. We find that below $`T_c750`$ K holes introduced by Sr doping lose mobility and enter a statically ordered charge glass phase with loosely correlated checkerboard arrangement of empty and occupied $`d_{3z^2r^2}`$ orbitals (Co<sup>3+</sup> and Co<sup>2+</sup>). The dynamics of the resultant mixed spin system is governed by the anisotropic nature of the crystal-field Hamiltonian and the peculiar exchange pattern produced by the orbital order. It undergoes a spin freezing transition at much a lower temperature, $`T_s30`$ K.
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The instability of the doped transition-metal oxides towards formation of cooperative charge and spin ordered phases is one of the central phenomena in the physics of colossal magnetoresistance (CMR) materials and high temperature superconductors. Charge segregation into lines which separate stripes of antiphase antiferromagnetic domains in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4+y</sub> (LSCO) cuprates at small $`x`$, is, according to some theories, key to their superconductivity . In perovskite manganates, built of essentially isostructural Mn-O layers , doping first destroys the cooperative Jahn-Teller (JT) distortion of the orbitally-ordered insulating state , and induces a transition to a ferromagnetic metal (FM) . The extremely strong response of transport properties in the FM phase to an applied magnetic field, termed CMR, results from the strong Hund’s coupling of charge carriers to the Mn<sup>4+</sup> core spins . Then, at doping $`x0.5`$ an instability towards another kind of charge-orbital order, with a checkerboard arrangement of Mn<sup>3+</sup>/Mn<sup>4+</sup> ions, results in an antiferromagnetic (AFM) insulating ground state. This doping level is argued to be of greatest technological relevance , since at half-doping transport properties usually show the strongest response to a magnetic field. To understand the physics of charge/spin ordered (CO/SO) phases, and, in particular, to answer the question of whether CO/SO are independent instabilities or closely coupled, the relative importance of the various interactions (superexchage, double exchange, Coulomb repulsion, JT distortion) needs to be clarified. Experimentally this can be done by studying different compounds in which the strengths of the above interactions vary.
The first complete characterization of the spin and charge order in the half-doped regime by means of neutron diffraction was reported for La<sub>1.5</sub>Sr<sub>0.5</sub>MnO<sub>4</sub> . Nuclear scattering accompanying the CO was found below T$`{}_{co}{}^{}217`$ K, where the steep rise in resistivity manifests the transition to an insulating state. A magnetic signal consistent with simple collinear antiferromagnetic (AFM) order appeared below T$`{}_{so}{}^{}110`$ K. CO was confirmed in a recent X-ray resonant scattering experiment , which has also provided the first direct evidence of the concomitant orbital order (OO). Similar charge, orbital and magnetic order was also observed in the ground state of half-doped pseudocubic manganates . In most cases, the charge orders at a somewhat higher temperature than the spins do, but whether the CO instability occurs independently or is driven by the magnetic/orbital fluctuations is a topic of continuing debate .
Here we report a detailed neutron scattering study of the closely related layered cobalt oxide La<sub>1.5</sub>Sr<sub>0.5</sub>CoO<sub>4</sub>. We argue that strong single-ion anisotropy, mediated by the relativistic spin-orbit coupling usually neglected in cuprates and manganates, effectively decouples charge ordering from low-energy spin fluctuations. This is reflected in the spectacular difference in CO and SO transition temperatures, $`T_{co}/T_{so}20`$ in this material. In fact, strong planar anisotropy leads to a quenching of the spin angular momentum on the Co<sup>3+</sup> sites, in full analogy with the well-known quenching of the orbital momentum by the crystal field. Combined with the checkerboard arrangement of the doped holes in the CoO<sub>2</sub>, planes this makes the spin system of La<sub>1.5</sub>Sr<sub>0.5</sub>CoO<sub>4</sub> a strongly frustrated square lattice antiferromagnet.
Electronic configurations of Co<sup>2+</sup> ($`3d^7`$) and Co<sup>3+</sup> ($`3d^6`$) ions are related to those of Mn<sup>4+</sup> and Mn<sup>3+</sup>, respectively, by virtue of electron-hole symmetry. However, in the case of cobalt, Hund’s rule is in close competition with the cubic crystal field, the latter splitting the $`3d`$ level into a lower-lying $`t_{2g}`$ triplet and an $`e_g`$ doublet of $`d_{x^2y^2}`$ and $`d_{3z^2r^2}`$ orbitals. As a result, even though Co<sup>2+</sup> is in the high-spin state ($`t_{2g}^5e_g^2`$, S=3/2), Co<sup>3+</sup> obtained by adding the fourth hole may be either in the high $`t_{2g}^4e_g^2`$ (HS, S=2), intermediate $`t_{2g}^5e_g^1`$ (IS, S=1), or low spin state $`t_{2g}^6`$ (LS, S=0) . Transitions between different spin states are not uncommon in Co compounds. If LS is the ground state, a decrease in free energy due to the higher magnetic entropy may drive transitions to IS and HS states with increasing temperature, as observed in LaCoO<sub>3</sub> . For the purpose of the present study, however, the spin state of the Co<sup>3+</sup> ions is not important at sufficiently low temperatures where spin order occurs. Indeed, any integer spin will be frozen in a singlet state by strong planar anisotropy, rendering ions effectively non-magnetic for low-energy fluctuations .
We studied a piece (0.48 g) of high quality single crystal of $`\mathrm{La}_{1.5}\mathrm{Sr}_{0.5}\mathrm{CoO}_4`$ grown by the floating-zone method . In the temperature range $`6`$ K$`T600`$ K covered in our experiment, the crystal remained in the tetragonal “HTT” phase (space group $`I4/mmm`$), with low $`T`$ lattice parameters $`a=3.83`$ Å and $`c=12.5`$ Å. However, to index the superlattice peaks it is convenient to choose a unit cell that is twice as large ($`\sqrt{2}a\times \sqrt{2}a\times c`$) corresponding to space group $`F4/mmm`$.
Most of the experiments were performed on the BT2 and BT4 3-axis thermal neutron spectrometers at the NIST Center for Neutron Research. Some preliminary measurements were done on the H8 spectrometer at the High Flux Beam Reactor at Brookhaven. In all cases, PG(002) reflections were used at the monochromator and analyser, supplemented by PG filters to suppress the higher-order contamination. The energy of the scattered neutrons was fixed at $`E_f=14.7`$ meV. Beam collimations before the monochromator and after the analyser were kept at $`60^{}`$ and $`100^{}`$, respectively, while around the sample they were set either to $`20^{}20^{}`$ or to $`42^{}62^{}`$ depending on the resolution required. Our sample is a cylinder with axis parallel to direction, $`D4`$ mm, $`L3`$ mm. It was mounted in the displex refrigerator with axis vertical, allowing wavevector transfers in the $`(h0l)`$ reciprocal lattice plane. Rocking curves about the direction reveal a mosaic broadening of less than $`0.25^{}`$. Normalization of the scattered intensity was done using the incoherent scattering from a vanadium sample.
Figure 1 shows a survey of the low-temperature elastic scattering in $`\mathrm{La}_{1.5}\mathrm{Sr}_{0.5}\mathrm{CoO}_4`$. Peaks of magnetic and structural origin are easily distinguished by their characteristic wavevector dependences. The magnetic intensity is proportional to the Fourier-transform of the density of unpaired electrons (the magnetic form factor) squared, $`|f(Q)|^2`$, which becomes weaker at large $`Q`$, while the intensity from small nuclear displacements increases as $`Q^2`$. This leads us to conclude that the peaks in Fig. 1(b), centered at slightly incommensurate positions $`𝝉\pm 𝐐_m`$, $`𝐐_m=(0.5+ϵ,0,1)`$ with $`ϵ0.017`$, are magnetic, while stronger feature in Fig. 1(c) and features in Fig. 1(d) result from the modulations of the atomic positions accompanying the charge order of Co<sup>2+</sup>/Co<sup>3+</sup> ions. None of the superstructure peaks are resolution-limited in $`Q`$, indicating finite-range correlated glass type modulations. They are static on the time scale $`\delta t1`$ ps as determined by the energy resolution of our experiment. The diffuse nature of the scattering is most apparent in scans along the $`𝐜^{}`$ () direction, which reflects the anisotropy of the correlation lengths: in-plane correlations are much better developed than inter-plane (see Table 1).
The incommensurability of the charge-order peaks, if any, is too small to be directly resolved, because of the short correlation length and the overlap of peaks at $`𝝉\pm 𝐐_c`$. However, the scattering is consistent with the modulation wave vectors $`𝐐_c=(1+2ϵ,0,l)`$ and $`(0,1+2ϵ,l)`$, with $`l=0`$ or 1. If $`ϵ`$ is zero, the charge order corresponds to a checkerboard superstructure within each CoO<sub>2</sub> plane, as indicated in Fig. 2(c). The magnetic modulation corresponds to almost antiferromagnetic order on either the Co<sup>2+</sup> or Co<sup>3+</sup> sublattice, with a superposed long-wavelength modulation along a (or b).
For the quantitative analysis of nuclear scattering accompanying the charge order, we use the following correlation function between the displacement $`𝜺_{\mu 𝐫}`$ of atom $`\mu `$ at the lattice point $`𝐫=r_a\widehat{𝐚}+r_b\widehat{𝐛}+r_c\widehat{𝐜}`$ and $`𝜺_{\mu ^{}\mathrm{𝟎}}`$ of atom $`\mu ^{}`$ at $`\mathrm{𝟎}`$:
$$\epsilon _{\mu 𝐫}^\alpha \epsilon _{\mu ^{}\mathrm{𝟎}}^\beta =\epsilon _\mu ^\alpha \epsilon _\mu ^{}^\beta \mathrm{cos}(𝐐_c𝐫)e^{(|r_a|+|r_b|)/\xi _{}^{co}|r_c|/\xi _{}^{co}},$$
(1)
describing a short-range harmonic modulation. Brackets denote an average over the crystal volume; $`\alpha ,\beta =x,y,z`$; and $`\xi _{}^{co}`$ and $`\xi _{}^{co}`$ are the charge-order correlation lengths parallel and perpendicular to the planes, respectively. The scattering cross-section obtained from Eq. (1) upon appropriate Fourier transformation has the form of factorized “lattice Lorentzians” in three directions, weighted with $`\left|_\mu (𝐐𝜺_\mu )b_\mu e^{i𝐐r_\mu }\right|^2`$, similar to that used to describe finite-range stripe correlations in LSCO . Here $`b_\mu `$ is the scattering length (including the Debye-Waller factor) of the nucleus at position $`𝐫_\mu `$ in the unit cell. The $`Q`$-depedence of the structural diffuse intensity evident in scans along $`𝐜^{}`$, Fig 1(d), suggests that it arises from breathing-type distortions of the oxygen octahedra surrounding the Co ions. Indeed, it appears sufficient to consider displacements of the in-plane oxygens O(1), $`\epsilon _{O(1)}^{x,y}`$, and apical oxygens O(2), $`\epsilon _{O(2)}^z`$, along the corresponding Co-O bonds, to obtain a good description of the observed “charge” scattering. Solid curves in Fig. 1(c),(d) present the neutron intensity calculated from Eq. (1) after appropriate normalization and correction for the spectrometer resolution. The parameters $`\epsilon _{O(1)}^{x,y}=0.042(4)`$ Å, $`\epsilon _{O(2)}^z=0.066(2)`$ Å, and the correlation lengths shown in Table 1 were obtained in a global fit of all scans around $`𝝉\pm 𝐐_c`$. The error bars shown do not include any possible systematic error from the vanadium normalization, which relied upon a knowledge of the vertical beam divergences (note the absence of arbitrary scaling factors between calculated and measured intensities in all panels of Fig. 1).
Figure 2 shows the temperature dependence of the peak intensity of the diffuse “charge” scattering and lattice thermal expansion in the a and c directions. Although we could not reach the CO melting temperature, an estimate of $`T_{co}750`$ K is obtained by extrapolation shown in Fig. 2(b). Strong nonlinear decrease of the $`c`$ lattice spacing which accompanies the charge order is indirect evidence of the concomitant orbital order. This is because the checkerboard ordering of empty and occupied out-of-plane $`d_{3z^2r^2}`$ orbitals, combined with the body-centered stacking of Co-O planes, $`𝜹`$$`=[\frac{1}{2}0\frac{1}{2}]`$, allows a reduction of the inter-plane spacing.
We analyze the magnetic scattering cross section starting from a spin-spin correlation function similar to Eq. (1). By simultaneously fitting all of the magnetic peaks measured at $`T=10(3)`$ K, we obtain a SO incommensurability $`ϵ=0.017(1)`$ and the correlation lengths shown in Table 1. We also find that the spin structure is planar, and isotropic within the $`ab`$ plane. This agrees with the static susceptibility data , which yield an estimated $`D400`$ K for the XY-type anisotropy energy. For $`TD`$ this restricts any integer spin to a singlet state . Further support for the conclusion that Co<sup>3+</sup> ions are effectively non-magnetic is provided by small value of the frozen magnetic moment, $`\mu =1.4(1)\mu _B`$ per Co site refined for a single domain model. Assuming two equivalent $`𝐐_m`$ domains with half of the sites occupied by magnetic Co<sup>2+</sup> ions, we obtain $`\mu _{\mathrm{Co}^{2+}}2.9\mu _B`$, within 20% of what is expected for S=3/2.
The main role of the ”non-magnetic” Co<sup>3+</sup> ions is to bridge the Co<sup>2+</sup> ions, providing effective antiferromagnetic coupling. As shown in Fig. 2(c), there are two Co<sup>2+</sup>–O–Co<sup>3+</sup>–O–Co<sup>2+</sup> exchange pathways between nearest neighbor (nn) Co<sup>2+</sup> ions, and one between next-nearest neighbors (nnn). As a result, the spin system appears to be a quasi-two-dimensional square-lattice antiferromagnet with nnn to nn exchange ratio $`J_2/J_10.5`$; i.e., it is in the critical region where frustration destroys
Néel order . This provides a natural explanation for the short spin-spin correlation length in the $`ab`$ plane and the peculiar spin freezing transition revealed by the temperature dependences shown in Fig. 3. Despite a somewhat “order parameter”-like dependence of the (quasi)elastic magnetic peak intensity, the correlation lengths do not diverge. In fact, the defining feature of the transition is the disappearance of the energy width $`\mathrm{\Gamma }_E`$ of the magnetic scattering, i.e. a divergence of the relaxation time for the spin fluctuations \[Fig. 3(b)\]. Upon appropriate correction for the instrumental resolution, we are able to refine the $`\mathrm{\Gamma }_E(T)`$ dependence down to $`0.05`$ meV. From the power-law fit $`\mathrm{\Gamma }_E(T)(TT_{so})^\zeta `$ shown in Fig. 3(b) we estimate $`\zeta =2.7(3)`$ and $`T_{so}30`$ K.
Finally, to underscore the significance of these results we compare the charge and spin order of $`\mathrm{La}_{1.5}\mathrm{Sr}_{0.5}\mathrm{CoO}_4`$ with that of La<sub>1.5</sub>Sr<sub>0.5</sub>MnO<sub>4</sub> and other manganates. Firstly, the charge order in $`\mathrm{La}_{1.5}\mathrm{Sr}_{0.5}\mathrm{CoO}_4`$ occurs independently of the magnetic order, arising at temperatures $`20`$ times higher than the characteristic energy scale of the cooperative spin fluctuations. Secondly, compared to the CO in the manganate, it has a much shorter correlation range, $`\xi _{co}(\mathrm{Mn})/\xi _{co}(\mathrm{Co})10`$, but, surprisingly, results in an even stronger charge localization with an activation behavior of electrical conductivity with $`E_a6000`$ K . Thirdly, as a result of the checkerboard order of Co<sup>2+</sup>/Co<sup>3+</sup> ions and the strong $`XY`$ anisotropy, the spin system is a stack of weakly interacting square antiferromagnets with nearly critical frustration, $`J_2/J_10.5`$. In contrast to La<sub>1.5</sub>Sr<sub>0.5</sub>MnO<sub>4</sub>, where simple collinear AFM order occurs, the cobaltate undergoes a freezing transition into a short-range correlated, incommensurate spin-glass state. In fact, this is a rare example where the evolution of the critical spin fluctuations in the course of glassification can be studied by neutron scattering over a dynamic range of about two decades. We intend to further investigate the dynamics of the spin freezing transition in the near future.
We thank NIST Center for neutron Research for hospitality during the experiments. This work was carried out under Contract NO. DE-AC02-98CH10886, Division of materials Sciences, US Department of Energy.
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# Energetic particle acceleration in a 3D magnetic field reconnection model: a role of MHD turbulence
## 1 Introduction
Particle acceleration processes in regions of magnetic field reconnection can play an important role in various astrophysical sites by providing high energy particles, heating thermal plasma and by changing the magnetic field configuration (e.g. Somov 1994). The first of these processes, accelerating particles to cosmic ray energies, acts due to electric fields occurring in the reconnection region. Two approaches were applied to describe the acceleration process. In the more basic one, analytic or semi-analytic models for the reconnection region provided background for the derivation of particle trajectories and discussion of the energy distribution of particles escaping from the reconnection volume. A simple discussion of particle trajectories presented by Speiser (1965; see also Sonnerup 1971) shows that even a tiny magnetic field component near the neutral current sheet efficiently removes ergetic particles and decreases its final energy. Further studies by Stern (1979) and Wagner et al. (1981) revealed the importance of the O-type neutral line regions for the acceleration, due to partial particle trapping in the accelerating volume. Deeg et al. (1991) and Somov & Kosugi (1997) discussed several features of the acceleration process in the X-type stationary reconnection. The time dependent effects of the acceleration process were introduced by considering the background magnetic and electric fields as derived from simple MHD simulations (cf. Sato et al. 1982, Scholer & Jamitzky 1987, Atkinson et al. 1989, Zelenyi et al. 1990).
Substantial progress in considering realistic reconnection processes came from the two dimensional ($``$ 2D; Matthaeus et al. 1984, Ambrosiano et al. 1988, Scholer & Jamitzky 1989, Veltri at al. 1998) and 3D (Birn & Hesse 1994, Schopper et al. 1999, Kliem et al. 1998) MHD modelling involving a perturbed magnetic field as the initial condition. This approach yields a complex structure with X-type and O-type null points and magnetic field perturbations of scales comparable to the macroscopic structure dimensions. Ambrosiano et al. (1988) show that such a turbulent neutral point mechanism influences the acceleration process in two ways. It enhances the reconnection magnetic field while producing a stochastic electric field that gives rise to momentum diffusion, and it also produces magnetic bubbles and other irregularities that can temporarily trap test particles in the strong reconnection field for times comparable to the magnetofluid characteristic time. As a result, the very flat particle spectra formed can extend to higher energies in comparison to the unperturbed conditions.
One should also note that a realistic description of the acceleration process requires considering the full 3D configuration of the reconnection region (e.g. Birn & Hesse 1994, Schopper et al. 1999, Veltri et al. 1998). Besides the role which the mean field structure plays, particle transport depends in a qualitative way on the turbulence dimensionality (cf. Giacalone & Jokipii 1994, Michałek & Ostrowski 1996) and the involved wave modes (e.g. Schlickeiser & Miller 1997, Michałek & Ostrowski 1998).
Besides providing a substantial improvement in understanding of the process, a general deficiency of the above mentioned papers discussing particle acceleration in the turbulent reconnection regions are the complicated interrelations between such factors as the assumed or derived mean magnetic field structure, its dimensionality, the assumed or derived form of turbulence, and the particle acceleration process. For example in the paper of Ambrosiano et al. (1988) it was difficult to separate the role of diffusive particle motions from trapping by the structures formed in/near the reconnection layer. In order to clarify the situation, with only one of these factors modifying the acceleration process in turbulent reconnection, we decided to use a simple model which allows us to consider the role of random particle trajectory perturbations separately. The process acting in the vicinity of the X-type reconnection does not provide a means for particle trapping in the mean (unperturbed) structure of the reconnecting field. As described in the next section, we use the 3D analytic model of stationary reconnection of Craig et al. (1995) to define the background model, being a reference for the model involving particle pitch-angle momentum scattering due to MHD turbulence. Perturbations of particle trajectories due to a turbulent magnetic field component were simulated using small-amplitude pitch-angle momentum scattering, enabling modelling of both small and large amplitude turbulence in a wide wave vector range. No second-order Fermi acceleration processes are allowed within this approach. In section 3 we present results of particle energy spectrum modelling for different scattering amplitudes. Comparison of the acceleration process to the unperturbed one confirms that the turbulence can substantially increase the acceleration efficiency, enabling particles to form flat high-energy spectra with much increased final energies. Then, in section 4, we briefly summarize these results.
In the present simulations we scale the parameters of the reconnection region to those characteristic for the solar flare (e.g. Miller et al. 1997). However, we do not aspire to present a flare acceleration model. We provide all numerical values in SI units.
## 2 Simulations of the acceleration process
### 2.1 The 3D reconnecting magnetic field structure
The considered reconnection region (Fig. 1) involves in-flowing plasma from the ‘top’ and the ‘bottom’ toward the $`x=0`$ plane, and outflowing to the sides. All physical quantities, the flow velocity, the magnetic field and the electric field are symmetric with respect to the origin of the reference frame, with the reconnection “plane” (magnetic field is strictly 0 only in a single point in the center) being slightly inclined with respect to the $`x=0`$ plane. To describe this structure analytically we use a 3D model of the stationary reconnection by Craig et al. (1995), which provides the analytic structure for the magnetic field $`\stackrel{}{B}`$, the electric field $`\stackrel{}{E}`$ and the plasma velocity $`\stackrel{}{V}`$. The following parameters are introduced in this model: $`\alpha `$, $`\lambda `$, $`\eta `$ and $`k`$ define topologic structure of the magnetic field; $`\mu `$ define the reconnection area thickness; $`\rho `$ is the plasma density; $`\xi `$ denotes the external electric field component; $`\mu _0`$ is the magnetic permeability of the vacuum. With the use of these parameters and the notation involving a radius vector $`\stackrel{}{r}=\{x`$ , $`y`$ , $`z\}`$ and the unit vectors $`\widehat{x}`$, $`\widehat{y}`$ and $`\widehat{z}`$ along the respective axes,
$$\stackrel{}{B}=\lambda \stackrel{}{P}(\stackrel{}{r})+\stackrel{}{Q}(x),$$
$`(2.1)`$
$$\stackrel{}{V}=\left[\stackrel{}{P}(\stackrel{}{r})+\lambda \stackrel{}{Q}(x)\right]\frac{1}{\sqrt{\rho \mu _0}},$$
$`(2.2)`$
$$\stackrel{}{E}=\left[\frac{\sqrt{\pi }\eta }{2\mu _0}Z^{}(0)e^{(\mu x)^2}\right]\widehat{y}+\xi \left[12x\mu \mathrm{daw}(\mu x)\right]\widehat{z},$$
$`(2.3)`$
where
$$\stackrel{}{P}=\alpha \left[x\widehat{x}+ky\widehat{y}+(1k)z\widehat{z}\right],$$
$`(2.4)`$
$$\stackrel{}{Q}=\frac{\xi }{\eta \mu }\sqrt{\frac{\mu _0}{\rho }}\mathrm{daw}(\mu x)\widehat{y}+\left[\frac{\sqrt{\pi }}{2\mu }Z^{}(0)\mathrm{erf}(\mu x)+Z(0)\right]\widehat{z},$$
$`(2.5)`$
and
$$\mathrm{daw}(x)_0^x\mathrm{exp}(t^2x^2)dt,\mathrm{erf}(x)_0^x\mathrm{exp}(t^2)dt,$$
$$\mu ^2=\frac{\alpha }{2\eta }1\lambda .$$
The function $`Z(x)`$ of Craig et al. (1995) satisfies the differential equation
$$\alpha (1\lambda ^2)[(1k)Z+xZ^{}]+\eta Z^{\prime \prime }=\alpha \gamma _2\lambda (2k)x,$$
$`(2.6)`$
where one requires $`\gamma _2k=0`$ ($`\gamma _2`$ or $`k`$ must be zero). With these expressions one can construct the reconnection region in two ways. It is possible to choose the parameters $`\xi `$ , $`Z^{}(0)`$, $`Z(0)`$, $`\alpha `$, $`\beta `$, $`\eta `$, $`\mu `$ in such a way that either the magnetic field component $`B_{\widehat{z}}`$ or $`B_{\widehat{y}}`$ vanishes in some area:
* i.) $`B_y`$ vanishes where $`\lambda \alpha ky=\frac{\xi \mu _0}{\eta \mu }\mathrm{daw}(\mu x)`$
* ii.) $`B_z`$ vanishes where $`\beta (1k)z=\frac{\sqrt{\pi }}{2\mu }Z^{}(0)\mathrm{erf}(\mu x)`$.
The maximum value of $`r`$ allowing for one of the above conditions to be satisfied gives the half length of the reconnection region $`L`$ (a scale for the layer with negligible magnetic field near the null point). For example, for the case (ii) one has
$$L=\frac{\pi Z^{}(0)}{4\beta (1k)\mu }.$$
$`(2.7)`$
In the case (i), where the function $`daw`$ determines the behaviour of the magnetic field, the gradient of the magnetic field is greater than in the case (ii). As a result the simulation time for the required accuracy is longer. Moreover, within the turbulent volume, at a large distance from the centre the function $`daw`$ grows slower resulting in less efficient acceleration. For these reasons in the present simulations we consider solutions determined with the condition (ii) satisfied. We take the following values of the model parameters: $`\alpha =1.010^9`$, $`\lambda =\frac{\beta }{\alpha }`$, where $`\beta =1.010^{10}`$, $`\eta =1.010^6`$, $`k=0.15`$ , $`\rho =M_p10^{16}`$ ($`M_p`$ proton mass), $`\xi =0`$, $`Z(0)=0`$, $`Z^{}(0)=710^4`$. Values of these parameters were derived from fitting the model magnetic, electric and velocity field to the mentioned solar flare conditions, with the plasma density equal to $`10^{16}`$ m<sup>-3</sup>.
Let us note that the simple analytic form of the considered solution limits the range of physically acceptable parameters to provide the possible boundary values of $`\stackrel{}{V}`$, $`\stackrel{}{B}`$, $`\stackrel{}{E}`$ . In the present simulations we choose such boundary values for the solution at the edge of our considered volume to provide physical parameters near the reconnection layer that are close to the ones estimated for the solar flares. In the vicinity of the reconnection area, just outside of the reconnecting current layer, the magnetic field induction is assumed to be $`B=1.6810^3`$ T and the respective gyroradius (gyroperiod) for a $`1`$ MeV proton is $`r_g=87`$ m ($`T_g=410^5`$ s), and $`r_g=3.4`$ km ($`T_g=8.110^5`$ s) for a $`1`$ GeV proton. The assumed linear dimensions of the full considered reconnection region are $`\mathrm{\Delta }x=210^4`$ km, $`\mathrm{\Delta }y=210^5`$ km and $`\mathrm{\Delta }z=1.510^5`$ km. The size of the volume of the perturbed magnetic field is in the centre $`\mathrm{\Delta }x=210^3`$ km, $`\mathrm{\Delta }y=610^3`$ km and $`\mathrm{\Delta }z=210^3`$ km. One should note that the selected particular conditions are not essential for the presented considerations.
### 2.2 Monte Carlo modelling procedure
In the simulations we used a Monte Carlo approach including a trajectory splitting technique to improve statistics at larger energies (cf. Ostrowski 1991). To derive particle spectra we recorded ‘weights’ of particles escaping from the simulation volume in a given energy range. For every escaping particle a randomly<sup>1</sup><sup>1</sup>1With probability of selecting a given particle being proportional to its weight. chosen particle still active in the simulations was split into two identical particles, each with half the weight of the original particle. Then the trajectories of both particles were followed, but due to the applied random momentum scattering they evolve in different ways. The maximum simulation time $`t_{max}`$ ($`t_{max}`$ is the upper limit in Fig. 4) was chosen to be sufficiently large to be unreachable by high weight particles. Thus any further increasing of $`t_{max}`$ does not influence the resulting spectra in a visible way.
In the simulations, we inject test particles with initial energy $`E_0`$ = $`2`$ MeV in the vicinity of the reconnection null point and we follow their trajectories by numerical integration of particle equations of motion. The integration is completed when a particle either crosses the boundary of the considered reconnection region (‘escape boundary’), or the time limit $`t_{max}=100`$ s is reached. Particles were scattered in a small region around the reconnection null point only. This region is shown in the two upper panels of Fig. 1 as black rectangles, and in expanded form in the lower panel. The trajectory computation times are much longer for low energy particles and this is the main reason why we start with initial proton energies substantially larger than the thermal energy; for the discussion of the injection problem of energetic solar flare particles one can consult the discussion in Miller et al. (1997). As the long integrations performed here require high accuracy we use a variable step fourth-order Runge-Kutta integration with the parameters chosen in such a way that any further accuracy increase does not affect the simulated trajectories. In the simulation we use only 100 particles because the required high accuracy of integration leads to extensive integration times. Let us remember that with the applied trajectory splitting technique the number of particles forming the spectrum at different energies is the same.
We performed a number of numerical tests of the code applied in the simulations. We checked by hand the derivation of the resulting particle phase space co-ordinates in a few randomly chosen individual integration steps of the code algorithm. Then we derived particle trajectories within a few simple uniform magnetic fields oriented randomly with respect to the chosen reference frame. The results coincided within the numerical accuracy with the derived analytic trajectories. Finally, for the actual considered magnetic field structure in the reconnection volume we positively checked conservation of particle energy with the plasma velocity set equal to zero.
In the simulations particles gain energy mostly in the vicinity of the reconnection layer when drifting in the $`\stackrel{}{V}\stackrel{}{B}`$ electric fields. Away from this region, while moving across the magnetic field gradient, particles can gain and lose energy, but the mean energy change is small. Somov & Kosugi (1997) estimated the accelerated particle energy as a product of the magnetic field, the plasma velocity and the reconnection area length. The energy gains derived in our unperturbed reconnection model are in agreement with this estimate if we take the distance traversed by a particle within the reconnection area as the required length. The effect of particle escape from the reconnection volume has been discussed by Speiser (1965) for a highly simplified reconnection model. We confirm his results showing that a small vertical magnetic field component within the reconnection layer increases the particle escape substantially.
In the simulations we consider the particle scattering process only in a small volume containing the reconnection null point, where the magnetic energy is dissipated (see Fig. 1). The real reconnection regions are expected to show analogous structures, with the turbulence amplitude growing toward the centre, including null points of the magnetic field. This choice also resulted from the fact that the Craig et al. (1995) three dimensional model of the stationary reconnection is not too realistic at large distances from the null point, as, for example, it involves the inflow velocity growing without a limit when increasing distance from the centre. For the perturbed trajectories model we use a simple pitch-angle scattering approach (e.g. Ostrowski 1991) intended to model the particle scattering at MHD waves. In this case integration of the particle equations of motion is performed in the electromagnetic field defined by the unperturbed analytic model, but trajectory perturbations are introduced every constant time interval $`\mathrm{\Delta }t`$, when the particle momentum vector is randomly scattered within a narrow cone along its original direction. The scattering is performed in the local plasma rest frame and it conserves particle energy in this frame. In the present simulations we consider the uniform momentum scattering within a cone with the half opening angle equal to $`11^{}`$. The perturbation intensity is controlled by changing $`\mathrm{\Delta }t`$ and it is characterized with $`\mathrm{}\kappa _{}/\kappa _{}`$, the ratio of the cross-field diffusion coefficient to the diffusion coefficient along the magnetic field. The value of $`\mathrm{}`$ was determined in auxiliary simulations performed in the uniform magnetic field with the value characteristic for a region close outside the reconnection current sheet (see section 2.3). One should also note that decrease of the magnetic field toward the reconnection site leads to increasing the effective turbulence amplitude (i.e. the value of $`\mathrm{}`$; in the limit of $`B=0`$ we have $`\mathrm{}=1.0`$), but in the present simulations the particle gyroradius is always larger than the reconnection layer thickness near the considered X-type null point.
As the mean scattering time $`\mathrm{\Delta }t`$ is assumed to be constant within a given simulation run, particle trajectories are affected by perturbations with intensity depending on particle energy and the background magnetic field. For non-relativistic particles with constant angular velocities of their gyration movements, assuming constant $`\mathrm{\Delta }t`$ is equivalent to introducing scattering acts at constant gyrophase steps. Thus the resulting value of $`\mathrm{}`$ does not depend on energy, as expected for the flat wave power spectrum $`F(k)k^1`$. This slightly unrealistic wave spectrum allows, on the other hand, evaluation of the role of diffusive effects for the same scattering amplitude at all considered particle energies. A discussion of a more realistic Kolmogorov wave spectrum within the finite wave vector range will be presented in the next paper (in preparation). However, as such a wave spectrum carries more energy in long waves, the results are expected to show a transition from our low $`\mathrm{}`$ results to the large $`\mathrm{}`$ ones.
In attempting to compare our simplified scattering model with the real turbulence with an amplitude of waves resonant for a particle of a given energy, $`\delta B_r`$ ($`\delta B_r^2/8\pi F(k_r)k_r`$, where $`k_r=2\pi /r_g`$), one can refer to a qualitative discussion of energetic particle diffusion presented by Drury (1983). With his scaling $`\kappa _{}(\delta B_r/B)^2`$ and $`\kappa _{}(\delta B_r/B)^2`$, the amplitude for resonance waves can be evaluated as $`\delta B_r/B\mathrm{}^{1/4}`$.
### 2.3 Derivation of the $`\mathrm{}\kappa _{}/\kappa _{}`$ parameter
The respective values of $`\mathrm{}\kappa _{}/\kappa _{}`$ were derived in auxiliary simulations involving the spatially uniform background magnetic field with the induction of $`1.510^3`$ T and particles with energies equal to the initial energy $`E_0=2`$ MeV. Trajectories of a large number of particles were followed with the imposed scattering process involving the momentum angular scattering, uniform within a cone of half opening angle equal to $`11^{}`$ and with the cone axis directed along the original momentum vector. The only parameter varying between the simulations was the time interval between successive scattering events, $`\mathrm{\Delta }t`$. The resulting diffusion coefficients were derived from growing particle dispersions along the background field ($`\kappa _{}`$) and along two orthogonal axes perpendicular to the background field ($``$ along the $`1`$\- or $`2`$-axis). The results of such computations are presented in Fig. 2. Presentation of two derived values of $`\kappa _1/\kappa _{}`$ and $`\kappa _2/\kappa _{}`$ allows one to evaluate the accuracy of these computations. The values used in the paper are fits to the asymptotic $`(T\mathrm{})`$ value of $`\mathrm{}=0.5(\kappa _1+\kappa _2)/\kappa _{}`$. For a sequence of scattering times $`\mathrm{\Delta }t`$ = $`10^6`$, $`10^5`$, $`10^4`$ and $`10^3`$ we derived the respective values of $`\mathrm{}`$ = $`1.210^2`$, $`610^3`$, $`610^5`$, and $`610^7`$. The results derived without applying any scattering are indicated by $`\mathrm{}=0`$.
## 3 Spectra of accelerated particles
In order to provide qualitative evaluations of turbulence effects in the volume of reconnecting magnetic field, but considering it only as a factor introducing random motion component to particle trajectories, we performed simulations of energetic proton spectra with a varying amount of turbulence ($``$ scattering). As explained above, this approach assumes the existence of short wave magnetic field perturbations to be present in the limited volume – the black rectangle in Fig. 1 – near the central neutral point, but we do not consider the influence of the turbulence on the reconnection process. Thus it is a complementary approach to that using MHD modelling of the turbulent reconnection including wave perturbations from a narrow wave vector range, as discussed in section 1.
We performed simulations of particle evolution starting at the same ‘injection’ energy $`E_0=2`$ MeV, avoiding consideration of the real injection process at much lower energies (cf. Miller et al. 1997). For each set of particles we derived the spectrum of particles escaping from the reconnection volume, as illustrated in Fig. 3. In the non-perturbed model ($`\mathrm{}=0`$, curve A) protons can increase their initial energy by approximately $`70`$% . One should note that the injected energetic particles can gain as well as loose energy. Introducing trajectory perturbations results in substantial modification of the acceleration process (curves B, C, D, E in Fig. 3). The spectrum energy cut-off shifts to higher values and the spectrum becomes harder when the amount of scattering (‘turbulence amplitude’) is increased. In our simulations the resulting flat spectra extend up to $`50`$ MeV for models with strong turbulence (Model D and E), and a steeper part of the spectrum is recorded at energies above $`100`$ MeV. This behaviour results from the fact that the diffusive component introduced in to particle trajectories by the scattering enables some particles to stay in the reconnection region much longer and diffuse back close to the null point from outside. As illustrated, this can have a pronounced influence on the acceleration process by substantially increasing the particle mean energy gain and providing much larger energies of individual particles. There is a general trend for the acceleration efficiency to increase with the perturbation amplitude in the considered range of values for the $`\mathrm{}`$ parameter.
Acceleration by the uniform electric field is characterized with the energy-independent rate of particle energy increase, leading to a linear relation between the particle acceleration time and its final energy. In contrast, a plot of correlation between the individual particle acceleration time ($``$ the physical simulation time) and its final energy presented in Fig. 4 shows that the acceleration rate depends on energy, being preferentially defined by particle drifts in the ‘$`V`$ x $`B`$’ electric fields. A large energy dispersion in a given time reflects variety of involved particle diffusive trajectories.
## 4 Final remarks
Consideration of energetic particle acceleration accompanying the magnetic field reconnection process requires knowledge of the electromagnetic field structure in the region of interest. Until now the available numerical models of such fields were oversimplified and fail to consistently include together the short and long MHD wave modes. Therefore, in the present simulations we apply an analytic model for a reconnecting field as the background, and the particle scattering process is imposed as the small amplitude uncorrelated angular momentum perturbations. A comparison of the acceleration process in such a model to the unperturbed model reveals the important result that inclusion of the turbulent field component into the reconnection volume can change the acceleration process in a qualitative way, enabling particles to reach much higher final energies and significantly increase their mean energy gain.
A serious limitation on the validity of our simplified modelling arises due to a non-self-consistent introduction of the MHD turbulent field. The trajectory perturbations in the turbulent medium are considered without taking into account the influence – in both directions – of these turbulent motions on the reconnection process itself. As mentioned in section 1, the presence of magnetic field turbulent structures increases mixing in the medium and is the source for anomalous resistivity leading to more effective reconnection as discussed recently by Lazarian & Vishniac 2000. Subsequently, it leads to more efficient particle acceleration (Matthaeus et al. 1984, Ambrosiano et al. 1988, Scholer & Jamitzky 1989, Veltri at al. 1998, Birn & Hesse 1994, Kliem et al. 1998, Schopper et al. 1999). As discussed in the present paper from a slightly different perspective it is expected to increase the cosmic ray acceleration efficiency, and influence the involved time scales and details of the resulting energy spectrum. However, for a given scattering conditions within the reconnection region the spectrum upper energy cut-off is limited by the ‘global’ perturbed structure of the reconnecting volume and not by the local conditions within the thin reconnection current sheet.
The present work was supported by the Komitet Badań Naukowych through the grants PB 179/P03/96/11 and PB 258/P03/99/17.
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# I Introduction
## I Introduction
The cosmic censorship hypothesis (CCH) is one of the most important open problems in general relativity since it plays an important role in theories of black hole physics . The CCH roughly states that singularities forming in gravitational collapse must be hidden behind event horizons and hence invisible to outside observers. Many types of gravitational collapse have been studied so far in the context of the CCH. Some of them produce globally naked singularities, but they cannot immediately be counterexamples to the CCH because the CCH requires suitable matter and appropriate initial conditions.
The well-known example of spacetime which describes the dynamical formation of naked singularity is the Lemaître-Tolman-Bondi (LTB) spacetime . The LTB spacetimes describe the gravitational collapse of a spherically symmetric dust cloud and have an explicit form of the entire spacetime metric. It has been proved that the LTB spacetimes have ingoing null naked singularity from the generic initial data . However this model is rather simplified because pressure is not taken into account.
Some aspects of the effect of pressure on the naked singularity formation have been investigated. Several important results have been obtained. Ori and Piran investigated the spherically symmetric collapse of a perfect fluid numerically under the assumption of self-similarity , and analytic discussions based on self-similarity followed it . Harada also solved numerically the non-self-similar spherically symmetric collapse of a perfect fluid . Their results are summarized as that naked singularity can occur in the spherically symmetric collapse of a perfect fluid for a sufficiently soft equation of state. In contrast to isotropic pressure, to include only tangential pressure is turned out to be more tractable . Among the solutions with vanishing radial pressure, there is a system of a spherical cloud of counterrotating particles, in which the physical origin of tangential pressure is clear . The each particle in the cluster has its angular momentum so that the average effect of all particles is a non-vanishing tangential pressure. This model is particularly interesting from a physical point of view because it gives insights into rotational effects on the gravitational collapse, without raising terrible difficulties. Harada, Iguchi and Nakao (HIN) have recently analyzed singularity occurrence in this model . Using the mass-area coordinates, which were first introduced by Ori and were applied in this system by Magli , HIN found a new exact solution that describes dynamical formation of massless naked singularity. The results show that the counter-rotation undresses the covered singularity, and in particular tangential pressure and rotation may induce the formation of naked singularity. Since the HIN solution is given in the mass-area coordinates, the motion of each shell and the global properties of the spacetime are not trivial. We study in this paper the HIN solution in detail to make them clear.
The plan of this paper is as follows. In the next section, we review the HIN solution. In Sec. III, we study null geodesics in the HIN spacetime which are necessary to determine the causal structure of spacetime. It is discussed in Sec. IV that the central singularity is timelike. In Sec. V, we will see that the collapse asymptotes to the singular static Einstein cluster. Section VI is devoted to conclusions. We use units with $`c=G=1`$ and follow the sign conventions of the textbook by Misner, Thorne and Wheeler about the metric, Riemann and Einstein tensors .
## II The HIN spacetime
### A The HIN solution
The HIN solution is the spherical cloud of counterrotating particles which is marginally bound. The specific angular momentum $`L(r)`$ of each particle at comoving radius $`r`$ equals to $`4F(r)`$ , where $`F(r)`$ is the conserved Misner-Sharp mass function . Using comoving coordinates, the line element for this system is reduced to
$$ds^2=e^{2\nu (t,r)}dt^2+R^{}{}_{}{}^{2}(t,r)(1+\frac{16F^2(r)}{R^2(t,r)})dr^2+R^2(t,r)(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(1)
where $`\nu (t,r)`$ and $`R(t,r)`$ satisfy the following set of coupled partial derivative equations:
$`\nu ^{}(t,r)`$ $`=`$ $`{\displaystyle \frac{16F^2}{R(R^2+16F^2)}}R^{}(t,r),`$ (2)
$`V`$ $``$ $`e^\nu \dot{R}(t,r)={\displaystyle \frac{|R4F|\sqrt{2F}}{\sqrt{R(R^2+16F^2)}}}.`$ (3)
The prime and overdot denote the partial derivatives with respect to $`r`$ and $`t`$, respectively. The energy density $`ϵ(t,r)T_t^t`$ and the tangential pressure $`\mathrm{\Pi }(t,r)T_\theta ^\theta =T_\varphi ^\varphi `$ are given by
$`ϵ(t,r)`$ $`=`$ $`{\displaystyle \frac{F^{}}{4\pi R^2R^{}}},`$ (4)
$`\mathrm{\Pi }(t,r)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{16F^2}{R^2+16F^2}}ϵ.`$ (5)
Using the coordinate transformation $`m=F(r)`$ and $`R=R(t,r)`$, the HIN solution is given in the mass-area coordinates $`(m,R)`$ by
$$ds^2=Adm^22BdRdmCdR^2+R^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(6)
where $`A`$, $`B`$ and $`C`$ are given by
$`A`$ $`=`$ $`H\left(1{\displaystyle \frac{2m}{R}}\right),`$ (7)
$`B`$ $`=`$ $`{\displaystyle \frac{R}{|R4m|}}\sqrt{{\displaystyle \frac{RH}{2m}}},`$ (8)
$`C`$ $`=`$ $`{\displaystyle \frac{1}{V^2}}={\displaystyle \frac{R(R^2+16m^2)}{2m(R4m)^2}},`$ (9)
$`\sqrt{H(m,R)}`$ $``$ $`\sqrt{B^2AC}={\displaystyle \frac{(F^1)_{,m}((F^1)^2+16m^2)}{|F^14m|\sqrt{2mF^1}}}`$ (12)
$`+\text{sign}(F^14m)[({\displaystyle \frac{R^216mR+144m^2}{3\sqrt{2}m(R4m)}}\sqrt{{\displaystyle \frac{R}{m}}}+4\sqrt{2}\mathrm{ln}{\displaystyle \frac{\sqrt{R}+2\sqrt{m}}{|\sqrt{R}2\sqrt{m}|}})`$
$`({\displaystyle \frac{(F^1)^216mF^1+144m^2}{3\sqrt{2}m(F^14m)}}\sqrt{{\displaystyle \frac{F^1}{m}}}+4\sqrt{2}\mathrm{ln}{\displaystyle \frac{\sqrt{F^1}+2\sqrt{m}}{|\sqrt{F^1}2\sqrt{m}|}})].`$
We denote the inverse function of $`F(r)`$ as $`F^1=F^1(m)`$ and have set $`R=F^1(m)`$ on the initial spacelike hypersurface, which corresponds to $`R(0,r)=r`$. The energy density in the mass-area coordinates is given by
$$ϵ(R,m)=\frac{\sqrt{R^2+16m^2}}{4\pi R^3|V|\sqrt{H}}.$$
(13)
We can assume that the metric functions are $`C^{\mathrm{}}`$ class with respect to the local Cartesian coordinates at least in the neighborhood of the center $`r=0`$ before encountering a central singularity. By this assumption, the metric variables in the comoving coordinates are expanded near the center as
$`\nu (t,r)`$ $`=`$ $`\nu _0(t)+\nu _2(t)r^2+\nu _4(t)r^4+O(r^6),`$ (14)
$`R(t,r)`$ $`=`$ $`R_1(t)r+R_3(t)r^3+R_5(t)r^5+O(r^7).`$ (15)
Then from Eq. (3), the arbitrary mass function $`F(r)`$ is also expanded as
$$F(r)=F_3r^3+F_5r^5+F_7r^7+O(r^9),$$
(16)
and the inverse function $`F^1(m)`$ is approximately given by
$$F^1(m)=r=\left(\frac{m}{F_3}\right)^{1/3}\frac{F_5}{3F_{3}^{}{}_{}{}^{2}}m+\frac{4F_{5}^{}{}_{}{}^{2}3F_3F_7}{9F_{3}^{}{}_{}{}^{11/3}}m^{5/3}+O(m^{7/3}),$$
(17)
where we expand $`F^1(m)`$ up to the order of $`m^{5/3}`$ because all terms up to this order are needed for later calculations. We can set $`\nu _0(t)=0`$ by using the rescaling freedom of the time coordinate, and Eq. (2) yields $`\nu _2(t)=0`$ . Thus the leading order of $`\nu (t,r)`$ is given by
$$\nu (t,r)=\frac{4F_{3}^{}{}_{}{}^{2}}{R_{1}^{}{}_{}{}^{2}(t)}r^4+O(r^6).$$
(18)
For $`R(t,r)`$, we obtain
$$R(t,r)=\left[1\frac{t}{t_0}\right]^{2/3}r+O(r^3),$$
(19)
where $`t_0`$ is given by
$$t_0=\frac{1}{3}\sqrt{\frac{2}{F_3}}.$$
(20)
The break down of these expansions means the appearance of central singularity. Since the break down is given by $`R_1(t)=0`$, the singularity formation time is $`t=t_0`$.
On the initial spacelike hypersurface, the regularity requires $`R(t=0,r)=r>4F=O(r^3)`$ in a sufficiently small, but finite region around $`r=0`$. From Eq. (3), $`R(t,r)`$ is a monotonically decreasing function with respect to $`t`$, and then each initially collapsing shell approaches $`R=4F`$. Using proper time $`\tau (t,r)=e^\nu 𝑑t`$ of each shell, the behavior of this approach is
$$R4F\mathrm{exp}\left(\frac{\tau }{8F}\right).$$
(21)
This behavior shows that $`R`$ approaches $`4F`$ asymptotically and thus $`R`$ is always $`R4F`$ if it initially holds. Since the trapped region is given by $`0R<2F`$, the region around the center is not trapped eternally. Hence the central singularity is globally naked.
### B Algebraic root equation method
The nakedness of the central singularity in the HIN spacetime is also confirmed by examining the algebraic root equation, which probes the existence of outgoing radial null geodesics from the center . To obtain the root equation, consider the radial null geodesic equation in mass-area coordinates
$$\frac{dR}{dm}=\frac{B\sqrt{H}}{C}=J_{}\sqrt{H(m,R)},$$
(22)
where $`J_{}`$ is given by
$$J_{}|V|\left(\frac{R}{\sqrt{R^2+16F^2}}|V|\right)=\frac{|z1|}{(z^2+1)\sqrt{2z}}\left(z\frac{|z1|}{\sqrt{2z}}\right),$$
(23)
and we have defined $`zR/4m(1)`$ for later convenience. The upper and lower signs refer to outgoing and ingoing null geodesics in the collapse phase, respectively. Hereafter we use this sign convention. To investigate the behavior of null geodesics near the center, we define
$$x(m)=\frac{R}{2m^\alpha },$$
(24)
where $`\alpha `$ is a constant. Applying l’Hospital theorem to Eq. (24), we obtain the root equation
$`x_0`$ $`=`$ $`\underset{m0}{lim}{\displaystyle \frac{m^{1\alpha }}{2\alpha }}{\displaystyle \frac{dR}{dm}}`$ (25)
$`=`$ $`\underset{m0}{lim}{\displaystyle \frac{m^{3(1\alpha )/2}(x_02m^{1\alpha })}{2\alpha x_0\left(x_{0}^{}{}_{}{}^{2}+4m^{2(1\alpha )}\right)}}\sqrt{H(m,2x_0m^\alpha )}\left[x_0^{3/2}m^{(1\alpha )/2}(x_02m^{1\alpha })\right],`$ (26)
where $`x_0x(0)`$ is introduced. If we find a consistent set of $`x_0`$ and $`\alpha `$, we have null geodesics which behave as Eq. (24) terminating at the center $`m=0`$. The root equation method picks up only the geodesics behaving as Eq. (24). To find the another possible null geodesics, we must solve the null geodesic equation.
## III Null geodesics in the HIN spacetime
From the detailed analysis of Eq.(26), we can prove that all possible values of $`\alpha `$ for positive finite $`x_0`$ are 1/3, 7/9 and 1. In reality, these three values of $`\alpha `$ correspond to three different types of null geodesics in the HIN spacetime, and it will be proved in this section that there are no other null geodesics in the HIN spacetime. We will see below the derivations and properties of these three types of null geodesics in detail.
### A $`\alpha =1/3`$ : regular null geodesics
The null geodesics with $`\alpha =1/3`$ correspond to regular null geodesics. For this case, the value of $`x_0`$ for $`\alpha =1/3`$ cannot be determined by the root equation method. We will see that the regular null geodesics actually have $`\alpha =1/3`$. From Eqs. (1), (18) and (19), the radial null geodesic equation for lowest order is
$$\frac{dt}{dr}\pm \left[1\frac{t}{t_0}\right]^{2/3}.$$
(27)
Inserting the solution $`t=t(r)`$ of Eq. (27) to Eq. (19), we obtain the regular null geodesics in the mass-area coordinates
$$R(m)\left(1\frac{t(0)}{t_0}\right)^{2/3}\left(\frac{m}{F_3}\right)^{1/3},$$
(28)
where $`t(0)<t_0`$ is the time when the null geodesics arrive at or emanate from the regular center $`r=0`$. The arbitrariness of $`x_0`$ for $`\alpha =1/3`$ comes from the arbitrary constant $`t(0)`$ of Eq. (28).
We can find that the scalar curvature is finite at the center along these null geodesics (see Appendix A). In fact, we can prove that all the null geodesics with $`\alpha =1/3`$ terminate at or emanate from the regular center. The integration of Eq. (3) by using the proper time $`\tau (t,r)`$ is
$$\tau (t,r)=c(m)4m_{R_0(m)/4m}^z\frac{\sqrt{2s(s^2+1)}}{s1}𝑑s,$$
(29)
where $`R_0(m)`$ is the initial area radius $`R(0,r)=F^1(m)`$ and $`c(m)`$ is an arbitrary function of $`m`$. For the regular spacetime, $`\nu (t,0)=\nu _0(t)=0`$ implies that we can set $`\tau =t`$ on $`m=0`$. Then the singularity appears on $`m=0`$ when the proper time is $`\tau =t_0`$. The elliptic integral of Eq. (29) is explicitly integrated by approximating the integrand. $`R_0(m)/4m1`$ is satisfied around the regular center $`m=0`$ and if we consider the condition $`z1`$, the integrand of Eq. (29) is approximately $`\sqrt{2s}`$. Then, $`\tau `$ for each shell becomes
$$\tau c(m)+t_0\frac{\sqrt{2}}{3}\left(\frac{R}{m^{1/3}}\right)^{3/2}.$$
(30)
We consider the null geodesics of Eq. (24) for $`\alpha =1/3`$. From Eq. (30), the proper time along these geodesics is
$$\tau =c(m)+t_0\frac{4}{3}x^{3/2}.$$
(31)
For the regular shell motion of Eq. (19), Eq. (30) becomes
$$\tau =c(m)+t.$$
(32)
Taking into account $`\tau =t`$ on $`m=0`$, we have $`c(0)=0`$ and thus Eq. (31) implies $`\tau <t_0`$ at $`m=0`$. Therefore all the null geodesics behaving as $`Rm^{1/3}`$ terminate at $`m=0`$ before the singularity appears.
### B $`\alpha =7/9`$ : the earliest singular null geodesic
The singular null geodesic with $`\alpha =7/9`$ for both ingoing and outgoing geodesics was found by Harada, Iguchi and Nakao . The behavior of the null geodesic as
$$R2x_0m^{7/9},$$
(33)
is not regular and thus it terminates at the central singularity. The coefficient $`x_0`$ is given by
$$x_0=\left(\frac{24F_3^2F_5}{4\sqrt{2}F_3^{13/6}}\right)^{2/3},$$
(34)
for $`F_5<24F_{3}^{}{}_{}{}^{2}`$. Note that $`F_5<24F_{3}^{}{}_{}{}^{2}`$ is the same as the requirement of no shell-crossing singularity and this condition holds if $`ϵ(0,r)`$ is a non-increasing function of $`r`$. It is found that Eq. (30) reduces to $`\tau =t_0`$ at $`m=0`$ for $`\frac{1}{3}<\alpha <1`$. Thus all the null geodesics of $`\frac{1}{3}<\alpha <1`$, particularly $`\alpha =7/9`$, terminate at the singularity. The scalar curvature diverges along the null geodesic with $`\alpha =7/9`$ (see Appendix A).
We will see below that there is only one ingoing or outgoing null geodesic with $`\alpha =7/9`$ and that this null geodesic is the first one which arrives at or emanates from the singularity at $`t=t_0`$. The null geodesic equation in mass-area coordinates has a singular point at $`(m,R)=(0,0)`$. To make the singular point tractable, we introduce new coordinates as
$`\chi `$ $``$ $`m^{1/9},`$ (35)
$`\vartheta `$ $``$ $`\left({\displaystyle \frac{R}{m^{7/9}}}\right)^{3/2},`$ (36)
and then Eq. (22) becomes
$$\frac{d\vartheta }{d\chi }+\frac{6}{\chi }(\vartheta \lambda )=\lambda \mathrm{\Psi }_{}(\chi ,\vartheta ),$$
(37)
where $`\mathrm{\Psi }_{}`$ is given by
$`\mathrm{\Psi }_{}(\chi ,\vartheta )`$ $`=`$ $`{\displaystyle \frac{1}{\chi }}\left({\displaystyle \frac{\psi _{}(\chi ,\vartheta )}{\lambda }}6\right),`$ (38)
$`\psi _{}(\chi ,\vartheta )`$ $`=`$ $`{\displaystyle \frac{9}{2}}\left(3\chi ^2\vartheta ^{1/3}\sqrt{H(\chi ,\vartheta )}J_{}(\chi ,\vartheta )\vartheta \right),`$ (39)
and we have introduced a parameter $`0<\lambda <\mathrm{}`$.
The form of Eq. (37) is similar to the form of the null geodesic equation of the LTB spacetime given in . One may now follow to Christodoulou’s argument to show the existence and uniqueness of a continuous solution of Eq. (37) . It is sufficient to consider the null geodesic equation in the neighborhood of the center $`\chi =0`$ since we are interested in the radial null geodesics terminating at $`\chi =0`$. There is no singular point at $`\chi >0`$ where $`\vartheta `$ is strictly positive by definition. Problems appear when we consider the center $`\chi =0`$. We expand $`J_{}`$ and $`\sqrt{H}`$ around $`\chi =0`$ using Eq. (17),
$`J_{}`$ $``$ $`\sqrt{2}\left({\displaystyle \frac{\chi }{\vartheta ^{1/3}}}\right)2\left({\displaystyle \frac{\chi }{\vartheta ^{1/3}}}\right)^24\sqrt{2}\left({\displaystyle \frac{\chi }{\vartheta ^{1/3}}}\right)^3+O\left(\left({\displaystyle \frac{\chi }{\vartheta ^{1/3}}}\right)^4\right),`$ (40)
$`\sqrt{H}`$ $``$ $`{\displaystyle \frac{8x_{0}^{}{}_{}{}^{3/2}}{9}}{\displaystyle \frac{1}{\chi ^3}}+{\displaystyle \frac{1}{3\sqrt{2}}}\left({\displaystyle \frac{\vartheta ^{1/3}}{\chi }}\right)^32\sqrt{2}\left({\displaystyle \frac{\vartheta ^{1/3}}{\chi }}\right)+O(\chi ^3)+O\left({\displaystyle \frac{\chi }{\vartheta ^{1/3}}}\right),`$ (41)
and from this expansion, $`\psi _{}`$ is
$$\psi _{}(\chi ,\vartheta )\psi _0+\psi _1(\vartheta )\chi +\psi _2(\vartheta )\chi ^2+O(\chi ^3),$$
(42)
where the coefficients of each order are
$`\psi _0`$ $`=`$ $`12\sqrt{2}x_0^{3/2},`$ (43)
$`\psi _1(\vartheta )`$ $`=`$ $`{\displaystyle \frac{\vartheta ^{2/3}}{\sqrt{2}}}\left(9+{\displaystyle \frac{2\psi _0}{\vartheta }}\right),`$ (44)
$`\psi _2(\vartheta )`$ $`=`$ $`12\vartheta ^{1/3}\left(6+{\displaystyle \frac{\psi _0}{3\vartheta }}\right).`$ (45)
If we choose the parameter $`\lambda `$ to be $`\lambda =\lambda _0(2x_0)^{3/2}`$, $`\mathrm{\Psi }_{}`$ is at least $`C^1`$ in $`\chi 0,\vartheta >0`$. We can apply the contraction mapping principle to Eq. (37) to find that there exists the solution satisfying $`\vartheta (0)=\lambda _0`$, and moreover that it is the unique solution to Eq. (37) which is continuous at $`\chi =0`$. This solution with $`\vartheta (0)=\lambda _0`$ exactly agrees with the geodesics of Eq. (33). The proof is presented in Appendix B. Therefore there is no other solution with $`0<\vartheta (0)<\mathrm{}`$. In other words, another possible solution must be $`\vartheta (0)=0`$ or $`\mathrm{}`$.
We consider possible solutions with $`\vartheta (0)=\mathrm{}`$. When $`\vartheta (0)=\mathrm{}`$, Eq. (37) is approximated around $`\chi =0`$ by using Eqs. (40) and (41) which are valid even in this limit as follows:
$$\frac{d\vartheta }{d\chi }\frac{6\vartheta }{\chi }.$$
(46)
The integration of this equation gives $`\vartheta (\chi )1/\chi ^6`$. This behavior coincides with the regular null geodesics of Eq. (28) because all the geodesics behaving as $`Rm^{1/3}`$ terminate at the regular center as we have shown before.
It is important that $`R(t,r)`$ is a monotonically decreasing function with respect to $`t`$, and thus that $`\vartheta `$ decreases as $`t`$ increases when $`\chi `$ is fixed. Because of this time direction in the mass-area coordinates and the fact that there are no other null geodesics with $`(2x_0)^{3/2}<\vartheta (0)<\mathrm{}`$, the geodesic with $`\vartheta (0)=(2x_0)^{3/2}`$ is the first null geodesic which arrives at or emanates from the appeared singularity. Hence we conclude that the arrival or emanational time in comoving coordinates is the singularity formation time $`t=t_0`$.
### C $`\alpha =1`$ : later singular null geodesics
There are also null geodesics with $`\alpha =1`$. $`x(m)`$ is expressed for $`\alpha =1`$ by a non-analytic function
$$x(m)2+2\mathrm{exp}\left(\frac{D}{\sqrt{m}}\right),$$
(47)
where $`D`$ is the positive constant which parameterizes the null geodesics. We find that the scalar curvature diverges at the center along these null geodesics as is seen in Appendix A. We will see these null geodesics in detail below.
In the coordinate $`\vartheta `$, these null geodesics are described by solutions with $`\vartheta (0)=0`$ if they exist. We first search the solution with $`z(0)=\mathrm{}`$ and $`\vartheta (0)=0`$. Under these conditions, Eq. (22) reduces around $`m=0`$ to
$$\frac{dz}{dm}=\frac{1}{4m}\left(\sqrt{H}J_{}4z\right)\frac{1}{m\sqrt{z}}\left(\frac{\sqrt{2}x_{0}^{}{}_{}{}^{3/2}}{9}m^{1/3}\frac{2}{3}z^{3/2}\right).$$
(48)
For $`z^{3/2}m^{1/3}`$, we can neglect the second term of Eq. (48) and immediately integrate it. However, the solution contradicts $`z^{3/2}m^{1/3}`$ because it is given by $`zm^{2/9}`$. For $`z^{3/2}m^{1/3}`$ or $`z^{3/2}m^{1/3}`$, there are consistent solutions. The solution of Eq. (48) for $`z^{3/2}m^{1/3}`$ is
$$z\frac{1}{m^{2/3}},$$
(49)
which corresponds regular null geodesics with $`\alpha =1/3`$. When $`z^{3/2}m^{1/3}`$, the integration gives
$$zm^{2/9},$$
(50)
which corresponds to the unique null geodesic with $`\alpha =7/9`$. These two solutions satisfy the condition $`z(0)=\mathrm{}`$, but do not satisfy $`\vartheta (0)=0`$. Consequently there is no solution which satisfies these two conditions.
We consider the geodesics with $`z(0)<\mathrm{}`$ and $`\vartheta (0)=0`$. Using l’ Hospital theorem, we obtain a restriction on $`J_{}\sqrt{H}`$ as
$$4z(0)=\underset{m0}{lim}\frac{dR}{dm}=\underset{m0}{lim}J_{}\sqrt{H}<\mathrm{}.$$
(51)
Though $`J_{}`$ is strictly positive at $`m=0`$ as long as $`z(0)>1`$, $`\sqrt{H}`$ diverges to positive infinity at $`m=0`$ as
$$\sqrt{H}=4\sqrt{2}\left(\frac{\beta }{m^{1/3}}+\gamma m^{1/3}+O(m)\right)+\frac{4\sqrt{2}}{3}\frac{\sqrt{z}(z^24z+9)}{z1}+4\sqrt{2}\mathrm{ln}\frac{\sqrt{z}+1}{\sqrt{z}1},$$
(52)
where we have defined
$`\beta `$ $``$ $`{\displaystyle \frac{\sqrt{2}}{9}}x_0^{3/2},`$ (53)
$`\gamma `$ $``$ $`{\displaystyle \frac{1440F_{3}^{}{}_{}{}^{4}+48F_{3}^{}{}_{}{}^{2}F_517F_{5}^{}{}_{}{}^{2}+12F_3F_7}{216F_{3}^{}{}_{}{}^{23/6}}}.`$ (54)
Thus, there exists no solution with the boundary condition $`1<z(0)<\mathrm{}`$. However it may be possible for null geodesics to satisfy Eq. (51) only if $`z(0)=1`$.
We introduce the new coordinate $`yz1`$ to study the solution $`y(m)`$ with $`y(0)=0`$. By expanding Eq. (22) around $`y=0`$, we obtain the consistent solution with $`y(0)=0`$ in the lowest order approximation. We give here the null geodesic equation which is expanded up to the enough order so that the difference between ingoing and outgoing null geodesics appears in the expansion;
$$\frac{dy}{dm}\frac{y}{2m}\left(\mathrm{ln}y+\frac{\beta }{m^{1/3}}+\delta \sqrt{2}+\gamma m^{1/3}\right).$$
(55)
It can be integrated to
$$y\mathrm{exp}\left(\frac{D}{\sqrt{m}}+\frac{3\beta }{m^{1/3}}+\delta \sqrt{2}+\frac{3\gamma }{5}m^{1/3}\right),$$
(56)
where $`D>0`$ is an integration constant and $`\delta `$ is defined as $`\delta 2\mathrm{ln}2\frac{8}{3}`$. We will see in Sec. V that the constant $`D`$ is related to the time when the null geodesics terminate at the center $`m=0`$. The geodesics of Eq. (56) satisfy $`y(0)=0`$ and they are classified into the geodesics with $`\alpha =1`$. To obtain the proper time when the null geodesics terminate at the singularity, we consider Eq. (29) again. By the approximation of the integrand, Eq. (29) reduces around $`m=0`$ to
$$\tau c(m)4m\left(2\mathrm{ln}y(m)\frac{t_0}{4m}\right).$$
(57)
Then the proper time when the null geodesics terminate at $`m=0`$ is $`\tau =t_0`$ and it means that the null geodesics terminate at the central singularity. As a result, it is concluded that all the solutions with $`\vartheta (0)=0`$ are the solutions which terminate at the singularity with $`y(0)=0`$.
## IV Causal structure of the HIN spacetime
### A Timelike singularity
It is important that Eq. (56) includes all the null geodesics emanating from the center after the singularity appears. We construct double null coordinates by using Eq. (56) and study the causal structure of the spacetime which is covered by these null geodesics. From Eq. (56), we introduce double null coordinates $`(u,v)`$ which satisfy around $`m=0`$
$`2\sqrt{2}u`$ $``$ $`\sqrt{m}\left(\mathrm{ln}y+{\displaystyle \frac{3\beta }{m^{1/3}}}+\delta \sqrt{2}+{\displaystyle \frac{3\gamma }{5}}m^{1/3}\right),`$ (58)
$`2\sqrt{2}v`$ $``$ $`\sqrt{m}\left(\mathrm{ln}y+{\displaystyle \frac{3\beta }{m^{1/3}}}+\delta +\sqrt{2}+{\displaystyle \frac{3\gamma }{5}}m^{1/3}\right).`$ (59)
Then, $`m`$ and $`y`$ are expressed in the null coordinates by
$`\sqrt{m}`$ $``$ $`vu,`$ (60)
$`y`$ $``$ $`\mathrm{exp}\left({\displaystyle \frac{\sqrt{2}(v+u)}{vu}}+{\displaystyle \frac{3\beta }{(vu)^{2/3}}}+\delta +{\displaystyle \frac{3\gamma }{5}}(vu)^{2/3}\right).`$ (61)
We are considering a sufficiently small, but finite region around $`m=0`$. Then Eq. (60) restricts $`u`$ and $`v`$ to $`uv`$. From Eqs. (60) and (61), $`dm`$ and $`dR`$ are given by
$`dm`$ $``$ $`2(vu)(dvdu),`$ (62)
$`dR4dm`$ $``$ $`4\sqrt{2}y\left(v+u\sqrt{2}\beta (vu)^{1/12}+(\sqrt{2}+1)(vu)+{\displaystyle \frac{\sqrt{2}}{5}}\gamma (vu)^{5/12}\right)du`$ (64)
$`+4\sqrt{2}y\left(v+u\sqrt{2}\beta (vu)^{1/12}+(\sqrt{2}1)(vu)+{\displaystyle \frac{\sqrt{2}}{5}}\gamma (vu)^{5/12}\right)dv.`$
Inserting Eqs. (62) and (64) into Eq. (6) and expanding the metric functions $`A,B`$ and $`C`$ around $`(m,y)=(0,0)`$, we obtain the line element in the double null coordinates as
$$ds^2512(vu)^2dvdu+R^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2).$$
(65)
The details of the calculation are given in Appendix C. In the null coordinates, the central singularity is represented at $`u=v`$ which corresponds to the center $`m=0`$ because of Eq. (60). It is found that the world line of $`u=v`$ is timelike. Hence the central singularity in the spacetime which is covered by the null geodesics given by Eq. (56) is timelike.
### B Penrose diagram
Here we summarize the obtained results. The whole of the spacetime is covered by three types of null geodesics which can be classified by the value of $`\alpha `$ into $`\alpha =1/3`$, $`7/9`$ and $`1`$. The null geodesics with $`\alpha =1/3`$ emanate from or terminate at the regular center and they are parametrized by one parameter. There is only one null geodesic with $`\alpha =7/9`$ and it is the earliest one which emanates from or terminates at the naked singularity. The null geodesics with $`\alpha =1`$ emanate from or terminate at the timelike singularity and they are parametrized by one parameter. From these results, it is now possible to draw the conformal diagram of the HIN spacetime (see Fig. 1). For comparison, we also present the conformal diagram of the LTB spacetime with naked singularity (see Fig. 2). It is found that the effect of counterrotation makes the singularity timelike in this model.
### C Curvature strength of the singularity
We should note the curvature strength of the singularity in the HIN spacetime. According to Tipler and Królak, we classify the curvature strength whether the singularity satisfies the strong curvature condition (SCC) or the limiting focusing condition (LFC) . Harada, Nakao and Iguchi studied the SCC and the LFC for spherically symmetric spacetimes with vanishing radial pressure . Applying theorems 1, 2 and 3 of their paper to the HIN spacetime, we can know the curvature strength along each null geodesic. For $`\alpha =7/9`$, not the SCC but only the LFC is satisfied along the null geodesics. For $`\alpha =1`$, the SCC is satisfied along the null geodesics. In conclusion, the naked singularity is relatively weak at the formation but becomes strong after that.
## V Asymptotic Behavior
It is worth while examining the asymptotic behavior of the spacetime. We concern with the null geodesic behavior in the comoving coordinates corresponding to Eq. (56). To obtain an insight, we consider the asymptotic behavior of the HIN solution in comoving coordinates. Since the time coordinate $`t`$ at the center does not proceed after the singularity formation, we must introduce new time coordinate which no longer agrees with the proper time at $`r=0`$. We denote the new time coordinate as $`T`$ and consider that the time $`t`$ in Eqs. (1), (2) and (3) is replaced by $`T`$.
We take the limit $`R=4F`$ because $`R(T,r)`$ approaches $`4F(r)`$ asymptotically. From Eq. (2), the metric in this limit is given by
$$ds^2=\frac{4F(r)}{K(T)}dT^2+2(4F^{}(r))^2dr^2+(4F(r))^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2).$$
(66)
The function $`K(T)`$ is an arbitrary function of $`T`$ which comes from the integration of Eq. (2). For simplicity we set $`K(T)=K=constant`$ by rescaling the time coordinate. The energy density in this limit is
$`ϵ`$ $`=`$ $`{\displaystyle \frac{1}{256\pi F^2}}.`$ (67)
This solution is a static system of counterrotating particles, which is called the Einstein cluster. In particular, the solution has timelike naked singularity at the center.
To study the asymptotic behavior of the shell motion for fixed $`r`$, we set perturbed quantities $`e(T,r)`$ and $`d(T,r)`$ for the coordinates $`(T,r)`$ around $`R=4F`$ as,
$`R(T,r)`$ $`=`$ $`4F(r)\left[1+e(T,r)\right],`$ (68)
$`e^{2\nu (T,r)}`$ $`=`$ $`e^{d(T,r)}{\displaystyle \frac{R(T,r)}{K}}{\displaystyle \frac{4F(r)}{K}}[1+e(T,r)+d(T,r)].`$ (69)
Inserting these quantities into Eqs. (2) and (3), we obtain following perturbed equations in linear order,
$`\dot{e}(T,r)`$ $`=`$ $`{\displaystyle \frac{e(T,r)}{4\sqrt{KF(r)}}},`$ (70)
$`d^{}(T,r)`$ $`=`$ $`{\displaystyle \frac{F^{}(r)}{F(r)}}e(T,r).`$ (71)
The solutions of these equations are
$`e(T,r)`$ $`=`$ $`E(r)\mathrm{exp}\left({\displaystyle \frac{T}{4\sqrt{KF(r)}}}\right),`$ (72)
$`d(T,r)`$ $`=`$ $`{\displaystyle e(T,r)\frac{F^{}}{F}𝑑r}+G(T),`$ (73)
where $`E(r)`$ and $`G(T)`$ are arbitrary functions of the comoving radius $`r`$ and time $`T`$ respectively. We have implicitly assumed that the time $`T`$ is very large compared to the singularity formation time. However, Eqs. (72) and (73) imply that this perturbation scheme is valid soon after the singularity formation time as long as we consider the region in which the radius $`r`$ is very small.
In the asymptotic region, the proper time becomes
$$\tau (T,r)=2\sqrt{\frac{F}{K}}[T+const.+O\left(e(T,r)\right)].$$
(74)
It is found that the asymptotic behavior of Eqs. (72) and (74) completely coincides with the already obtained behavior of Eq. (21). Using these results of perturbation, the null geodesic equation in the comoving coordinates is
$$\frac{dT}{dr}=\pm \sqrt{\frac{2K}{F}}\left(2F^{}+2Fe^{}F^{}d\right),$$
(75)
and it is immediately integrated to
$$T(r)=T(0)\pm \sqrt{2K}\left(4\sqrt{F(r)}+\frac{dr}{\sqrt{F(r)}}\left[2Fe^{}F^{}d\right]\right).$$
(76)
Integration constant $`T(0)`$ is the time when these null geodesics terminate at or emanate from the singularity, and this parameterization of the null geodesic family corresponds to the fact that the singularity is timelike. From Eqs. (68), (72) and (76), we obtain the null geodesics in the mass-area coordinates corresponding to Eq. (76),
$$y=\frac{R}{4F}1=E(r)\mathrm{exp}\left(\frac{T(0)}{4\sqrt{KF(r)}}\sqrt{2}\frac{1}{2\sqrt{2F(r)}}\frac{dr}{\sqrt{F(r)}}[2Fe^{}F^{}d]\right).$$
(77)
Comparing this result to Eq. (56), the parameter $`D`$ relates to the arrival time $`T(0)`$ as $`T(0)=4D\sqrt{K}`$.
## VI Conclusions
We have studied the causal structure of the HIN spacetime and it was shown that the central massless singularity of this spacetime is timelike. To show this fact, we have investigated the radial null geodesics in detail. The null geodesics are classified to three types, $`\alpha =`$ 1/3, 7/9 and 1, by their dependence of $`R`$ on $`m`$ near the center. One is regular and the other two are singular. The classification of the singular geodesics corresponds to their arrival or emanational time at the central singularity. The $`\alpha =7/9`$ type is the earliest singular null geodesic which arrives at or emanates from the singularity at its formation time $`t_0`$, while the geodesics with $`\alpha =1`$ arrive at or emanate from the singularity after its appearance.
We have shown that singular null geodesics with $`\alpha =1`$ exactly parametrized by their arrival or emanational time and that there is only one set of ingoing and outgoing geodesics for each parameter. This fact shows that the central singularity has timelike property. We have also constructed double null coordinates around the central singularity from the null geodesics with $`\alpha =1`$. The line element in this double null coordinates shows that there is timelike singularity in this spacetime. We have considered the asymptotic behavior of this spacetime after the singularity appeared in comoving coordinates. It gives us understanding of the null geodesic behavior in mass-area coordinates, and of the parametrization of the null geodesic family. The curvature strength of this singularity was also investigated. The LFC is satisfied along the null geodesics with $`\alpha =7/9`$ and the SCC is satisfied for $`\alpha =1`$. The curvature strength of the naked singularity is relatively weak at the formation and it becomes strong after that.
In summary, the HIN solution describes a dynamical formation of timelike naked singularity, that is, the birth of timelike singularity. The solution dynamically tends to the static singular Einstein cluster as time proceeds. Though the HIN system is simply composed of collisionless particles, the collapse leads to the nontrivial causal structure. It implies that the effects of rotation and tangential pressure play important roles in the final stage of collapse, particularly the singularity formation.
###### Acknowledgements.
We are grateful to T. Nakamura, H. Kodama, T.P. Singh, K. Nakao, A. Ishibashi and S.S. Deshingkar for helpful discussions. This work was partly supported by the Grant-in-Aid for Scientific Research (No. 05540) from the Japanese Ministry of Education, Science, Sports and Culture.
## A Scalar Curvature along the null geodesics
Singularities are boundary points of spacetime where the normal differentiability breaks down. If the energy density or the curvature invariant diverges at boundary points, the points are singularities. In this appendix, we give the scalar curvature $`R_\mu ^\mu `$ in the HIN spacetime along the radial null geodesics terminating at the center $`r=0`$, and show that the center is singular along the geodesics with $`\alpha =7/9`$ and $`1`$.
The scalar curvature $`R_\mu ^\mu `$ of the HIN spacetime is given by
$$R_\mu ^\mu =\frac{8\pi R^2}{R^2+16F^2}ϵ=\frac{\sqrt{2}}{(R4m)\sqrt{mR}\sqrt{H}}.$$
(A1)
Along the geodesics $`R=2xm^{1/3}`$, which are shown to be regular null geodesics, $`R_\mu ^\mu `$ is given by
$$R_\mu ^\mu =\frac{3}{4x^3}+O(m^{2/3}),$$
(A2)
and it is finite at $`m=0`$. On the other hand, when we consider $`R_\mu ^\mu `$ along the geodesics, $`R=2x_0m^{7/9}`$ and $`R=4m(1+y)`$ of Eq. (56), the scalar curvatures are given by
$$R_\mu ^\mu =\frac{9}{28x_{0}^{}{}_{}{}^{3}m^{4/3}}+O(m^{10/9}),$$
(A3)
and
$$R_\mu ^\mu =\frac{1}{64m^2}+O(m^{5/2}y),$$
(A4)
respectively. The scalar curvature diverges as $`m0`$ in these cases, and therefore the center is definitely singular.
## B contraction mapping principle
We prove the existence of the unique continuous solution of Eq. (37). It is based on the discussion given by Christodoulou for the LTB model . Consider the differential equation obtained from Eq. (37) by replacing $`\vartheta `$ in $`\mathrm{\Psi }_{}(\chi ,\vartheta )`$ by a given continuous function $`\overline{\vartheta }>0`$:
$$\frac{d\vartheta }{d\chi }+\frac{6}{\chi }(\vartheta \lambda _0)=\lambda _0\mathrm{\Psi }_{}(\chi ,\overline{\vartheta }(\chi )).$$
(B1)
Since we have chosen $`\lambda =\lambda _0`$, $`\mathrm{\Psi }_{}(\chi ,\overline{\vartheta })`$ is at least $`C^1`$ in the strip $`\chi [0,\chi _1]`$ and $`\overline{\vartheta }(0,\mu _1]`$, where we assume $`\chi _1`$ is sufficiently small. The continuous solution of Eq. (B1) is only the solution with $`\vartheta (0)=\lambda _0(>0)`$. This solution is given by
$$\vartheta (\chi )=\lambda _0\left(1+\chi _0^1s^6\mathrm{\Psi }_{}(s\chi ,\overline{\vartheta }(s\chi ))𝑑s\right).$$
(B2)
We consider the nonlinear map $`T_{\lambda _0}`$ defined as $`\vartheta =T_{\lambda _0}(\overline{\vartheta })`$. Since we are considering the function $`\overline{\vartheta }`$ with $`\overline{\vartheta }(0)>0`$, we obtain $`T_{\lambda _0}(\overline{\vartheta })(0)=\lambda _0`$ from expansions of Eqs. (40) and (41). Thus the possible fixed point of $`T_{\lambda _0}`$ is only one which satisfies $`\vartheta _{\mathrm{FP}}(0)=\lambda _0`$. To prove the existence of the fixed point, we consider the set $`V_\mu `$ consisting of all $`\overline{\vartheta }`$ such that for $`\chi [0,\chi _2]`$,
$$\mu _{}\overline{\vartheta }(\chi )\mu _+,$$
(B3)
where the lower and the upper bound are defined by $`\mu _{}\lambda _0\mu _2`$ and $`\mu _+\lambda _0+\mu _2`$ for a sufficiently small $`\mu _2`$. From these bounds, the nonlinear map is also restricted in
$$\mu _{}T_{\lambda _0}(\overline{\vartheta })(\chi )\mu _+,$$
(B4)
if we choose a $`\chi `$ such that
$$\chi \chi _3\frac{7\mu _2}{\lambda _0\mathrm{\Delta }_1},$$
(B5)
where
$$\mathrm{\Delta }_1=\underset{0\chi \chi _2}{sup}\underset{\mu _{}\overline{\vartheta }\mu _+}{sup}|\mathrm{\Psi }_{}|.$$
(B6)
The map $`T_{\lambda _0}`$ sends $`V_\mu `$ into itself for all $`\chi _4\mathrm{min}\{\chi _2,\chi _3\}`$. Let
$$\mathrm{\Delta }_2=\underset{0\chi \chi _2}{sup}\underset{\mu _{}\overline{\vartheta }\mu _+}{sup}\left|\frac{\mathrm{\Psi }_{}}{\overline{\vartheta }}\right|.$$
(B7)
Then we obtain from Eq. (B2) for $`\overline{\vartheta }_1,\overline{\vartheta }_2V_\mu `$,
$$T_{\lambda _0}(\overline{\vartheta }_1)T_{\lambda _0}(\overline{\vartheta }_2)\frac{\chi _4\lambda _0\mathrm{\Delta }_2}{7}\overline{\vartheta }_1\overline{\vartheta }_2,$$
(B8)
where $``$ $``$ denotes the supremum norm. If we choose $`\chi _4<7/(\lambda _0\mathrm{\Delta }_2)`$, the map $`T_{\lambda _0}`$ is contractive in $`V_\mu `$. Hence $`T_{\lambda _0}`$ has a unique fixed point $`\vartheta _{\mathrm{FP}}V_\mu `$ which is given by
$$\vartheta _{\mathrm{FP}}(\chi )=T_{\lambda _0}(\vartheta _{\mathrm{FP}})(\chi )=\lambda _0\left(1+\chi _0^1s^6\mathrm{\Psi }_{}(s\chi ,\vartheta _{\mathrm{FP}}(s\chi ))𝑑s\right).$$
(B9)
Therefore we have the unique continuous solution with $`\vartheta (0)=\lambda _0`$.
## C derivation of eq. (65)
In this appendix, we construct double null coordinates by using Eq. (56). We define double null coordinates $`(u,v)`$ as Eqs. (58) and (59). From this definition, $`dm`$ and $`dR`$ are given by
$`dm`$ $`=`$ $`p_udu+p_vdv,`$ (C1)
$`dR`$ $`=`$ $`4dm+(q_udu+q_vdv),`$ (C2)
where
$`p_u`$ $``$ $`p_v2(uv),`$ (C3)
$`q_u`$ $``$ $`4\sqrt{2}y\left(v+u\sqrt{2}\beta (vu)^{1/12}+(\sqrt{2}+1)(vu)+{\displaystyle \frac{\sqrt{2}}{5}}\gamma (vu)^{5/12}\right)du,`$ (C4)
$`q_v`$ $``$ $`4\sqrt{2}y\left(v+u\sqrt{2}\beta (vu)^{1/12}+(\sqrt{2}1)(vu)+{\displaystyle \frac{\sqrt{2}}{5}}\gamma (vu)^{5/12}\right)dv.`$ (C5)
The line element of Eq. (1) is rewritten by the coordinate transformation from $`(m,R)`$ to $`(u,v)`$
$`ds^2`$ $`=`$ $`Udu^2+Vdv^2+Wdudv+R^2d\mathrm{\Omega }^2,`$ (C6)
where the metric functions $`U,V`$ and $`W`$ are given by
$`U`$ $`=`$ $`(A+8B+16C)p_u^2Cq_u^22(B+4C)p_uq_u,`$ (C7)
$`V`$ $`=`$ $`(A+8B+16C)p_v^2Cq_v^22(B+4C)p_vq_v,`$ (C8)
$`W`$ $`=`$ $`2(A+8B+16C)p_up_v2Cq_uq_v2(B+4C)(p_uq_v+p_vq_u).`$ (C9)
From Eq. (C3), we obtain
$$WUVC(q_u+q_v)^2,$$
(C10)
and then Eq. (C6) becomes
$$ds^2C(q_u+q_v)^2dudv+U(du^2dudv)+V(dv^2dudv)+R^2d\mathrm{\Omega }^2.$$
(C11)
If the condition
$$C(q_u+q_v)^2U,V$$
(C12)
is satisfied, the double null coordinates would be constructed approximately.
We expand $`\sqrt{H}`$ of Eq. (52) around $`y=0`$,
$$\sqrt{H}=h(m)+\sqrt{2}\left(\frac{8}{y}4\mathrm{ln}y+4(3+\delta )+y+\frac{3}{4}y^2+O(y^3)\right),$$
(C13)
where $`h(m)`$ is defined as
$$h(m)=4\sqrt{2}\left(\frac{\beta }{m^{1/3}}+\gamma m^{1/3}+O(m)\right).$$
(C14)
Inserting Eq. (C13) into $`A`$, $`B`$ and $`C`$, metric functions $`U`$, $`V`$ and $`W`$ are calculated. As a result we obtain
$`\underset{m0}{lim}{\displaystyle \frac{1}{m}}C(q_u+q_v)^2`$ $`=`$ $`512,`$ (C15)
$`\underset{m0}{lim}{\displaystyle \frac{U}{m}}`$ $`=`$ $`0,`$ (C16)
$`\underset{m0}{lim}{\displaystyle \frac{V}{m}}`$ $`=`$ $`0,`$ (C17)
and it satisfies the condition of Eq. (C12). Hence the line element around the center is given by
$$ds^2512(vu)^2dudv+R^2d\mathrm{\Omega }^2.$$
(C18)
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# 1 Introduction
## 1 Introduction
The latest atmospheric neutrino results based on 1117 days of data from Super Kamiokande are still consistent with a standard two neutrino oscillation $`\nu _\mu \nu _\tau `$ with a near maximal mixing angle $`\mathrm{sin}^22\theta _{23}>0.88`$ and a mass square splitting $`\mathrm{\Delta }m_{23}^2`$ from $`1.5\times 10^3`$ to $`5\times 10^3eV^2`$ at 90% CL . The sterile neutrino oscillation hypothesis $`\nu _\mu \nu _s`$ is excluded at 99% CL. Super Kamiokande is also beginning to provide important clues concerning the correct solution to the solar neutrino problem. The latest results from 1117 days of data from Super Kamiokande sees a one sigma day-night asymmetry, and a flat energy spectrum, which together disfavour the small mixing angle (SMA) MSW solution , the just-so vacuum oscillation hypothesis and the sterile neutrino hypotheses. All three possibilities are now excluded at 95% CL. The results allow much of the large mixing angle (LMA) MSW region, which now looks like the leading candidate for the solution to the solar neutrino problem. For example a typical point in the LMA MSW region is $`\mathrm{sin}^22\theta _{12}0.75`$ and $`\mathrm{\Delta }m_{12}^22.5\times 10^5eV^2`$.
The see-saw mechanism implies that the three light neutrino masses arise from some large mass scales corresponding to the Majorana masses of some heavy “right-handed neutrinos” $`N_R^p`$ $`M_{RR}^{pq}`$ ($`p,q=1,\mathrm{},Z`$) whose entries take values which extend from $`10^{14}`$ GeV down to perhaps several orders of magnitude lower. The presence of electroweak scale Dirac mass terms $`m_{LR}^{ip}`$ (a $`3\times Z`$ matrix) connecting the left-handed neutrinos $`\nu _L^i`$ ($`i=1,\mathrm{}3`$) to the right-handed neutrinos $`N_R^p`$ then results in a very light see-saw suppressed effective $`3\times 3`$ Majorana mass matrix
$$m_{LL}=m_{LR}M_{RR}^1m_{LR}^T$$
(1)
for the left-handed neutrinos $`\nu _L^i`$, which are the light physical degrees of freedom observed by experiment. If the neutrino masses arise from the see-saw mechanism then it is natural to assume the existence of a physical neutrino mass hierarchy $`m_{\nu _1}m_{\nu _2}m_{\nu _3}`$, which implies $`\mathrm{\Delta }m_{23}^2m_{\nu _3}^2`$, and $`\mathrm{\Delta }m_{12}^2m_{\nu _2}^2`$, which fixes $`m_{\nu _3}5.6\times 10^2eV`$, and (assuming the LMA MSW solution) $`m_{\nu _2}5.2\times 10^3eV`$, with rather large errors. Thus $`m_{\nu _2}/m_{\nu _3}0.1`$. In view of such a 23 mass hierarchy the presence of a large 23 mixing angle looks a bit surprising at first sight, especially given our experience with small quark mixing angles. Several explanations have been proposed , but the simplest idea is that the contributions to the 23 block of the light effective Majorana matrix come predominantly from a single right-handed neutrino, which causes the 23 subdeterminant to approximately vanish. This mechanism, called single right-handed neutrino dominance (SRHND), was proposed in , and developed for bi-maximal mixing in .
In this paper we shall be concerned with the effect of radiative corrections in models based on SRHND and $`U(1)`$ family symmetry , and in particular on the case of bi-maximal mixing in these models . We choose these models because analytic estimates suggest that the 23 hierarchy arises from a physical mechanism which accounts for the smallness of the 23 subdeterminant, and so this hierarchy should be stable under radiative corrections. Although the radiative corrections to atmospheric mixing are only a modest $`10\%`$, this can nevertheless play an important role in achieving near maximal atmospheric mixing. Our results may be compared to the several RG studies of various models which already exist in the literature , . Many of the existing studies do not take proper account of the heavy right-handed neutrino mass thresholds, often only including the effects of running due to the low energy dimension 5 see-saw operator from the see-saw scale down to low energies , or running from the high energy scale down to low energies including a single right-handed neutrino mass threshold , assuming that all the right-handed neutrinos are degenerate. Our results indicate that the corrections in running from the high energy scale, through the heavy neutrino threshold region, down to the lightest right-handed neutrino mass scale are just as important (and in some cases more so) than the traditionally calculated radiative corrections in running from the lightest right-handed neutrino mass scale down to low energies.
The paper is organised as follows. In section 2 we define the MSSM with $`Z`$ right-handed neutrinos, where the heavy neutrino mass matrix arises from the vacuum expectation value (VEV) of a singlet field $`\mathrm{\Sigma }`$, and write down the renormalisation group equations (RGEs) relevant for running from the unification scale $`M_U2\times 10^{16}GeV`$ down to low energies. In section 3 we review the analytic discussion of the phenomenological conditions for SRHND and LMA MSW involving the different types of heavy Majorana neutrino textures, and show how a $`U(1)`$ family symmetry may be used to satisfy them . In section 4 we discuss three explicit examples of this kind, and then perform our numerical RG analysis of these cases. Section 5 concludes the paper.
## 2 The MSSM with $`Z`$ Right-handed neutrinos
We consider the Yukawa terms with two Higgs doublets augmented by $`Z`$ right-handed neutrinos, which, are given by
$`_{yuk}`$ $`=`$ $`ϵ_{ab}[Y_{ij}^uH_u^aQ_i^bU_j^c+Y_{ij}^dH_d^aQ_i^bD_j^c+Y_{ij}^eH_d^aL_i^bE_j^cY_{ip}^\nu H_u^aL_i^bN_p^c`$ (2)
$`+`$ $`{\displaystyle \frac{1}{2}}Y_{RR}^{pq}\mathrm{\Sigma }N_p^cN_q^c]+H.c.`$
where $`ϵ_{ab}=ϵ_{ba}`$, $`ϵ_{12}=1`$, and the remaining notation is standard except that the $`Z`$ right-handed neutrinos $`N_R^p`$ have been replaced by their CP conjugates $`N_p^c`$ with $`p,q=1,\mathrm{},Z`$ and we have introduced a singlet field $`\mathrm{\Sigma }`$ whose vacuum expectation value (VEV) induces a heavy Majorana matrix $`M_{RR}=<\mathrm{\Sigma }>Y_{RR}`$. When the two Higgs doublets get their VEVS $`<H_u^2>=v_2`$, $`<H_d^1>=v_1`$ with $`\mathrm{tan}\beta v_2/v_1`$ we find the terms
$$_{yuk}=v_2Y_{ij}^uU_iU_j^c+v_1Y_{ij}^dD_iD_j^c+v_1Y_{ij}^eE_iE_j^c+v_2Y_{ip}^\nu N_iN_p^c+\frac{1}{2}M_{RR}^{pq}N_p^cN_q^c+H.c.$$
(3)
Replacing CP conjugate fields we can write in a matrix notation
$$_{yuk}=\overline{U}_Lv_2Y^uU_R+\overline{D}_Lv_1Y^dD_R+\overline{E}_Lv_1Y^eE_R+\overline{N}_Lv_2Y^\nu N_R+\frac{1}{2}N_R^TM_{RR}N_R+H.c.$$
(4)
where we have assumed that all the masses and Yukawa couplings are real and written $`Y^{}=Y`$. The diagonal mass matrices are given by the following unitary transformations
$`v_2Y_{diag}^u=V_{uL}v_2Y^uV_{uR}^{}=\mathrm{diag}(\mathrm{m}_\mathrm{u},\mathrm{m}_\mathrm{c},\mathrm{m}_\mathrm{t}),`$
$`v_1Y_{diag}^d=V_{dL}v_1Y^dV_{dR}^{}=\mathrm{diag}(\mathrm{m}_\mathrm{d},\mathrm{m}_\mathrm{s},\mathrm{m}_\mathrm{b}),`$
$`v_1Y_{diag}^e=V_{eL}v_1Y^eV_{eR}^{}=\mathrm{diag}(\mathrm{m}_\mathrm{e},\mathrm{m}_\mu ,\mathrm{m}_\tau ),`$
$`M_{RR}^{diag}=\mathrm{\Omega }_{RR}M_{RR}\mathrm{\Omega }_{RR}^{}=\mathrm{diag}(\mathrm{M}_{\mathrm{R1}},\mathrm{},\mathrm{M}_{\mathrm{RZ}}).`$ (5)
Below the mass of the lightest right-handed neutrino $`M_{R1}`$ the right-handed neutrino masses may be integrated out of the theory, which corresponds to replacing the last two terms in Eq.2 by a dimension 5 operator
$$ϵ_{ab}Y_{ip}^\nu H_u^aL_i^bN_p^c+\frac{1}{2}M_{RR}^{pq}N_p^cN_q^c+H.c.\frac{1}{2}\kappa _{ij}(ϵ_{ab}H_u^aL_i^b)(ϵ_{a^{}b^{}}H_u^a^{}L_j^b^{})+H.c.$$
(6)
where $`\kappa =Y_\nu M_{RR}^1Y_\nu ^T`$ is simply related to the see-saw mass matrix in Eq.1 when the Higgs fields are replaced by their VEVs
$$m_{LL}=v_2^2\kappa .$$
(7)
Having constructed the light Majorana mass matrix it must then be diagonalised by unitary transformations,
$$m_{LL}^{diag}=V_{\nu L}m_{LL}V_{\nu L}^{}=\mathrm{diag}(\mathrm{m}_{\nu _1},\mathrm{m}_{\nu _2},\mathrm{m}_{\nu _3}).$$
(8)
The CKM matrix is given by
$$V_{CKM}=V_{uL}V_{dL}^{}$$
(9)
The leptonic analogue of the CKM matrix is the MNS matrix defined as
$$V_{MNS}=V_{eL}V_{\nu L}^{}.$$
(10)
which may be parametrised by a sequence of three rotations about the 1,2 and 3 axes, as in the standard CKM parametrisation,
$$V_{MNS}=R_{23}R_{13}R_{12}$$
(11)
where
$$R_{23}=\left(\begin{array}{ccc}\hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill c_{23}& \hfill s_{23}\\ \hfill 0& \hfill s_{23}& \hfill c_{23}\end{array}\right),R_{13}=\left(\begin{array}{ccc}\hfill c_{13}& \hfill 0& \hfill s_{13}\\ \hfill 0& \hfill 1& \hfill 0\\ \hfill s_{13}& \hfill 0& \hfill c_{13}\end{array}\right),R_{12}=\left(\begin{array}{ccc}\hfill c_{12}& \hfill s_{12}& \hfill 0\\ \hfill s_{12}& \hfill c_{12}& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1\end{array}\right)$$
(12)
where $`s_{ij}=\mathrm{sin}\theta _{ij}`$, $`c_{ij}=\mathrm{cos}\theta _{ij}`$, and $`\theta _{ij}`$ refer to lepton mixing angles. Note that we completely ignore CP violating phases in this paper. From the unitarity conditions of the MNS matrix elements in Eq.11, and parametrisation in Eq.12, the mixing angles can be expressed in terms of elements of $`V_{MNS}`$ as
$$S_{sol}=\mathrm{sin}^22\theta _{12}=\frac{4V_{e2}^2V_{e1}^2}{(V_{e2}^2+V_{e1}^2)^2},$$
(13)
$$S_{at}=\mathrm{sin}^22\theta _{23}=\frac{4V_{\mu 3}^2V_{\tau 3}^2}{(V_{\mu 3}^2+V_{\tau 3}^2)^2}$$
(14)
In principle these mixing angles may not be exactly the same as those obtained from two generation analysis of the experimental results.
The renormalisation group equations (RGEs) to one-loop order are:
$`{\displaystyle \frac{dY^u}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}[N_q.Y^u+Y^u.N_u+(N_{H_u})Y_u]`$
$`{\displaystyle \frac{dY^d}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}[N_q.Y^d+Y^d.N_d+(N_{H_d})Y_d]`$
$`{\displaystyle \frac{dY^\nu }{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}[N_l.Y^\nu +Y^\nu .N_\nu +(N_{H_u})Y_\nu ]`$
$`{\displaystyle \frac{dY^e}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}[N_l.Y^e+Y^e.N_e+(N_{H_d})Y_e]`$
$`{\displaystyle \frac{dY_{RR}}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}[N_\nu .Y_{RR}+Y_{RR}.N_\nu +(N_\mathrm{\Sigma })Y_{RR}]`$
(15)
where the wavefunction anomalous dimensions are
$`N_q`$ $`=`$ $`({\displaystyle \frac{8}{3}}g_3^2+{\displaystyle \frac{3}{2}}g_2^2+{\displaystyle \frac{1}{30}}g_1^2)IY^uY_{}^{u}{}_{}{}^{}Y^dY_{}^{d}{}_{}{}^{}`$
$`N_u`$ $`=`$ $`({\displaystyle \frac{8}{3}}g_3^2+{\displaystyle \frac{8}{15}}g_1^2)I2Y_{}^{u}{}_{}{}^{}Y^u`$
$`N_d`$ $`=`$ $`({\displaystyle \frac{8}{3}}g_3^2+{\displaystyle \frac{2}{15}}g_1^2)I2Y_{}^{d}{}_{}{}^{}Y^d`$
$`N_l`$ $`=`$ $`({\displaystyle \frac{3}{2}}g_2^2+{\displaystyle \frac{3}{10}}g_1^2)IY^eY_{}^{e}{}_{}{}^{}Y^\nu Y_{}^{\nu }{}_{}{}^{}`$
$`N_e`$ $`=`$ $`({\displaystyle \frac{6}{5}}g_1^2)I2Y_{}^{e}{}_{}{}^{}Y^e`$
$`N_\nu `$ $`=`$ $`2Y_{}^{\nu }{}_{}{}^{}Y^\nu Y_{RR}^{}Y_{RR}`$
$`N_{H_u}`$ $`=`$ $`({\displaystyle \frac{3}{2}}g_2^2+{\displaystyle \frac{3}{10}}g_1^2)3Tr(Y_{}^{u}{}_{}{}^{}Y^u)Tr(Y_{}^{\nu }{}_{}{}^{}Y^\nu )`$
$`N_{H_d}`$ $`=`$ $`({\displaystyle \frac{3}{2}}g_2^2+{\displaystyle \frac{3}{10}}g_1^2)3Tr(Y_{}^{d}{}_{}{}^{}Y^d)Tr(Y_{}^{e}{}_{}{}^{}Y^e)`$
$`N_\mathrm{\Sigma }`$ $`=`$ $`Tr(Y_{RR}^{}Y_{RR})`$ (16)
where $`t=\mathrm{ln}\mu `$ ($`\mu `$ is the $`\overline{MS}`$ scale) and $`I`$ is the unit matrix. The RGEs for the gauge couplings are
$$\frac{dg_i}{dt}=\frac{1}{16\pi ^2}b_ig_i^3$$
(17)
where $`b_i=(\frac{33}{5},1,3)`$.
The RGEs above are used to run the Yukawa matrices down from high energies (say the unification or string scale) down to the heaviest right-handed neutrino mass, $`M_{RZ}`$, which we assume is equal to the VEV of the $`\mathrm{\Sigma }`$ field. At this mass scale we perform a rotation of the right-handed neutrino fields to the basis in which $`M_{RR}`$ is diagonal, according to Eq.5, so that $`Y^\nu `$ is replaced by $`Y^{}_{}{}^{}\nu =Y^\nu \mathrm{\Omega }_{RR}^{}`$. In this basis one then runs the remaining RGEs (apart from $`Y_{RR}`$) down through the right-handed neutrino thresholds $`\mathrm{diag}(\mathrm{M}_{\mathrm{RZ}},\mathrm{},\mathrm{M}_{\mathrm{R1}})`$, decoupling each right-handed neutrino at its mass threshold from the $`Y^{}_{}{}^{}\nu `$ contributions which appear on the right-hand side of the RGEs. To be explicit, we replace on the right-hand sides of the RGEs,
$$Y_{ip}^{}_{}{}^{}\nu Y_{ip}^{}_{}{}^{}\nu \theta _p$$
(18)
where $`\theta _p=\theta (\mathrm{ln}\mu \mathrm{ln}M_{Rp})`$. For the (diagonal) $`Y_{RR}`$ on the right-hand side we replace it by
$$Y_{RRpp}Y_{RRpp}\theta _p.$$
(19)
In the next section we shall see that for the cases of interest it is not necessary to diagonalise $`M_{RR}`$ in order to implement decoupling of the right-handed neutrinos, and one may remain in the basis defined by the $`U(1)`$ family charges, since decoupling is facilitated by the simple structures of $`M_{RR}`$.
Below the mass of the lightest right-handed neutrino one may use the RGE for $`\kappa `$, the coefficient of the dimension 5 neutrino mass operator,
$$\frac{d\kappa }{dt}=\frac{1}{16\pi ^2}[(6g_2^2+\frac{6}{5}g_1^2)\kappa 6\kappa Tr(Y^uY_{}^{u}{}_{}{}^{})(Y^eY_{}^{e}{}_{}{}^{})\kappa \kappa (Y^eY_{}^{e}{}_{}{}^{})^T]$$
(20)
In solving Eq.20, it is convenient to diagonalise $`Y^u`$ and $`Y^e`$ at the scale of the lightest right-handed neutrino, using Eqs.5, then make the approximation of keeping only the largest third family Yukawa eigenvalues. In the diagonal charged lepton basis, $`\kappa `$ must be transformed to $`\kappa ^{}`$ given by
$$\kappa ^{}=V_{eL}\kappa V_{eL}^{}$$
(21)
then, retaining the third family $`t`$ and $`\tau `$ Yukawa couplings only, the RGEs for the elements of $`\kappa ^{}`$ are given by
$$\frac{d\kappa _{ij}^{}}{dt}=\frac{1}{16\pi ^2}\kappa _{ij}^{}[6g_2^2+\frac{6}{5}g_1^26h_t^2\delta _{i3}h_\tau ^2\delta _{3j}h_\tau ^2]$$
(22)
Following from Eq.22 we see that the elements of $`m_{LL}^{}(M_{R1})=v_2^2\kappa ^{}(M_{R1})`$ at high energy are renormalised down to $`m_{LL}^{}(m_t)=v_2^2\kappa ^{}(m_t)`$ at low energy, ignoring the running of $`v_2`$, according to
$$m_{LL}^{}(m_t)=e^{\frac{6}{5}I_{g_1}}e^{6I_{g_2}}e^{6I_t}\left(\begin{array}{ccc}m_{LL11}^{}(M_{R1})& m_{LL12}^{}(M_{R1})& m_{LL13}^{}(M_{R1})e^{I_\tau }\\ m_{LL21}^{}(M_{R1})& m_{LL22}^{}(M_{R1})& m_{LL23}^{}(M_{R1})e^{I_\tau }\\ m_{LL31}^{}(M_{R1})e^{I_\tau }& m_{LL32}^{}(M_{R1})e^{I_\tau }& m_{LL33}^{}(M_{R1})e^{2I_\tau }\end{array}\right)$$
(23)
where
$$I_f=\frac{1}{16\pi ^2}_{\mathrm{ln}m_t}^{\mathrm{ln}M_{R1}}h_f^2(t)𝑑t,I_{g_i}=\frac{1}{16\pi ^2}_{\mathrm{ln}m_t}^{\mathrm{ln}M_{R1}}g_i^2(t)𝑑t$$
(24)
where $`f=t,\tau `$ and $`i=1,2,3`$. In running down to $`m_t`$ the charged lepton matrix will remain diagonal to good approximation, so that the low energy MNS matrix is simply given by
$$V_{MNS}=V_{\nu L}^{{}_{}{}^{}}.$$
(25)
where $`V_{\nu L}^{}`$ is the matrix which diagonalises $`m_{LL}^{}(m_t)`$,
$$m_{LL}^{{}_{}{}^{}diag}(m_t)=V_{\nu L}^{}m_{LL}^{}(m_t)V_{\nu L}^{{}_{}{}^{}}=\mathrm{diag}(\mathrm{m}_{\nu _1},\mathrm{m}_{\nu _2},\mathrm{m}_{\nu _3}).$$
(26)
The mixing matrix $`V_{MNS}`$ and hence neutrino mixing angles, are the running quantities which can be computed at different energy scales. For example, the running of the neutrino mixing angle relevant to the atmospheric neutrino deficit, $`\theta _{23}`$, can be understood from the evolution equation
$$16\pi ^2\frac{d}{dt}\mathrm{sin}^22\theta _{23}=2\mathrm{sin}^22\theta _{23}(1\mathrm{sin}^2\theta _{23})(h_\tau ^2h_\mu ^2)\frac{m_{LL}^{33^{}}+m_{LL}^{22^{}}}{m_{LL}^{33^{}}m_{LL}^{22^{}}}$$
(27)
This equation (see Babu et al in ref.) describes the evolution of the physical 23 mixing angle assuming that we are in the diagonal charged lepton mass basis.
## 3 Three right-handed neutrinos and SRHND
We now specialise to three right-handed neutrinos and review the conditions for achieving SRHND and the LMA MSW solution . The statement of SRHND is that, of the three right-handed neutrinos, one of them, $`N_{R3}`$, makes the dominant contribution to the 23 block of $`m_{LL}`$. This ensures that the 23 sub-determinant approximately vanishes, and a 23 mass hierarchy therefore naturally results.
We first write the neutrino Yukawa matrix in general as
$$Y_\nu =\left(\begin{array}{ccc}a^{}& a& d\\ b^{}& b& e\\ c^{}& c& f\end{array}\right)$$
(28)
There are now three distinct textures for the heavy Majorana neutrino matrix which maintain the isolation of the dominant right-handed neutrino $`N_{R3}`$, namely the diagonal, democratic and off-diagonal textures introduced previously. We consider each of them in turn. Note that, assuming SRHND, the contribution to the lepton 23 and 13 mixing angles from the neutrino sector are approximately
$$\mathrm{tan}\theta _{23}\frac{e}{f},\mathrm{tan}\theta _{13}\frac{d}{\sqrt{e^2+f^2}},$$
(29)
so that Super-Kamiokande and CHOOZ imply
$$def$$
(30)
The condition on the 12 mixing angle such that it is relevant for the LMA MSW solution is discussed separately for each case below.
### 3.1 Diagonal Texture
$$M_{RR}=\left(\begin{array}{ccc}X^{}& 0& 0\\ 0& X& 0\\ 0& 0& Y\end{array}\right)$$
(31)
We may invert the heavy Majorana matrix and construct the light Majorana matrix using the see-saw mechanism,
$$m_{LL}=\left(\begin{array}{ccc}\frac{d^2}{Y}+\frac{a^2}{X}+\frac{a^2}{X^{}}& \frac{de}{Y}+\frac{ab}{X}+\frac{a^{}b^{}}{X^{}}& \frac{df}{Y}+\frac{ac}{X}+\frac{a^{}c^{}}{X^{}}\\ .& \frac{e^2}{Y}+\frac{b^2}{X}+\frac{b^2}{X^{}}& \frac{ef}{Y}+\frac{bc}{X}+\frac{b^{}c^{}}{X^{}}\\ .& .& \frac{f^2}{Y}+\frac{c^2}{X}+\frac{c^2}{X^{}}\end{array}\right)v_2^2$$
(32)
The SRHND condition is
$$\frac{e^2}{Y}\frac{ef}{Y}\frac{f^2}{Y}\frac{xy}{X},\frac{x^{}y^{}}{X^{}}$$
(33)
where $`x,ya,b,c`$ and $`x^{},y^{}a^{},b^{},c^{}`$. <sup>2</sup><sup>2</sup>2The beauty of SRHND is that it automatically implies $`m_{\nu _2}m_{\nu _3}`$ due to the approximately vanishing 23 subdeterminant, without the need for appeal to cancellations. To understand this simply drop the $`1/X`$ terms and observe that the 23 subdeterminant vanishes which implies a massless eigenvalue which is a rather extreme case of a hierarchy! The 12 mixing angle determines whether we have the LMA MSW or SMA MSW solution, and this depends on the relative magnitude of the sub-dominant entries of $`m_{LL}`$, as discussed in . The condition for LMA MSW is
$$max(\frac{ab}{X},\frac{ac}{X},\frac{a^{}b^{}}{X^{}},\frac{a^{}c^{}}{X^{}})max(\frac{b^2}{X},\frac{bc}{X},\frac{c^2}{X},\frac{b^2}{X^{}},\frac{b^{}c^{}}{X^{}},\frac{c^2}{X^{}})$$
(34)
### 3.2 Off-Diagonal Texture
This is defined by:
$$M_{RR}=\left(\begin{array}{ccc}0& X& 0\\ X& 0& 0\\ 0& 0& Y\end{array}\right)$$
(35)
The off-diagonal case is qualitatively different from the other two cases and gives
$$m_{LL}=\left(\begin{array}{ccc}\frac{d^2}{Y}+\frac{2aa^{}}{X}& \frac{de}{Y}+\frac{a^{}b}{X}+\frac{ab^{}}{X}& \frac{df}{Y}+\frac{a^{}c}{X}+\frac{ac^{}}{X}\\ .& \frac{e^2}{Y}+\frac{2bb^{}}{X}& \frac{ef}{Y}+\frac{b^{}c}{X}+\frac{bc^{}}{X}\\ .& .& \frac{f^2}{Y}+\frac{2cc^{}}{X}\end{array}\right)v_2^2$$
(36)
SRHND is now defined by the conditions
$$\frac{e^2}{Y}\frac{ef}{Y}\frac{f^2}{Y}\frac{xx^{}}{X}$$
(37)
where $`xa,b,c`$ and $`x^{}a^{},b^{},c^{}`$, The LMA MSW solution condition is :
$$max(\frac{a^{}b}{X},\frac{ab^{}}{X},\frac{a^{}c}{X},\frac{ac^{}}{X})max(\frac{bb^{}}{X},\frac{b^{}c}{X},\frac{bc^{}}{X},\frac{cc^{}}{X})$$
(38)
### 3.3 Democratic Texture
The democratic case (assuming the Majorana masses in the upper block are of the same order but are not exactly equal) is defined by:
$$M_{RR}=\left(\begin{array}{ccc}X& X& 0\\ X& X& 0\\ 0& 0& Y\end{array}\right)$$
(39)
The order of magnitude of $`m_{LL}`$ is:
$$m_{LL}=\left(\begin{array}{ccc}\frac{d^2}{Y}+O(\frac{a^2}{X})+O(\frac{a^2}{X})& \frac{de}{Y}+O(\frac{ab}{X})+O(\frac{a^{}b^{}}{X})& \frac{df}{Y}+O(\frac{ac}{X})+O(\frac{a^{}c^{}}{X})\\ .& \frac{e^2}{Y}+O(\frac{b^2}{X})+O(\frac{b^2}{X})& \frac{ef}{Y}+O(\frac{bc}{X})+O(\frac{b^{}c^{}}{X})\\ .& .& \frac{f^2}{Y}+O(\frac{c^2}{X})+O(\frac{c^2}{X})\end{array}\right)v_2^2$$
(40)
In this case the SRHND conditions are:
$$\frac{e^2}{Y}\frac{ef}{Y}\frac{f^2}{Y}\frac{xy}{X}\frac{x^{}y^{}}{X}$$
(41)
where $`x,ya,b,c`$ and $`x^{},y^{}a^{},b^{},c^{}`$. The LMA MSW condition for a large 12 angle is
$$max(\frac{ab}{X},\frac{ac}{X},\frac{a^{}b^{}}{X},\frac{a^{}c^{}}{X})max(\frac{b^2}{X},\frac{bc}{X},\frac{c^2}{X},\frac{b^2}{X},\frac{b^{}c^{}}{X},\frac{c^2}{X})$$
(42)
### 3.4 $`U(1)`$ Family Symmetry
Introducing a $`U(1)`$ family symmetry , , , provides a convenient way to organise the hierarchies within the various Yukawa matrices. For definiteness we shall focus on a particular class of model based on a single pseudo-anomalous $`U(1)`$ gauged family symmetry . We assume that the $`U(1)`$ is broken by the equal VEVs of two singlets $`\theta ,\overline{\theta }`$ which have vector-like charges $`\pm 1`$ . The $`U(1)`$ breaking scale is set by $`<\theta >=<\overline{\theta }>`$ where the VEVs arise from a Green-Schwartz mechanism with computable Fayet-Illiopoulos $`D`$-term which determines these VEVs to be one or two orders of magnitude below $`M_U`$. Additional exotic vector matter with mass $`M_V`$ allows the Wolfenstein parameter to be generated by the ratio
$$\frac{<\theta >}{M_V}=\frac{<\overline{\theta }>}{M_V}=\lambda 0.22$$
(43)
The idea is that at tree-level the $`U(1)`$ family symmetry only permits third family Yukawa couplings (e.g. the top quark Yukawa coupling). Smaller Yukawa couplings are generated effectively from higher dimension non-renormalisable operators corresponding to insertions of $`\theta `$ and $`\overline{\theta }`$ fields and hence to powers of the expansion parameter in Eq.43, which we have identified with the Wolfenstein parameter. The number of powers of the expansion parameter is controlled by the $`U(1)`$ charge of the particular operator. The fields relevant to neutrino masses $`L_i`$, $`N_p^c`$, $`H_u`$, $`\mathrm{\Sigma }`$ are assigned $`U(1)`$ charges $`l_i`$, $`n_p`$, $`h_u=0`$, $`\sigma `$. From Eqs.43, the neutrino Yukawa couplings and Majorana mass terms may then be expanded in powers of the Wolfenstein parameter,
$$M_{RR}\left(\begin{array}{ccc}\lambda ^{|2n_1+\sigma |}& \lambda ^{|n_1+n_2+\sigma |}& \lambda ^{|n_1+n_3+\sigma |}\\ .& \lambda ^{|2n_2+\sigma |}& \lambda ^{|n_2+n_3+\sigma |}\\ .& .& \lambda ^{|2n_3+\sigma |}\end{array}\right)<\mathrm{\Sigma }>$$
(44)
The conditions which ensure that the third dominant neutrino is isolated require that the elements $`\lambda ^{|n_1+n_3+\sigma |}`$, $`\lambda ^{|n_2+n_3+\sigma |}`$ be sufficiently small. The diagonal, off-diagonal and democratic textures then emerge as approximate cases . The neutrino Yukawa matrix is explicitly
$$Y_\nu \left(\begin{array}{ccc}\lambda ^{|l_1+n_1|}& \lambda ^{|l_1+n_2|}& \lambda ^{|l_1+n_3|}\\ \lambda ^{|l_2+n_1|}& \lambda ^{|l_2+n_2|}& \lambda ^{|l_2+n_3|}\\ \lambda ^{|l_3+n_1|}& \lambda ^{|l_3+n_2|}& \lambda ^{|l_3+n_3|}\end{array}\right)$$
(45)
which may be compared to the notation in Eq.28. The requirement of large 23 mixing and small 13 mixing expressed in Eq.30 becomes
$$|n_3+l_2|=|n_3+l_3|,|n_3+l_1||n_3+l_3|=1or2$$
(46)
The remaining conditions for the $`U(1)`$ charges depend on the specific heavy Majorana texture under consideration .
The charged lepton Yukawa matrix is given by
$$Y_e\left(\begin{array}{ccc}\lambda ^{|l_1+e_1|}& \lambda ^{|l_1+e_2|}& \lambda ^{|l_1+e_3|}\\ \lambda ^{|l_2+e_1|}& \lambda ^{|l_2+e_2|}& \lambda ^{|l_2+e_3|}\\ \lambda ^{|l_3+e_1|}& \lambda ^{|l_3+e_2|}& \lambda ^{|l_3+e_3|}\end{array}\right)$$
(47)
where $`e_i`$ are the $`U(1)`$ charges of the charged lepton singlet fields.
For the quarks we shall assume a common form for the textures of $`Y^u`$ and $`Y^d`$
$$Y^u\left(\begin{array}{ccc}\lambda ^8& \lambda ^5& \lambda ^3\\ \lambda ^7& \lambda ^4& \lambda ^2\\ \lambda ^5& \lambda ^2& 1\end{array}\right),Y^d\left(\begin{array}{ccc}\lambda ^4& \lambda ^3& \lambda ^3\\ \lambda ^3& \lambda ^2& \lambda ^2\\ \lambda & 1& 1\end{array}\right)\lambda ^n$$
(48)
## 4 Renormalisation Group Analysis of SRHND
In we tabulated the simplest charges which satisfy all the conditions given above, and so provide a natural account of the atmospheric and solar neutrinos via the LMA MSW effect. In Table 1 we consider one example of each of the three cases, namely case A (Diagonal $`M_{RR}`$), case B (Off-diagonal $`M_{RR}`$) and case C (Democratic $`M_{RR}`$). The $`U(1)`$ charges along with corresponding $`M_{RR}`$, $`Y^\nu `$, obtained from Eqs.44, 45, and other relevant parameters are outlined for each case, including the charged lepton charges, are also shown. Note that the zeroes in $`M_{RR}`$ appear after small angle rotations on the right-handed neutrino fields, which will not affect the perturbative expansion in powers of $`\lambda `$ in $`Y^\nu `$. The order unity coefficients $`b_{ij}`$, which are always present in $`U(1)`$ models, are defined in this basis.
In case A, before rotating to the diagonal charged lepton mass basis, we may easily estimate the order of the entries in $`m_{LL}`$ in Eq.32 from the matrices $`Y^\nu `$ and $`M_{RR}`$ in Table 1, and hence verify the SRHND conditions Eq.33, and the LMA MSW condition Eq.34. Using Eq.32 we find
$$m_{LL}\left(\begin{array}{ccc}\lambda +\lambda +\lambda & 1+\lambda +\lambda ^2& 1+\lambda +\lambda ^2\\ .& \frac{1}{\lambda }+\lambda +\lambda ^3& \frac{1}{\lambda }+\lambda +\lambda ^3\\ .& .& \frac{1}{\lambda }+\lambda +\lambda ^3\end{array}\right)\frac{v_2^2}{<\mathrm{\Sigma }>}$$
(49)
where the first entry in each element corresponds to the $`1/Y`$ contributions from $`N_{3R}`$, which clearly dominate the 23 block by a factor of $`\lambda ^2`$, and the 12 and 13 elements by a factor of $`\lambda `$. The 13 element is smaller than the elements in the 23 block by a factor of $`\lambda `$ leading to a 13 CHOOZ angle of this order. The subdominant entries in the 12,13,22,23 elements are the same order, leading to a large 12 angle suitable for LMA MSW.
Similarly in case B, we can show that the SRHND condition Eq.37 and LMA MSW condition Eq.38 are satisfied. We estimate the order of the entries in $`m_{LL}`$ in Eq.36 from the matrices $`Y^\nu `$ and $`M_{RR}`$ in Table 1 as,
$$m_{LL}\left(\begin{array}{ccc}\lambda ^5+2\lambda ^5& \lambda ^3+\lambda ^5+\lambda ^3& \lambda ^3+\lambda ^5+\lambda ^3\\ .& \lambda +2\lambda ^3& \lambda +\lambda ^3+\lambda ^3\\ .& .& \lambda +2\lambda ^3\end{array}\right)\frac{v_2^2}{<\mathrm{\Sigma }>}$$
(50)
where the first entry in each element corresponds to the $`1/Y`$ contributions coming from the right-handed neutrino $`N_{3R}`$, and clearly dominates the 23 block by a factor of $`\lambda ^2`$. In this case it does not dominate the other elements outside the 23 block. The 13 element from $`N_{3R}`$ is suppressed relative to the 23 block elements by a factor of $`\lambda ^2`$, leading to a CHOOZ angle of this order. The subdominant entries in the 12,13,22,23 elements are again the same order, leading to a large 12 angle suitable for LMA MSW.
Finally in case C, we verify that the SRHND condition Eq.41 and LMA MSW condition Eq.42 are satisfied, by constructing the entries in $`m_{LL}`$ in Eq.40 from the matrices $`Y^\nu `$ and $`M_{RR}`$ in Table 1,
$$m_{LL}\left(\begin{array}{ccc}\lambda +\lambda +\lambda & 1+\lambda +\lambda & 1+\lambda +\lambda \\ .& \frac{1}{\lambda }+\lambda +\lambda & \frac{1}{\lambda }+\lambda +\lambda \\ .& .& \frac{1}{\lambda }+\lambda +\lambda \end{array}\right)\frac{v_2^2}{<\mathrm{\Sigma }>}$$
(51)
where the first entry in each element corresponds to the $`1/Y`$ contributions coming from the dominant right-handed neutrino, which dominates in the 23 block by a factor of $`\lambda ^2`$ and in the 12 and 13 elements by a factor of $`\lambda `$. The 13 element is smaller than the 23 elements by a factor of $`\lambda `$ leading to a CHOOZ angle of this order. The subdominant 12,13,22,23 elements are all the same order leading to LMA MSW.
We now turn to a numerical treatment of these three cases. In integrating the RGEs down to low energies the heavy thresholds for the diagonal texture in Eq.31 are dealt with exactly as described in the previous section, except that no rotation is required to get to the diagonal $`M_{RR}`$ basis since we already begin with that form. For the democratic texture in Eq.39 there are essentially only two mass thresholds to consider $`X`$ and $`Y`$ since $`N_{R1}`$ and $`N_{R2}`$ are approximately degenerate with mass $`X`$ and we can ignore any small mass difference between them to leading order. Similarly for the off-diagonal texture in Eq.35 we also have only two mass thresholds to consider $`X`$ and $`Y`$ since $`N_{R1}`$ and $`N_{R2}`$ are now exactly degenerate with mass $`X`$. Thus for both democratic and off-diagonal textures we can replace on the right-hand sides of the RGEs by
$$Y_{i1}^\nu Y_{i1}^\nu \theta _1,Y_{i2}^\nu Y_{i2}^\nu \theta _2,Y_{i3}^\nu Y_{i3}^\nu \theta _3$$
(52)
where $`\theta _{1,2}=\theta (\mathrm{ln}\mu \mathrm{ln}X)`$, $`\theta _3=\theta (\mathrm{ln}\mu \mathrm{ln}Y)`$. We replace $`Y_{RR}`$ on the right-hand side of the RGEs by
$$Y_{RRij}Y_{RRij}\theta _i\theta _j.$$
(53)
Tables 2-4 give the numerical values of quantities at three different energy scales: the GUT scale $`M_U=2.0\times 10^{16}`$ GeV, the lightest right-handed neutrino mass $`M_{R_1}10^{14}`$ GeV, and the low energy scale $`m_t=175`$ GeV. To begin with the left-handed Majorana neutrino mass matrix <sup>3</sup><sup>3</sup>3Strictly speaking the see-saw mechanism does not operate at energy scales higher than the right-handed neutrino mass scales, but such results do provide a meaningful measure of the effects of radiative corrections from the GUT scale down to low energies. $`m_{LL}(M_X)`$ arising from the sum of dominant and subdominant contributions from $`N_{R3}`$ and $`N_{R1}`$,$`N_{R2}`$ respectively, are numerically computed at the energy scale $`\mu =M_U=2.0\times 10^{16}`$GeV, as seen in tables 2,3,4 for the three cases (A,B,C) for large $`\mathrm{tan}\beta `$. The corresponding MNS mixing matrices which in turn give $`S_{sol}`$ and $`S_{at}`$ along with neutrino masses ($`m_{\nu 1},m_{\nu 2},m_{\nu 3}`$) are also estimated. Next, as already discussed in section 2, we calculate the radiative corrections to $`m_{LL}^{}`$ and $`V_{MNS}`$ through the running of the Yukawa couplings $`Y^u`$, $`Y^d`$, $`Y^\nu `$, $`Y^e`$, $`Y_{RR}`$, and the gauge couplings $`g_1`$, $`g_2`$, $`g_3`$ from high energy scale $`M_U`$ through the heavy neutrino threshold region at successive steps down to the lightest right-handed neutrino mass $`(M_{R_1})`$. The corresponding $`m_{LL}^{}(M_{R_1})`$ (in the diagonal charged lepton basis) and $`V_{MNS}(M_{R_1})`$ are again estimated as shown in tables 2,3 and 4 for three cases. Further, the radiative corrections to $`m_{LL}^{}`$ and $`V_{MNS}`$ from the lightest right-handed neutrino down to low energies ( say top-quark mass scale $`m_t`$) are taken in the usual way through the running of the coefficient of the dimension 5 operator $`\kappa ^{}`$ in the diagonal charged lepton basis. These low energy results in tables 2,3,4 can be compared with observational data, and $`S_{sol},S_{at}`$ defined in Eqs.13,14 and the neutrino masses given as the diagonal elements of $`m_{LL}^{diag}`$ are seen to be in good agreement with atmospheric data and the LMA MSW solutions. The CHOOZ constraint is also satisfied in all cases, with the 13 element of $`V_{MNS}`$ being larger for cases A and C than for case B, as expected from the analaytical estimates above. The charged lepton mass ratios are $`m_e/m_\mu 0.0066`$ and $`m_\mu /m_\tau 0.07`$ at low energy scale, compared to the experimental values 0.005 and 0.06, respectively.
It is instructive to examine the RG evolution of the neutrino masses and mixing angles. The numerical results in Tables 1-4 show that neutrino masses $`m_{\nu _2}`$ and $`m_{\nu _3}`$ are decreasing from high energy scale $`M_U`$ to low energy scale $`m_t`$ by about $`40\%`$, but the ratios $`m_{\nu _2}/m_{\nu _3}`$ increase by about ($`24\%`$, $`17\%`$, $`55\%`$) corresponding to low energy ratios of (0.05, 0.10, 0.08) for cases (A,B,C), respectively. The atmospheric mixing quantity $`S_{at}`$ increases by about $`10\%`$ in the three cases, and approaches maximal mixing in all cases. This can be understood from Eq.27, which shows that in the diagonal charged lepton basis, the sign of the RGE evolution is determined by the sign of $`m_{LL}^{33^{}}m_{LL}^{22^{}}`$, which in the examples in Tables 2,3,4 is always positive. This means that the overall sign of the RGE is negative, which implies an increasing $`S_{at}`$ as the energy scale is reduced. Note that the largest contribution to the atmospheric mixing angle arises from the charged lepton sector, for this choice of parameters, as is clear from examining $`V_{eL}`$ in Tables 2,3,4. The solar mixing quantity $`S_{sol}`$ also increases, but the increases of $`15\%`$ are mild in comparison to $`S_{at}`$. In cases A,C $`S_{sol}`$ approaches maximal mixing, while in case B its low energy value is about 0.71, which are well near the best fit $`sin^22\theta _{12}0.76`$.
In Figures 1-3 we display the RG running of the mixing angles $`S_{sol}`$ and $`S_{at}`$ and the ratio of neutrino masses $`m_{\nu 2}/m_{\nu 3}`$ as a function of $`t=\mathrm{ln}\mu `$, for the three cases A,B,C. We prefer to show the variation in the neutrino mass ratio, rather than the absolute values of the two neutrino mass eigenvalues since they will depend to some extent on the running vacuum expectation value, although the effect on the present analysis is not large since we assume high $`\mathrm{tan}\beta `$. We also do not consider the ratio $`m_{\nu 1}/m_{\nu 2}`$ since this is not experimentally measurable in the hierarchical case.
In Figure 1(a) we show the variations of $`S_{sol}`$ and $`S_{at}`$ with energy scale for case A. As already stated, the atmospheric mixing angle runs more rapidly than the solar mixing angle, and although $`S_{at}`$ starts out smaller then $`S_{sol}`$ at $`M_U`$, it quickly grows larger. Note the effect of the heavy right-handed neutrino mass thresholds which change the slope of the curves, which grow steeply over the heavy threshold region. In Figure 1(b) we give the variation of the neutrino mass ratio $`m_{\nu 2}/m_{\nu 3}`$ with $`t=\mathrm{ln}\mu `$ for case A. Here the effects of the three heavy right-handed neutrino thresholds is clearly seen, with again a steep rise in this mass ratio over the heavy threshold region.
In Figure 2(a) we show the variations of $`S_{sol}`$ and $`S_{at}`$ with energy scale for case B. Here the atmospheric angle starts out larger than the solar angle, but still grows more rapidly. In Figure 2(b) we give the variation of the neutrino mass ratio $`m_{\nu 2}/m_{\nu 3}`$ with $`t=\mathrm{ln}\mu `$ for case B. The qualitative shape of this curve is similar to Figure 1(b), but there are only two heavy neutrino mass thresholds in this case, and also the mass ratio is larger throughout.
In Figure 3(a) we show the variations of $`S_{sol}`$ and $`S_{at}`$ with energy scale for case B. Here the solar angle starts out much larger than the atmospheric angle, and as before the atmospheric angle grows more rapidly and approaches the solar angle. In Figure 3(b) we give the variation of the neutrino mass ratio $`m_{\nu 2}/m_{\nu 3}`$ with $`t=\mathrm{ln}\mu `$ for case B. Because of the choice of $`c_{12}`$ in Table 1, the two heavy right-handed neutrino mass thresholds are very close together in this case.
Fig.1(a) Case A (Diagonal $`M_{RR}`$): Variation of $`S_{sol}`$ and $`S_{at}`$ with energy scale $`t=\mathrm{ln}\mu `$, which are represented by solid-line and dotted-line respectively.
Fig.1(b) Case A (Diagonal $`M_{RR}`$): Variation of the neutrino mass ratio $`m_{\nu 2}/m_{\nu 3}`$ with $`t=\mathrm{ln}\mu `$.
Fig.2(a) Case B (Off-diagonal $`M_{RR}`$): Variation of $`S_{sol}`$ and $`S_{at}`$ with $`t=\mathrm{ln}\mu `$, which are represented by solid-line and dotted-line respectively.
Fig.2(b) Case B (Off-diagonal $`M_{RR}`$): Variation of the neutrino mass ratio $`m_{\nu 2}/m_{\nu 3}`$ with $`t=\mathrm{ln}\mu `$.
Fig.3(a) Case C (Democratic $`M_{RR}`$): Variation of $`S_{sol}`$ and $`S_{at}`$ with energy scale $`t=\mathrm{ln}\mu `$, which are represented by solid-line and dotted-line respectively.
Fig.3(b) Case C (Democratic $`M_{RR}`$): Variation of the neutrino mass ratio $`m_{\nu 2}/m_{\nu 3}`$ with $`t=\mathrm{ln}\mu `$.
It is important to emphasise that more than $`50\%`$ of the radiative corrections to all the physical quantities $`m_{\nu _2}/m_{\nu _3}`$, $`S_{at}`$, $`S_{sol}`$ arises from the RG evolution over the range from the GUT scale $`M_U=2.0\times 10^{16}`$GeV through the heavy right-handed neutrino threshold region down to the lightest right-handed neutrino mass scale $`M_{R_1}10^{14}`$GeV. Therefore even though this range covers only two orders of magnitude in energy, it is just as important as the RG evolution effects from $`M_{R_1}10^{14}`$GeV down to low energies which covers over 12 orders of magnitude in energy, and is the region commonly considered in the literature.
## 5 Conclusion
We have studied the effects of radiative corrections on neutrino masses and mixing angles, in evolving the theory defined at the GUT scale down to low energies. Our approach is to run down all the Yukawa matrices from the GUT scale down through the heavy right-handed neutrino mass thesholds to low energy, replacing the neutrino Yukawa matrices by the dimension 5 neutrino mass operator at the lowest right-handed neutrino mass threshold. We have found that in the realistic cases considered, the atmospheric and solar neutrino mixing parameters receive radiative corrections of order $`10\%`$ and $`5\%`$, respectively, while the ratios of the neutrino masses change by up about $`50\%`$ in going from $`M_U`$ to low energy. Importantly, more than $`50\%`$ of the overall radiative corrections arises from running through the threshold region, even though it only accounts for two orders of magnitude in energy. Thus many of the existing analyses in the literature which ignore this threshold region could endanger a significant error.
We have considered a realistic class of models known as SRHND , in which the contribution to the 23 block of the light effective Majorana matrix $`m_{LL}`$ is dominated by a single right-handed neutrino. The small neutrino mass hierarchy $`m_{\nu _2}/m_{\nu _3}0.1`$ originates from the smallness of the 23 subdeterminant of $`m_{LL}`$ (which is zero in the limit that a single right-handed neutrino is the only contribution). We have focussed on cases with two large mixing angles $`\theta _{23}`$ and $`\theta _{12}`$, with the CHOOZ angle $`\theta _{13}`$ being small - the so-called bi-maximal mixing scenario. The couplings of the dominant right-hand neutrino controls the 23 and 13 mixing angles, and the subdominant right-handed neutrino couplings control the 12 angle. A $`U(1)`$ family symmetry is used to generate a controlled expansion of all the Yukawa couplings in powers of the Wolfenstein parameter $`\lambda `$, and with a suitable choice of $`U(1)`$ charges SRHND may be achieved with subdominant right-handed neutrino contributions being of the correct order of magnitude to generate the desired neutrino mixing angles and mass hierarchy.
The $`U(1)`$ charges also control the charged lepton Yukawa matrix, resulting in large contributions to the physical lepton mixing angles from the charged lepton sector. In general one would expect the two contributions to the 23 mixing angle (from the charged leptons and the neutrinos) not to cancel exactly due to order unity coefficients in the Yukawa matrices which are not predicted, and our numerical results include these effects. Similarly it would be surprising if the physical mixing angles turned out to be exactly maximal at the GUT scale. Therefore we have considered examples in which the 23 and 12 angles start out large but not maximal, and we have also assumed large $`\mathrm{tan}\beta `$, where the $`\tau `$ Yukawa coupling is large and the effect of radiative corrections is maximised. Although the effect of radiative corrections on the mixing angles is always $`10\%`$, showing that the models are quite stable, we have shown that the effects may play an important role in driving the initially large (but not maximal) mixing angle towards its maximal value. This is true both of the atmospheric mixing angle and the solar mixing angle, although in the latter case the effects are milder, which helps to explain why the atmospheric angle is larger than the solar angle. In principle, a different choice of parameters could have caused the neutrino angles to have grown smaller at low energies. However it is significant that for the cases considered both mixing angles become magnified showing that the low energy approximate bi-maximal scenario could partly result from radiative corrections in SRHND models.
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# 1 Introduction
## 1 Introduction
The aim of our study here is to continue the search for better understanding on the physics in the infrared of non-Abelian asymptotic free gauge theories (e.g., QCD). The questions of central interests are: i) the mechanism of confinement; ii) the mechanism of flavor (chiral) symmetry breaking; and the relation between the two; iii) the existence of other phases (CFT, oblique confinement, etc.), and iv) the $`\theta `$ dependence, CP properties, etc. There has been substantial progress recently in this field, coming from the study of supersymmetric models. For instance, an interesting hint came from the study of $`N=2`$ supersymmetric QCD with $`SU(2)`$ gauge group, in which supersymmetry is softly broken to $`N=1`$: in some of the vacua, condensation of magnetic monopoles leads to confinement and flavor symmetry breaking simultaneously.
We study here more general classes of $`N=2`$ supersymmetric $`SU(n_c)`$, $`USp(2n_c)`$ and $`SO(n_c)`$ gauge theories with $`n_f`$ quarks, with a small adjoint mass breaking supersymmetry to $`N=1`$. The generalizations turn out to be highly nontrivial, and the resulting variety of dynamical possibilities much richer than might be expected from the $`SU(2)`$ cases studied by Seiberg and Witten, or from the pure ($`n_f=0`$) $`SU(n_c)`$ theory discussed by Douglas and Shenker. For the exact Seiberg-Witten curves for $`N=2`$ supersymmetric $`SU(n_c)`$, $`USp(2n_c)`$ and $`SO(n_c)`$ gauge groups see .
Before presenting the model and discuss our main results, let us mention a recent work in which the interrelation between confinement and chiral symmetry breaking in the standard (non-supersymmetric) QCD with $`SU(2)`$ gauge group, was studied , by using the Faddeev-Niemi gauge field decomposition. It is argued that in the ground state of ($`SU(2)`$) QCD there are two complementary, competing configurations which are important: one - meronlike configurations or regularized Wu-Yang monopoles - responsible for confinement but in itself has nothing to do with chiral symmetry breaking, while the instantonlike configurations are fundamental for the chiral symmetry breaking but are themselves unrelated to confinement.
### 1.1 The model
The models discussed here are described by the Lagrangian,
$$L=L^{(N=2)}(W,\mathrm{\Phi },\stackrel{~}{Q}_i,Q^i)+m_i\stackrel{~}{Q}_iQ^i|_F+\mu \mathrm{\Phi }^2|_F,$$
(1.1)
with $`m_i,\mu \mathrm{\Lambda }`$, where the first term of the standard $`N=2`$ supersymmetric Lagrangian with massless hypermultiplets (quarks) in the fundamental representation of the gauge group, $`m_i`$ is the bare quark mass of the $`i`$-th flavor, and the adjoint mass $`\mu `$ breaks supersymmetry to $`N=1`$. The models have no flat directions so that there are finite number of isolated $`N=1`$ vacua, keeping track of which provides us with a quite nontrivial check of our analyses. Also, only those theories are considered in which the interactions become strong in the infrared. The global symmetry of the models are:
$`SU(n_c):G_F`$ $`=`$ $`U(n_f)(m_im,\text{or}\mathrm{\hspace{0.17em}\hspace{0.17em}0});`$
$`USp(2n_c):G_F`$ $`=`$ $`SO(2n_f)(m_i0);`$
$`SO(n_c):G_F`$ $`=`$ $`USp(2n_f)(m_i0).`$ (1.2)
Also, the global discrete symmetries such as $`Z_{2n_cn_f}`$ in $`SU(n_c)`$ play important roles.
### 1.2 The results
The most striking results of our analysis, summarized in Table 1 and Table 2 for $`SU(n_c)`$ and $`USp(2n_c)`$ theories, are the following.
The ’t Hooft - Mandelstam picture of confinement, caused by the condensation of monopoles in the maximal Abelian subgroup $`U(1)^k`$, ($`k=`$ Rank of the gauge group), is in fact realized only in some of the vacua. In a more “typical” vacuum of $`SU(n_c)`$ gauge theory, the effective, infrared degrees of freedom involve are a set of dual quarks, interacting with low-energy effective non-Abelian $`SU(r)`$ gauge fields. The condensation of these magnetic quarks as well as of certain Abelian monopoles also present in the theory, upon $`\mu `$ perturbation, lead to confinement and dynamical symmetry breaking. The semi-classical monopoles may be interpreted as baryonic composites made of these magnetic quarks and monopoles, which break up into their constituents before they become massless, as we move from the semiclassical region of the space of $`N=2`$ vacua (parametrized by a set of gauge invariant VEVS) towards the relevant singularity.
The second most interesting result is that the special vacua in $`SU(n_c)`$ theory as well as the entire first group of vacua in $`USp(2n_c)`$ or $`SO(n_c)`$ theory correspond to various nontrivial infrared fixed points (SCFT). The low-energy effective degrees of freedom in general contain relatively nonlocal states and there is no local effective Lagrangian description of these theories, though the symmetry breaking pattern can be found from the analysis at large adjoint mass $`\mu `$.
Finally, in both type of gauge theories, for large number of flavors, there is a second group of vacua in free-magnetic phase, with no confinement and no spontaneous flavor symmetry breaking. In these vacua the low energy degrees of freedom are weakly interacting non-Abelian dual quarks and gauge particles, as well as some monopoles of products of $`U(1)`$ groups. In $`SO(n_c)`$ theories, the situation is qualitatively similar to $`USp(2n_c)`$ cases; however, the effective gauge group and the unbroken global group in the vacua in free-magnetic phase are given by $`SO(\stackrel{~}{n}_c)=SO(2n_fn_c+4)`$ and $`USp(2n_f)`$, respectively, in these theories.
## 2 Analyses
Our analyses leading to these results consist of several independent steps :
i) Semi-classical analysis, yielding the number of the vacua, $`𝒩`$;
ii) Determination of dynamical symmetry breaking pattern at $`\mu \mathrm{\Lambda }`$ fixed and $`m_i0`$;
iii) Check of the correct decoupling of the adjoint fields in the limit, $`\mu \mathrm{}`$, with $`m_i`$ and $`N=1`$ scale factor $`\mathrm{\Lambda }_1\mu ^{\frac{n_c}{3n_cn_f}}\mathrm{\Lambda }^{\frac{2n_cn_f}{3n_cn_f}}`$ (for $`SU(n_c)`$ for instance) fixed;
iv) Study of $`N=1`$ vacua at $`m_i,\mu \mathrm{\Lambda }`$, from the Seiberg-Witten curves, through the mass perturbation around CFT singularities ;
v) Study of $`N=1`$ vacua at $`m_i,\mu \mathrm{\Lambda }`$, by using the low-energy effective action , leading to a clear microscopic picture of infrared physics;
vi) Numerical study of the maximal singularities of the Seiberg-Witten curves in the cases of rank $`2`$ gauge groups ($`SU(3)`$, $`USp(4)`$, $`SO(4)`$, $`SO(5)`$);
vii) Study of monodromy around the singularities of the curves;
viii) Semi-classical analysis of the monopole flavor multiplet structure à la Jackiw-Rebbi.
These independent analyses lead, in a remarkably subtle way, to mutually consistent answers as regards the number of the vacua and physical properties of each of them. We shall discuss the two among the most important aspects of our analysis, ii) and v), below, but let us mention here the results of the semiclassical study, i). As the models have no flat directions, one can simply minimize the classical potential and get, after taking into account the appropriate Witten’s index in case part of the gauge group remains unbroken by classical VEVS, the number of vacua. For $`SU(n_c)`$ theory with $`n_f`$ flavors, there are
$$𝒩=\underset{r=0}{\overset{\text{min}\{n_f,n_c1\}}{}}(n_cr)\left(\begin{array}{c}n_f\\ r\end{array}\right)$$
(2.1)
classical solutions. Note that when $`n_f`$ is equal to or less than $`n_c`$ the sum over $`r`$ is done readily, and Eq. (2.1) is equal to
$$𝒩_1=(2n_cn_f)\mathrm{\hspace{0.17em}2}^{n_f1},(n_fn_c).$$
(2.2)
Similarly, we find for $`USp(2n_c)`$
$$𝒩=\underset{r=0}{\overset{min\{n_c,n_f\}}{}}(n_cr+1)\left(\begin{array}{c}n_f\\ r\end{array}\right).$$
(2.3)
vacua, which reduces for smaller values of $`n_f`$ to a closed expression,
$$𝒩=(2n_c+2n_f)\mathrm{\hspace{0.17em}2}^{n_f1},(n_fn_c).$$
(2.4)
For $`SO(n_c)`$ theories the result is
$$𝒩=\underset{r=0}{\overset{min\{[n_c/2],n_f\}}{}}w(n_c2r)\left(\begin{array}{c}n_f\\ r\end{array}\right)+\left(\genfrac{}{}{0pt}{}{n_f}{n_c/2}\right),$$
(2.5)
where
$$w(N)=N2,N5,$$
(2.6)
and $`w(N)=4,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}1},`$ for $`N=4,\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}0},`$ respectively, and the last term is present only for $`2n_fn_c,n_c=\text{even}.`$ The formulas Eq.(2.1)-Eq.(2.4) correctly reduce to the well-known result
$$𝒩=n_f+2,$$
(2.7)
in the case of the $`SU(2)`$ theory. Note that the generalization is nontrivial: e.g., $`𝒩n_f+n_c`$ in $`SU(n_c)`$ theory! The complexity of the formulae Eq.(2.1) - Eq.(2.5) as compared to Eq.(2.7) signals indeed the presence of a rich variety of dynamical possibilities in general $`SU(n_c)`$, $`USp(2n_c),`$ and $`SO(n_c)`$ theories, some of which might well be important in the understanding of the standard QCD.
## 3 Determination of dynamical symmetry breaking pattern at large $`\mu `$
At large $`\mu `$ ($`\mu \mathrm{\Lambda }`$), the effective superpotential can be read off from the bare Lagrangian by integrating out the heavy, adjoint fields and by adding to it the known exact instanton–induced superpotentials of the corresponding $`N=1`$ theories. $`N=1`$ supersymmetry guarantees that physics depend on $`\mu `$ holomorphically, so there cannot be any phase transition as $`\mu `$ is varied to smaller values.
### 3.1 $`SU(n_c)`$: $`n_fn_c+1`$
When the number of flavors is relatively small, the effective superpotential takes the form: <sup>1</sup><sup>1</sup>1To be precise, this is the form of the superpotential for generic $`n_f<n_c`$. For special cases $`n_f=n_c`$ and $`n_f=n_c+1,`$ one must use appropriate superpotentials involving baryonlike composites as well as mesons. See .
$$W=\frac{1}{2\mu }\left[\text{Tr}M^2\frac{1}{n_c}(\text{Tr}M)^2\right]+\text{Tr}(Mm)+(n_cn_f)\frac{\mathrm{\Lambda }_1^{(3n_cn_f)/(n_cn_f)}}{(detM)^{1/(n_cn_f)}},$$
(3.1)
where $`M_j^i\stackrel{~}{Q}_jQ^i`$ are the $`n_f\times n_f`$ meson matrix, and
$$\mathrm{\Lambda }_1\mu ^{\frac{n_c}{3n_cn_f}}\mathrm{\Lambda }^{\frac{2n_cn_f}{3n_cn_f}}$$
(3.2)
is the invariant mass scale of the $`N=1`$ SQCD (without the adjoint fields). The last term of Eq.(3.1) is the Affleck-Dine-Seiberg instanton-induced superpotential, the first arises from integrating out the adjoint fields $`\mathrm{\Phi }`$, $`m=\text{diag}(m_1,m_2,\mathrm{},m_f)`$ is the bare quark mass matrix. After making an Ansatz, $`M=\text{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _{n_f}),`$ one can straightforwardly find the minima of the potential. We find $`(2n_cn_f){}_{n_f}{}^{}C_{r}^{}`$ vacua in which the global symmetry is spontaneously broken (in $`m_i0`$ limit) as
$$U(n_f)U(r)\times U(n_fr).$$
(3.3)
By summing over $`r`$, $`r=0,1,2,\mathrm{}[\frac{n_f}{2}]`$, one finds $`𝒩_1`$ of Eq.(2.2).
### 3.2 $`SU(n_c)`$: $`n_fn_c+2`$
When the number of flavor exceeds $`n_c+1`$, physics at low energies is described by the effective superpotential,
$$𝒲=\stackrel{~}{q}Mq+\text{Tr}(mM)\frac{1}{2\mu }\left[\text{Tr}M^2\frac{1}{n_c}(\text{Tr}M)^2\right],$$
(3.4)
where $`q`$ stands for the $`n_f`$ dual quarks in the fundamental representation of the dual gauge group $`SU(\stackrel{~}{n}_c)`$, $`M`$ is the meson matrix as in Eq.(3.1). By first minimizing the potential with respect to $`q`$ and $`M`$, one finds
$$𝒩_2=\underset{r=0}{\overset{\stackrel{~}{n}_c1}{}}{}_{n_f}{}^{}C_{r}^{}(\stackrel{~}{n}_cr)$$
(3.5)
of vacua, in which VEVS behave as
$$q0,M0,$$
(3.6)
in the limit, $`m_i0`$. In other words, in these vacua the global $`SU(n_f)\times U(1)`$ symmetry remains unbroken.
One seems to encounter a puzzle though: the number of the vacua found here $`𝒩_2`$ is always less than the known total number of vacua, $`𝒩`$ (Eq.(2.1)). Where are other vacua?
Actually, we have tacitly assumed $`\text{Rank}M<n_f`$ above, for otherwise the dual quarks are all massive and the theory reduces to the pure $`SU(\stackrel{~}{n}_c)`$ Yang-Mills in the infrared: its strong interaction dynamics must be taken into account in order to get information about its ground state.<sup>2</sup><sup>2</sup>2 In fact, a related puzzle is how Seiberg’s dual Lagrangian \- the first two terms of Eq. (3.4) - can give the right number of vacua for the massive $`N=1`$ SQCD with $`n_f>n_c+1`$. By following the same method as below but with $`\mu =\mathrm{}`$, we do find the correct number ($`n_c`$) of vacua. In order to retrieve these vacua, we must first integrate out the dual quark fields. The instanton effects in the dual gauge group $`SU(\stackrel{~}{n}_c)`$ leads to a superpotential, which is identical to (actually continuation of ) $`W`$ in Eq.(3.1)! The minimization of such a potential yields $`𝒩_1=(2n_cn_f)\mathrm{\hspace{0.17em}2}^{n_f1}`$ vacua as before.
But then, for consistency, the sum of $`𝒩_1`$ and $`𝒩_2`$ must be equal to $`𝒩`$ of Eq.(2.1). As one can show easily by changing the dummy variable and by using the known identities among the binomial coefficients, this is indeed so.
To sum up, in $`SU(n_c)`$ theories the exact global $`U(n_f)`$ symmetry in the equal mass (or massless) limit, which is spontaneously broken to $`U(r)\times U(n_fr)`$ in $`(2n_cn_f){}_{n_f}{}^{}C_{r}^{}`$ vacua, $`r=0,1,\mathrm{},[n_f/2]`$. When the number of the flavor is larger ($`n_f>n_c+1`$), we find another class of vacua, with no global symmetry breaking. We shall see that these match the vacua in the free-magnetic phase at small $`\mu `$.
### 3.3 $`USp(2n_c)`$, $`SO(n_c)`$
The analysis in the cases of $`USp(2n_c)`$ or $`SO(n_c)`$ theories is similar, although the results are qualitatively different from the case of $`SU(n_c)`$ theory.
In $`USp(2n_c)`$ (or $`SO(n_c)`$) theories, for small numbers of flavors, the chiral $`SO(2n_f)`$ (or $`USp(2n_f)`$) symmetry in the massless limit is always spontaneously broken to $`U(n_f)`$. This result nicely agrees with what is expected generally from bi-fermion condensate of the standard form in non supersymmetric theories, and forms a result in closest analogy with what is supposed to occur in the standard QCD with small number of flavors.
Finally, in the cases of $`USp(2n_c)`$ or $`SO(n_c)`$ theories too, there exist also other vacua without dynamical symmetry breaking, when the number of flavor is greater ($`n_f>n_c+2`$ or $`2n_f>n_c4`$, respectively).
## 4 Quantum vacua at $`\mu \mathrm{\Lambda }`$
At small $`\mu `$, the infrared properties of the theory are described by certain singularities of Seiberg-Witten curves . To be concrete take the case of $`SU(n_c)`$ gauge theory with $`n_f`$ flavors. $`N=1`$ supersymmetric vacua are found by requiring that the curve
$$y^2=\underset{k=1}{\overset{n_c}{}}(x\varphi _k)^2+4\mathrm{\Lambda }^{2n_cn_f}\underset{j=1}{\overset{n_f}{}}(x+m_j),$$
(4.1)
where
$$\varphi =\text{diag}(\varphi _1,\varphi _2,\mathrm{}),$$
(4.2)
describe the gauge invariant VEVS,
$$u=\text{Tr}\mathrm{\Phi }^2=\underset{i<j}{}\varphi _i\varphi _j,u_3=\text{Tr}\mathrm{\Phi }^3=\underset{i<j<k}{}\varphi _i\varphi _j\varphi _k,$$
(4.3)
etc., to be maximally singular. Since there are $`n_c1`$ free parameters, up to $`n_c1`$ pairs of branch points can be made to coincide by appropriate choices of $`\{\varphi \}`$, and this corresponds to the condition that there are maximal number of massless monopoles of Abelian subgroup $`U(1)^{n_c1}SU(n_c).`$
Such a connection follows from the by now well understood association of monopole masses with the integral over canonical cycles of meromorphic differentials on the curve such as Eq.(4.1) . Also one can show that, upon $`\mu `$ perturbation, only these singular points of QMS (quantum space of vacua) lead to supersymmetric ground states.
Once the low-energy degrees of freedoms (monopoles) are identified and their quantum numbers known, it is in principle straightforward to analyze the properties of the vacua.
It turns out that the limit $`m_i0`$ is highly nontrivial, and physics in the infrared is far richer than one might have expected from the knowledge of $`SU(2)`$ gauge theory . In particular, one finds that the low-energy degrees of freedom are not always the monopoles of the maximal Abelian subgroup envisaged in the Nambu-’t Hooft-Mandelstam mechanism.
The crucial steps of our analysis are the points iv) and v) of Sec. 2. The first of these steps shows how all the $`N=1`$ vacua, selected out by the adjoint mass perturbation, are associated with the various universality classes of SCFT , and allows us to relate the quantum vacua at small $`\mu `$ to those at large $`\mu `$; the second step leads to the microscopic picture of confinement and dynamical symmetry breaking, summarized below.
## 5 Microscopic Picture of Dynamical Symmetry Breaking
In $`SU(n_c)`$ theories with $`n_f`$ flavors, there are two group of vacua. In the first group of vacua with finite meson or dual quark vacuum expectation values (VEVS), labeled by an integer $`r`$, $`r[n_f/2]`$, the system is in confinement phase. The nature of the actual carrier of the flavor quantum numbers depends on $`r`$. In vacua with $`r=0`$, magnetic monopoles are singlets of the global $`U(n_f)`$ group, hence no global symmetry breaking accompanies confinement.
In vacua with $`r=1`$, the light particles are magnetic monopoles in the fundamental representation of $`U(n_f)`$ flavor group. Their condensation leads to the confinement and flavor symmetry breaking, simultaneously.
In vacua labeled by $`r`$, $`2r<n_f/2`$ ($`rn_fn_c`$), the grouping of the associated singularities on the Coulomb branch, with multiplicity, $`{}_{n_f}{}^{}C_{r}^{}`$, at first sight suggests the condensation of monopoles in the rank-$`r`$ anti-symmetric tensor representation of the global $`SU(n_f)`$ group. Actually, this does not occur. The low-energy degrees of freedom of these theories are $`n_f`$ magnetic quarks (in $`\underset{¯}{r}`$) plus a number of singlet monopoles of a non-Abelian effective $`SU(r)\times U(1)^{n_cr}`$ gauge theory .
Monopoles in higher representations of $`SU(n_f)`$ flavor group, even if they exist semi-classically, break up into magnetic quarks before they become massless at singularities on the Coulomb branch. It is the condensation of the latter that induces confinement and flavor symmetry breaking, $`U(n_f)U(r)\times U(n_fr)`$, in these vacua. The system thus realizes the global symmetry of the theory in a Nambu-Goldstone mode, without having unusually many Nambu-Goldstone bosons. It is a novel mechanism for confinement and dynamical symmetry breaking.
In the special cases with $`r=n_f/2`$, still another dynamical scenario takes place. In these cases, the interactions among the monopoles and dyons become so strong that the low-energy theory describing them is a nontrivial SCFT, with conformal invariance explicitly broken by the adjoint or quark masses. Although the symmetry breaking pattern is known $`U(n_f)U(n_f/2)\times U(n_f/2)`$, the low-energy degrees of freedom in general involve relatively nonlocal fields and their interactions cannot be described in terms of a local action.
Finally, there are vacua in which magnetic-quarks do not condense and remain as physically observable particles at long distances, interacting with non-Abelian dual gauge particles (free magnetic phase), when the number of the flavor $`n_f`$ exceeds $`n_c+1`$. The global $`U(n_f)`$ symmetry remains unbroken in these vacua. These precisely match the second group of vacua found at large $`\mu `$ in which no condensate forms.
In $`USp(2n_c)`$ theories, again, we find two groups of vacua, whose properties are shown in Table 2. The most salient difference as compared to the $`SU(n_c)`$ theory is that here the entire first group of vacua corresponds to a SCFT. It is a nontrivial superconformal theory: one does not have a local effective Lagrangian description for those theories.<sup>3</sup><sup>3</sup>3Except for $`n_f=3`$ and $`n_f=2`$, in which cases we expect a local description to be valid, see Sec. 8.2 of Carlino et. al. . Nonetheless, the symmetry breaking pattern can be deduced, from the analysis done at large $`\mu `$: $`SO(2n_f)`$ symmetry is always spontaneously to $`U(n_f)`$.
To see better what is going on, it is instructive to consider the equal but nonvanishing quark mass case first. See Table 3. The flavor symmetry group of the underlying theory is now explicitly broken to $`U(n_f)`$. The first group of vacua split into various subgroups of vacua labeled by $`r`$, $`r=0,1,2,\mathrm{},[\frac{n_f1}{2}]`$, each of which is described by a local effective gauge theory of Argyres-Plesser-Seiberg for $`SU(n_c)`$ theory (!), with gauge group $`SU(r)\times U(1)^{n_cr+1}`$ and $`n_f`$ (dual) quarks in the fundamental representation of $`SU(r)`$.<sup>4</sup><sup>4</sup>4The key fact that some of the SCFT in $`USp(2n_c)`$ or in $`SO(n_c)`$ theories at $`m_i=m0`$ are in the same universality classes as those occurring in the $`SU(n_c)`$ theory, was first noted by Eguchi et. al. . Our perturbation theory in masses around the SCFT and the effective action analysis show that these SCFT are indeed the $`N=1`$ vacua which survive the adjoint mass perturbation (missed in the analysis of ) . Indeed, the gauge invariant composite VEVS characterizing these vacua differ by some positive powers of $`m`$, so that the validity of each effective theory is limited to small fluctuations of order of $`m^a,`$ $`a>0,`$ around each vacuum.
In the limit $`m0`$ these points in the quantum moduli space (QMS) collapse into one single point and accordingly the range of the validity of each local effective action shrinks to zero, implying that we have a nontrivial SCFT here, with mutually nonlocal massless states. In the example of $`USp(4)`$ theory with $`n_f=4`$, we have explicitly verified this by determining the singularities and branch points at finite equal mass $`m`$ and then by studying the limit $`m0`$.
The first group of vacua in the $`SO(n_c)`$ gauge theory has similar characteristics as the one in $`USp(2n_c)`$ theory just discussed. These cases, together with the special $`r=n_f/2`$ vacua for the $`SU(n_c)`$ theory, reveal another new mechanism for dynamical symmetry breaking: although the global symmetry breaking pattern deduced indirectly looks familiar enough, the low-energy degrees of freedom are relatively nonlocal dual quarks and dyons. It would be interesting to get a better understanding of this phenomenon.
For large numbers of flavor, there are also vacua, just as in large $`n_f`$ $`SU(n_c)`$ theories, with no confinement and no dynamical flavor symmetry breaking. Physics around these vacua can be explicitly studied by use of the effective low-energy action at the ”special points” of . The low-energy particles are solitonlike magnetic quarks which weakly interact with dual (in general) non-Abelian gauge fields: the system is in the free magnetic phase. These are smoothly connected to the second group of vacua, with no symmetry breaking, found at large $`\mu `$.
Properties of various vacua in $`SU(n_c)`$ and $`USp(2n_c)`$ theories are illustrated schematically in Fig. 1 and Fig. 2.
## Acknowledgments
The author thanks G. Carlino, P. Kumar and H. Murayama for fruitful and enjoyable collaborations, and the organizers of the workshop “Continuous Advances in QCD” (Univ. Minnesota, May 2000) for inviting him to present and discuss these recent results in such a stimulating atmosphere.
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# Cluster Temperature Evolution: The Mass-Temperature Relation
## 1 Introduction
Surveys of distant clusters of galaxies are now realizing their promise as cosmological indicators (e.g., Henry 1997; Eke et al. 1998; Borgani et al. 1999a; Donahue & Voit 1999). Because clusters are the largest virialized objects in the universe, and the latest objects to form in hierarchical models of structure formation, their rate of evolution is quite sensitive to cosmological parameters. However, because cluster masses are difficult to measure directly, a surrogate for cluster mass is usually used when comparing cluster observations to structure-formation models. In the X-ray regime, the simplest cluster observables are X-ray luminosity ($`L_X`$) and emissivity-weighted temperature ($`T_X`$). Temperature is more directly related to a cluster’s mass, but luminosity can also be mapped to mass via an $`L_XT_X`$ relation. In either case, the cruical link between models and observations is the mass-temperature relation.
Recent analyses of high-redshift cluster temperature functions from the Einstein Extended Medium-Sensitivity Survey (EMSS) have shown that the matter density of the universe probably lies in the range $`0.2<\mathrm{\Omega }_\mathrm{M}<0.7`$ (Henry 1997; Donahue et al. 1998; Bahcall & Fan 1998; Eke et al. 1998; Donahue & Voit 1999; but see Blanchard & Bartlett 1998; Viana & Liddle 1999). Even though these conclusions are based on rather few clusters, the systematic errors in these measurements of $`\mathrm{\Omega }_\mathrm{M}`$ are comparable to the statistical errors (Donahue & Voit 1999). With the flood of cluster temperature measurements expected over the next few years from Chandra and XMM, we will have the opportunity to measure $`\mathrm{\Omega }_\mathrm{M}`$ much more precisely. If we are to take full advantage of these measurements, we will need to reduce the systematic errors that currently exist in the modeling. Here we concentrate on the uncertainties in the mass-temperature relation itself.
This paper analyzes the role of the mass-temperature relation in characterizing cluster temperature evolution. Section 2 outlines the formalism used to decribe the evolution of the cluster mass function. Section 3 investigates the physics underlying the mass-temperature relation, showing that the standard derivation of this relation from spherical top-hat collapse is flawed and suggesting a new context for understanding the physical effects that govern this relation. Section 4 discusses how to normalize the mass-temperature relation and evaluates how severely uncertainties in this normalization affect measurements of cosmological parameters. Section 5 summarizes the paper.
## 2 Mass-Function Evolution
Numerical simulations have shown that the Press-Schechter formalism (Press & Schechter 1974) and its various extensions (e.g., Lacey & Cole 1993) characterize gravitationally driven structure formation with surprising fidelity, particularly on cluster scales (e.g., Lacey & Cole 1994; Eke, Cole, & Frenk 1996; Bryan & Norman 1998; Borgani et al. 1999a). Because of its successes, this formalism is often used to relate observed cluster temperature and luminosity functions to cosmological models (e.g., Henry & Arnaud 1991; Eke et al. 1996; Borgani et al. 1999a). While some recent large-scale simulations have shown that the Press-Schechter formula might be somewhat less successful at predicting the number density of the most massive clusters (e.g., Governato et al. 1999), we will assume in this paper that it is an exact description of the cluster mass function, because here we are more concerned with analyzing systematic problems with the mass-temperature relation.
This section briefly outlines the formalism we will employ for expressing the cluster mass function. We derive expressions for mass-function evolution in both open and flat universes, and we assess the effects of a cosmological constant on cluster evolution. Many of these results have been derived elsewhere; we compile them here as background for subsequent sections.
### 2.1 Press-Schechter Formalism
The Press-Schechter formalism for structure formation and its extensions describe how virialized objects grow from a field of initial density perturbations. One defines $`\delta (𝐱,t;M)`$ to be the local fractional overdensity of the universe smoothed on mass scale $`M`$ and centered on comoving point $`𝐱`$ at time $`t`$. While these fluctuation amplitudes are linear, they grow in proportion to the function $`D(t)`$, which depends on $`\mathrm{\Omega }_\mathrm{M}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$, and their rms amplitude on scale $`M`$ can be expressed as $`\sigma (M)D(t)/D(t_0)`$. Ultimately, some of these fluctuations grow non-linear, and they are assumed to virialize when their amplitudes, extrapolated from the linear regime according to $`D(t)`$, exceed some critical threshold for virialization $`\delta _c(t)`$, which also depends on $`\mathrm{\Omega }_\mathrm{M}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$. One can then trace the merger history of a mass parcel beginning at location $`𝐱`$ from time $`t_1`$ to the present by keeping track of the largest $`M`$ for which $`\delta (𝐱,t_1;M)D(t)/D(t_1)>\delta _c(t)`$.
Assuming that the perturbations are Gaussian, we can assess the number density of virialized objects with mass $`>M`$ by evaluating the quantity $`\nu _c(M,t)[\delta _c(t)/\sigma (M)][D(t_0)/D(t)]`$, which is the critical virialization threshold in units of the characteristic fluctuation amplitude. The probability that a given mass parcel is part of a virialized structure of mass $`>M`$ is then equal to $`\mathrm{erfc}[\nu _c(M,t)/\sqrt{2}]`$, where $`\mathrm{erfc}(q)`$ is the complementary error function. Thus, the overall mass density in virialized objects exceeding mass $`M`$ is
$$\rho (>M)=\rho _0\mathrm{erfc}\left(\frac{\nu _c}{\sqrt{2}}\right)=\frac{2\rho _0}{\sqrt{\pi }}_{\nu _c/\sqrt{2}}^{\mathrm{}}e^{x^2}𝑑x,$$
(1)
where $`\rho _0`$ is the mean mass density of the universe. Differentiating this expression with respect to $`M`$ and dividing the result by $`M`$ yields the familiar Press-Schechter formula for the comoving differential number density $`dn`$ of virialized objects within mass interval $`dM`$:
$$\frac{dn}{dM}(M,t)=\left(\frac{2}{\pi }\right)^{1/2}\frac{\mathrm{\Omega }_\mathrm{M}\rho _{\mathrm{cr},0}}{M^2}\left|\frac{d\mathrm{ln}\sigma }{d\mathrm{ln}M}\right|\nu _c(M,t)\mathrm{exp}[\nu _c^2(M,t)/2],$$
(2)
where $`\rho _{\mathrm{cr},0}=3H_0^2/8\pi G`$ represents the present-day critical mass density.
The evolution of $`dn/dM`$ depends solely on $`\nu _c(M,t)`$. At present, we have $`\nu _c(M,t_0)=\delta _c(t_0)/\sigma (M)`$, and in principle, we can determine $`\sigma (M)`$ by fitting equation (2) to the current distribution of cluster masses. The function $`\sigma (M)`$ can be approximated by a power law with index $`\alpha =(n+3)/6`$ on cluster scales, so that $`\sigma (M)=\sigma _8(M/M_8)^\alpha `$ with $`M_8=(H_0^2\mathrm{\Omega }_\mathrm{M}/2G)(8h^1\mathrm{Mpc})^3=6.0\times 10^{14}\mathrm{\Omega }_\mathrm{M}h^1M_{}`$. Holding $`\sigma (M)`$ fixed, we can project the current cluster mass distribution backward in time as long as we know the functions $`\delta _c(t)`$ and $`D(t)`$. More specific expressions for $`\sigma (M)`$ describe how $`n`$ changes with the mass scale in CDM-like models, but the analytical simplicity of the power-law form better serves the illustrative purposes of this paper.
### 2.2 Cluster Evolution and $`\mathrm{\Omega }_\mathrm{M}`$
Massive cluster evolution is very sensitive to $`\mathrm{\Omega }_\mathrm{M}`$ because the number density of large clusters depends exponentially on $`\nu _c^2`$ which is $`1`$. Slight differences in the rate at which $`\nu _c`$ evolves therefore translate into large differences in cluster evolution (e.g., Oukbir & Blanchard 1992; Eke et al. 1996; Viana & Liddle 1996). To illustrate how dramatic these differences can be, we will briefly outline the case of cluster evolution in an open universe with no cosmological constant.
Following Lacey & Cole’s (1993) treatment of perturbation growth when $`\mathrm{\Omega }_\mathrm{M}<1`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ (see also the Appendix), we have $`\nu _c(M,t)=\omega (t)/\sigma (M)`$ with
$$\omega (t)\delta _c(t)D(t_0)/D(t)=\frac{3}{2}D(t_0)[1+(t_\mathrm{\Omega }/t)^{2/3}],$$
(3)
where $`t_\mathrm{\Omega }=\pi \mathrm{\Omega }_\mathrm{M}/H_0(1\mathrm{\Omega }_\mathrm{M})^{3/2}`$. Normalizing the present value of $`\nu _c`$ at some fiducial mass scale $`M_0`$, we can write
$$\nu _c(M,t)=\nu _c(M_0,t_0)\left(\frac{M}{M_0}\right)^\alpha \frac{1+(t_\mathrm{\Omega }/t)^{2/3}}{1+(t_\mathrm{\Omega }/t_0)^{2/3}}.$$
(4)
Because the cluster number density at mass scale $`M_0`$ obeys $`dn/dM\nu _c\mathrm{exp}(\nu _c^2/2)`$, we can use this equation to gauge how rapidly clusters at this mass scale evolve.
Figure 1 illustrates how sensitively the number-density evolution of massive clusters ($`M_05\times 10^{14}h^1M_{}`$) depends on $`\mathrm{\Omega }_\mathrm{M}`$. For the purposes of this illustration, we have adopted the $`\sigma _8`$ fitting formulae of Eke et al. (1996) and have assumed that $`n=1.5`$ on this mass scale. In this example, the number density of clusters in a universe with $`\mathrm{\Omega }_\mathrm{M}=0.8`$ grows by about five orders of magnitude from $`z=1`$ to the present, a rapid rate of evolution that contrasts sharply with the single order of magnitude expected in a universe with $`\mathrm{\Omega }_\mathrm{M}=0.2`$. Changing the perturbation spectrum within observationally allowed bounds changes the quantitative predictions somewhat, but the qualitative conclusion remains the same: the number density of massive clusters evolves much more rapidly for $`\mathrm{\Omega }_\mathrm{M}1`$ than for $`\mathrm{\Omega }_\mathrm{M}1`$.
### 2.3 Cluster Evolution and $`\mathrm{\Lambda }`$
Clusters evolve slightly more rapidly in a flat, $`\mathrm{\Omega }_\mathrm{\Lambda }>0`$ universe than they do in an open, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ universe with the same value of $`\mathrm{\Omega }_\mathrm{M}`$ (e.g., Eke et al. 1996). When a cosmological constant is operating, the universe’s density remains close to critical later in time, promoting perturbation growth at lower redshifts. However, cluster evolution is considerably less sensitive to $`\mathrm{\Omega }_\mathrm{\Lambda }`$ than it is to $`\mathrm{\Omega }_\mathrm{M}`$.
In order to characterize cluster evolution in a flat universe with $`\mathrm{\Omega }_\mathrm{\Lambda }>0`$ we require expressions for $`D(t)`$ and $`\delta _c(t)`$, which are derived in the Appendix. From these expressions we can construct the threshold function
$$\omega (t)=9\xi _c(t)D(t_0),$$
(5)
where $`\xi _c(t)`$, defined in the the Appendix, is proportional to the specific energy of a perturbation that collapses at time $`t`$. Normalizing $`\nu _c`$ as before at the mass scale $`M_0`$, we can write
$$\nu _c(M,t)=\nu _c(M_0,t_0)\left(\frac{M}{M_0}\right)^\alpha \frac{\xi _c(t)}{\xi _c(t_0)}.$$
(6)
Plugging this expression into equation (2) then yields the desired formulae for cluster evolution.
The dotted lines in Figure 1 show how the number density of clusters at mass scale $`M_0`$ evolves when $`\mathrm{\Omega }_\mathrm{\Lambda }>0`$, for the $`\sigma _8`$ normalization of Eke et al. (1996) and $`n=1.5`$. Note that the rate of cluster evolution is quite insensitive to $`\mathrm{\Omega }_\mathrm{\Lambda }`$. In this particular case, the best-fitting $`\sigma _8`$ for a flat universe is slightly higher than that for an open universe, which almost compensates for the slightly more rapid rate of evolution owing to $`\mathrm{\Lambda }`$. In general, the best-fitting $`\mathrm{\Omega }_\mathrm{M}`$ in an open universe is $`0.1`$ higher than in a flat universe (e.g., Donahue & Voit 1999).
## 3 Theoretical Mass-Temperature Relations
The Press-Schechter formalism conveniently describes the rate at which virialized objects of mass $`M`$ accumulate in the universe. If we could observe cluster masses directly, then comparisions between Press-Schechter predictions and observed cluster evolution would be simple. Several types of observables, such as X-ray temperature, cluster velocity dispersion, and weak lensing, are related to cluster masses, but linking these quantities with the proper Press-Schechter $`M`$ values requires careful attention.
This section focuses on the relation between cluster mass and cluster temperature, the crucial relationship for linking X-ray observations of clusters to models of structure formation. We first outline the observational evidence for a well-behaved mass-temperature relationship. Then we analyze the standard derivation of the mass-temperature relation, which is based on collapse of a spherical top-hat perturbation to an isothermal sphere. This derivation yields a relation similar to the observed relation, but it fails to conserve energy, indicating that it omits important physical effects. In an effort to understand this relation more deeply, we present a model for cluster formation, drawn from the merging-halo formalism of Lacey & Cole (1993), which accounts for the fact that massive clusters accrete matter quasi-contiuously. Analyzing clusters in this context enables us to identify the physical effects that make up for the lack of energy conservation. The primary advantage of the continuous formation model is that it more naturally reproduces the late-time evolution of clusters, and the section concludes by comparing predictions of cluster temperature evolution drawn from the continuous formation model with those from the spherical top-hat model.
### 3.1 Observational Evidence
Simple scaling arguments suggest that the X-ray temperatures of clusters ($`T_X`$) should be directly related to their masses. One way to define a cluster’s mass is to specify a characteristic radius $`r_\mathrm{\Delta }`$ within which the mean density is $`\mathrm{\Delta }`$ times the critical density $`\rho _{\mathrm{cr}}`$, so that $`M_\mathrm{\Delta }=4\pi r_\mathrm{\Delta }^3\rho _{\mathrm{cr}}\mathrm{\Delta }/3`$. If all cluster potentials share the same density distribution, $`\rho (r/r_\mathrm{\Delta })`$, and the X-ray gas is isothermal, then $`T_XM_\mathrm{\Delta }/r_\mathrm{\Delta }M_\mathrm{\Delta }^{2/3}`$. Numerical simulations of cluster formation demonstrate that this scaling ought to be remarkably tight, with a scatter of only $`1520`$% (Evrard, Metzler, & Navarro 1996; Bryan & Norman 1998). These simulations also provide normalizations for the mass-temperature relation that can be compared with actual clusters.
A recent observational investigation of the cluster mass-temperature relation at $`z0.1`$ by Horner, Mushotsky, & Scharf (1999) supports the results of the simulations. They show that cluster masses derived from velocity dispersions (Girardi et al. 1998) agree well with those inferred from ASCA temperatures (Fukazawa 1997), using the scaling law from Evrard et al. (1996) at $`\mathrm{\Delta }=200`$:
$$M_{200}=(1.4\times 10^{15}h^1M_{})\left(\frac{T_X}{10\mathrm{keV}}\right)^{3/2}.$$
(7)
The scatter in the observed mass-temperature relation is $`30`$%, but it decreases by a factor of 2 for the clusters with the highest numbers of measured galaxy redshifts, suggesting that the scatter intrinsic to the mass-temperature relation is probably quite small. However, some concerns remain: a handful of outliers deviate from the standard relation by up to 50% and the mass normalization one finds from hydrostatic modeling of a subset of these clusters is 40% smaller (Horner et al. 1999).
Similar comparisions at higher redshifts are more difficult, but the available data indicate that the mass-temperature relation remains well-behaved. Hjorth, Oukbir, & van Kampen (1998) have compared masses derived from gravitational lensing analyses for 8 clusters at $`0.17z0.54`$ with the X-ray temperatures of these clusters. Their best-fit mass-temperature relation agrees with the Evrard et al. (1996) scaling law within the observational errors, and they conclude that this scaling law can be used to measure masses to within 27%. For clusters at even higher redshifts ($`0.53z0.83`$), Donahue et al. (1999) show that the observed relation between X-ray temperatures and cluster velocity dispersions remains consistent with the low-$`z`$ relation.
### 3.2 Virial Mass and the Late-Formation Approximation
The seemingly good behavior of the cluster mass-temperature relation is fortunate for those who wish to study cluster evolution with X-ray telescopes, but care must be taken when relating these temperature-derived masses to the virial masses demanded by the Press-Schechter formalism. From the simulations and observations, we know that the mass within a specified density contrast is straightforwardly related to temperature. However, the density contrast $`\mathrm{\Delta }_{\mathrm{vir}}`$ corresponding to the virial radius depends in general on $`\mathrm{\Omega }_\mathrm{M}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$. Thus, in order to characterize cluster evolution properly, we need to know how $`\mathrm{\Delta }_{\mathrm{vir}}(t;\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$ changes with time.
The usual approach to defining a cluster’s virial mass is to approximate cluster formation with the evolution of a spherical top-hat perturbation (e.g., Peebles 1993). Such a perturbation formally collapses to the origin at a particular moment ($`t_c`$) which is taken to be the moment of virialization. The virialization time thus equals twice the time required for the perturbation to reach its turnaround radius ($`r_{\mathrm{ta}}`$). A naive application of the virial theorem, assuming that the perturbation is cold at maximum expansion, dictates that the cluster’s final potential energy ought to be twice its potential energy at turnaround. Hence, the cluster’s virial radius is assumed to be half its turnaround radius ($`r_{\mathrm{vir}}=r_{\mathrm{ta}}/2`$). According to this prescription, $`\mathrm{\Delta }_{\mathrm{vir}}`$ is a well-defined function of cosmic time and the parameters $`\mathrm{\Omega }_\mathrm{M}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ (Lacey & Cole 1993; Kitayama & Suto 1996; Oukbir & Blanchard 1997). In the case of $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, this function can be concisely expressed as $`\mathrm{\Delta }_{\mathrm{vir}}=8\pi ^2/(Ht)^2`$, where $`H`$ is the Hubble constant at time $`t`$. If we additionally assume that each cluster we see at a given redshift $`z`$ has just reached the moment of virialization, an assumption known as the late-formation approximation, then $`M_{\mathrm{vir}}T_X^{3/2}\rho _{\mathrm{cr}}^{1/2}\mathrm{\Delta }_{\mathrm{vir}}^{1/2}`$.
In a critical $`\mathrm{\Omega }_\mathrm{M}=1`$ universe, the late-formation approximation is valid because massive clusters develop rapidly at all redshifts; the effective moment of virialization is always close to the moment of observation. However, in a universe with $`\mathrm{\Omega }_\mathrm{M}<1`$, cluster formation is currently shutting down, and one must account for differences between the moment of virialization and the moment of observation. This problem grows most severe at late times in a $`\mathrm{\Omega }_\mathrm{M}1`$ universe, because the quantity $`\rho _{\mathrm{cr}}\mathrm{\Delta }_{\mathrm{vir}}`$ as determined via the late-formation approximation declines indefinitely. The $`M_{\mathrm{vir}}`$ associated with a given $`T_X`$ therefore rises steadily, even though cluster evolution has essentially stopped. This spurious late-time evolution of the $`M_{\mathrm{vir}}T_X`$ relation is an undesirable artifact of the late-formation approximation.
One approach to solving this problem is to account explicitly for the difference between the moment of virialization and the moment of observation in the context of a merging-halo formalism for cluster growth (Viana & Liddle 1996; Kitayama & Suto 1996). Another equivalent but mathematically simpler approach, which § 3.3 describes in detail, is to consider how the $`M_{\mathrm{vir}}T_X`$ relation should evolve in a population of clusters that gradually accrete their matter over an extended period of time, a more realistic scenario for the growth of very massive clusters (Voit & Donahue 1998).
Calculating the normalization of the $`M_{\mathrm{vir}}T_X`$ relation under the late-formation approximation is also somewhat problematic. Because a virialized cluster’s potential is approximately isothermal, one would like to approximate it with a singular isothermal sphere, truncated at radius $`r_{\mathrm{vir}}`$, within which the mean density is $`\rho _{\mathrm{cr}}\mathrm{\Delta }_{\mathrm{vir}}`$. The one-dimensional velocity dipersion within such a potential is $`\sigma _{1\mathrm{D}}^2=GM/2r_{\mathrm{vir}}`$ (Binney & Tremaine 1987), which leads to the following relation between virial mass and temperature:
$`kT_X`$ $`=`$ $`{\displaystyle \frac{GM^{2/3}\mu m_p}{2\beta }}\left[{\displaystyle \frac{4\pi }{3}}\rho _{\mathrm{cr}}\mathrm{\Delta }_{\mathrm{vir}}\right]^{1/3}`$ (8)
$`=`$ $`(1.38\mathrm{keV})\beta ^1h^{2/3}M_{15}^{2/3}\mathrm{\Delta }_{\mathrm{vir}}^{1/3}\left[{\displaystyle \frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_M(z)}}\right]^{1/3}(1+z)`$
where $`\beta =\mu m_p\sigma _{1\mathrm{D}}^2/kT_X`$ and $`\mu m_p=1\times 10^{24}\mathrm{g}`$ is the mean mass per gas particle.
Comparisions of the $`M_{\mathrm{vir}}T_X`$ relation in equation (8) with the masses and temperatures of simulated clusters indicate that $`\beta ^10.81`$ (e.g., Bryan & Norman 1998). This nearness of $`\beta `$ to unity appears to validate the assumptions governing the derivation of equation (8), but the approximate agreement between this equation and the simulations turns out to be something of a coincidence. The total energy of a collapsing spherical top-hat perturbation is $`3GM^2/5r_{\mathrm{ta}}`$. After the collapsed perturbation virializes into an isothermal sphere, a naive application of the virial theorem that disregards boundary effects would place the total kinetic energy of the system at $`3GM^2/5r_{\mathrm{ta}}`$, corresponding to $`\sigma _{1\mathrm{D}}^2=2GM/5r_{\mathrm{ta}}`$. The virial radius of the relaxed system would then be $`5r_{\mathrm{ta}}/4`$, a factor of 2.5 larger than assumed in the derivation of equation (8), and its temperature would correspondingly be 2.5 times lower.
In fact, truncation of a virialized system at some $`r_{\mathrm{vir}}`$ implies the existence of a confining pressure, unaccounted for in the top-hat collapse model, that alters the usual virial relationship between potential and kinetic energy (e.g., Carlberg, Yee, & Ellingson 1997). In the case of a singular isothermal sphere, the total kinetic energy is three times the absolute value of the total energy. Thus, energy-conserving collapse of a spherical top-hat perturbation into a pressure-truncated singular isothermal sphere should yield $`\sigma _{1\mathrm{D}}^2=6GM/5r_{\mathrm{ta}}`$, implying $`r_{\mathrm{vir}}=5r_{\mathrm{ta}}/12`$ (e.g., Shapiro, Iliev, & Raga 1999). This result is close to the naive assumption that $`r_{\mathrm{vir}}=r_{\mathrm{ta}}/2`$, but it is valid only if a confining pressure is applied at the virial radius.
### 3.3 Continuously Forming Clusters
The inconsistencies in the top-hat, late-formation derivation of the $`M_{\mathrm{vir}}T_X`$ relation outlined above indicate, perhaps unsurprisingly, that the top-hat collapse model excludes important physical effects that contribute to the normalization of the relation. In particular, the top-hat model accounts for neither the energy and mass accumulation during the early stages of cluster formation nor the confining effects of matter that continues to fall in, both of which significantly increase the temperature associated with a given mass. This section shows how these missing effects can be addressed in the context of a simple model in which massive clusters are allowed to form gradually, rather than instantaneously.
In hierarchical models for structure formation, the growth of the largest clusters is quasi-continuous. The most massive clusters are so rare that they almost never merge with another cluster of similar size (e.g., Lacey & Cole 1993). Rather, they grow by continually accumulating much smaller virialized objects. In the notation of § 2, their masses grow like $`M\omega ^{3/(n+3)}`$ (Lacey & Cole 1993; Voit & Donahue 1998). Because each bit of infalling matter carries with it a specific energy $`ϵ`$, we can compute the virial energy $`E`$ of the cluster by integrating $`E=ϵ𝑑M`$. The cluster temperature itself is proportional to $`E/M`$, so this integral also leads to a relation between virial mass and temperature.
Voit & Donahue (1998) treat the case of continuous cluster growth when $`\mathrm{\Omega }_\mathrm{M}<1`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, finding that $`Mx^{3m/5}`$, where $`x=1+(t_\mathrm{\Omega }/t)^{2/3}`$ and $`m=5/(n+3)`$. Here we extend that calculation to include the constant of proportionality between energy and mass. Drawing on the Appendix, we express the specific energy of infalling matter at time $`t`$ as
$$ϵ(t)=\frac{1}{2}\left(\frac{2\pi GM}{t_\mathrm{\Omega }}\right)^{2/3}(x1).$$
(9)
Thus, we obtain
$$\frac{E}{M}=\frac{3}{10}\frac{m}{m1}\left(\frac{2\pi G}{t_\mathrm{\Omega }}\right)^{2/3}M^{2/3}\left[\left(\frac{t_\mathrm{\Omega }}{t}\right)^{2/3}+\frac{1}{m}\right].$$
(10)
In the limit of large $`m`$, which corresponds to the late-formation approximation, this expression reduces to
$$\frac{E}{M}=\frac{3}{5}ϵ(t),$$
(11)
which is identical to the $`E/M`$ ratio for a spherical top-hat perturbation of mass $`M`$ that virializes at time $`t`$. A similar procedure yields the mass-temperature relation in a flat $`\mathrm{\Omega }_\mathrm{\Lambda }>0`$ universe. From the Appendix, we have $`ϵ(t)M^{2/3}\xi _c(t)`$ and $`\omega (t)\xi _c(t)`$, giving
$$\frac{E}{M}=\frac{3}{5}\frac{m}{m1}ϵ(t),$$
(12)
which again reduces to $`3ϵ(t)/5`$ in the limit of large $`m`$.
Two factors in equation (10) drive $`E/M`$ higher than the late-formation value. The $`m/(m1)`$ factor, also present in equation (12), accounts for the effects of early infall; continuous cluster formation tends to create hotter clusters than top-hat formation because more of the mass is assembled early, at a higher mean density. For values of $`n`$ typical of cluster scales ($`2n1`$), this factor ranges from 1.2 to 1.7. The $`1/m`$ term in the bracketed factor of equation (10) accounts for the cessation of cluster formation when $`tt_\mathrm{\Omega }`$. At late times in an open universe, $`E/M`$ should remain constant, but in the late-formation approximation $`E/M`$ falls indefinitely because the fiducial density scale never stops dropping.
Relating $`E/M`$ to temperature requires an expression for the relationship between the total virial energy $`E`$ and the total kinetic energy $`E_K`$. When an external pressure $`P`$ confines the boundary of a spherically symmetric virialized system, the appropriate form of the virial theorem can be written
$$E_K=E+4\pi Pr_{\mathrm{vir}}^3.$$
(13)
If we take the velocity dispersion to be isothermal ($`\sigma _{1\mathrm{D}}=\mathrm{const}.`$), then $`P=\rho (r_{\mathrm{vir}})\sigma _{1\mathrm{D}}^2`$, and
$$E_K=\frac{\overline{\rho }}{\overline{\rho }2\rho (r_{\mathrm{vir}})}E,$$
(14)
where $`\overline{\rho }`$ is the mean density within the virial radius. In this formulation, the ratio $`E_K/E`$ depends on the shape of the potential within $`r_{\mathrm{vir}}`$. If the local density is negligible at $`r_{\mathrm{vir}}`$, then the confining pressure is effectively zero and $`E_K=E`$. If the potential strictly obeys $`\rho r^2`$, then $`E_K=3E`$.
Because we wish to derive an approximate mass-temperature relation for comparison with observations and simulations giving the temperature and mass of a cluster within $`r_{200}`$, let us compute $`E_K/E`$ for the “universal” density profile of Navarro, Frenk, & White (1997), truncated at $`r_{200}`$, under the assumption that $`\sigma _{1\mathrm{D}}`$ is constant. The ratio of mean density to local density within $`r_{200}`$ is then
$$\frac{\overline{\rho }}{\rho (r_{200})}=3\frac{(1+c)^2}{c^2}\mathrm{ln}\left[(1+c)\frac{c}{1+c}\right],$$
(15)
where $`c`$ is a parameter that quantifies the concentration of matter toward the cluster’s center. Simulations by Eke, Navarro, & Frenk (1998) show that $`c46.5`$ for clusters at $`0<z<1`$ in a $`\mathrm{\Omega }_\mathrm{M}=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ universe, and simulations of cluster formation by Navarro et al. (1997) show that the most massive clusters in critical universes exhibit similar levels of concentration. Values of $`c5`$ lead to density contrast factors $`4`$ and $`E_K/E2`$.
An intriguing alternative approach by Shapiro et al. (1999) investigates the post-collapse structures of clusters by seeking the minimum-energy solution among a family of non-singular truncated isothermal spheres. The minimum-energy solution turns out to closely resemble the self-similar spherical infall solution of Bertschinger (1985). Because the truncation radius of this minimum-energy isothermal sphere is nearly equal to the accretion-shock radius of the infall solution, Shapiro et al. (1999) suggest that continual infall naturally maintains the confining pressure on the virialized isothermal sphere. In this model, the density contrast factor at the truncation radius is 3.73, and $`E_K/E=2.17`$.
Taken together, the effects of early accretion and pressure confinement make up for the lack of energy conservation in the top-hat, late-formation derivation of equation (8). In the early-time limit ($`tt_\mathrm{\Omega }`$), equation (10) yields the following mass-temperature relation:
$$kT_X=\left[\frac{2}{5}\frac{m}{m1}\frac{E_K}{E}\right]\frac{GM^{2/3}\mu m_p}{2\beta }\left(\frac{4\pi }{3}\rho _{\mathrm{cr}}\mathrm{\Delta }_{\mathrm{vir}}\right)^{1/3}.$$
(16)
When $`n2`$ and $`E_K/E2`$, the prefactor in brackets is close to unity, making this expression nearly identical to equation (8).
The lesson here is that the assumptions underlying equation (8) are physically unsound. The approximate agreement between the $`M_{\mathrm{vir}}T_X`$ normalization derived via the top-hat collapse model and those derived from simulations and observations is largely coincidental. As long as $`\mathrm{\Omega }_\mathrm{M}`$ is not very small, the time-dependent factors in equations (8) and (10) do not differ by a large factor. However, because equation (10) more faithfully reflects the behavior of cluster formation in all the appropriate limits, we prefer to base the $`M_{\mathrm{vir}}T_X`$ relation on the continuous-formation model.
### 3.4 Late-Formation vs. Continuous-Formation
The $`M_{\mathrm{vir}}T_X`$ relations in equation (8), derived using the late-formation approximation, and equation (10), derived using the continuous-formation approximation, differ in both normalization and time-dependent behavior. The following section will discuss the importance of properly normalizing the $`M_{\mathrm{vir}}T_X`$ relation. Here we wish to examine how differences in time-dependence alone translate into different predictions for cluster evolution. In order to isolate the time-dependent behavior, we can identically normalize both $`M_{\mathrm{vir}}T_X`$ relations to equation (7) at $`z=0`$ and compare the resulting cluster temperature functions.
Figure 2 shows the evolution of the temperature function for 8 keV clusters, given the $`\sigma _8`$ normalization of Eke et al. (1996). Because the $`M_{\mathrm{vir}}T_X`$ relation evolves less strongly in the continuous-formation case, the rise in cluster temperature at a given mass as $`z`$ increases does not compensate as fully for the drop in the number of clusters at that mass. The evolution of the temperature function at a given value of $`\mathrm{\Omega }_\mathrm{M}`$ is therefore stronger in the continuous-formation case (see also Viana & Liddle 1999). Correspondingly, the best-fitting $`\mathrm{\Omega }_\mathrm{M}`$ to a given observed amount of cluster evolution will be lower. In this particular case, the difference amounts to $`0.1`$ in $`\mathrm{\Omega }_\mathrm{M}`$ for a best-fitting $`\mathrm{\Omega }_\mathrm{M}0.3`$.
Because the statistical errors in $`\mathrm{\Omega }_\mathrm{M}`$ derived from cluster temperature evolution are also $`0.1`$, this discrepancy between the late-formation and continuous-formation approximations will need to be resolved if we are to take full advantage of the cluster temperature measurements expected from Chandra and XMM. The best way to proceed will be to test how well these temperature-function predictions represent the results of large-scale structure simulations. However, cluster temperature functions will have to be extracted directly from the simulated data, without resorting to a mass-temperature conversion step, presumably by using the cluster-particle velocity dispersion as a surrogate for cluster temperature.
## 4 Normalization of the $`M_{\mathrm{vir}}T_X`$ Relation
Both the top-hat and continuous-formation derivations of the $`M_{\mathrm{vir}}T_X`$ relation given in the previous section have holes which must be plugged with knowledge gained from simulations. In fact, any such spherically symmetric representation glosses over aspects of cluster formation that are inherently three-dimensional. Thus, it seems wise to normalize these analytical expressions to the results of numerical simulations. This procedure appears simple enough, but one must bear in mind that the normalization depends on $`\mathrm{\Omega }_\mathrm{M}`$ and that simulations have been done only for a few particular values of $`\mathrm{\Omega }_\mathrm{M}`$. Here we explain how we choose to normalize the $`M_{\mathrm{vir}}T_X`$ relation then investigate the consequences of an offset in the normalization.
### 4.1 Normalizing to Simulations
Because of the good agreement between the observational compilation of Horner et al. (1999) and the simulations of Evrard et al. (1996), we would like to normalize the $`M_{\mathrm{vir}}T_X`$ relation accordingly. Applying the time-dependence factors derived for continuously forming clusters to the empirical mass-temperature relation in equation (7) thus gives
$$kT_X=(8.0\mathrm{keV})\left(\frac{M}{10^{15}h^1M_{}}\right)^{2/3}\left[\frac{(t_\mathrm{\Omega }/t)^{2/3}+1/m}{(t_\mathrm{\Omega }/t_0)^{2/3}+1/m}\right]$$
(17)
and
$$kT_X=(8.0\mathrm{keV})\left(\frac{M}{10^{15}h^1M_{}}\right)^{2/3}\frac{\xi _c(t)}{\xi _c(t_0)}.$$
(18)
for open and flat universes, respectively.
Figure 3 compares this empirical normalization with the normalizations of the $`M_{\mathrm{vir}}T_X`$ relations derived by Eke et al. (1996) and Voit & Donahue (1998) at $`z=0`$, assuming $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$. The dashed line indicates the normalization given in equation (7), which is presumed to be independent of $`\mathrm{\Omega }_\mathrm{M}`$. The curve labeled “late formation” shows the normalization of equation (8), derived from the top-hat, late-formation model. This normalization drops steadily with decreasing $`\mathrm{\Omega }_\mathrm{M}`$ because the density contrast factor $`\mathrm{\Delta }_{\mathrm{vir}}`$ grows smaller as $`\mathrm{\Omega }_\mathrm{M}`$ declines. Clusters modelled in this way are therefore less compact and cooler than one would expect from the critical density alone. When $`\mathrm{\Omega }_\mathrm{M}=1`$, this normalization is only 4% below the empirical value, but if $`\mathrm{\Omega }_\mathrm{M}=0.2`$, it lies 20% below this value, corresponding to mass discrepancies of 6% and 30%, respectively.
The behavior of the normalizations derived from continuous-formation models is more complicated. Voit & Donahue (1998) normalized these relations to the Eke et al. (1996) relation at $`\mathrm{\Omega }_\mathrm{M}=1`$ to simplify comparisons. For $`\mathrm{\Omega }_\mathrm{M}1`$, they are lower than at $`\mathrm{\Omega }_\mathrm{M}=1`$ for the same reason as in the late-formation model. However, if $`\mathrm{\Omega }_\mathrm{M}1`$, cluster formation happened long before $`z0`$, when the universe was considerably denser. Clusters are therefore denser and hotter than one would expect from the current critical density. As a result, the temperature normalization of the $`n=2`$ case deviates by less than 10% over the range $`0.2<\mathrm{\Omega }_\mathrm{M}<1.0`$. In the $`n=1`$ case the normalization is actually 18% higher than the empirical value (27% in mass) at $`\mathrm{\Omega }_\mathrm{M}=0.2`$.
This comparison illustrates why the procedure of normalizing the $`M_{\mathrm{vir}}T_X`$ relation to simulations is imperfect. In general, we expect this normalization to vary with $`\mathrm{\Omega }_\mathrm{M}`$ in a way that depends on $`n`$. Given that we have simulations for only a handful of cosmological models, how do we unambiguously normalize these relations? Furthermore, different choices for extrapolating this normalization to other values of $`\mathrm{\Omega }_\mathrm{M}`$ can lead to normalizations that differ by as much as 40% in temperature (60% in mass) at a given value of $`\mathrm{\Omega }_\mathrm{M}`$. Because of these uncertainties in the normalization of the $`M_{\mathrm{vir}}T_X`$ relation, it is important to understand how offsets in the normalization affect cosmological parameters derived from the cluster temperature function.
### 4.2 Normalization Offset and $`\sigma _8`$
The $`M_{\mathrm{vir}}T_X`$ relation is invoked twice in the usual derivation of $`\sigma _8`$ from the low-redshift temperature function. In both instances, an overestimate of the mass associated with a given temperature drives the best-fitting value of $`\sigma _8`$ higher. For example, a 50% offset in the mass normalization changes $`\sigma _8`$ by about 15%. Systematic uncertainties in the $`M_{\mathrm{vir}}T_X`$ relation therefore lead to systematic uncertainties in $`\sigma _8`$, limiting the usefulness of the temperature function as a tool to measure the perturbation amplitude.
The first place the $`M_{\mathrm{vir}}T_X`$ relation enters is in the conversion of the theoretical cluster mass function $`dn/dM`$ in equation (2) to the cluster temperature function
$$\frac{dn}{dT}(T,t)=\frac{3}{2}\left(\frac{2}{\pi }\right)^{1/2}\frac{\mathrm{\Omega }_\mathrm{M}\rho _{\mathrm{cr},0}}{TM(T,t)}\left|\frac{d\mathrm{ln}\sigma }{d\mathrm{ln}M}\right|\nu _c[M(T,t),t]\mathrm{exp}\{\nu _c^2[M(T,t),t)]/2\}.$$
(19)
An $`M_{\mathrm{vir}}T_X`$ relation that overestimates $`M(T)`$ by a fractional amount $`\delta _M`$ will underpredict the density $`dn/dT`$ by the same fractional amount. This source of error drives the best-fitting value of $`\nu _c`$ lower by a fractional amount $`\delta _\nu \delta _M/(\nu _c^21)`$.
Once $`\nu _c(T)`$ has been derived over a given range of temperatures, one can determine $`\sigma (T)=\delta _c(t_0)/\nu _c(T)`$. These $`\sigma `$ values will be too high by a fraction $`\delta _M/(\nu _c^21)`$ if there is a normalization offset. Conversion of $`\sigma (T)`$ to $`\sigma (M)`$ contributes another term to the systematic error budget. If $`M(T)`$ is overestimated, the mismapping of temperature to mass inflates $`\sigma _8`$ by a fractional amount $`\alpha \delta _M`$.
As an example of these effects, consider the derivation of $`\sigma _8`$ from the abundance of $`>5`$ keV clusters. Markevitch (1998) finds that the number density of such clusters at $`z0`$ is $`7.0\times 10^7h^3\mathrm{Mpc}^3`$. According to the mass-temperature relation in equation (7), the mass of a 5 keV cluster is $`5.0\times 10^{14}h^1M_{}`$, leading to an overall mass density $`\rho (>5\mathrm{keV})2.3\times 10^{32}h^2\mathrm{g}\mathrm{cm}^3`$ in such objects. Plugging this value into equation (1) and solving for $`\nu _c`$ yields $`\nu _c(5\mathrm{keV})3.2`$ for $`\mathrm{\Omega }_\mathrm{M}=1`$ and $`\nu _c(5\mathrm{keV})2.8`$ for $`\mathrm{\Omega }_\mathrm{M}=0.3`$, numbers that are consistent with more rigorous fits to cluster temperature data using a similar mass-temperature relation (Donahue & Voit 1999).
Conversion of these $`\nu _c`$ values to $`\sigma _8`$ values depends on the shape of the initial perturbation spectrum and the value of the virialization threshold $`\delta _c`$. For the purposes of this analysis, we will assume that $`\delta _c`$ is known perfectly, implying that $`\sigma (5.0\times 10^{14}h^1M_{})0.53`$ for $`\mathrm{\Omega }_\mathrm{M}=1.0`$ and $`\sigma (5.0\times 10^{14}h^1M_{})0.59`$ for $`\mathrm{\Omega }_\mathrm{M}=0.3`$. Extrapolating along the mass spectrum assuming, for example, that $`n=1.5`$ then leads to $`\sigma _80.51`$ for $`\mathrm{\Omega }_\mathrm{M}=1.0`$ and $`\sigma _80.76`$ for $`\mathrm{\Omega }_\mathrm{M}=0.3`$.
Now let us inflate the normalization of the mass-temperature relation by 50% in mass. The overall mass density in objects $`>5`$ keV rises to $`\rho (>5\mathrm{keV})3.5\times 10^{32}h^2\mathrm{g}\mathrm{cm}^3`$, yielding $`\nu _c(5\mathrm{keV})3.1`$ for $`\mathrm{\Omega }_\mathrm{M}=1`$ and $`\nu _c(5\mathrm{keV})2.7`$ for $`\mathrm{\Omega }_\mathrm{M}=0.3`$. Because 5 keV corresponds to $`7.5\times 10^{14}h^1M_{}`$ under the alternative normalization, we now have $`\sigma (7.5\times 10^{14}h^1M_{})0.54`$ for $`\mathrm{\Omega }_\mathrm{M}=1.0`$ and $`\sigma (7.5\times 10^{14}h^1M_{})0.61`$ for $`\mathrm{\Omega }_\mathrm{M}=0.3`$. Extrapolation to $`\sigma _8`$, again assuming $`n=1.5`$, gives $`\sigma _80.58`$ for $`\mathrm{\Omega }_\mathrm{M}=1.0`$ and $`\sigma _80.87`$ for $`\mathrm{\Omega }_\mathrm{M}=0.3`$. Note that these systematic changes, of order 15%, exceed the typically quoted measurement errors for $`\sigma _8`$ at a fixed value of $`\mathrm{\Omega }_\mathrm{M}`$.
Additional uncertainty in the derived value of $`\sigma _8`$ can arise from uncertainty in the slope $`n`$ of the perturbation spectrum. For example, the $`n2.3`$ slope derived from the Markevitch cluster sample (Markevitch 1998; Donahue & Voit 1999) leads to a considerably lower derived value of $`\sigma _8`$ when $`\mathrm{\Omega }_\mathrm{M}1`$. In the case of $`\mathrm{\Omega }_\mathrm{M}=0.3`$, a measurement of $`\sigma (5.0\times 10^{14}h^1M_{})0.59`$ extrapolates to $`\sigma _80.66`$, corresponding to clusters of temperature $`2.5\mathrm{keV}`$, below the temperature limit of the sample. Blanchard et al. (1999) have recently argued that the Markevitch sample is incomplete at low temperatures, implying $`n>2.3`$. That is probably why the best-fitting $`\sigma _8`$ values of Donahue & Voit (1999), based primarily on the Markevitch clusters, seem unusually low. (Errors on $`\sigma _8`$ quoted in that paper refer only to the statistical errors in $`\sigma _8`$ at the best-fitting value of $`\mathrm{\Omega }_\mathrm{M}`$.) Ideally, one would like to measure $`\sigma _8`$ from the number density of clusters at temperatures corresponding to the appropriate mass scale, but to do this, one first needs a reliable $`M_{\mathrm{vir}}T_X`$ relation, in addition to a well constrained value of $`\mathrm{\Omega }_\mathrm{M}`$.
The upshot of this analysis is that $`\sigma _8`$ values derived from the cluster temperature function contain systematic errors that depend on the mass-temperature relation. These systematic errors are currently comparable to the measurement errors. Until the $`M_{\mathrm{vir}}T_X`$ relation is better understood, $`\sigma _8`$ values derived from the cluster temperature function will have to be treated with caution. Conversely, predictions of cluster temperature functions that invoke $`\sigma _8`$ values derived from other kinds of data will also contain systematic errors owing to normalization uncertainties in the $`M_{\mathrm{vir}}T_X`$ relation.
### 4.3 Normalization Offset and $`\mathrm{\Omega }_\mathrm{M}`$
The systematic problems discussed above in relating $`\sigma _8`$ to $`dn/dT`$ need not lead to unwarranted pessimism about deriving $`\mathrm{\Omega }_\mathrm{M}`$ from the evolution of $`dn/dT`$. The key quantity in establishing the rate of cluster temperature evolution is not $`\sigma _8`$, but rather $`\nu _c(T,t_0)`$, whose systematic errors are considerably smaller, $`5`$% instead of $`15`$%. In order to evaluate the impact of a $`M_{\mathrm{vir}}T_X`$ normalization offset on predictions of temperature evolution, we will first analyze the case of $`\mathrm{\Omega }_\mathrm{M}=1`$, then consider how a normalization offset affects measurements of $`\mathrm{\Omega }_\mathrm{M}`$.
When $`\mathrm{\Omega }_\mathrm{M}=1`$, the time-dependent parts of $`M(T,t)`$ and $`\nu (T,t)`$ simplify to $`M(1+z)^{3/2}`$ and $`\nu _c(1+z)^{(23\alpha )/2}`$. The amount of cluster evolution at a fixed temperature $`T`$ can therefore be written as
$$C(T,z)=\frac{\frac{dn}{dT}(T,z)}{\frac{dn}{dT}(T,0)}=(1+z)^{(53\alpha )/2}\mathrm{exp}\{\frac{\nu _c^2(T,0)}{2}[(1+z)^{(23\alpha )}1]\}.$$
(20)
A fractional overestimate $`\delta _M`$ of cluster masses thus leads to an underestimate of the amount of evolution by a factor $`\mathrm{exp}\{\delta _M[(1+z)^{23\alpha }1]\}`$. If the overestimate of cluster masses is 50%, this factor amounts to a 20% underestimate of evolution at $`z=0.3`$ and a 70% underestimate at $`z=0.8`$ for $`n=1.5`$.
These uncertainties are relatively minor compared to the expected amount of evolution. For example, if $`n=1.5`$ and $`\nu _c(T,z=0)3.2`$, we expect $`C(T,z=0.3)0.3`$ and $`C(T,z=0.8)0.03`$. Larger values of $`\nu _c`$, characteristic of hotter clusters, lead to even more evolution. Because these evolution predictions are over an order of magnitude larger than the systematic errors, $`M_{\mathrm{vir}}T_X`$ normalization discrepancies do not seriously affect the conclusion that cluster temperature evolution rules out $`\mathrm{\Omega }_\mathrm{M}=1`$, particularly when clusters at $`z>0.5`$ are included.
Partially because of the potentially significant uncertainty in $`\sigma _8`$, certain authors have been cautious about the conclusions that can be drawn about $`\mathrm{\Omega }_\mathrm{M}`$ from cluster temperature evolution (Colafranceso, Mazzotta, & Vittorio 1997; Viana & Liddle 1999; Borgani et al. 1999b). However, maximum likelihood methods of determining $`\mathrm{\Omega }_\mathrm{M}`$ that compare the cluster temperature function at $`z0`$ directly with the cluster temperature function at higher redshifts (e.g., Henry 1997; Donahue & Voit 1999) can obtain stronger constraints on $`\mathrm{\Omega }_\mathrm{M}`$ because they are differential measurements in which much of the uncertainty in $`\sigma _8`$ and $`\delta _c`$ cancels. Figure 4 shows the cluster evolution predictions that result when a representative range of $`\sigma _8`$ values is considered. Here we allow $`0.5\sigma _8\mathrm{\Omega }_\mathrm{M}^{0.470.1\mathrm{\Omega }_\mathrm{M}}0.6`$. At the $`z0.3`$ redshift of the Henry (1997) clusters the prediction of the high-$`\sigma _8`$, $`\mathrm{\Omega }_\mathrm{M}=1`$ model is only a factor of two lower than the low-$`\sigma _8`$, $`\mathrm{\Omega }_\mathrm{M}=0.5`$ model, underscoring the need to identify systematic sources of uncertainty in $`\sigma _8`$ before deriving evolutionary predictions for clusters from it. However, Figure 5 paints a somewhat rosier picture. Here we allow $`\nu _{c0}=\nu _c(5\mathrm{keV},z=0)`$ to span a range that corresponds to a factor of two range in the mass normalization at 5 keV, or equivalently, a factor of two range in the number density of 5 keV clusters at $`z=0`$. The resulting systematic uncertainty in the best-fitting $`\mathrm{\Omega }_\mathrm{M}`$ is $`0.1`$.
### 4.4 Toward Greater Precision
When using Press-Schechter methods to model cluster evolution, one should always keep in mind that they are useful because they efficiently approximate numerical simulations. Our confidence in these methods is rooted in the fact that they reproduce the mass function of simulated clusters reasonably well. Less work has been done on comparisons of simulated cluster temperature functions with temperature functions derived from Press-Schechter mass functions using an $`M_{\mathrm{vir}}T_X`$ relation (e.g., Bryan & Norman 1998, Pen 1998).
Given the ambiguities surrounding the $`M_{\mathrm{vir}}T_X`$ relation and the very definition of a cluster’s mass, the most robust way to model the evolution of $`dn/dT`$ would seem to be with a version of the Press-Schechter formalism that describes the temperature function directly. In such a scheme, one would separate the function $`\nu _c(T,t)`$ that determines the evolution of $`dn/dT`$ into a temperature-dependent part $`\nu _c(T,0)`$ and a time-dependent part $`g(t;\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$. In principle, $`g(t;\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$ could be derived from a grid of simulated cluster temperature functions. However, creating such a simulation set would be very expensive. The simulation volume would need to be extremely large to obtain adequate statistics on rare, high-temperature clusters.
Maximum likelihood fits to cluster temperature surveys are essentially seeking the best-fitting $`\nu _c(T,z)=\nu _c(T,0)g(z;\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$. Figure 5 shows that this technique is fairly robust with respect to systematic uncertainties in the $`M_{\mathrm{vir}}T_X`$ normalization. Statistical uncertainties in the normalization and slope of $`\nu _c(T,0)`$ are handled naturally by the maximum likelihood method. Insofar as the dependence of $`\nu _c(T,z)`$ on cosmological parameters is accurate, this technique currently has the potential to deliver values of $`\mathrm{\Omega }_\mathrm{M}`$ that are accurate to $`0.1`$. However, it remains to be seen how accurately our assumed forms for $`g(z;\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$ reproduce the results of large-scale clustering simulations.
## 5 Summary
X-ray surveys of distant clusters are placing increasingly more stringent constraints on $`\mathrm{\Omega }_\mathrm{M}`$. The lack of extreme evolution in the cluster temperature function strongly indicates that $`\mathrm{\Omega }_\mathrm{M}<1`$. One of the crucial ingredients in placing such constraints on $`\mathrm{\Omega }_\mathrm{M}`$ is the $`M_{\mathrm{vir}}T_X`$ relation that converts cluster temperatures to cluster masses, enabling us to relate X-ray temperature surveys to theoretical models for cluster formation. If we are to extract accurate values of $`\mathrm{\Omega }_\mathrm{M}`$ from the larger cluster temperature surveys expected from Chandra and ASCA, we need to ensure that our $`M_{\mathrm{vir}}T_X`$ relation faithfully describes cluster evolution when coupled with Press-Schechter analysis. To that end, this paper has analyzed our current understanding of the cluster mass-temperature relation in an effort to identify the systematic errors it introduces into measurements of cosmological parameters.
We find that the usual derivation of the $`M_{\mathrm{vir}}T_X`$ relation, which assumes that clusters form by spherical top-hat collapse and that we are observing them immediately after they formed, is physically inconsistent. The rough agreement between the $`M_{\mathrm{vir}}T_X`$ normalization derived in this way and the normalization determined from numerical models of clusters is therefore somewhat coincidental. To obtain the proper normalization, one needs to account both for the fact that much of a cluster’s mass accreted well before the moment we are observing it and for the non-zero density at $`r_{200}`$, which requires a surface pressure term to be included in the virial theorem.
Because of these shortcomings of the spherical top-hat picture, we advocate a more realistic scenario for deriving the $`M_{\mathrm{vir}}T_X`$ relation in which clusters form quasi-contiuously. An expression for the $`M_{\mathrm{vir}}T_X`$ relation can be derived in the context of hierarchical merging, but the normalization of this relation depends on the concentration parameter $`c`$ of the cluster, which must be obtained from simulations. The primary advantage of this form for the $`M_{\mathrm{vir}}T_X`$ relation is that, unlike the spherical top-hat model, it properly reproduces the cessation of cluster evolution at late times if $`\mathrm{\Omega }_\mathrm{M}<1`$.
Given the systematic uncertainties in setting the normalization of the $`M_{\mathrm{vir}}T_X`$ relation, we have investigated their impact on the derivation of cosmological parameters from the cluster temperature function. Because two applications of the $`M_{\mathrm{vir}}T_X`$ relation are needed to extract $`\sigma _8`$, this parameter is particularly susceptible to uncertainties in the mass-temperature normalization: a mass-normalization uncertainty of 50% leads to a 15% uncertainty in $`\sigma _8`$. However, only a single application of the $`M_{\mathrm{vir}}T_X`$ relation is needed to extract $`\mathrm{\Omega }_\mathrm{M}`$, making it less vulnerable to normalization uncertainties. The systematic error in $`\mathrm{\Omega }_\mathrm{M}`$ owing to uncertainties in the $`M_{\mathrm{vir}}T_X`$ relation is $`0.1`$.
Improvements in our understanding of the $`M_{\mathrm{vir}}T_X`$ relation await comparisions of theoretically-derived cluster temperature functions with structure-formation simulations large enough to contain many hot clusters. In essence, the $`dn/dT`$ expression derived from the $`M_{\mathrm{vir}}T_X`$ relation and the Press-Schechter mass function is no more than an elaborate fitting formula for representing the results of simulations. The ideal mass-temperature relation will therefore be the one that reproduces simulated temperature functions or velocity-dispersion functions most accurately. Even better would be a fitting formula that gives $`dn/dT`$ directly without passing through murky intermediate steps involving ill-defined cluster masses. Until we have resolved these systematic uncertainties in deriving $`dn/dT`$, our constraints on $`\mathrm{\Omega }_\mathrm{M}`$ from cluster temperatures will not grow appreciably tighter.
Megan Donahue’s support, encouragement, and advice has been invaluable to the author, who would also like to acknowledge Pat Henry, Stefano Borgani, and the referee for helpful comments and NASA grants NAG5-3257 and NAG5-3208 for partial support. APPENDIX CONTINUOUS CLUSTER GROWTH IN $`\mathrm{\Omega }_\mathrm{M}<1`$ UNIVERSES In the spirit of Gunn & Gott (1972), we can idealize continuous cluster growth as occurring through the sequential collapse and virialization of an infinite series of concentric shells, each obeying the equation of motion
$$\ddot{R}=\frac{GM}{R^2}+\frac{\mathrm{\Lambda }R}{3},$$
(1)
where $`R`$ is the radius of the shell encompassing mass $`M`$, and $`\mathrm{\Lambda }`$ is the cosmological constant. The specific energy of the matter in the shell is then
$$ϵ=\frac{\dot{R}^2}{2}\frac{GM}{R}\frac{\mathrm{\Lambda }R^2}{6},$$
(2)
which remains constant until the shell virializes. If the shell ever reaches the critical radius $`R_{\mathrm{cr}}=(3GM/\mathrm{\Lambda })^{1/3}`$, cosmic repulsion dominates gravity from then on, and the shell never collapses.
We can cast the equation of motion for the shell into dimensionless form by defining $`x=R/R_{\mathrm{cr}}`$, $`\theta =\mathrm{\Lambda }^{1/2}t`$, and $`\xi =ϵR_{\mathrm{cr}}^2\mathrm{\Lambda }^1`$, so that
$$\frac{dx}{d\theta }=\left[\frac{2}{3x}+2\xi +\frac{x^2}{3}\right]^{1/2}.$$
(3)
A particular shell will collapse if $`x_0^3+6\xi x_0+2=0`$ for some $`x_0`$ in the range $`0x_0<1`$. Solving this cubic equation, we find
$$x_0=2^{3/2}|\xi |^{1/2}\mathrm{cos}\left(\frac{\alpha _\xi }{3}\frac{2\pi }{3}\right),$$
(4)
where $`\alpha _\xi `$ is defined by $`\mathrm{cos}\alpha _\xi =(8|\xi |^3)^{1/2}`$ with $`\pi /2\alpha _\xi \pi `$. The shell therefore reaches its maximum radius $`R_{\mathrm{cr}}x_0(\xi )`$ at a time $`\mathrm{\Lambda }^{1/2}\theta _0(\xi )`$, where
$$\theta _0(\xi )=_0^{x_0(\xi )}\frac{x^{1/2}dx}{\left(\frac{x^3}{3}+2\xi x+\frac{2}{3}\right)^{1/2}}.$$
(5)
Because of the symmetry of the motion, the shell collapses to the origin at $`t_c(\xi )=2\mathrm{\Lambda }^{1/2}\theta _0(\xi )`$, which we will take to be the time of virialization.
The overall scale factor of the universe obeys a similar equation of motion, for which the specific energy of the background matter is
$`ϵ_b`$ $`=`$ $`{\displaystyle \frac{\dot{a}^2}{2}}{\displaystyle \frac{H_0^2}{2}}\mathrm{\Omega }_\mathrm{M}a^1{\displaystyle \frac{H_0^2}{2}}\mathrm{\Omega }_\mathrm{\Lambda }a^2`$ (6)
$`=`$ $`{\displaystyle \frac{H_0^2}{2}}(1\mathrm{\Omega }_\mathrm{M}\mathrm{\Omega }_\mathrm{\Lambda }).`$ (7)
The universe changes from decelerating to accelerating when $`a=a_{\mathrm{cr}}(\mathrm{\Omega }_\mathrm{M}/2\mathrm{\Omega }_\mathrm{\Lambda })^{1/3}`$, and with the definition $`w=a/a_{\mathrm{cr}}`$, the dimensionless equation of motion for the background expansion becomes
$$\frac{dw}{d\theta }=\left[\frac{2}{3w}+2\xi _b+\frac{w^2}{3}\right]^{1/2},$$
(8)
where $`\xi _b=ϵ_ba_{\mathrm{cr}}^2\mathrm{\Lambda }^1`$.
Taking advantage of the formal similarities between these equations of motion, we can derive an expression for perturbation growth in an open universe. In the linear regime, $`\delta \rho /\rho =3\delta w/w`$, where $`\delta w=xw`$, and at any moment in time we have
$$_0^x\frac{y^{1/2}dy}{\left(y^3+6\xi y+2\right)^{1/2}}=_0^w\frac{y^{1/2}dy}{\left(y^3+6\xi _by+2\right)^{1/2}},$$
(9)
At early times, when $`x^3+2|6\xi x|`$ and $`w^3+2|6\xi _bw|`$, we obtain
$$\frac{\delta \rho }{\rho }=9(\xi _b\xi )\frac{(w^3+2)^{1/2}}{w^{3/2}}_0^w\frac{y^{3/2}dy}{(y^3+2)^{3/2}},$$
(10)
and at the earliest times ($`w1`$), this expression reduces to
$$\frac{\delta \rho }{\rho }=\frac{9}{5}(\xi _b\xi )w.$$
(11)
Perturbations at early times grow like $`1/(1+z)`$, as expected, and their amplitudes are proportional to the specific energy difference between the perturbation and the background. Note that if the universe is flat, the background specific energy $`\xi _b`$ vanishes, and the perturbation amplitude within a shell is directly proportional to the shell’s specific energy.
In the case of a vanishing cosmological constant, we can relate the perturbation amplitudes explicitly to their collapse times $`t_c`$. When the cosmological constant is very small, we have $`|\xi |1`$ and $`t_c(\pi /3\sqrt{2})|\xi |^{3/2}\mathrm{\Lambda }^{1/2}`$, so a shell that collapses and virializes at time $`t_c`$ carries with it a specific energy
$$ϵ=\frac{1}{2}\left(\frac{2\pi GM}{t_c}\right)^{2/3}.$$
(12)
If we define $`t_\mathrm{\Omega }=\pi \mathrm{\Omega }_\mathrm{M}/H_0(1\mathrm{\Omega }_\mathrm{M}\mathrm{\Omega }_\mathrm{\Lambda })^{3/2}`$, then at very early times,
$$\frac{\delta \rho }{\rho }\frac{3}{2}\frac{(12\pi )^{2/3}}{10}\left(\frac{t}{t_\mathrm{\Omega }}\right)^{2/3}\left[1+(t_\mathrm{\Omega }/t_c)^{2/3}\right].$$
(13)
According to Lacey & Cole (1993), the growth rate for linear perturbations in this limit is $`D(t)[(12\pi )^{2/3}/10](t/t_\mathrm{\Omega })^{2/3}`$, and we retrieve their expression for $`\delta _c(t)`$:
$$\delta _c(t)=\frac{3}{2}D(t)\left[1+(t_\mathrm{\Omega }/t)^{2/3}\right].$$
(14)
Time therefore enters the Press-Schechter formula via the parameter $`\nu _c\delta _c(t)/D(t)[1+(t_\mathrm{\Omega }/t)^{2/3}]`$.
In the case of a flat universe with $`\mathrm{\Omega }_\mathrm{M}<1`$, we have
$$\frac{\delta \rho }{\rho }=9\xi D(t)$$
(15)
where
$$D(t)=D[w(t)]=\frac{(w^3+2)^{1/2}}{w^{3/2}}_0^w\frac{y^{3/2}dy}{(y^3+2)^{3/2}},$$
(16)
in agreement with Eke et al. (1996). Inverting the function $`t_c(\xi )`$ yields the function $`\xi _c(t)`$ giving the specific energy $`ϵ(t)=R_{cr}^2\mathrm{\Lambda }\xi _c(t)`$ of a shell that collapses to the origin at time $`t`$. The collapse threshold becomes
$$\delta _c(t)=9\xi _c(t)D(t),$$
(17)
and time enters the Press-Schechter formula via $`\nu _c\delta _c(t)/D(t)\xi _c(t)`$.
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# Non local parity transformations and anomalies
## 1 Introduction
Because of the presence of ultraviolet divergences, most quantum field theory models require the use of a regularization scheme at some intermediate step in the renormalization program. The regularization procedure, to be effective, will modify the large momentum behaviour of the theory. In some cases, there are symmetries which are particularly sensitive to that modification, leading to the existence of anomalies, namely, symmetries which remain broken even after the regularization is removed. In other words, performing a symmetry transformation and renormalizing are non-commuting operations.
In the context of the functional integral quantization, anomalies are traditionally attributed to the non invariance of the integration measure under a classical symmetry transformation. The non-invariance of the integration measure shows up in that the Jacobian corresponding to the classical symmetry transformation is non trivial and, in general,ill-defined. When a proper regularization is introduced, this Jacobian gives rise to an anomalous term in the effective action .
In reference , an alternative procedure to study anomalies within the functional integral approach was presented. It amounts to considering the functional integral for a regularized action, and to restrict the study to symmetry transformations that leave that action invariant. For the cases of the conformal and chiral anomalies, we have shown that those symmetry transformations do exist, and that they tend to the usual ones when the cutoff is removed. The transformations are necessarily non local, but on a length scale of the order of the inverse of the cutoff. The outcome of performing those symmetry transformations is that, while the regularized classical action is indeed invariant (even for a finite cutoff), the integration measure is not. Moreover, this non invariance is explicit, since the corresponding Jacobian determinant is well defined and different from $`1`$, for any value of the cutoff. Contact with the usual anomalies is made by taking the infinite cutoff limit, what reproduces the known results.
In this letter, we will point out that this procedure can also be applied to the case of the parity anomaly in $`2+1`$ dimensions, since the regularized action does have a non local symmetry, which coincides with the classical parity symmetry in the infinite cutoff limit. Moreover, the associated Jacobian for this discrete transformation is finite, and reproduces the parity anomaly, a Chern-Simons term in the effective action -. This shows that the procedure we have developed in ref. can also be applied to the case of a discrete symmetry.
The organization of this letter is as follows: In section 2, we discuss the classical parity transformations and its generalization to the (quantum) regularized case. The generalized transformations are then applied to the calculation of the parity anomaly in section 3. Section 4 contains our conclusions.
## 2 Generalized parity transformations
Let us consider massive Dirac fermions coupled to an external Abelian gauge field $`A_\mu `$, in $`2+1`$ dimensions. The (Minkowski spacetime) action for the system is
$$S_f[A]=d^3x_f$$
(1)
where the Lagrangian density $`_f`$ is given by
$$_f=\overline{\psi }(x)(i/e/A(x)m)\psi (x).$$
(2)
In our conventions, $`g_{\mu \nu }=\mathrm{diag}(1,1,1)`$, and the Dirac matrices verify $`\{\gamma _\mu ,\gamma _\nu \}=2g_{\mu \nu }`$, and $`\gamma _\mu ^{}=\gamma ^\mu `$.
We recall that, in $`2+1`$ dimensions, parity transformations are, in fact, tantamount to spatial reflections:
$$x=(x^0,x^1,x^2)\stackrel{𝒫}{}x^P=(x^0,x^1,x^2)$$
(3)
since a spatial inversion is a rotation in $`\pi `$ when the number of space coordinates is even.
The gauge and spinor fields transform according to the rules:
$$A(x)=(A^0(x),A^1(x),A^2(x))A^P(x^P)=(A^0(x),A^1(x),A^2(x))$$
$$\psi (x)\psi ^P(x^P)=\gamma _1\psi (x)\overline{\psi }(x)\overline{\psi ^P}(x^P)=\overline{\psi }(x)\gamma _1.$$
(4)
Hence, under $`𝒫`$, the classical action transforms as follows:
$$S[A]S_f^P[A]=d^3x^P\overline{\psi }_P(x_P)(i/_Pe/A_P(x_P)m)\psi _P(x_P)$$
$$=d^3x\overline{\psi }(x)(i/e/A(x)+m)\psi (x).$$
(5)
Writing the mass dependence of the action explicitly, we have the relation:
$$S_f^P(m)=S_f(m).$$
(6)
This indicates, of course, that the massless classical theory is parity invariant.
Regarding the regularized action, although not every regularization method can be implemented in terms of an action, there are many important cases where this can be done. Examples of those are the Euclidean cutoff, lattice, and Pauli-Villars (PV) regularization methods. We shall use the PV method, since it maintains most of the symmetries, except parity (see for a thorough discussion in odd-dimensional spaces). This greatly simplifies the discussion. In our $`2+1`$ dimensional example, the functional integral is rendered convergent by the addition of just one bosonic regulator spinor field $`\varphi `$, whose mass $`\mathrm{\Lambda }`$ plays the role of a cutoff. The regularized action $`S_f^r`$ is:
$$S_f^r=id^3x[\overline{\psi }(x)(/+ie/A)\psi (x)+\overline{\varphi }(x)(/+ie/A+i\mathrm{\Lambda })\varphi (x)].$$
(7)
It is possible to write the regularized action (7) in terms of just one fermionic field $`\mathrm{\Psi }`$, although at the expense of equipping that field with a non local action. To that end, we consider the regularized functional integral $`𝒵^r[A]`$ corresponding to (7):
$$𝒵^r[A]=𝒟\psi 𝒟\overline{\psi }𝒟\varphi 𝒟\overline{\varphi }\mathrm{exp}\{d^3x[\overline{\psi }(x)/D\psi (x)+\overline{\varphi }(x)(/D+i\mathrm{\Lambda })\varphi (x)]\}$$
(8)
where $`/D(/+ie/A)`$.
We first integrate out the regulator in (8):
$$𝒵^r[A]=𝒟\psi 𝒟\overline{\psi }\stackrel{1}{det}[1\frac{i/D}{\mathrm{\Lambda }}]\mathrm{exp}\left(d^3x\overline{\psi }/D\psi \right)$$
(9)
(we have neglected an irrelevant constant factor). We then make the change of variables:
$`\psi (x)`$ $`=`$ $`\left[1{\displaystyle \frac{i/D}{\mathrm{\Lambda }}}\right]^{\frac{1}{2}}\mathrm{\Psi }(x)`$ (10)
$`\overline{\psi }(x)`$ $`=`$ $`\overline{\mathrm{\Psi }}(x)\left[1{\displaystyle \frac{i/D}{\mathrm{\Lambda }}}\right]^{\frac{1}{2}}`$ (11)
under which the measure transforms according to:
$$𝒟\psi 𝒟\overline{\psi }=𝒟\mathrm{\Psi }𝒟\overline{\mathrm{\Psi }}det[1\frac{i/D}{\mathrm{\Lambda }}].$$
(12)
The result of the two previous steps is that we may rewrite (9) in the equivalent form:
$$𝒵^r[A]=𝒟\mathrm{\Psi }𝒟\overline{\mathrm{\Psi }}\mathrm{exp}(S_f^{nl})$$
(13)
where $`S_f^{nl}`$ denotes a non-local form of the regularized action,
$$S_f^{nl}=d^3xd^3y\overline{\mathrm{\Psi }}(x)𝒟(x,y)\mathrm{\Psi }(y).$$
(14)
with
$$𝒟(x,y)=\left[\frac{/D}{1{\displaystyle \frac{i/D}{\mathrm{\Lambda }}}}\right](x,y).$$
(15)
We have found it convenient to adopt a ‘bracket’ like notation to write the action (14),
$$S_f^{nl}=\mathrm{\Psi }|𝐃|\mathrm{\Psi },𝐃=\frac{/D}{1{\displaystyle \frac{i/D}{\mathrm{\Lambda }}}},$$
(16)
since it avoids writing operators kernels and integrations explicitly. Note that the ‘bra’ includes the $`\gamma _0`$ factor for the adjoint field.
To understand the behaviour of (14) under parity transformations, we previously need to know the $`𝐃`$ operator transformation properties. Denoting by $`/D(x,y)`$ the kernel of $`/D`$, we see that
$$\gamma _1/D^P(x^P,y^P)\gamma _1=/D(x,y).$$
(17)
We can also express this result as:
$$\gamma _1/D^P\gamma _1=/D,$$
(18)
what in turn implies for $`𝐃`$:
$$\gamma _1𝐃^P(\mathrm{\Lambda })\gamma _1=𝐃(\mathrm{\Lambda }).$$
(19)
The parity transformed non local regularized action becomes
$$S_f^{nlP}=\mathrm{\Psi }^P|𝐃^P(\mathrm{\Lambda })|\mathrm{\Psi }^P=\mathrm{\Psi }^P|\gamma _1𝐃(\mathrm{\Lambda })\gamma _1|\mathrm{\Psi }^P,$$
(20)
where we need to plug in the parity transformed of $`\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}`$. The transformation rules for the new fields, can be obtained as follows: from the relation between $`\mathrm{\Psi }`$ and $`\psi `$, we learn that
$$\psi ^P(x^P)=\left[1\frac{i/D^P}{\mathrm{\Lambda }}\right]^{\frac{1}{2}}\mathrm{\Psi }^P(x^P)$$
(21)
and
$$\psi ^P(x^P)=\gamma _1\left[1+\frac{i/D}{\mathrm{\Lambda }}\right]^{\frac{1}{2}}\gamma _1\mathrm{\Psi }^P(x^P).$$
(22)
On the other hand, we have the relation
$$\psi ^P(x^P)=\gamma _1\psi (x)=\gamma _1\left(1\frac{i/D}{\mathrm{\Lambda }}\right)^{\frac{1}{2}}\mathrm{\Psi }(x).$$
(23)
Then,
$$\mathrm{\Psi }^P(x^P)=\gamma _1\frac{\sqrt{1+\frac{i/D}{\mathrm{\Lambda }}}}{\sqrt{1\frac{i/D}{\mathrm{\Lambda }}}}\mathrm{\Psi }(x)=\gamma _1x|\frac{𝐃(\mathrm{\Lambda })}{\sqrt{𝐃(\mathrm{\Lambda })𝐃(\mathrm{\Lambda })}}|\mathrm{\Psi }.$$
(24)
We obtain an analogous relation for $`\overline{\mathrm{\Psi }}(x)`$. Using the compact notation, we see that
$$|\mathrm{\Psi }^P=\gamma _1\frac{𝐃(\mathrm{\Lambda })}{\sqrt{𝐃(\mathrm{\Lambda })𝐃(\mathrm{\Lambda })}}|\mathrm{\Psi }$$
(25)
$$\mathrm{\Psi }^P|=\mathrm{\Psi }|\frac{𝐃(\mathrm{\Lambda })}{\sqrt{𝐃(\mathrm{\Lambda })𝐃(\mathrm{\Lambda })}}\gamma _1,$$
(26)
and these are the generalized parity transformations we were looking for. We note that they tend to the standard parity transformations when the cutoff is removed ($`\mathrm{\Lambda }\mathrm{}`$).
It is immediate to verify that, with these transformation rules, the non-local form of the regularized action remains invariant:
$$S_f^{nlP}=S_f^{nl}.$$
(27)
Being the Pauli-Villars fields integrated out exactly, the effective action (14) incorporates the complete effect of the regularization, in the sense that the contribution to $`S_f^{nl}`$ of the determinant associated with the Pauli-Villars field is totally included. For example, were one to use this effective action to compute the vacuum polarization, the corresponding diagrams will automatically be convergent to all orders.
## 3 Parity anomaly
Although the regularized action is invariant under the generalized parity transformations, the fermionic measure acquires a non trivial Jacobian:
$$𝒟\mathrm{\Psi }_P𝒟\overline{\mathrm{\Psi }}_P=𝒟\mathrm{\Psi }𝒟\overline{\mathrm{\Psi }}𝒥$$
(28)
where
$$𝒥=det\left[\left(\frac{\sqrt{D(\mathrm{\Lambda })D(\mathrm{\Lambda })}}{D(\mathrm{\Lambda })}\right)^2\right]=det\left(\frac{i/D\mathrm{\Lambda }}{i/D+\mathrm{\Lambda }}\right).$$
(29)
The parity anomaly is obtained from the infinite-$`\mathrm{\Lambda }`$ limit of this Jacobian. This result is known, and coincides of course with (twice) the leading term in a derivative expansion of the effective action $`I_f^{\mathrm{eff}}[A,\mathrm{\Lambda }]`$: We see that:
$$\mathrm{exp}(iI_{\mathrm{eff}}[A,\mathrm{\Lambda }]iI_{\mathrm{eff}}[A,\mathrm{\Lambda }])=det\left(\frac{/D+i\mathrm{\Lambda }}{/Di\mathrm{\Lambda }}\right)=𝒥.$$
(30)
$$I_{\mathrm{eff}}[A,\mathrm{\Lambda }]=i\mathrm{ln}det(/D+i\mathrm{\Lambda }).$$
(31)
Hence the Jacobian is expressed as a function af the effective action. In the $`\mathrm{\Lambda }\mathrm{}`$ one obtains -:
$$I_{\mathrm{eff}}[A,\mathrm{\Lambda }]\frac{\mathrm{\Lambda }}{|\mathrm{\Lambda }|}S_{CS},$$
(32)
where the $`CS`$ action $`S_{CS}`$ is defined by:
$$S_{CS}=\frac{e^2}{4\pi }d^3xϵ^{\mu \nu \rho }A_\mu (x)_\nu A_\rho (x).$$
(33)
Thus we can write for the Jacobian:
$$𝒥=\mathrm{exp}\left\{\pm 2iS_{CS}[A]\right\},$$
(34)
which is the parity anomaly result. It is not renormalized by higher order correction, as a consequence of the Coleman-Hill theorem .
For the case of a continuous symmetry transformation, the Fujikawa Jacobian resulting from the associated infinitesimal change of variables can be used to calculate the potentially anomalous divergence of the corresponding Noether current. This procedure does not, of course, have a direct parallel in the present, discrete symmetry case. Nevertheless, from the knowledge of the Jacobian (34), one can make explicit the anomalous behavior of the fermion current under a parity transformation. Indeed, from eq.(13), we can write
$$𝒵^r[A^P]\mathrm{exp}\left(I_{eff}^r[A^P]\right)=𝒟\mathrm{\Psi }𝒟\overline{\mathrm{\Psi }}\mathrm{exp}(S_f^{nl}[A^P,\overline{\mathrm{\Psi }},\mathrm{\Psi }]).$$
(35)
Now, changing the fermion variables as in (25)-(26), and using the invariance (27) of the non local action, we see that
$$𝒵^r[A^P]=𝒥[A^P]𝒟\mathrm{\Psi }^P𝒟\overline{\mathrm{\Psi }}^P\mathrm{exp}(S_f^{nl}[A^P,\overline{\mathrm{\Psi }}^P,\mathrm{\Psi }^P])=𝒥[A^P]𝒵^r[A]$$
(36)
or
$$\mathrm{exp}(I_{eff}^r[A^P]+I_{eff}^r[A]))=𝒥[A^P].$$
(37)
Differentiation of the regularized effective action with respect to the gauge field $`A`$ yields the ground state current $`j_\mu [A]`$ so that one has
$$ej_\mu [A^P]+ej_\mu [A]=\frac{\delta \mathrm{log}𝒥[A^P]}{\delta A^\mu }.$$
(38)
Using relations (4) for the l.h.s. and (34) for the r.h.s., we can then write the anomalous, parity-odd vacuum current:
$$j_\mu [A]^{odd}=\frac{1}{2e}\frac{\delta \mathrm{log}𝒥[A^P]}{\delta A^\mu }=\pm i\frac{e}{4\pi }d^3xϵ^{\mu \nu \rho }_\nu A_\rho (x).$$
(39)
which is the well-known result first obtained in .
It is immediate to verify that all the steps leading the parity anomaly for the Abelian case can be generalized to the non-Abelian case as well, with the only modification that the Jacobian is now related to the non-Abelian massive determinant. The result is of course that this Jacobian is the exponential of ($`2i`$ times) the non-Abelian Chern-Simons action.
## 4 Conclusions
We have shown that the parity anomaly in 2+1 dimensions in a theory of massless fermions coupled to an external gauge field can be obtained from the Jacobian for a generalized symmetry transformation of the regularized theory. This procedure has the virtue of disentangling the symmetries and infinities, by looking for transformations which are symmetries of the regularized action. This avoids involving modes with an arbitrarily high momentum into the symmetry transformation, and the bonus is that the resulting Jacobian automatically takes care of the anomaly. It is worth noting that it is precisely the fact that the regularization breaks the symmetry what renders the symmetry transformations of the regularized action non-local. Non-anomalous local symmetries have, in this setting, a distinctive property: they are always local, even when acting on the regularized action.
Our calculations required the use of a non-local form of the regularized action and we used ‘regularized’ fermion fields $`\mathrm{\Psi }`$. The non-locality of the action, however, is hidden in the non-physical region of momenta above the cutoff. For example, when $`A=0`$, the $`𝐃`$ operator has a kernel:
$$/D(x,y)=/(1+\frac{i/}{\mathrm{\Lambda }})\frac{d^3k}{(2\pi )^3}\frac{e^{ik(xy)}}{1(k/\mathrm{\Lambda })^2}$$
(40)
which is non-local on a scale $`1/\mathrm{\Lambda }`$.
## Acknowledgements
C. D. F. would like to acknwledge Prof. F. D. Mazzitelli, for useful conversations and correspondence.
M. L. C. is supported by a CNEA fellowship. C. D. F. is a member of CONICET, and is partially supported by ANPCyT (PICT 97/1040) and Fundación Antorchas grants. F. A. S is partially supported by CICBA as Investigador, and through grants CONICET (PIP 4330/96), and ANPCyT (PICT 97/2285).
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# Resumming the large-𝑁 approximation for time evolving quantum systems
## I Introduction
The need to understand quantum systems in real time in a quantum field theoretic setting arose from attempts to understand various early universe scenarios. These scenarios are based on the evolution of scalar fields either through their role as inflation fields or as topological defect forming fields. One would like to understand the quantum evolution of these fields rather than rely on unjustified treatments based on studying their classical evolution. The study of the “slow rollover” transition in an upside down harmonic approximation by Guth and Piref:Guth was the first attempt to understand whether classical approximations could be justified. However, one really needed to go beyond the harmonic approximation to address the nonlinear aspects of double well (and Mexican hat) potentials. These non-linear aspects effect production of topological defects as well as the nature of the oscillation at the bottom of the well which causes reheating.
Our ultimate goal is to be able to describe accurately over relevant time periods the nonlinear aspects of quantum field theory evolutions. Although in one-dimensional quantum mechanics, one can rely on a numerical solution of the Schrödinger equation to understand the time evolution of the system accurately over long time periods, in field theory contexts the numerical solution of the functional Schrödinger equation is presently beyond the reach of the largest computers. One important question is how to decrease the number of degrees of freedom in a manner consistent with certain physical requirements such as conservation of energy, preservation of positivity and boundedness of expectation values. Although this is guaranteed in variational approximations, approximations based on various truncation schemes, whether perturbative or non-perturbative in nature often fail to preserve these physical requirements. For example, naive truncations of the coupled Green functions equations beyond the truncation at the two-point function level lead to secular behavior (unboundedness at late times). This is also true for the $`1/N`$ expansion which is derivable from an effective action. The second question is, after guaranteeing these properties, how accurately have we described the time evolution.
The simplest truncations of the field theory have been based on gaussian variational methodsref:Hartree ; ref:GV , or the related leading order in large-$`N`$ approximation (LOLN) ref:LOLN ; ref:ctpN . These two approximations can be shown to be equivalent to a classical Hamiltonian dynamics for the variational parameters (or equivalently the Green functions) which leads to probability conservation at the quantum level so that the results always lead to conserved energy, and positive and bounded expectation values. Unfortunately, hard scatterings which lead to thermalization are ignored so that important physics is left out. The approximation also is numerically inaccurate after a few oscillations in quantum mechanical applications, unless the anharmonic coupling constant $`g`$ divided by the number of fields $`N`$ in an $`O(N)`$ model is quite small. In this paper we will be comparing our methods of going beyond mean field theory (Hartree or large-$`N`$) with exact numerical simulations of a quantum mechanical $`O(N)`$ model. In this way we can see how accurate the approximations are as a function of $`N`$ as well as study numerically if the approximation maintains the various physical requirements we posit, such as boundedness and positive definiteness of expectation values. The reason for using this quantum mechanical model is that exact simulations can be done at *all* $`N`$, so that accuracy of the method as a function of the parameter $`1/N`$ can be studied. By restricting ourselves to a quantum mechanics problem we unfortunately will not be able to study questions of thermalization. A complementary approach has been undertaken by Aarts, Bonini, and Wetterichabw where they consider *classical* 1+1 dimensional $`\varphi ^4`$ field theory (for $`N=1`$). There one can look at some aspects of classical thermalization (as long as one keeps a cutoff because of the Raleigh-Jeans divergence) but one is restricted to low values of $`N`$ so one cannot study the $`N`$ dependence of the result. Also one cannot study the *quantum* aspects of the problem. In the above paper, Aarts *et.al.* study a truncation of the Green functions at the four-point level, which is known to lead to unboundedness and secularity in quantum mechanical (as well as classical) applications. It will be interesting in the future to apply the approximations we are using here to classical 1+1 dimensional $`\varphi ^4`$ to see if, and how well, they describe the thermalization.
There are several ways of approaching the problem of thinning the degrees of freedom of the quantum field theory. One of the earliest was based on making a variational approximation to the functional Schrödinger equation. The variational approach has the advantage of leading to a Hamiltonian dynamical system for the variational parameters as well as to a density matrix which has positivity properties. Energy conservation and positivity and boundedness of expectation values are automatically guaranteed. However, even for the simple problem of the quantum roll, the gaussian, or time dependent Hartree approximation, studied by Cooper, Pi and Stancioffref:Hartree , and improvements which are based on trial wave functions of the form of a polynomial times a gaussianref:CC , were found to be only accurate for relatively short time periods (one or a few oscillations) when compared to the exact numerical solution of the Schrödinger equation. In quantum mechanics, except for exceptional situations, the wave function in multiwell situations gets very complicated very quickly and is not easily described by a small number of variational parameters.
A second approach has been a direct $`1/N`$ expansion of the path integral in the Schwinger-Keldysh-Bakshi-Mahanthappa closed time path formalismref:CTP . In this approach the connected Green functions have the property that they start at order $`G_{2n}1/N^{n1}`$. Thus if we retain only a certain order in the expansion, there is a truncation in the order of Green functions retained. This approach was applied recently to the quantum roll problemref:paper1 and was found to suffer from the secularity problem — although the short time behavior of the result was improved by including $`1/N`$ corrections, an exact reexpansion in terms of $`1/N`$ leads to corrections in the Green functions that are of the form $`\pm t/N`$ and so the individual corrections become unbounded as well as non positive definite. In this approach, although energy is conserved, individual contributions are not positive definite and unphysical behavior is found.
A third approach has been to consider the complete set of equal time Green functions. These obey first order local equations in time, as in the Schrödinger approach. This approach has been nicely systematized and an equation for the generating functional obtained by Wetterich and collaborators in a series of papers ref:EQT . However, naive truncations of the equal time Green function hierarchy again have the problem that although there is a conserved energy, one cannot show that this truncation (except at the two-point level) corresponds to a positive definite probability so that expectation values are not necessarily bounded or positive definite. Truncated at the two point function level, this approach is identical to the Hartree approximation. However, simulations based on truncations assuming 6th order or 8th order 1-PI graphs, could be set to zero, were carried out for the $`O(N)`$ $`x^4`$ th oscillator problem, and secularity was found for many choices of initial conditionsref:bett . So we, as quantum field theorists, having entered the domain of nonequilibrium phenomena, are now beset with all the problems faced by our plasma and condensed matter brethren more than 40 years ago!
In both quantum and classical many-body systems, the dynamical equations are an infinite hierarchy of coupled equations which relate given ensemble averages, whether classical or quantum, to successively more complicated ones. To make the solution of this hierarchy possible, some truncation scheme is necessary. Most naive truncation schemes which, for example, just truncate the hierarchy of coupled correlators at a particular order, do not preserve various physical properties required of the system — such as positivity of the spectral components of the Green function and conservation of probability. A corollary of this is that in most perturbation schemes, secularity arises quickly with each term in the perturbation series, growing with higher powers of the time $`t`$. In his seminal paper of 1961, Robert Kraichnanref:Robert discussed in detail the key issues and obtained a partial solution to the problem by demanding that the approximations one should use should correspond to some physically realizable dynamical system. This would guarantee positivity and secularity would be avoided. The reason why variational approximations avoid these problems is exactly because they lead to a Hamiltonian dynamical system for the variational parameters (which are related to equal time correlation functions). He also discussed scenarios where particular classes of graphs, which contained the relevant dynamics, are summed and he suggested some physically motivated approximations which did not suffer from any diseases. In field theory one rarely has the parameter control to make such guesses, however some progress in QCD has been made by summing hard thermal loopsref:SDQCD , which already tells us some of the graphs that we want to include. In plasma physics, one wants to make sure that the approximation to the dynamics is robust enough so that the photon propagator includes polarization effects, which give Debye screening. This is related to the hard thermal loop summation in QCD.
To find resummation schemes that avoid the secularity problem we will rely on the experience of our many-body and plasma physics friends. To calculate the conductivity of a non-relativistic plasma, it is known what graphs are necessary to sum in order to get agreement with experimental resultsref:plasma1 ; ref:plasma2 . Basically the conductivity is found from the vertex function which must satisfy an integral equation which sums ladders of the Debye screened photon propagator. The two approximations we will discuss here will differ on whether the equivalent of the Debye screened photon propagator for the anharmonic oscillator is treated in lowest order in mean field theory, or is self-consistently determined. In studying the conductivity of a relativistic plasma the first approximation has the advantage of obeying the correct Ward identities (but violating energy conservation to order $`1/N`$) whereas the second preserves energy conservation but violates Ward identities (to order $`1/N^2`$). Here we are not studying QED, and the Ward identities of the $`O(N)`$ model for the quantum mechanics problem are much simpler than those of QED and energy conservation is a more important constraint on the accuracy of the answer. We will include both approximations here mainly because of the recent interest in the gauge invariant approximation for the relativistic plasmaref:emil , and also because in truncations of Schwinger-Dyson equations, it is often too difficult to solve for the photon propagator self consistently, and so one is often forced to try the more drastic approximation of using the mean field propagator in the resummation scheme. By studying this approximation in a quantum mechanics problem we will see the shortcomings of such an approach.
In what follows we will discuss two approaches to obtaining the above two truncations of the exact Schwinger-Dyson equation and apply them to the problem of the quantum roll — the long time behavior of $`N`$ coupled anharmonic oscillators with “radial” symmetry in an $`N`$-dimensional space. This particular problem has been studied by us previouslyref:paper1 exactly and in the next to leading order in the large-$`N`$ approximation (NLOLN) and is interesting because exact numerical solutions can be found for arbitrary $`N`$. What we found previously, is that for the parameter set studied ($`g1,M^2=2`$), the next to leading order in large-$`N`$ contributions became unbounded for $`N<21`$. For larger $`N`$, where the approximation was physical, it had the failing that it was unable to track the spreading of the exact wave function which led to the envelope of the oscillations found for $`\widehat{x}^2(t)`$ contracting at late times and then reexpanding. A related study of large-$`N`$ for quantum mechanics in the context of the equal time correlators by Bettencourt and Wetterichref:bett , also displayed growing modes for various initial conditions.
The resummation presented here will allow one to track the contraction for some period, but at later times it also fails in that it leads to small oscillations about a fixed point value. In field theory settings, where one hopes that this approximation will lead to thermalization, optimistically this fixed point behavior will become physical and be related to thermal equilibration. Whether this is true or not can be checked by studying this approximation for classical evolutions averaged over a distribution of initial conditions described by a an initial probability distribution in phase space.
In what follows we will present numerical solutions for the quantum roll problem for the $`O(N)`$ model, and compare them to these two different approximations to the Schwinger-Dyson equations, which sum infinite numbers of leading order and next to leading order in $`1/N`$ graphs. Our approach will be to introduce a composite “field”
$$\chi =\frac{g}{2}\left(\underset{i=1}{\overset{N}{}}x_i^2r_0^2\right),$$
which is treated on equal footing to the field $`x`$. By doing that, the Schwinger-Dyson equations for the theory will have the same topology as those of QED with $`x`$ playing the role of the electron and $`\chi `$ the role of the photon. At leading order in large $`N`$ in N-flavor QED, one sums all the fermion loop vacuum polarization corrections to the photon propagator which gives the Debye screening. Here the bare photon propagator is replaced by a local interaction in the graphs for the $`\chi `$ propagator in LOLN. The next consideration, important for charged plasmas, is that to obtain reasonable agreement with experiments on the conductivity of the plasma, the vertex function must sum all the ladders with the Debye screened propagator as the kernel in the integral equation. The two resummation schemes which we discuss in this paper both have this property.
The approximation which we call the bare vertex approximation (BVA), uses the full Green function for $`x`$ as well as the full Green function for $`\chi `$ in a 2-PI Hartree graph contribution to the effective action. This is in contrast to an earlier scheme for going beyond $`1/N`$ref:Hu using the 2-PI formalism which is based only on the $`x`$ Green functions. The BVA approximation sums an infinite Geometric series of 2-PI graphs of the single field formalism. Recent simulations in a toy 1+1 dimensional scalar field theoryref:berges show that the approximation described inref:Hu already has the ability to thermalize arbitrary initial conditions, so we are confident that the BVA approximation will also have that feature when applied to a field theory problem. The BVA can also be obtained by setting the full vertex function to unity in the Schwinger-Dyson equations for the one- and two-point functions with external sources hence the origin of its name. The second approximation we will study, which we call the dynamic Debye screening approximation (DDSA), makes the further assumption that the full $`\chi `$ propagator can be replaced by the lowest order in $`1/N`$ composite field propagator in all the integral equations. The main interest in the DDSA results from it being the lowest order resummation scheme that *exactly* preserves QED Ward identities. Both these approximations are free from the difficulties found in the perturbative $`1/N`$ expansion, which we display for comparison. We find that the BVA is accurate at modest times $`25`$ oscillations when $`N>10`$. At later times it settles down to oscillating about an unphysical fixed point. The DDSA approximation violates energy conservation at order $`1/N`$ and as a result becomes inaccurate after several oscillations. In spite of this, it is numerically more accurate for a longer period of time than the Hartree approximation at small and modest values of $`N`$.
It should be kept in mind that quantum mechanics and quantum field theory are very different. For example, in the quantum mechanics application discussed here, the graphs of the $`O(1/N)`$ corrections do not correspond to interparticle collisions (as they do in field theory) since we are restricting ourselves to one-particle quantum mechanics. Nevertheless quantum mechanical examples provide excellent test beds for key issues such as positivity violation, boundedness, and late time accuracy of the approximations. It is precisely these questions that we are hoping to understand in this paper.
## II The $`O(N)`$ model
The classical Lagrangian for the $`O(N)`$ model of $`N`$ non-linear oscillators is given by:
$$L(x,\dot{x})=\frac{1}{2}\underset{i=1}{\overset{N}{}}\dot{x}_i^2\frac{g}{8N}\left(\underset{i=1}{\overset{N}{}}x_i^2r_0^2\right)^2.$$
(1)
The Schrödinger equation for this problem is given by:
$$\left\{\frac{1}{2}\underset{i=1}{\overset{N}{}}\frac{^2}{x_i^2}+V(r)\right\}\psi (x,t)=i\frac{\psi (x,t)}{t},$$
(2)
where $`V(r)`$ is a potential of the form
$$V(r)=\frac{g}{8N}\left(r^2r_0^2\right)^2,r^2=\underset{i=1}{\overset{N}{}}x_i^2.$$
(3)
For the quantum roll problem there is spherical symmetry. This means that we can assume a solution of the form $`\psi (r,t)=\varphi (r,t)/r^{(N1)/2}`$, in which case the time dependent Schrödinger equation for $`\varphi (r,t)`$ reduces toref:BlazotRipka :
$$\left\{\frac{1}{2}\frac{^2}{r^2}+U(r)\right\}\varphi (r,t)=i\frac{\varphi (r,t)}{t},$$
(4)
with an effective one dimensional potential $`U(r)`$ given by
$$U(r)=\frac{(N1)(N3)}{8r^2}+\frac{g}{8N}\left(r^2r_0^2\right)^2.$$
(5)
It is this equation that we will solve numerically to obtain exact numerical solutions as a function of $`N`$. $`U(r)`$ has a minimum at $`r=r_{\text{min}}`$. In our simulations, we have fixed our mass scale $`M^2`$, defined as the second derivative of $`U(r)`$ at the minimum, to have a value of 2, independent of $`N`$.
Returning to the Lagrangian formulation, it is useful for the purposes of obtaining a large-$`N`$ expansion to introduce scaled variables:
$`x_i`$ $`\sqrt{N}x_i,`$ $`r_0`$ $`\sqrt{N}r_0.`$ (6)
Then the Lagrangian scales by a factor of $`N`$:
$$L/N=L_N(x,\dot{x})=\frac{1}{2}\underset{i=1}{\overset{N}{}}\dot{x}_i^2\frac{g}{8}\left(\underset{i=1}{\overset{N}{}}x_i^2r_0^2\right)^2.$$
(7)
We use these scaled variables in this paper, so that the rescaled $`r_01`$. Next we introduce a composite coordinate $`\chi `$ by adding to (7) a term:
$$\frac{1}{2g}\left(\chi \frac{g}{2}\left(\underset{i=1}{\overset{N}{}}x_i^2r_0^2\right)\right)^2.$$
(8)
The Lagrangian (7) then becomes:
$$\begin{array}{c}L_N(x,\chi ;\dot{x},\dot{\chi })=\underset{i=1}{\overset{N}{}}\left[\frac{1}{2}(\dot{x}_i^2\chi x_i^2)+j_ix_i\right]\hfill \\ \hfill +\frac{r_0^2\chi }{2}+\frac{\chi ^2}{2g}+J\chi ,\end{array}$$
(9)
where we have also added sources $`j_i`$ and $`J`$ coupling to $`x_i`$ and $`\chi `$ respectively. From this Lagrangian we get the Heisenberg equations of motion for the operators $`\widehat{x}_i(t)`$ and $`\widehat{\chi }(t)`$:
$$\widehat{\ddot{x}}_i(t)+\widehat{\chi }(t)\widehat{x}_i(t)=j_i(t),$$
$$\frac{\widehat{\chi }(t)}{g}=\frac{1}{2}\left(\underset{i=1}{\overset{N}{}}\widehat{x}_i^2(t)r_0^2\right)J(t).$$
(10)
Here, and in the following, we indicate operators by “hats.” Taking expectation values with respect to an initial density matrix we obtain the c-number equations:
$$\widehat{\ddot{x}}_i(t)+\widehat{\chi }(t)\widehat{x}_i(t)=j_i,$$
$$\frac{\widehat{\chi }(t)}{g}=\frac{1}{2}\left(\underset{i=1}{\overset{N}{}}\widehat{x}_i^2(t)r_0^2\right)J(t).$$
(11)
By rewriting the quartic interaction in terms of the composite field $`\chi `$, the induced interaction of the form $`\chi x_i^2`$ is reminiscent of $`N`$ flavor QED with interaction $`A_\mu \overline{\psi }_i\gamma ^\mu \psi _i`$. The fact that these two theories have the same topological structure will allow us to use the intuition gained in classical plasmas to make appropriate approximations.
To simplify notation we include all independent coordinates in one vector. We define:
$`x_\alpha (t)`$ $`=[\chi (t),x_1(t),x_2(t),\mathrm{},x_N(t)],`$
$`j_\alpha (t)`$ $`=[\stackrel{~}{J}(t),j_1(t),j_2(t),\mathrm{},j_N(t)].`$ (12)
for $`\alpha =0,1,\mathrm{},N`$, and where $`\stackrel{~}{J}(t)=J(t)r_0^2/2`$. Absorbing the factor $`r_0^2/2`$ into the current means that $`\stackrel{~}{J}(t)`$ is not zero when $`J(t)`$ is set to zero. Greek indices run from $`0`$ to $`N`$, whereas Latin indices go from $`1`$ to $`N`$. Using this extended notation, the generating functional $`Z[j]`$ and connected generator $`W[j]`$ is given by the path integral:
$$Z[j]=e^{iNW[j]}=\underset{\alpha =0}{\overset{N}{}}dx_\alpha \mathrm{exp}\left\{iNS_N[x;j]\right\}$$
(13)
where the action $`S_N[x;j]`$ is given by:
$$\begin{array}{c}S_N[x;j]=\frac{1}{2}\underset{\alpha ,\beta }{}_𝒞dt_𝒞dt^{}x_\alpha (t)\mathrm{\Delta }_{\alpha ,\beta }^1[x](t,t^{})x_\beta (t^{})\hfill \\ \hfill +\underset{\alpha }{}_𝒞dtx_\alpha (t)j_\alpha (t),\end{array}$$
(14)
and where $`\mathrm{\Delta }_{\alpha ,\beta }^1[x](t,t^{})`$ is given by:
$$\mathrm{\Delta }_{\alpha ,\beta }^1[x](t,t^{})=\left(\begin{array}{cc}D^1(t,t^{})& 0\\ 0& G_{ij}^1[\chi ](t,t^{})\end{array}\right),$$
(15)
with
$`D^1(t,t^{})`$ $`={\displaystyle \frac{1}{g}}\delta _𝒞(t,t^{}),`$
$`G_{ij}^1[\chi ](t,t^{})`$ $`=\left\{{\displaystyle \frac{\mathrm{d}^2}{\mathrm{d}t^2}}+\chi (t)\right\}\delta _{ij}\delta _𝒞(t,t^{}).`$ (16)
In what follows it will be useful to introduce another matrix inverse Green function $`G_{\alpha \beta }^1[x](t,t^{})`$ as follows:
$`G_{\alpha ,\beta }^1[x](t,t^{})={\displaystyle \frac{\delta ^2S_N[x;j]}{\delta x_\alpha (t)\delta x_\beta (t^{})}}`$
$`=\left(\begin{array}{cc}D^1(t,t^{})& \overline{K}_j^1(t,t^{})\\ K_i^1(t,t^{})& G_{i,j}^1(t,t^{})\end{array}\right),`$ (17)
with $`D^1(t,t^{})`$ and $`G_{i,j}^1(t,t^{})`$ given by Eq. (16), and $`K_i^1[x](t,t^{})=\overline{K}_i^1[x](t,t^{})=x_i(t)\delta _𝒞(t,t^{})`$.
## III The Schwinger-Dyson equations
The Schwinger-Dyson equations are integral equations for the Green functions. The Green functions can be obtained by functional differentiation of the path integral for the generating function in the presence of external sources. After setting the external sources to zero, one obtains an infinitely coupled hierarchy of coupled equations for the Green functions. For an initial value problem, the boundary conditions on the Green functions can be implemented by using a time ordered product where the time ordering refers to the closed time path contour of the Schwinger-Keldysh-Bakshi-Mahanthappa formalismref:CTP . A detailed discussion of that formalism as applied to implementing the $`1/N`$ expansion for this particular problem is described in ref. ref:ctpN . One way to generate the equations is to consider the identityref:IZ :
$$\underset{\beta }{}dx_\beta \frac{\delta }{\delta x_\alpha (t)}e^{iNS_N[x;j]}=0,$$
(18)
from which we find:
$$\frac{1}{g}\chi (t)+\frac{1}{2}\left\{\underset{i}{}\left[x_i^2(t)+\frac{1}{N}𝒢_{ii}(t,t)/i\right]r_0^2\right\}=J(t),$$
$$\left\{\frac{\mathrm{d}^2}{\mathrm{d}t^2}+\chi (t)\right\}x_i(t)+\frac{1}{N}𝒦_i(t,t)/i=j_i(t),$$
(19)
where $`x_i(t)`$ and $`\chi (t)`$ are *average values* of the operators,
$`x_i(t)`$ $`{\displaystyle \frac{\delta W[J,j]/i}{\delta j_i(t)}}=\widehat{x}_i(t),`$
$`\chi (t)`$ $`{\displaystyle \frac{\delta W[J,j]/i}{\delta J(t)}}=\widehat{\chi }(t),`$
and where the Green functions $`𝒢_{\alpha ,\beta }[j](t,t^{})`$ are defined by:
$`𝒢_{\alpha ,\beta }[j](t,t^{})`$ $`={\displaystyle \frac{\delta x_\alpha (t)}{\delta j_\beta (t^{})}}={\displaystyle \frac{\delta ^2W[j]}{\delta j_\alpha (t)\delta j_\beta (t^{})}}`$
$`=\left(\begin{array}{cc}𝒟(t,t^{})& 𝒦_j(t,t^{})\\ \overline{𝒦}_i(t,t^{})& 𝒢_{i,j}(t,t^{})\end{array}\right).`$ (20)
Eq. (19) is identical to Eq. (11). In this equation and in what follows, $`x_i`$ and $`\chi `$ now correspond to the expectation values:
The Green functions are explicitly given by
$`𝒟(t,t^{})`$ $`={\displaystyle \frac{\delta ^2W[J,j]}{\delta J(t)\delta J(t^{})}}`$ $`𝒦_i(t,t^{})`$ $`={\displaystyle \frac{\delta ^2W[J,j]}{\delta J(t)\delta j_i(t^{})}}`$
$`\overline{𝒦}_i(t,t^{})`$ $`={\displaystyle \frac{\delta ^2W[J,j]}{\delta j_i(t)\delta J(t^{})}}`$ $`𝒢_{i,j}(t,t^{})`$ $`={\displaystyle \frac{\delta ^2W[J,j]}{\delta j_i(t)\delta j_j(t^{})}}.`$
The integrability conditions require that $`\overline{𝒦}_i(t,t^{})=𝒦_i(t^{},t)`$. To obtain the Schwinger-Dyson equations it is advantageous to Legendre transform to the expectation value of the coordinate variables $`x_\alpha (t)`$, as the independent variable instead of the currents. The effective action generating functional of 1-PI graphs is given by a Legendre transformation:
$$\mathrm{\Gamma }[x]=W[j]_𝒞dt\underset{\alpha }{}\{x_\alpha (t)j_\alpha (t)\}.$$
(21)
So since $`j_\alpha (t)=\delta \mathrm{\Gamma }[x]/\delta x_\alpha (t)`$, the equations of motion (19) give values for derivatives of $`\mathrm{\Gamma }[x]`$:
$`{\displaystyle \frac{\delta \mathrm{\Gamma }[x]}{\delta \chi (t)}}`$ $`={\displaystyle \frac{1}{g}}\chi (t)`$
$`+{\displaystyle \frac{1}{2}}\left\{{\displaystyle \underset{i}{}}\left[x_i^2(t)+{\displaystyle \frac{1}{N}}𝒢_{ii}(t,t)/i\right]r_0^2\right\}`$ (22)
$`{\displaystyle \frac{\delta \mathrm{\Gamma }[x]}{\delta x_i(t)}}`$ $`=\left\{{\displaystyle \frac{\mathrm{d}^2}{\mathrm{d}t^2}}+\chi (t)\right\}x_i(t)+{\displaystyle \frac{1}{N}}𝒦_i(t,t)/i.`$ (23)
However the Green functions here, $`𝒢_{ii}(t,t)`$ and $`𝒦_i(t,t)`$ are defined in Eq. (20) as functionals of the currents $`j_\alpha (t)`$. These must be expressed as functionals of $`x_\alpha (t)`$ by inverse relations. We define these inverse Green functions, which are functionals of $`x_\alpha (t)`$, by:
$`𝒢_{\alpha ,\beta }^1[x](t,t^{})`$ $`={\displaystyle \frac{\delta j_\alpha (t)}{\delta x_\beta (t^{})}}={\displaystyle \frac{\delta ^2\mathrm{\Gamma }[x]}{\delta x_\alpha (t)\delta x_\beta (t^{})}}`$
$`=\left(\begin{array}{cc}𝒟^1(t,t^{})& \overline{𝒦}_j^1(t,t^{})\\ 𝒦_i^1(t,t^{})& 𝒢_{i,j}^1(t,t^{})\end{array}\right),`$
where explicitly
$`𝒟^1(t,t^{})`$ $`={\displaystyle \frac{\delta ^2\mathrm{\Gamma }[\chi ,x]}{\delta \chi (t)\delta \chi (t^{})}},`$
$`\overline{𝒦}_i^1(t,t^{})`$ $`={\displaystyle \frac{\delta ^2\mathrm{\Gamma }[\chi ,x]}{\delta \chi (t)\delta x_i(t^{})}},`$
$`𝒦_i^1(t,t^{})`$ $`={\displaystyle \frac{\delta ^2\mathrm{\Gamma }[\chi ,x]}{\delta x_i(t)\delta \chi (t^{})}},`$
$`𝒢_{i,j}^1(t,t^{})`$ $`={\displaystyle \frac{\delta ^2\mathrm{\Gamma }[\chi ,x]}{\delta x_i(t)\delta x_j(t^{})}}.`$
Again we have $`\overline{𝒦}_i^1(t,t^{})=𝒦_i^1(t^{},t)`$. The inverse Green functions are given by differentiating the equations of motion, Eqs. (22) and (23), with respect to the coordinates. Using
$$_𝒞dt^{}\underset{\beta }{}𝒢_{\alpha ,\beta }^1[x](t,t^{})𝒢_{\beta ,\gamma }[j](t^{},t^{\prime \prime })=\delta _{\alpha ,\gamma }\delta _𝒞(t,t^{\prime \prime }),$$
we find:
$`{\displaystyle \frac{\delta 𝒢_{\alpha ,\beta }[j](t_1,t_2)}{\delta x_\gamma (t_3)}}`$ $`={\displaystyle _𝒞}dt_4{\displaystyle _𝒞}dt_5{\displaystyle \underset{\delta ,ϵ}{}}𝒢_{\alpha ,\delta }[j](t_1,t_4)`$
$`\times \mathrm{\Gamma }_{\delta ,ϵ,\gamma }[x](t_4,t_5,t_3)𝒢_{ϵ,\beta }[j](t_5,t_2),`$ (24)
where $`\mathrm{\Gamma }_{\alpha ,\beta ,\gamma }[x](t_1,t_2,t_3)`$ is the three-point vertex function, defined by:
$`\mathrm{\Gamma }_{\alpha ,\beta ,\gamma }[x](t_1,t_2,t_3)`$ $`={\displaystyle \frac{\delta 𝒢_{\alpha ,\beta }^1[x](t_1,t_2)}{\delta x_\gamma (t_3)}}`$
$`={\displaystyle \frac{\delta ^3\mathrm{\Gamma }[x]}{\delta x_\alpha (t_1)\delta x_\beta (t_2)\delta x_\gamma (t_3)}}.`$ (25)
Explicitly, we find an equation of the form:
$$𝒢_{\alpha ,\beta }^1(t,t^{})=G_{\alpha ,\beta }^1(t,t^{})+\mathrm{\Sigma }_{\alpha ,\beta }(t,t^{}),$$
(26)
where $`G_{\alpha ,\beta }^1(t,t^{})`$ is given by Eq. (17). The generalized self energy $`\mathrm{\Sigma }_{\alpha ,\beta }(t,t^{})`$ is given by:
$$\mathrm{\Sigma }_{\alpha ,\beta }(t,t^{})=\left(\begin{array}{cc}\mathrm{\Pi }(t,t^{})& \mathrm{\Omega }_j(t,t^{})\\ \overline{\mathrm{\Omega }}_i(t,t^{})& \mathrm{\Sigma }_{ij}(t,t^{})\end{array}\right),$$
(27)
and where the polarization $`\mathrm{\Pi }(t,t^{})`$, self energy $`\mathrm{\Sigma }_{ij}(t,t^{})`$, and the off diagonal terms $`\mathrm{\Omega }_i(t,t^{})`$ and $`\overline{\mathrm{\Omega }}_i(t,t^{})`$ are given by:
$`\mathrm{\Pi }(t,t^{})`$ $`={\displaystyle \frac{i}{2N}}{\displaystyle \underset{i,\alpha ,\beta }{}}{\displaystyle _𝒞}dt_1{\displaystyle _𝒞}dt_2𝒢_{i,\alpha }(t,t_1)\mathrm{\Gamma }_{\alpha ,\beta ,0}(t_1,t_2,t^{})𝒢_{\beta ,i}(t_2,t),`$
$`\mathrm{\Sigma }_{ij}(t,t^{})`$ $`={\displaystyle \frac{i}{N}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle _𝒞}dt_1{\displaystyle _𝒞}dt_2𝒢_{i,\alpha }(t,t_1)\mathrm{\Gamma }_{\alpha ,\beta ,j}(t_1,t_2,t^{})𝒢_{\beta ,0}(t_2,t)`$
$`\mathrm{\Omega }_i(t,t^{})`$ $`={\displaystyle \frac{i}{2N}}{\displaystyle \underset{j,\alpha ,\beta }{}}{\displaystyle _𝒞}dt_1{\displaystyle _𝒞}dt_2𝒢_{j,\alpha }(t,t_1)\mathrm{\Gamma }_{\alpha ,\beta ,i}(t_1,t_2,t^{})𝒢_{\beta ,j}(t_2,t)`$
$`\overline{\mathrm{\Omega }}_i(t,t^{})`$ $`={\displaystyle \frac{i}{N}}{\displaystyle \underset{\alpha \beta }{}}{\displaystyle _𝒞}dt_1{\displaystyle _𝒞}dt_2𝒢_{i,\alpha }(t,t_1)\mathrm{\Gamma }_{\alpha ,\beta ,0}(t_1,t_2,t^{})𝒢_{\beta ,0}(t_2,t).`$ (28)
In order to solve the equation for the two point function, Eq. (26), one requires knowledge of the three point function, defined by Eq. (25). This in turn requires knowledge of the four-point function, *ad infinitum*. It is this infinite hierarchy of coupled Green function equations that corresponds to solving exactly the Schrödinger equation.
The matrix inversion of Eq. (26) gives the set of coupled equations,
$$\begin{array}{c}𝒢_{\alpha ,\beta }(t,t^{})=G_{\alpha ,\beta }(t,t^{})\underset{\gamma ,\delta }{}_𝒞dt_1_𝒞dt_2G_{\alpha ,\gamma }(t,t_1)\hfill \\ \hfill \times \mathrm{\Sigma }_{\gamma ,\delta }(t_1,t_2)𝒢_{\delta ,\beta }(t_2,t^{}),\end{array}$$
(29)
where
$$G_{\alpha ,\beta }(t,t^{})=\left(\begin{array}{cc}D(t,t^{})& K_i(t,t^{})\\ \overline{K}_i(t,t^{})& G_{ij}(t,t^{})\end{array}\right).$$
(30)
with
$$\begin{array}{cc}\hfill \underset{j}{}\{[\frac{\mathrm{d}^2}{\mathrm{d}t^2}+\chi (t)]\delta _{ij}& +gx_i(t)x_j(t)\}G_{jk}(t,t^{})\hfill \\ & =\delta _{ik}\delta _𝒞(t,t^{}),\hfill \end{array}$$
(31)
$$D(t,t^{})=g\delta _𝒞(t,t^{})+g^2\underset{ij}{}x_i(t)G_{ij}(t,t^{})x_j(t^{}),$$
(32)
$$\overline{K}_j(t,t^{})=K_j(t^{},t)=g\underset{i}{}G_{ji}(t,t^{})x_i(t^{}).$$
(33)
When $`x_i(t)0`$, one notes that $`D(t,t^{})`$ is not the inverse of $`D^1(t,t^{})`$.
The vertex function $`\mathrm{\Gamma }_{\alpha ,\beta ,\gamma }[x](t_1,t_2,t_3)`$ defined in (25) is obtained by differentiation of Eq. (26) with respect to $`x_\gamma (t)`$. We find:
$$\begin{array}{c}\mathrm{\Gamma }_{\alpha ,\beta ,\gamma }[x](t_1,t_2,t_3)=\frac{\delta 𝒢_{\alpha ,\beta }^1[x](t_1,t_2)}{\delta x_\gamma (t_3)}\hfill \\ \hfill =f_{\alpha ,\beta ,\gamma }\delta _𝒞(t_1,t_2)\delta _𝒞(t_1,t_3)+\mathrm{\Phi }_{\alpha ,\beta ,\gamma }[x](t_1,t_2,t_3).\end{array}$$
(34)
Here $`f_{i,j,0}=f_{0,i,j}=f_{i,0,j}=\delta _{ij}`$, otherwise $`f`$ is zero. $`\mathrm{\Phi }_{\alpha ,\beta ,\gamma }[x](t_1,t_2,t_3)`$ is given by derivatives of the self-energy matrix:
$$\mathrm{\Phi }_{\alpha ,\beta ,\gamma }[x](t_1,t_2,t_3)=\frac{\delta \mathrm{\Sigma }_{\alpha ,\beta }[x](t_1,t_2)}{\delta x_\gamma (t_3)},$$
(35)
and is of order $`1/N`$.
We are interested in resummation schemes that are exact to order $`1/N`$ for $`x_i^2`$. We see from Eqs. (34) and (35) that it is consistent to replace $`\mathrm{\Gamma }_{\alpha ,\beta ,\gamma }[x](t_1,t_2,t_3)`$ in Eq. (29) by the first term in Eq. (34) to obtain a resummation which is exact to order $`1/N`$. To simplify our discussion of the exact Schwinger-Dyson equation for the vertex function, we will only consider the case of the quantum roll where $`x_i(t)=0`$.
Following the treatment of ref. ref:CGHT , we have for the 3-$`\chi `$ vertex:
$$\mathrm{\Lambda }(t_1,t_2,t_3)=\frac{\delta 𝒟^1(t_1,t_2)}{\delta \chi (t_3)}=\underset{ijk}{}_𝒞dt_4_𝒞dt_5𝒢_{ij}(t_3,t_4)𝒢_{ik}(t_3,t_5)M_{jk}(t_4,t_2;t_5,t_1),$$
where $`M_{jk}(t_4,t_2;t_5,t_1)`$ is 1-PI in the channel $`x+x\chi +\chi `$. The lowest order in $`1/N`$ contribution to $`M(t_4,t_5;t_2,t_3)`$ is:
$$M_{jk}(t_4,t_2;t_5,t_1)=\delta _𝒞(t_4,t_2)\delta _𝒞(t_5,t_1)𝒢_{jk}(t_2,t_1).$$
(36)
When $`x_i(t)=0`$, the exact Schwinger-Dyson equation for the $`\chi `$-$`x`$-$`x`$ vertex is
$$\begin{array}{c}\mathrm{\Gamma }_{ij}(t_1,t_2,t_3)=\frac{\delta 𝒢^1(t_1,t_2)}{\delta \chi (t_3)}=\delta _{ij}\delta _𝒞(t_1,t_2)\delta _𝒞(t_1,t_3)_𝒞\mathrm{d}t_4_𝒞\mathrm{d}t_5_𝒞\mathrm{d}t_6_𝒞\mathrm{d}t_7\times \hfill \\ \hfill \left\{\underset{klmn}{}\mathrm{\Gamma }_{kl}(t_4,t_5,t_3)𝒢_{km}(t_4,t_6)𝒢_{ln}(t_5,t_7)𝒦_{1mn}(t_5,t_2;t_7,t_1)+\mathrm{\Lambda }(t_4,t_5,t_3)𝒟(t_4,t_5)𝒟(t_6,t_7)𝒦_{2ij}(t_5,t_2;t_7,t_1)\right\}.\end{array}$$
(37)
where $`𝒦_1`$ and $`𝒦_2`$ are the s-channel 2-PI scattering amplitudes for the reactions: $`x+xx+x`$ and $`\chi +\chi x+x`$, respectively.
This is shown pictorially in Fig. 1. In general one then has to obtain equations for the 2-PI scattering amplitudes as well as for $`\mathrm{\Lambda }`$. These will depend on even higher $`n`$-point functions, *ad infinitum*. In our approximations made at the two-point function level, the 2-PI s-channel scattering amplitudes $`K_1`$ and $`K_2`$, used in the equations for the vertex function, will turn out to be graphs for one-particle exchange in the t-channel of the $`\chi `$\- and $`x`$-particles respectively.
In our truncations of the Schwinger-Dyson equations, we will always replace the full three-point vertex function by the bare one in the equations for $`x`$ and $`𝒢`$ in the *presence* of external sources. Once this truncation is made, then for the problem we are addressing here (the approximate time evolution of $`N`$ quantum anharmonic oscillators) one never needs any of the $`N`$ point functions beyond the 1 and 2 point function equations. What will distinguish a further approximation we will call the DDSA is that we will also further approximate the $`\chi `$ propagator to be that of the LOLN approximation.
By making this bare vertex approximation in the equations for the one- and two-point Green functions, we have *not* relinquished our ability to calculate in this approximation all the higher connected Green functions. These are obtainable by further functional differentiation of the effective action. In particular if we wanted to use linear response theory (the Kubo formula) to obtain the electrical conductivity for a QED plasma, one would functionally differentiate the equation for the inverse two-point for the electron function with respect to $`A_\mu `$. In our problem the photon is replaced by the composite field $`\chi `$, and the electron by $`x_i`$.
Because of recent interest in studying plasma conductivity in both QED and QCD, we will spend extra time on comparing the equations obtained for the vertex function in the three approximations considered here. In conductivity calculations, it is necessary to sum all the ladder graphs in the equation for the vertex function to get good results for dilute plasmas. We will find that in NLOLN the vertex function is *not* an integral equation but is rather the sum of a few diagrams whereas the other two approximations lead to integral equations that sum an infinite number of diagrams. Another issue is in preserving Ward identities. One of the reasons the large-$`N`$ expansion was so interesting is that it is a complete reexpansion of the field theory which preserves Ward identities at each order. The QED plasma conductivity problem people ref:emil became interested in the DDSA because it *exactly* obeyed the Ward identities, whereas the BVA approximation violates Ward identities at order $`1/N^2`$. It is for this reasons we thought it appropriate to study the DDSA approximation, even though it violated energy conservation already at order $`1/N`$, hoping that at least at large $`N`$ it would be numerically accurate *and* satisfy Ward identities in QED applications.
The exact formula for the energy is given by:
$$E/N=\frac{1}{2}\underset{i}{}\left\{\widehat{\dot{x}}_i^2(t)+\widehat{\chi }(t)\widehat{x}_i^2(t)r_0^2\widehat{\chi }(t)\widehat{\chi }^2(t)/g\right\}.$$
(38)
When $`x_i(t)=\widehat{x}_i(t)=0`$ and $`\dot{x}_i(t)=0`$, one obtains:
$`E/N`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\{{\displaystyle \frac{^2𝒢(t,t^{})/i}{tt^{}}}|_{t^{}=t}+\chi (t)𝒢(t,t)/ir_0^2\chi (t){\displaystyle \frac{1}{g}}[\chi ^2(t)+{\displaystyle \frac{1}{N}}𝒟(t,t)/i]`$
$`+{\displaystyle \frac{1}{N}}{\displaystyle \underset{ijk}{}}{\displaystyle _𝒞}\mathrm{d}t_1{\displaystyle _𝒞}\mathrm{d}t_2{\displaystyle _𝒞}\mathrm{d}t_3𝒟(t_1,t_2)𝒢_{ij}(t_1,t)𝒢_{ik}(t,t_1)\mathrm{\Gamma }_{jk}(t_1,t_3,t_2)\},`$ (39)
where $`\mathrm{\Gamma }_{jk}(t_1,t_3,t_2)`$ is the full vertex function given in Eq. (37).
## IV Effective action for two-particle irreducible graphs
Since the approximations we are going to consider have a simple interpretation in terms of keeping a particular 2-PI vacuum graph in the generating functional of the 2-PI graphs, we would like to review this formalism following the approach of Cornwall, Jackiw, and Tomboulis (CJT)ref:CJT .
The first Legendre transform of the generating functional $`W[j]`$ of connected Green functions is widely known and used and is called the “effective action.” The higher Legendre transforms (second, third, etc.) were introduced by De Dominicis and Martinref:DM in quantum statistics. Dahmen and Jona-Lasinioref:DJL , and later Visil’ev and Kazanskiiref:VK , extended these ideas to quantum field theory. These methods were then used by Cornwall, Jackiw, and Tomboulis to discuss dynamical symmetry breaking in Hartree type approximations which later led to the second Legendre transformation formalism being called the CJT formalism. These higher order Legendre transformed actions have the advantage of being able to treat higher order Green functions on the same footing as the coordinates.
We will first summarize the general results of that paper before proceeding to the specific approximations we consider in this paper. The method of CJT is to introduce one- and two-body sources for the coordinates $`x_\alpha (t)`$ and the Green functions $`𝒢_{\alpha ,\beta }(t,t^{})`$ in the action, and then make a Legendre transformation to the one- and two-point functions. The resulting action, as a function of $`x`$ and $`𝒢`$, contains a term which is the sum of all two-particle irreducible vacuum graphs. This term can be written using the vertices of the interaction and $`𝒢`$. We use the extended notation for the coordinates and one-body sources, given in Eq. (12).
Thus, the generating functional $`Z[j,k]`$ for the CJT action is given by:
$$Z[j,k]=e^{iNW[j,k]}=\underset{\alpha =0}{\overset{N}{}}dx_\alpha \mathrm{exp}\left\{iNS_N[x;j,k]\right\},$$
with
$$S_N[x;j,k]=S_{\text{class}}[x]+\underset{\alpha }{}_𝒞dtx_\alpha (t)j_\alpha (t)+\frac{1}{2}\underset{\alpha ,\beta }{}_𝒞dt_𝒞dt^{}x_\alpha (t)k_{\alpha ,\beta }(t,t^{})x_\beta (t^{}),$$
(40)
where
$`S_{\text{class}}[x]`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle _𝒞}dt{\displaystyle _𝒞}dt^{}x_\alpha (t)\mathrm{\Delta }_{\alpha ,\beta }^1[x](t,t^{})x_\beta (t^{})=S_0+S_{\text{int}}[x],`$ (41)
$`S_0`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle _𝒞}dt{\displaystyle _𝒞}dt^{}x_\alpha (t)\mathrm{\Delta }_{0\alpha ,\beta }^1(t,t^{})x_\beta (t^{}),`$ (42)
$`S_{\text{int}}[x]`$ $`={\displaystyle \frac{1}{2}}{\displaystyle _𝒞}dt\chi (t){\displaystyle \underset{i}{}}x_i^2(t),`$ (43)
and where $`\mathrm{\Delta }_{0\alpha ,\beta }^1(t,t^{})`$ is given by:
$`\mathrm{\Delta }_{0\alpha ,\beta }^1(t,t^{})`$ $`=\left(\begin{array}{cc}D^1(t,t^{})& 0\\ 0& G_{0ij}^1(t,t^{})\end{array}\right),`$
$`G_{0ij}^1(t,t^{})`$ $`=\left\{{\displaystyle \frac{\mathrm{d}^2}{\mathrm{d}t^2}}\right\}\delta _{ij}\delta _𝒞(t,t^{}).`$
with $`D^1(t,t^{})`$ given by Eq. (16). In this formalism, we have separated out an “interaction” term, Eq. (43), which depends on the coordinates $`x_\alpha (t)`$, from a bare Green function $`G_{0ij}^1(t,t^{})`$, *which is independent of the coordinates $`x_\alpha (t)`$*, in contrast to our previous definitions in Eq. (16). The term $`r_0^2\chi (t)/2`$ has been absorbed into the definition of the current $`\stackrel{~}{J}(t)`$ in Eq. (12).
The second Legendre transform of $`W[j,k]`$ is the CJT effective action:
$$\begin{array}{c}\mathrm{\Gamma }[x,𝒢]=W[j,k]\underset{\alpha }{}_𝒞dtx_\alpha (t)j_\alpha (t)\hfill \\ \hfill +\frac{1}{2}\underset{\alpha ,\beta }{}_𝒞dt_𝒞dt^{}k_{\alpha ,\beta }(t,t^{})\left\{x_\alpha (t)x_\beta (t^{})+𝒢_{\alpha ,\beta }(t,t^{})\right\}\end{array}$$
CJT showed that $`\mathrm{\Gamma }[x,𝒢]`$ can be obtained as a series expansion in terms of 2-PI graphs. That is, introducing the functional operator,
$`G_{\alpha ,\beta }^1[x](t,t^{})`$ $`={\displaystyle \frac{\delta ^2S_0[x]}{\delta x_\alpha (t)\delta x_\beta (t^{})}}`$
$`=\left(\begin{array}{cc}D^1(t,t^{})& \overline{K}_j^1[x](t,t^{})\\ K_i^1[x](t,t^{})& G_{i,j}^1[x](t,t^{})\end{array}\right),`$ (44)
which is the same as the $`G_{\alpha ,\beta }^1[x](t,t^{})`$ as defined in Eq. (17), one can write the effective action in the form:
$$\begin{array}{c}\mathrm{\Gamma }[x,𝒢]=S_{\text{class}}[x]+\frac{i}{2}\mathrm{Tr}\{\mathrm{ln}[𝒢^1]\}\hfill \\ \hfill +\frac{i}{2}\mathrm{Tr}\{G^1[x]𝒢1\}+\mathrm{\Gamma }_2[x,𝒢].\end{array}$$
(45)
The quantity $`\mathrm{\Gamma }_2[x,𝒢]`$ has a simple graphical interpretation in terms of all the 2-PI vacuum graphs using vertices from the interaction term. The Hartree and leading order in large-$`N`$ approximation for the $`x^4`$ potential was obtained by CJT using a single two-loop vacuum graph in the $`O(N)`$ theory written in terms of only the coordinates $`x_i`$. Our strategy for obtaining a resummation of the large-$`N`$ approximation is to first rewrite the theory in terms of the composite field $`\chi `$, and the equivalent Lagrangian given in Eq. (9). Using these new variables, we then choose for $`\mathrm{\Gamma }_2[x,𝒢]`$ the 2-PI graphs shown in Fig. 3, which is now written in terms of the full $`\chi `$ and $`x`$ propagators and the trilinear coupling $`\chi (t)x_i^2(t)/2`$.
## V Bare Vertex Approximation
The bare vertex approximation (BVA) is obtained by setting the vertex function equal to its bare value in the exact equations for the one and two point functions. This is an energy conserving approximation which leads to integral equations for the three-$`\chi `$ vertex function as well as for the $`x`$-$`x`$-$`\chi `$ vertex function. The bare vertex approximation consists of making the replacement
$$\mathrm{\Gamma }_{\alpha ,\beta ,\gamma }[x](t_1,t_2,t_3)=f_{\alpha ,\beta ,\gamma }\delta _𝒞(t_1,t_2)\delta _𝒞(t_1,t_3).$$
(46)
in the exact Schwinger-Dyson equations for the self-energies, Eqs. (28). This gives for the BVA:
$`\mathrm{\Pi }(t,t^{})`$ $`={\displaystyle \frac{i}{2N}}{\displaystyle \underset{ij}{}}𝒢_{ij}(t,t^{})𝒢_{ji}(t^{},t),`$ (47)
$`\mathrm{\Omega }_i(t,t^{})`$ $`={\displaystyle \frac{i}{N}}{\displaystyle \underset{j}{}}\overline{𝒦}_j(t,t^{})𝒢_{ji}(t,t^{}),`$
$`\overline{\mathrm{\Omega }}_i(t,t^{})`$ $`={\displaystyle \frac{i}{N}}{\displaystyle \underset{j}{}}𝒦_j(t^{},t)𝒢_{ji}(t^{},t),`$
$`\mathrm{\Sigma }_{ij}(t,t^{})`$ $`={\displaystyle \frac{i}{N}}\left\{\overline{𝒦}_i(t,t^{})𝒦_j(t^{},t)+𝒢_{ij}(t,t^{})𝒟(t^{},t)\right\},`$
where we have used the symmetry property, $`𝒢_{ij}(t,t^{})=𝒢_{ji}(t^{},t)`$ and $`𝒦_i(t,t^{})=\overline{𝒦}_i(t^{},t)`$. Thus we find $`\overline{\mathrm{\Omega }}_i(t,t^{})=\mathrm{\Omega }_i(t^{},t)`$. The self-energies (47) are then used in Eqs. (26) to find the one- and two-point functions. For the Green functions, we find:
$$𝒢_{\alpha ,\beta }^1(t,t^{})=G_{\alpha ,\beta }^1(t,t^{})+\mathrm{\Sigma }_{\text{BVA}\alpha ,\beta }(t,t^{}),$$
(48)
with $`\mathrm{\Sigma }_{\text{BVA}\alpha ,\beta }(t,t^{})`$ given by Eq. (47). The inversion of Eq. (48) is given by Eq. (29), which is a set of four coupled integral equations for the four BVA Green functions, which must be solved simultaneously.
From Eqs. (22) and (23), the equations of motion for $`x_i(t)`$ and the gap equation for $`\chi (t)`$ is then given by:
$$\left\{\frac{\mathrm{d}^2}{\mathrm{d}t^2}+\chi (t)\right\}x_i(t)+\frac{1}{N}𝒦_i(t,t)/i=0,$$
(49)
$$\chi (t)=\frac{g}{2}\left\{\underset{i}{}\left[x_i^2(t)+\frac{1}{N}𝒢_{ii}(t,t)/i\right]r_0^2\right\}.$$
(50)
For the quantum roll, we further set $`x_i(t)=0`$. This means that $`𝒦_i(t,t)=\overline{𝒦}_i(t,t)=0`$, so that $`G_{\alpha \beta }(t,t^{})`$ is diagonal, and results in the following set of equations for the Green functions:
$`𝒟(t,t^{})`$ $`=D(t,t^{})`$
$`{\displaystyle _𝒞}dt_1{\displaystyle _𝒞}dt_2D(t,t_1)\mathrm{\Pi }(t_1,t_2)𝒟(t_2,t^{}),`$ (51)
$`𝒢_{ij}(t,t^{})`$ $`=G_{ij}(t,t^{})`$
$`{\displaystyle \underset{kl}{}}{\displaystyle _𝒞}dt_1{\displaystyle _𝒞}dt_2G_{ik}(t,t_1)\mathrm{\Sigma }_{kl}(t_1,t_2)𝒢_{lj}(t_2,t^{}),`$ (52)
where
$`\mathrm{\Pi }(t,t^{})`$ $`={\displaystyle \frac{i}{2N}}{\displaystyle \underset{ij}{}}𝒢_{ij}(t,t^{})𝒢_{ji}(t^{},t),`$
$`\mathrm{\Sigma }_{ij}(t,t^{})`$ $`={\displaystyle \frac{i}{N}}𝒢_{ij}(t,t^{})𝒟(t^{},t).`$ (53)
The gap equation for $`\chi (t)`$ becomes:
$$\chi (t)=\frac{g}{2}\left\{\frac{1}{N}\underset{i}{}𝒢_{ii}(t,t)/ir_0^2\right\}.$$
(54)
In addition, for this case, the initial conditions imply that we can take $`G_{ij}(t,t^{})`$ and $`𝒢_{ij}(t,t^{})`$ to be diagonal, which greatly simplify the integral equations. The BVA for the quantum roll requires that we solve equations (51), (52), (53), and (54) simultaneously using the numerical methods described in refs. ref:MDC and ref:BMIM .
Because of the interest in using the BVA approximation in QED (and QCD) plasma conductivity problems, we will discuss the integral equation one obtains for the vertex function in what follows. It was precisely because this approximation gives the sum of the graphs used in non-relativistic plasmas (see Fig. 2) in conductivity calculations which gave both accurate results as well as giving physical answers that initially interested us in this approximation.
The three-point vertex functions for the BVA are given by functional differentiation of the inverse two point functions:
$`\mathrm{\Lambda }(t_1,t_2,t_3)`$ $`\mathrm{\Gamma }_{000}(t_1,t_2,t_3)={\displaystyle \frac{\delta 𝒟^1(t_1,t_2)}{\delta \chi (t_3)}}`$ (55)
$`\mathrm{\Gamma }_{ij}(t_1,t_2,t_3)`$ $`\mathrm{\Gamma }_{ij0}(t_1,t_2,t_3)={\displaystyle \frac{\delta 𝒢_{ij}^1(t_1,t_2)}{\delta \chi (t_3)}},`$ (56)
and obtain the coupled integral equations:
$$\mathrm{\Lambda }(t_1,t_2,t_3)=\frac{i}{N}_𝒞dt_4_𝒞dt_5\underset{ijkl}{}𝒢_{ik}(t_1,t_4)\mathrm{\Gamma }_{kl}(t_4,t_5,t_3)𝒢_{lj}(t_5,t_2)𝒢_{ji}(t_2,t_1)$$
(57)
and
$$\begin{array}{c}\mathrm{\Gamma }_{ij}(t_1,t_2,t_3)=\delta _{ij}\delta _𝒞(t_1,t_2)\delta _𝒞(t_1,t_3)\hfill \\ \hfill _𝒞dt_4_𝒞dt_5\left\{\underset{kl}{}𝒢_{ik}(t_1,t_4)\mathrm{\Gamma }_{kl}(t_4,t_5,t_3)𝒢_{lj}(t_5,t_2)𝒟(t_2,t_1)+𝒢_{ij}(t_1,t_2)𝒟(t_2,t_4)\mathrm{\Lambda }(t_4,t_5,t_3)𝒟(t_5,t_1)\right\}.\end{array}$$
(58)
This is shown diagrammatically in Fig. 2. Looking at the diagrams, if we iterate these equations, we sum all the “rainbow” diagrams. As advertised, comparing these graphs with those shown in Fig. 1, $`𝒦_1`$ is approximated in the BVA by $`\chi `$ exchange and $`𝒦_2`$ by $`x`$ exchange in the $`t`$-channel.
Let us show that this approximation is easy to obtain from the CJT formalism once we treat $`𝒢`$ and $`𝒟`$ and $`𝒦`$ on exactly the same footing. We choose for our approximation to $`\mathrm{\Gamma }_2[𝒢]`$ the 2-PI graphs shown in Fig. 3. This gives:
$$\mathrm{\Gamma }_2[𝒢]=\frac{1}{4N}\underset{ij}{}_𝒞dt_1_𝒞dt_2𝒟(t_1,t_2)𝒢_{ij}(t_1,t_2)𝒢_{ji}(t_2,t_1)\frac{1}{2N}\underset{ij}{}_𝒞dt_1_𝒞dt_2\overline{𝒦}_i(t_1,t_2)𝒢_{ij}(t_1,t_2)𝒦_j(t_2,t_1).$$
(59)
Since the $`𝒟`$ propagator sums the contact term plus all the polarization bubbles $`\mathrm{\Pi }`$ of the original quartic interaction $`gx^4`$, if we reexpand $`𝒟`$ in a power series in $`\mathrm{\Pi }`$ then the first two terms in the series give the graphs used in the approximation ofref:Hu andref:berges . The CJT action is given by Eq. (45). The stationary condition for $`𝒢_{\alpha ,\beta }(t,t^{})`$ gives:
$$\frac{\delta \mathrm{\Gamma }[x,𝒢]}{\delta 𝒢_{\alpha \beta }}=\frac{i}{2}\left\{G_{\alpha \beta }^1[x]𝒢_{\alpha \beta }^1\right\}+\frac{\delta \mathrm{\Gamma }_2[𝒢]}{\delta 𝒢_{\alpha \beta }}=0,$$
or
$$𝒢_{\alpha ,\beta }^1(t,t^{})=G_{\alpha ,\beta }^1(t,t^{})+\mathrm{\Sigma }_{\text{BVA}\alpha ,\beta }[𝒢](t,t^{}),$$
where:
$$\mathrm{\Sigma }_{\text{BVA}\alpha ,\beta }[𝒢](t,t^{})=2i\frac{\delta \mathrm{\Gamma }_2[𝒢]}{\delta 𝒢_{\alpha \beta }(t,t^{})}.$$
(60)
Carrying out the derivatives of $`\mathrm{\Gamma }_2[𝒢]`$ given in Eq. (59), we find that $`\mathrm{\Sigma }_{\text{BVA}\alpha ,\beta }(t,t^{})`$ is exactly the same as found in Eq. (47) using the Schwinger-Dyson equations in the BVA approximation. The stationary condition for $`x_\alpha `$ also gives the same equations of motion for $`x_i(t)`$ and gap equation for $`\chi (t)`$ as found in Eqs. (49) and (50) using the Schwinger-Dyson equations in the BVA. Thus we conclude that the CJT action, as given in Eqs. (45) and (59), gives exactly the same set of equations as in the Schwinger-Dyson BVA truncation.
The energy for the BVA is obtained from (39) by using (46) for the vertex function. We find:
$`E/N`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\{{\displaystyle \frac{^2𝒢(t,t^{})/i}{tt^{}}}|_{t^{}=t}+\chi (t)𝒢(t,t)/ir_0^2\chi (t){\displaystyle \frac{1}{g}}[\chi ^2(t)+{\displaystyle \frac{1}{N}}𝒟(t,t)/i]`$
$`+{\displaystyle \frac{1}{N}}{\displaystyle \underset{ij}{}}{\displaystyle _𝒞}\mathrm{d}t_1𝒟(t_1,t)𝒢_{ij}(t_1,t)𝒢_{ji}(t,t_1)\}.`$ (61)
where, for our case, we have set $`x_i(t)=\dot{x}_i(t)=0`$. Since the BVA equations are derived from an effective action, energy is conserved.
## VI Dynamical Debye Screening Approximation
In plasma studies of the electric conductivity of fully ionized plasmas ref:plasma1 ; ref:plasma2 , it was found that in order to correctly determine the conductivity it was necessary to have an approximation where the photon propagator included the effects of dynamical Debye screening in the random phase approximation. This improved propagator was then used in a scattering kernel in the kinetic equations. In our model, the $`\chi `$ field plays the roll of the photon in the dynamics of the $`x_i`$ oscillators. The lowest approximation that includes the polarization effects in $`𝒟`$ is precisely the leading order in large-$`N`$ approximation to $`𝒟`$, namely $`𝒟_0`$ (see Eq. 69) which is discussed below in our derivation of the NLOLN approximation . The leading order in large-$`N`$ approximation is similar in spirit to the random phase approximation. The equation for $`𝒟^1(t,t^{})`$ in leading order in large-$`N`$ is given by:
$$𝒟_0^1(t,t^{})=\frac{1}{g}\delta _𝒞(t,t^{})+\mathrm{\Pi }_0(t,t^{}),$$
(62)
where
$$\begin{array}{c}\mathrm{\Pi }_0(t,t^{})=\frac{i}{2N}\underset{i,j}{}G_{ij}(t,t^{})G_{ji}(t^{},t)\hfill \\ \hfill +\underset{i,j}{}x_i(t)G_{ij}(t,t^{})x_j(t^{}).\end{array}$$
In the QED plasma problem, the $`\chi `$ propagator becomes the photon propagator and the delta function in $`𝒟_0`$ is replaced by the bare photon propagator. It is the bubble in $`\mathrm{\Pi }_0`$ that leads to the Debye screening of the photon. It is because of our interest in QED that we call this approximation the DDSA.
Let us now specialize to the case when $`x_i(t)=0`$. The equation for the full $`x`$ propagator $`𝒢`$ is:
$$\begin{array}{c}𝒢_{ij}(t,t^{})=G_{ij}(t,t^{})\hfill \\ \hfill \underset{k,l}{}_𝒞dt_1_𝒞dt_2G_{ik}(t,t_1)\mathrm{\Sigma }_{kl}(t_1,t_2)𝒢_{lj}(t_2,t^{}),\end{array}$$
(63)
with the self energy depending on the full $`𝒢`$ and the leading order in $`1/N`$ approximation to $`𝒟`$ given by Eq. (62):
$$\mathrm{\Sigma }_{kl}(t,t^{})=\frac{i}{N}𝒢_{kl}(t,t^{})D(t,t^{}).$$
(64)
The gap equation is:
$$\chi (t)=\frac{g}{2}\left\{\underset{i}{}\frac{1}{N}𝒢_{ii}(t,t)/ir_0^2\right\}.$$
(65)
There is a nontrivial vertex function in this approximation given by:
$`\mathrm{\Gamma }_{ij}(t_1,t_2,t_3)`$ $`={\displaystyle \frac{\delta 𝒢_{ij}^1[\chi ](t_1,t_2)}{\delta \chi (t_3)}}`$
$`=\delta _𝒞(t_1,t_2)\delta _𝒞(t_3,t_2)\delta _{ij}{\displaystyle \underset{kl}{}}{\displaystyle _𝒞}dt_4{\displaystyle _𝒞}dt_5\mathrm{\Gamma }_{kl}(t_4,t_5,t_3)𝒢_{ki}(t_4,t_1)D(t_1,t_2)𝒢_{jl}(t_2,t_5)`$
$`{\displaystyle _𝒞}dt_4{\displaystyle _𝒞}dt_5\mathrm{\Lambda }(t_4,t_5,t_3)D(t_4,t_1)𝒢_{ij}(t_1,t_2)D(t_2,t_5).`$ (66)
This equation can be obtained from the exact integral equation for $`\mathrm{\Gamma }`$ shown pictorially in Fig. 1 by making two approximations. The first is to approximate the exact three-$`\chi `$ vertex function by the triangle graph, which is the leading term in the $`1/N`$ expansion of this function. The second is to replace the scattering kernels, $`K_1`$ and $`K_2`$ by single particle exchange in the t-channel. The reason for our studying this approximation is that, the same approximation made in QED can be shown to be the lowest order resummation scheme that preserves Ward identities (ref:emil ).
The DDSA approximation can be derived from an effective action by modifying slightly the approach of Cornwall, Jackiw and Tomboulis (CJT)ref:CJT . The discussion that follows here is due to Emil Mottola and Luis Bettencourtref:emil . Thinking of the fields $`x`$ and $`\chi `$ as part of an $`N+1`$ component field, and considering the case that $`\widehat{x}(t)=0`$ where there is no mixed propagator, one can write a CJT like action for the generating functional of the twice Legendre transformed effective action as:
$$\begin{array}{c}\mathrm{\Gamma }[\chi ,𝒢,𝒟]=S_{\text{class}}[\chi ]+\frac{i}{2}\mathrm{Tr}\{\mathrm{ln}[𝒟^1]\}\hfill \\ \hfill +\frac{i}{2}\mathrm{Tr}\{\mathrm{ln}[𝒢^1]\}+\frac{i}{2}\mathrm{Tr}\{𝒟_0^1𝒟+G^1[\chi ]𝒢1\}+\mathrm{\Gamma }_2[𝒢].\end{array}$$
(67)
here $`G^1(t,t^{})`$ is defined by (16) and $`𝒟_0(t,t^{})`$ by Eq. (62). $`𝒟_0(t,t^{})`$ is considered an *external* parameter, and is not varied to obtain the equations of motion. In the DDSA, the 2-PI contribution to the action, $`\mathrm{\Gamma }_2[𝒢]`$, for the case when $`x_i(t)=0`$, is given by Eq. (59) with $`𝒟(t,t^{})`$ set equal to its LOLN value $`𝒟_0(t,t^{})`$:
$$\begin{array}{c}\mathrm{\Gamma }_2[𝒢]=\hfill \\ \hfill \frac{1}{4N}\underset{ij}{}_𝒞dt_1_𝒞dt_2𝒟_0(t_1,t_2)𝒢_{ij}(t_1,t_2)𝒢_{ji}(t_2,t_1).\end{array}$$
(68)
By varying the action (67), we reproduce Eqs. (63) and (65). Although there is an effective action for the DDSA approximation, since $`𝒟_0`$ is treated as an external time-dependent propagator, energy conservation is violated at order $`1/N`$. At modest $`N`$ we will find that this causes this approximation to become inaccurate after several oscillations. However, it is more accurate at these modest values of $`N`$ than the LOLN approximation, as well avoiding the unboundedness of the NLOLN approximation we discuss next.
## VII The large-$`N`$ approximation
The large-$`N`$ expansion is obtained from Eq. (13) by first integrating over all the $`x_i`$ and then evaluating the remaining functional integral over $`\chi `$ by steepest descent. The effective action, as a power series in $`1/N`$, is obtained from the first Legendre transform of the generating functional. In a previous paperref:ctpN , we obtained equations for the next to leading order large-$`N`$ approximation (NLOLN) to the action, and gave numerical results for the quantum roll. For completeness, we review those equations here. To order $`1/N`$, we obtain:
$$\begin{array}{c}\mathrm{\Gamma }_{\text{Large-N}}[x]=S_{\text{class}}[x]\hfill \\ \hfill +_𝒞dt\left\{\frac{i}{2}\underset{i}{}\mathrm{ln}[G_{ii}^1(t,t)]+\frac{i}{2N}\mathrm{ln}[𝒟_0^1(t,t)]\right\},\end{array}$$
where $`S_{\text{class}}[x]`$ is given by Eq. (41), and $`𝒟_0^1(t,t^{})`$ is the inverse propagators for $`\chi `$ in lowest order in the $`1/N`$ expansion, given by
$$𝒟_0^1(t,t^{})=D^1(t,t^{})+\mathrm{\Pi }_0(t,t^{}),$$
(69)
with
$`\mathrm{\Pi }_0(t,t^{})`$ $`={\displaystyle \frac{i}{2N}}{\displaystyle \underset{i,j}{}}G_{ij}(t,t^{})G_{ji}(t^{},t)`$
$`{\displaystyle \underset{i,j}{}}x_i(t)G_{ij}(t,t^{})x_j(t^{}).`$ (70)
Here $`D^1(t,t^{})`$ and $`G_{ij}^1(t,t^{})`$ are the same as Eqs. (16) that we defined earlier.
The equations of motion for the classical fields $`x_i(t)`$, to next to leading order in $`1/N`$, are given by:
$$\begin{array}{c}\left\{\frac{\mathrm{d}^2}{\mathrm{d}t^2}+\chi (t)\right\}x_i(t)\hfill \\ \hfill +i\underset{j}{}_𝒞dt^{}G_{ij}(t,t^{})𝒟_0(t,t^{})x_j(t^{})=0,\end{array}$$
(71)
with the gap equation for $`\chi (t)`$ given by
$$\chi (t)=\frac{g}{2}\left\{\underset{i}{}\left(x_i^2(t)+\frac{1}{N}\underset{i}{}𝒢_{ii}^{(2)}(t,t)/i\right)r_0^2\right\},$$
(72)
and where the second order $`x_i`$ propagator $`𝒢_{ij}^{(2)}(t,t)`$ and self energy $`\mathrm{\Sigma }_{ij}(t,t^{})`$ to order $`1/N`$ is given by:
$$\begin{array}{c}𝒢_{ij}^{(2)}(t,t^{})=G_{ij}(t,t^{})\hfill \\ \hfill \underset{k,l}{}_𝒞dt_1_𝒞dt_2G_{ik}(t,t_1)\mathrm{\Sigma }_{kl}(t_1,t_2)G_{lj}(t_2,t^{}),\end{array}$$
(73)
where
$$\mathrm{\Sigma }_{ij}(t,t^{})=\frac{i}{N}G_{ij}(t,t^{})𝒟_0(t,t^{})x_i(t)𝒟_0(t,t^{})x_j(t^{}).$$
We see here that the equation for $`𝒢`$ is the expansion of the BVA equation in a series of $`1/N`$, truncated at first order.
Let us now specialize to the case of the quantum roll problem where $`x_i(t)=0`$. In that case the two point inverse propagator for the $`x`$ field is
$`𝒢_{ij}^1[\chi ](t_1,t_2)`$ $`={\displaystyle \frac{\delta ^2\mathrm{\Gamma }_{\text{Large-N}}[x,\chi ]}{\delta x_i(t_1)\delta x_j(t_2)}}`$
$`=G_{ij}^1[\chi ](t_1,t_2)+\mathrm{\Sigma }_{ij}[\chi ](t_1,t_2),`$
with
$$\mathrm{\Sigma }_{ij}[\chi ](t,t^{})=\frac{i}{N}G_{ij}(t,t^{})𝒟_0(t,t^{})$$
However it is $`𝒢_{ij}^{(2)}(t,t^{})`$ which enters into Eq. (72) and not $`𝒢_{ij}(t,t^{})`$. Thus the solution for $`𝒢_{ij}(t,t^{})`$, which we might interpret as $`\widehat{x}_i(t)\widehat{x}_j(t^{})`$, does not enter into the dynamics of the solution! *This* $`𝒢_{ii}(t,t)`$ is positive definite, but quickly blows up.
The vertex function $`\mathrm{\Gamma }_{ij}(t_1,t_2,t_3)`$ is given by:
$`\mathrm{\Gamma }_{ij}(t_1,t_2,t_3)`$ $`={\displaystyle \frac{\delta 𝒢_{ij}^1[\chi ](t_1,t_2)}{\delta \chi (t_3)}}`$ (74)
$`=\delta _𝒞(t_1,t_2)\delta _𝒞(t_2,t_3)\delta _{ij}{\displaystyle \frac{i}{N}}G_{ij}(t_1,t_3)G_{ji}(t_3,t_2)𝒟_0(t_2,t_1)`$
$`{\displaystyle \frac{i}{N}}{\displaystyle _𝒞}dt_4{\displaystyle _𝒞}dt_5G_{ij}(t_1,t_2)𝒟_0(t_1,t_4)\mathrm{\Lambda }_0(t_4,t_5,t_3)𝒟_0(t_5,t_2),`$
where the lowest order in $`1/N`$ 3-$`\chi `$ vertex is given by
$`\mathrm{\Lambda }_0(t_4,t_5,t_3)`$ $`={\displaystyle \frac{\delta 𝒟_0^1(t_4,t_5)}{\delta \chi (t_3)}}`$
$`={\displaystyle \frac{i}{N}}{\displaystyle \underset{ijk}{}}G_{ij}(t_4,t_3)G_{kl}(t_3,t_5)G_{li}(t_5,t_4).`$
We immediately see that this is not an integral equation but again, is the lowest order in $`1/N`$ contribution to Eq. (57).
The inverse $`\chi `$ propagator gets $`1/N`$ corrections which are of two types, one is a self energy correction to the $`x`$ propagator and the other is a new three loop graph containing two lowest order $`\chi `$ propagators. We find
$`𝒟^1(1,2)`$ $`={\displaystyle \frac{\delta ^2\mathrm{\Gamma }_{\text{Large-N}}[x,\chi ]}{\delta \chi (t_1)\delta \chi (t_2)}}`$
$`={\displaystyle \frac{1}{g}}\delta _𝒞(1,2)\mathrm{\Pi }_0(1,2){\displaystyle \underset{ijkl}{}}{\displaystyle _𝒞}dt_3{\displaystyle _𝒞}dt_4G_{ij}(t_1,t_3)\mathrm{\Sigma }_{jk}(t_3,t_4)G_{kl}(t_4,t_2)G_{li}(t_2,t_1)`$
$`+{\displaystyle _𝒞}dt_3{\displaystyle _𝒞}dt_4{\displaystyle _𝒞}dt_5{\displaystyle _𝒞}dt_6\mathrm{\Lambda }_0(t_4,t_1,t_3)𝒟_0(t_3,t_5)\mathrm{\Lambda }_0(t_5,t_2,t_6)𝒟_0(t_6,t_4).`$
The last term in this equation is a $`1/N`$ correction to the vertex function. However, it is $`𝒟_0`$ and not $`𝒟`$ which enters Eq. (73), so that the BVA and the $`1/N`$ expansion will differ only by terms of order $`1/N^2`$. The BVA approximation treats $`x`$ and $`\chi `$ on exactly the same footing, whereas the large-$`N`$ expansion treats $`x`$ exactly, but then expands in loops of $`\chi `$. So at order $`1/N^2`$, the large-$`N`$ approximation will contain graphs omitted from the BVA approximation, and vice-versa.
## VIII Results and Conclusions
In this section we present the results of exact numerical simulations of the quantum roll, using initial conditions described in our previous paper on the large-$`N`$ approximationref:paper1 . We choose as our dimensional mass scale the second derivative of $`U(r)`$ at the minimum of the effective one dimensional potential $`U(r)`$. This mass scale was chosen to have value $`M^2=2`$. In terms of this mass scale, the coupling constant as well as the rescaled $`r_0`$ are of order one for all $`N`$. The exact manner in which $`g`$ and $`r_0`$ runs with $`N`$ is described in ref. ref:paper1 .
As $`N\mathrm{}`$ the Hartree and leading order large-$`N`$ approximation become exact and an initially gaussian wave packet remains gaussian with width equal to $`x^2(t)`$ oscillating in a known manner. At modest $`N`$, $`10<N<20`$ an initially Gaussian wave function develops a large number of nodes and so the wave function even at modest times is of the form Gaussian time a high order polynomial. In spite of this, $`x^2(t)`$ shows rather simple behavior. It oscillates with a constant amplitude for a reasonable period of time with an envelope that oscillates with a much longer time constant which increases with $`N`$. The Hartree and leading order large-$`N`$ approximations just oscillate with fixed amplitude. The NLOLN blows up in this regime. BVA attempts to track the contraction of the envelope but then contracts to a fixed point. The DDSA violates energy conservation at order $`1/N`$ so it becomes numerically inaccurate when $`1/N`$ effects become important which is at a time $`tN`$. Both BVA and DDSA do however stay bounded and positive definite during the time period of our numerical simulations. Higher order correlation functions show more complicated behavior and the approximations presented here are only accurate for a few oscillations in the regime $`3<N<20`$ consistent with the increasingly complicated evolving structure of the wave function.
In Figs. 4 through 6, we show the results for $`x^2(t)`$ as a function of $`t`$, comparing the bare vertex, the dynamic Debye screening, and the large-$`N`$ approximations to the exact solution, for $`N=3`$, $`10`$, and $`21`$. In Figs. 7 to 8, we show the same results for $`\chi (t)`$ as a function of $`t`$, and in Figs. 9 through 11, we give the results for $`\chi ^2(t)`$ \[For detailed views of these figures in color, see our web site at: http://www.theory.unh.edu/resum\].
In our previous studiesref:paper1 of the large-$`N`$ approximation, we found that the next to leading order large-$`N`$ approximation had the feature that the effective potential was not defined at small $`x`$ for $`N20`$, for our parameter set, and it was not until $`N`$ was greater than about $`20`$ that the large-$`N`$ expansion produced bounded values for $`x^2(t)`$. This result is reproduced here. For the limit $`N\mathrm{}`$ the quantity $`x^2(t)`$ corresponds to harmonic oscillations. At finite $`N`$, however, the exact solution for $`N21`$ has the property that the envelope of these oscillations contracts. As noted in the figures, only the bare vertex approximation attempts to follow this contraction. At $`N=21`$, the BVA is accurate up to a $`t130`$ before overshooting and then oscillating about a fixed point. This fixed point behavior shows that this approximation still neglects some important quantum phase information present in the exact solution.
In contrast to the NLOLN approximation, which breaks down for $`N<21`$, both the BVA and the DDSA have the feature that $`x^2(t)`$ remains positive definite, as well as being bounded at all $`N`$. This is true for all the expectation values that contribute to the energy. This conclusion is purely based on numerical evidence. We do not have a proof that this approximation corresponds to a positive definite probability distribution. However, all the moments we have studied (a total of five, as shown in Fig. 12), are all bounded.
The DDSA is more accurate than the second order large-$`N`$ approximation for $`N`$ less that $`20`$, but for $`N`$ greater than $`20`$, the reverse becomes true. However, neither approximation captures the true nonlinear shrinking of the envelope of the oscillations, even for $`N`$ greater than $`20`$.
Energy is conserved for the bare vertex and the second order large-$`N`$ approximations, but not for the dynamic Debye screening approximations, as pointed out in section VI. This is a serious drawback to the dynamic Debye screening approximation.
In all these figures, one can see that the bare vertex approximation tries to follow the envelope of the exact curve, whereas the dynamic Debye screening approximation does not do so. This is particularly striking for the cases when $`N`$ is less than $`21`$, where the dynamic Debye screening approximation yield unphysically large values for the expectation values.
In the BVA approximation we observe that $`x^2(t)`$ at late times has an envelope of decreasing oscillations about a fixed point. In fact as seen in Fig. 12 all the contributions to the energy in the BVA have the same feature that they asymptote to a fixed point. In Fig. 12 we display all five contributions to the energy at $`N=10`$ to demonstrate this fact. In contrast, as seen in the very long time run shown in Fig. 11, the exact solutions exhibit “recurrence” patterns of motion which are not captured in the BVA. In the $`1+1`$ dimensional field theory simulations of ref. ref:berges , all the Fourier components of the two particle correlation function showed this behavior which was given as evidence for thermalization. So one hopes that this “defect” of the BVA approximation in a quantum mechanics setting, will instead have the correct physics of thermalization in a field theory application where Poincaré recurrence times are expected to become very large. To see if this is true, we intend to study the BVA in classical 1+1 dimensional field theory where again exact simulations can be performedabw .
In summary we have found that both resummation methods described here, the BVA and the DDSA, produce positive definite and apparently bounded results for expectation values at all values of $`N`$. The bare vertex approximation appears to provide the best description of the motion, but cannot describe recurrences of the motion. Still, it provides an energy conserving and reasonably accurate description, and is a dramatic improvement over the next to leading order large-$`N`$ approximation when $`N<N_{\text{crit}}=21`$. As mentioned earlier, in the single particle quantum mechanics problem we studied here, the graphs do not correspond to particle collisions, so there is no possibility of studying thermalization. Thermalization questions need to be addressed in field theory applications. It will be important to show that the BVA approximation will lead to thermalization of arbitrary initial data as found in the 3-loop approximation of ref. ref:berges when applied to 1+1 dimensional quantum field theory. We would also like to study the analogue of the BVA approximation for a gaussian ensemble of initial conditions for a 1+1 dimensional classical field theory since that can also be studied exactly numericallyabw . These authors have shown that the classical field theory indeed thermalizes and we would like to know how accurately the classical version of our approximation captures this physics. This will be the subject of a future publication.
###### Acknowledgements.
We wish to thank Salman Habib for helpful discussions on understanding the numerical simulations and for continued advice. We wish to thank Prof. Gabor Kalman for explaining relevant plasma conductivity approximations, and Emil Mottola for suggesting our study of the “dynamic Debye screening” approximation and explaining its derivation from the CJT formalism. We also thank Juergen Berges for discussing with us his recent results on thermalization in a related approximation. JFD is supported in part by the U.S. Department of Energy under grant DE-FG02-88ER40410. He would like to thank the T-8 theory group at LANL, and the Institute for Nuclear Theory at the University of Washington, for hospitality during the course of this work. FC would like to thank Boston College and UNH for hospitality during the course of this work.
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# 1 INTRODUCTION
## 1 INTRODUCTION
One of the central results that has emerged from the studies of the AdS/CFT correspondence between $`5`$-dimensional gravity in the AdS bulk and $`4`$-dimensional Yang-Mills theory on its boundary is the identification of the renormalization scale of the latter with the radial coordinate of the AdS bulk and the radial evolution of the $`5`$-dimensional fields with RG flows of the couplings in the $`4`$-dimensional Yang-Mills theory . A conjectured generalization of this correspondence to $`4`$-dimensional boundary theories which include gravity has also recently been studied in the context of Randall-Sundram type brane-world scenarios and the cosmological constant problem .
The existence of a connection between RG flows in a D-dimensional theory, which includes gravity, and gravitational equations in (D+1)-dimensions, was recognized and pointed out quite sometime back in the perturbative studies of noncritical string theory . Perturbative noncritical string theory is formulated as D-dimensional matter coupled to 2-d quantum gravity . As is well-known, in this formulation of string theory the extra coordinate of the (D+1)-dimensional space is related to the conformal degree of freedom of the world-sheet metric and world-sheet gravitational dressing of the various $`\sigma `$-model couplings gives rise to their dependence on this extra coordinate. A connection between the RG flows of the D-dimensional fields and gravitational equations in (D+1)-dimensions arises because the dependence of the $`\sigma `$-model couplings on the conformal mode of the world-sheet metric is determined by gravitational equations in (D+1)-dimensions. In recent years this connection has been made precise in the context of AdS/CFT correspondence. The purpose of this note is to reexamine and further expand on the world-sheet approach of noncritical string theory to holographic RG in the light of these recent advances. The main advantage of this approach is that it provides a natural setting for the discussion of a generic holographic RG connection between D-dimensional boundary theories that contain gravity and (D+1)-dimensional gravitational dynamics. The world-sheet approach also provides a systematic handle on stringy (i.e. $`\alpha ^{}`$) as well as string loop corrections.
The key organizing principle in the first quantized approach to noncritical string theory is the requirement of world-sheet reparametrization invariance. In the background gauge-fixing method, a prescription for integrating over the 2-d metric which ensures this requirement of reparametrization invariance automatically ensures Weyl invariance with respect to the 2-d fiducial metric. In this approach, therefore, Weyl invariance with respect to the 2-d fiducial metric emerges as the principal consistency requirement which is needed to ensure world-sheet reparametrization invariance. For example, it is this requirement that determines the gravitational dressing of the $`\sigma `$-model couplings.
In this note we will consider the $`\sigma `$-model partition function of noncritical strings propagating in background fields. This partition function is in general not well-defined because of divergent contributions to it arising from correlators of microscopic loop operators whose liouville wavefunctions are not normalizable. Although this is an ultraviolet (small area) divergence on the world-sheet, from the (D+1)-dimensional target space point of view it is an infrared (large volume) divergence. It can be regularized by introducing a cut-off on the integration over the liouville zero mode. This is very much like the infrared regulator needed in the radial direction to evaluate the on-shell gravity action in AdS space. This way of regularizing the partition function introduces a “boundary” in the liouville direction. The regularized partition function depends on the location of the liouville boundary only implicitly through the values of the dressed couplings at the boundary. We will show here that a change in the location of the boundary gives rise to an RG flow equation for the partition function which looks exactly like the Hamilton-Jacobi constraint which an on-shell boundary gravitational action is expected to satisfy .
The plan of this paper is as follows. In the next section we first summarize the main results from the first quantized approach to noncritical string theory as D-dimensional matter coupled to 2-d gravity. We then discuss in detail the interpretation of the dependence of the $`\sigma `$-model couplings on the extra coordinate as RG dependence in the D-dimensional theory. In Sec. 3 we argue that the $`\sigma `$-model couplings should more correctly be interpreted as defining a boundary value problem in a (D+1)-dimensional gravity theory. We explain how the boundary arises from the need to regularize world-sheet ultraviolet divergences in the calculation of the partition function. We then show that, consistent with its interpretation as an on-shell boundary action, the regularized partition function staisfies a flow equation which looks very much like a Hamilton-Jacobi constraint equation which an on-shell boundary gravitational action is expected to satisfy. We end in Sec. 4 with some concluding remarks.
## 2 NONCRITICAL STRINGS AND HOLOGRAPHIC RG
In this section we first briefly summarize some old and rather well-known results from noncritical string theory. We then discuss the RG scale dependence interpretation of the gravitational dressing of the couplings. For simplicity we restrict the discussion to bosonic string, but extension to superstring is straightforward.
The starting point of the first quantized approach to noncritical strings in background fields is the world-sheet reparametrization invariant action
$`S`$ $`=`$ $`{\displaystyle \frac{1}{8\pi \alpha ^{}}}{\displaystyle }d^2\xi \sqrt{g}[_\alpha X^\mu _\beta X^\nu (g^{\alpha \beta }G_{\mu \nu }(X(\xi ))+ϵ^{\alpha \beta }B_{\mu \nu }(X(\xi )))`$ (2.1)
$`\text{ }+\alpha ^{}R^{(2)}\mathrm{\Phi }(X(\xi ))+T(X(\xi ))+\mathrm{}].`$
Here $`x^\mu `$’s, which are the zero modes of $`X^\mu (\xi )`$’s, parametrize a D-dimensional space with metric $`G_{\mu \nu }(x)`$ and other fields. The Polyakov path integral formally defines the partition function
$`Z[G_{\mu \nu },\mathrm{\Phi },B_{\mu \nu },\mathrm{}]={\displaystyle [𝒟g_{\alpha \beta }][𝒟X^\mu ]e^S}`$ (2.2)
which is a functional of the D-dimensional couplings $`G_{\mu \nu }`$, $`\mathrm{\Phi }`$, $`B_{\mu \nu }`$, etc.
Gravitational Dressing
In the quantum theory the various $`\sigma `$-model couplings get dressed by 2-d gravity. A reparametrization invariant prescription for determining these gravitational dressings is the following. One first fixes the conformal gauge
$$g_{\alpha \beta }=e^{\varphi (\xi )}\widehat{g}_{\alpha \beta }.$$
Here $`\varphi (\xi )`$ is the liouville mode and $`\widehat{g}_{\alpha \beta }`$ is a fiducial metric that depends on the moduli of the Riemann surface over which the action in (2.1) is defined. One then makes a transformation in the functional integral from the liouville mode $`\varphi (\xi )`$ to a field $`\eta (\xi )`$ with gaussian measure which, in the absence of the background fields, has the following action
$`{\displaystyle \frac{1}{8\pi }}{\displaystyle d^2\xi \sqrt{\widehat{g}}\left(\widehat{g}^{\alpha \beta }_\alpha \eta _\beta \eta +Q\widehat{R}^{(2)}\eta \right)}`$ (2.3)
where $`Q=\sqrt{(25D)/3}`$. When background fields are switched on, in the presence of 2-d gravity they get dressed, that is they become functions of $`\eta `$ . Thus
$$G_{\mu \nu }(x)G_{\mu \nu }(x,\eta ),\mathrm{\Phi }(x)\mathrm{\Phi }(x,\eta ),\mathrm{}$$
The $`\eta `$-dependence of the various fields is fixed by demanding that the above procedure preserve world-sheet reparametrization invariance. In particular, this means that the final results should be invariant under Weyl transformations of the fiducial metric $`\widehat{g}_{\alpha \beta }`$. This leads to the familiar beta-function equations for the dressed fields
$`0`$ $`=`$ $`R_{MN}+2_M_N\mathrm{\Phi }{\displaystyle \frac{1}{4}}H_{MPL}H_{N}^{}{}_{}{}^{PL}+\mathrm{}`$ (2.4)
$`0`$ $`=`$ $`^PH_{PMN}2^P\mathrm{\Phi }H_{PMN}+\mathrm{}`$ (2.5)
$`0`$ $`=`$ $`{\displaystyle \frac{D25}{3\alpha ^{}}}R^{(D+1)}4^P_P\mathrm{\Phi }+4_P\mathrm{\Phi }^P\mathrm{\Phi }`$ (2.6)
$`+{\displaystyle \frac{1}{12}}H_{MNP}H^{MNP}+\mathrm{}`$
etc. Here the indices $`M,N`$ etc. run over $`\mu ,\eta `$ and the dots represent $`\alpha ^{}`$ and string loop corrections. These equations have (D+1)-dimensional general covariance.
RG Flows
Although we considered D-dimensional matter in the above discussion, these considerations actually apply to any matter coupled to 2-d gravity. A term in the action of the form
$$\underset{i}{}d^2\xi \sqrt{g}\lambda ^iO_i(X(\xi ),g_{\alpha \beta }(\xi )),$$
where $`O_i(X(\xi ),g_{\alpha \beta }(\xi ))`$ is a local operator constructed from the matter fields $`\{X(\xi )\}`$, gets dressed to
$$\underset{i}{}d^2\xi \sqrt{\widehat{g}}\lambda ^i(\eta (\xi ))O_i(X(\xi ),\widehat{g}_{\alpha \beta }(\xi ))$$
and the $`\eta `$-dependence of the dressed coupling, $`\lambda ^i(\eta )`$, is determined by an appropriate vanishing beta-function condition. In case the couplings $`\lambda ^i`$ correspond to a CFT coupled to 2-d gravity, the $`\lambda ^i(\eta )`$ are independent of $`\eta `$. This identifies CFT’s as special points in the space spanned by the set of all possible couplings $`\lambda ^i`$, the so-called theory space. The more general case in which the couplings get dressed can be interpreted as giving RG flows between these special points corresponding to CFT’s.
It is important to emphasize here that the RG flow that we are talking about is not due to changes of cut-off in the 2-d QFT of $`X(\xi )`$’s and $`\eta (\xi )`$. Although this cut-off is needed to do computations, the vanishing beta-function conditions ensure that the couplings do not depend on it. The RG flows that we are talking about are similar to finite size scaling. The size is here provided by the invariant area of the world-sheet, or its conjugate, the world-sheet cosmological constant . In the conformal gauge, the flows thus correspond to the response of the couplings to changes of the physical scale brought about by shifts of the liouville mode, and hence of $`\eta `$. The vanishing beta-function equations describe how the dressed couplings change with precisely these shifts of $`\eta `$. In general, shifts of $`\eta `$ do not produce any simple scalings of the couplings, except near the points in theory space described by CFT’s. The trajectories in theory space given by the dressed couplings describe RG flows between two such points. An example of such an RG flow between two $`c<1`$ minimal models has been discussed in detail in where an explicit kink solution is given which describes RG flow due to a nearly marginal perturbation. <sup>1</sup><sup>1</sup>1RG flows in two-dimensional theories coupled to gravity have also been considered in .
For D-dimensional matter coupled to 2-d gravity, constant shifts of $`\eta `$ describe RG flows in a D-dimensional effective theory of gravity. To see this, consider the (D+1)-dimensional gravity theory, obtained after dressing by 2-d gravity, in the gauge <sup>2</sup><sup>2</sup>2As we shall discuss in the next section, it is natural to guess that the partition function of noncritical strings in background fields evaluates a (D+1)-dimensional gravitational action on-shell in this gauge.
$`G_{\eta \eta }=1,G_{\eta \mu }=0.`$ (2.7)
In this gauge, the (D+1)-dimensional metric is given by
$$ds_{D+1}^2=d\eta ^2+ds_D^2,$$
where the metric in a constant $`\eta `$ D-dimensional slice is given by
$$ds_D^2=G_{\mu \nu }(x,\eta )dx^\mu dx^\nu .$$
A shift in $`\eta `$ produces a change in the D-dimensional metric which is dictated by the (D+1)-dimensional gravitational equations. In general, such a change generates a local change of scale in the D-dimensional world and hence corresponds to a local generalization of the usual RG flows. In the particular case that
$$G_{\mu \nu }(x,\eta )=\mathrm{\Omega }(\eta )G_{\mu \nu }(x)$$
and $`\mathrm{\Omega }`$ is a monotonic function of $`\eta `$, shifts in $`\eta `$ give global changes of scale in the D-dimensional world and hence in this case we recover the standard RG flows in the D-dimensional theory.
## 3 THE FLOW EQUATION
In this section we will argue that the partition function for noncritical strings evaluates a (D+1)-dimensional gravitational action on-shell, in the gauge (2.7), for solutions with a boundary. Consistent with this interpretation, we will derive a flow equation for the partition function, which looks just like the Hamilton-Jacobi constraint that the on-shell gauge-fixed (D+1)-dimensional action must satisfy for ensuring full (D+1)-dimensional general covariance.
The partition function is formally given by the functional integral in (2.2) where $`S`$ is the action in (2.1). Now, it is well-known that the equations in (2.4)-(2.6) can be derived from a (D+1)-dimensional gravitational action. Since these equations also determine the 2-d gravitational dressing of the $`\sigma `$-model couplings, it would seem consistent to identify the partition function in (2.2) with the (D+1)-dimensional gravitational action evaluated on-shell. There is, however, a problem in this identification, which we will now discuss.
The set of equations (2.4)-(2.6) are partial differential equations of second or higher order, depending on the order to which terms are retained in $`\alpha ^{}`$ expansion. At the lowest order in $`\alpha ^{}`$, the equations are second order differential equations in $`\eta `$ for the dressed couplings. A general solution for these equations depends on two independent functions of $`x`$ for each of the background fields. For example, for the dilaton the general solution depends on $`\mathrm{\Phi }_0(x)`$ and $`\mathrm{\Phi }_0^{}(x)`$, which may be taken to be respectively the value of the dressed field $`\mathrm{\Phi }(x,\eta )`$ at some point $`\eta =\eta _0`$ and the value of its derivative $`_\eta \mathrm{\Phi }(x,\eta )`$ at $`\eta _0`$. Thus, at the lowest order in $`\alpha ^{}`$, the (D+1)-dimensional on-shell action should depend on two independent functions of $`x`$ for each of the background fields. The partition function in (2.2), however, depends only on one function of $`x`$ for each of the background fields since its dependence on the background fields is inherited from the action in (2.1) which has this property <sup>3</sup><sup>3</sup>3The situation becomes worse when higher order corrections in $`\alpha ^{}`$ are included in the differential equations in (2.4)-(2.6) since then the order in $`\eta `$ derivatives in these equations increases with a corresponding increase in the number of independent functions of $`x`$, for each background field, characterizing the general solution..
One might worry that the functional intergral in (2.2) defining the partition function for the noncritical string is only formal and that a more rigorous definition of this functional integral would show some subtility in counting the number of independent couplings. That in fact there is no proliferation of couplings in a more rigorous setting can be easily seen by reformulating (2.2) in the framework of dynamically triangulated random surfaces. <sup>4</sup><sup>4</sup>4See, for example, . Thus the partition function in (2.2) cannot be identified with a (D+1)-dimensional gravitational action evaluated on-shell for general solutions.
Actually, there are solutions to the differential equations (2.4)-(2.6) which depend on only one function of $`x`$ for each background field. These solutions involve (D+1)-dimensional spaces with a D-dimensional boundary <sup>5</sup><sup>5</sup>5We will assume the D-dimensional boundary metric to have a Euclidean signature, and a regularity condition in the bulk generally picks up only one of the two possible (at the lowest order in $`\alpha ^{}`$) solutions which evolves the boundary data into the bulk. The (D+1)-dimensional action evaluated on-shell for such solutions depends on only one function of $`x`$ for each of the background fields, which together constitute the boundary data.
A simple example of the above is provided by the coupling corresponding to the tachyon field in the noncritical string theory corresponding to the flat space linear dilaton solution of the equations (2.4)-(2.6). The gravitationally dressed tachyon coupling satisfies the following equation in this background :
$`(_\eta ^2Q_\eta +_x^2)T(x,\eta ){\displaystyle \frac{V}{T}}=0,`$ (3.1)
where $`V=T^2+O(T^3)`$ is the tachyon potential. Ignoring the cubic and higher order terms in the potential, and going to the momentum space conjugate to $`x`$, we get the two solutions
$`\stackrel{~}{T}_\pm (k,\eta )=e^{\frac{Q}{2}\eta }\psi _{\pm ,k}(\eta )`$ (3.2)
where $`\psi _{\pm ,k}(\eta )`$ is the liouville wavefunction,
$`\psi _{\pm ,k}(\eta )=e^{\pm \eta \sqrt{(k^2\frac{D1}{12})}}\psi _{\pm ,k}.`$ (3.3)
If we place a D-dimensional “boundary” in the liouville direction at $`\eta _0`$, we see that only the solution $`\psi _{}(\eta )`$ is regular at $`\eta +\mathrm{}`$. The other solution, $`\psi _+(\eta )`$, is regular at $`\eta \mathrm{}`$. The liouville wavefunctions for both these solutions are not normalizable. <sup>6</sup><sup>6</sup>6For $`D1`$ this is true only if $`k^2\frac{D1}{12}`$, and then these wavefunctions correspond to microscopic states. Wavefunctions corresponding to operators with $`k^2<\frac{D1}{12}`$ are normalizable and hence these correspond to macroscopic states. See, for example, for a review of liouville theory.
A more nontrivial example is provided by AdS gravity, which has been extensively studied recently in the context of the AdS/CFT correspondence. In this case the Fefferman-Graham theorem guarantees that there is a unique regular solution that evolves boundary data into the AdS bulk. The regular solution is, however, not normalizable , just like the liouville wavefunctions for microscopic states in the above example of tachyon in flat space and linear dilaton background. For example, the regular solution to linearized equations for a scalar field in the AdS background behaves at large $`r`$ as $`\lambda (x)r^{\mathrm{\Delta }D},`$ where $`\mathrm{\Delta }`$ is the dimension of the corresponding operator in the dual CFT description and $`D`$ is the dimension of the boundary. <sup>7</sup><sup>7</sup>7The other solution, which behaves at large $`r`$ as $`r^\mathrm{\Delta }`$, is normalizable.
We propose to identify the partition function in (2.2) with a “boundary” action evaluated on regular solutions obtained as described above. The boundary in $`\eta `$ is the D-dimensional space parametrized by the $`x^\mu `$’s and the background fields appearing in the $`\sigma `$-model action $`S`$, (2.1), essentially account for the boundary values of the corresponding dressed fields. <sup>8</sup><sup>8</sup>8It is important to note here that the $`\sigma `$-model couplings are not equal to the values of the dressed couplings at the boundary. The former can, however, be traded for the latter once the dressings are known. The example of the dressing of the tachyon in noncritical string theory considered above provides a good illustration of this. The tachyon coupling that enters the $`\sigma `$-model is $`\psi _k`$ which appears on the right hand side of (3.3). This can clearly be traded for $`\psi _k(\eta _0)`$ using (3.3). Moreover, since no independent functions of $`x`$ corresponding to the components $`G_{\eta \eta }`$ and $`G_{\eta \mu }`$ of the (D+1)-dimensional metric appear in $`S`$, we propose that the partition function actually evaluates the boundary action in the gauge (2.7). As evidence for this we point out that, by construction, the partition function has D-dimensional general covariance, as is required by the second of the conditions in (2.7). A more non-trivial check for our proposal is provided by the first gauge-fixing condition which requires the partition function to satisfy a Hamilton-Jacobi type of constraint equation. Later in this section we will derive a flow equation for the partition function which has a remarkable resemblance to such an equation.
The Liouville Boundary
Let us first try to understand how a boundary arises in the liouville direction. As is well-known, in critical string theory, where the 2-d metric is non-dynamical, the background fields appearing in the $`\sigma `$-model action must satisfy the beta function equations for conformal invariance. In the noncritical formulation of string theory, however, the D-dimensional background fields appearing in (2.1) are completely arbitrary, since it is the integration over the liouville mode that now enforces conformal invariance. To make the discussion more general, let us rewrite (2.2) as follows:
$`Z[\lambda ]={\displaystyle [𝒟g_{\alpha \beta }]Z_g[\lambda ]},`$ (3.4)
where
$`Z_g[\lambda ]={\displaystyle [𝒟X^\mu ]exp\left(\underset{i}{}d^2\xi \sqrt{g}\lambda ^i(X(\xi ))O_i(X(\xi ),g_{\alpha \beta }(\xi ))\right)}`$ (3.5)
Here the set $`\{O_i\}`$ forms a complete basis of closed string operators and so the set of couplings $`\{\lambda ^i\}`$ includes all the closed string modes. Now, from the point of view of the D-dimensional matter functional integral in (3.5), the 2-d metric is just an external fiducial metric. As is well-known, in this case for generic couplings $`\lambda ^i`$, $`Z_g[\lambda ]`$ in (3.5) has a conformal anomaly. This anomaly has the effect that in the generic case all the couplings $`\lambda ^i`$ get dressed by 2-d gravity. In perturbation theory, the liouville wavefunctions which give the dressings of the couplings are of two types, microscopic and macroscopic , depending on the operator they couple to. Since the wavefunctions which dress microscopic operators are not normalizable, integration over the liouville zero mode of correlation functions involving a sufficient number of these operators is generically divergent. This is a source of divergent contributions to the partition function, which thus needs a regulator to make it well-defined.
The dressing of the coupling corresponding to the dilaton operator plays a somewhat special role since this determines the effective string coupling. This can bring in other problems, so let us discuss the dressing of the dilaton coupling in some detail.
Let us first separate out the integral over the zero mode of $`x`$ in (3.5) and write
$`Z_g[\lambda ]={\displaystyle d^Dx𝒵_g[\lambda ;x]}.`$ (3.6)
Now, let us consider the dilaton coupling. The dressing of this coupling is controlled by the corresponding beta function, which is essentially the matter central charge, and is determined by the condition that the total central charge of the matter plus 2-d gravity system should vanish. As a result when we fix the conformal gauge and change over from the liouville mode to the variable $`\eta (\xi )`$ in the functional integral of $`𝒵_g[\lambda ;x]`$ over the 2-d metric, we pick up a linear term in the action for $`\eta (\xi )`$, similar to that in (2.3), but now with the coefficient $`Q`$ given by
$`Q`$ $`=`$ $`[{\displaystyle \frac{25D}{3}}+\alpha ^{}(R^{(D)}+4_\mu ^\mu \mathrm{\Phi }4_\mu \mathrm{\Phi }^\mu \mathrm{\Phi }{\displaystyle \frac{1}{12}}H_{\mu \nu \lambda }H^{\mu \nu \lambda })`$ (3.7)
$`+O(\alpha _{}^{}{}_{}{}^{2})]^{\frac{1}{2}}.`$
For flat space we recover the result in (2.3). In general there is a linear term in $`\eta (\xi )`$ even in critical dimensions. Also, since the other beta functions do not vanish for generic couplings, $`Q`$ is in general a function of $`x`$. The generic case is, therefore, difficult to deal with. However, since the couplings are arbitrary, we can choose them to be such that $`Q`$ is a real constant, <sup>9</sup><sup>9</sup>9$`Q`$ real is needed to ensure a space-like interpretation for $`\eta (\xi )`$. independent of $`x`$, or at least sufficiently slowly varying with $`x`$ so that its dependence on $`x`$ can be ignored to the first approximation. Assuming this simplifying choice, we then have an effective string coupling that grows at one end of the $`\eta `$ direction. This is true even in critical dimensions. String perturbation theory, therefore, breaks down because of this strong-coupling singularity.
In some cases this strong-coupling singularity can be removed by generating a potential for $`\eta (\xi )`$ by switching on some additional backgrounds. For example, in the case of flat space in noncritical dimensions, a potential for $`\eta (\xi )`$ is generated by a 2-d cosmological constant term. We will assume here that we are dealing with such a case and an appropriate coupling has been switched on to remove the strong coupling singularity. However, even in cases where the strong-coupling singularity can be removed in this way, generically there is a divergence in the partition function (3.4) which comes from integration over the opposite end of the $`\eta `$ direction where the effective string coupling becomes arbitrarily weak. As we have already remarked, this is because of the contributions to the $`\sigma `$-model partition function of correlators involving microscopic loop operators whose liouville wavefunctions are not normalizable. Thus the partition function in (3.5) diverges for generic couplings, $`\{\lambda ^i\}`$. <sup>10</sup><sup>10</sup>10This happens even when $`Q`$ in (3.7) vanishes. From the world-sheet point of view this divergence is ultraviolet in nature because it comes from 2-d surfaces of small area. However, from the (D+1)-dimensional point of view, this is an infrared divergence since it arises from the infinite volume in the (noncompact) liouville direction. One way of regulating this divergence is by introducing an appropriate cut-off on the integration over the zero mode of $`\eta (\xi )`$. This is how a “boundary” gets introduced in the $`\eta `$ direction, its location being at the value of the cut-off.
Once a boundary has been introduced in this way, we can trade-off the couplings appearing in (3.4) for the boundary values of the dressed couplings. Now, as discussed in the previous section, a shift in the liouville mode generates a local scale transformation on the boundary. Therefore, a shift of the cut-off generates an RG flow in the regularized partition function through the boundary values of the dressed couplings, leading to a flow equation which we will now derive.
The Equation
After fixing the conformal gauge and transforming from the liouville mode to the variable $`\eta (\xi )`$, the partition function in (3.4) may be written as
$`Z[\lambda ;\eta _0]={\displaystyle d^Dx𝒵[\lambda ;\eta _0;x]}`$ (3.8)
where
$`𝒵[\lambda ;\eta _0;x]={\displaystyle _{\eta _0}^{\mathrm{}}}𝑑\eta [\lambda (\eta );x],`$ (3.9)
and $`\eta `$ is the zero mode of $`\eta (\xi )`$ with $`\eta _0`$ the cut-off or the boundary value. On the right hand side of (3.9) we have made it explicit that the $`\eta `$-dependence comes entirely from the dressings of the couplings.
The flow equation can now be derived by making an $`x`$-dependent change in $`\eta _0`$, namely $`\eta _0\eta _0+ϵ(x)`$. <sup>11</sup><sup>11</sup>11This $`x`$-dependent change in $`\eta _0`$ is made possible by the fact that we could have chosen an $`x`$-dependent cut-off on $`\eta `$. This possibility of a local cut-off on $`\eta `$ is compatible with the requirement of world-sheet reparametrization invariance. Denoting the value of the dressed coupling at the boundary by $`\lambda _0^i(x)`$, the flow equation is
$`[\lambda (\eta _0);x]=_{\eta _0}\lambda _0^i(x){\displaystyle \frac{\delta Z}{\delta \lambda _0^i(x)}}`$ (3.10)
This equation follows from the fact that a change in $`𝒵[\lambda ;\eta _0;x]`$ produced by a shift in $`\eta _0`$ can be computed in two different ways. One is directly from the way $`\eta _0`$ appears as an integration limit in (3.9). This gives the left hand side of the equation. The other is by recognizing that $`𝒵[\lambda ;\eta _0;x]`$ depends on $`\eta _0`$ only through the boundary values of the dressed couplings since the $`\eta `$ dependence in $`[\lambda (\eta );x]`$ comes entirely from the liouville dressings of the couplings and the subsequent transformation to the gaussian variable $`\eta `$. This gives the right hand side of the equation.
At the lowest order in $`\alpha ^{}`$, we expect $`[\lambda (\eta );x]`$ to be only quadratic in $`\eta `$-derivatives of the dressed couplings. Let us write this out explicitly as <sup>12</sup><sup>12</sup>12We may assume the standard normalization for the “kinetic” term without any loss of generality. A possible linear term in $`\lambda ^i(x,\eta )`$ can be removed by a field redefinition of the original sigma-model couplings $`\lambda ^i`$. We will assume that this has been done and that the couplings $`\lambda ^i`$ have been chosen accordingly. Also, note that the right hand side of (3.11) is evaluated on-shell in the sense that the $`\eta `$-dressing of the various couplings is determined by the requirement of reparametrization invariance. We also mention that the form of $``$ assumed in (3.11) can be derived close to a CFT point in theory space.
$`[\lambda (\eta );x]={\displaystyle \frac{1}{2}}𝒢_{ij}_\eta \lambda ^i(x,\eta )_\eta \lambda ^j(x,\eta )+V[\lambda (\eta );x]`$ (3.11)
where $`\lambda ^i(x,\eta )`$ is the dressed coupling, $`𝒢^{ij}`$ is
the metric on the space of the couplings and $`V`$ is assumed to have a local expansion in $`x`$-derivatives of $`\lambda ^i(x,\eta )`$. We have made the reasonable assumption that $``$ has a low energy expansion in derivatives of $`\eta `$ and $`x`$. Once (3.11) is given, one can show that the variation of the partition function with respect to the boundary values of the dressed couplings is related to the “velocities” in the standard way,
$`{\displaystyle \frac{\delta Z}{\delta \lambda _0^i(x)}}=𝒢_{ij}_{\eta _0}\lambda _0^j(x).`$ (3.12)
The flow equation (3.10) may then be rewritten as
$`{\displaystyle \frac{1}{2}}𝒢^{ij}{\displaystyle \frac{\delta Z}{\delta \lambda _0^i(x)}}{\displaystyle \frac{\delta Z}{\delta \lambda _0^j(x)}}=V[\lambda _0;x]`$ (3.13)
which is the advertized Hamilton-Jacobi type of constraint equation that the regularized partition function must satisfy.
## 4 CONCLUDING REMARKS
In this note we have presented noncritical string theory as a boundary value problem, based on the observation that the liouville or conformal mode gives rise to an additional dimension. As we have argued, the boundary arises from the cut-off needed to regulate world-sheet ultraviolet divergences. We have shown that, under some reasonable assumptions, the partition function of the noncritical string $`\sigma `$-model action satisfies a Hamilton-Jacobi type of constraint equation as a functional of the boundary values of the $`\sigma `$-model couplings. The dependence of the couplings on the additional dimension is determined by the first order local RG flow equations (3.12). These equations were obtained for the bosonic string, but extension to the superstring is straightforward when RR backgrounds are absent. Since RR backgrounds couple to bilinears of space-time fermions in the $`\sigma `$-model, the analysis becomes complicated when these backgrounds are switched on . For this reason it is difficult to demonstrate explicitly that a Hamilton-Jacobi type of constraint equation continues to be satisfied in the presence of RR backgrounds, although we expect this to be the case. Finally we mention that the structure of the solution space of the RG flow equations (3.12) is presently not known . In order to address this issue it would be worthwhile to discuss the global topology of the RG flows along the lines presented in where the global topology of a class of $`c<1`$ models was exactly calculated using methods of Morse theory.
Acknowledgements
One of us (SRW) would like to thank Theory Division, CERN, for hospitality during a visit when part of this work was done.
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# Theory of interlayer tunneling in bi-layer quantum Hall ferromagnets
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## Abstract
Spielman et al. have recently observed a large zero-bias peak in the tunnel conductance of a bi-layer system in a quantum Hall ferromagnet state. We argue that disorder-induced topological defects in the pseudospin order parameter limit the peak size and destroy the predicted Josephson effect. We predict that the peak would be split and shifted by an in-plane magnetic field in a way that maps the dispersion relation of the ferromagnet’s Goldstone mode. We also predict resonant structures in the DC I-V characteristic under bias by an ac electric field.
\]
Exotic effects induced by inter-layer Coulomb interactions have made strongly coupled bi-layer quantum Hall systems at total Landau level filling factor $`\nu =1`$ the subject of numerous theoretical and experimental studies. When the layers are widely separated they behave as two weakly coupled $`\nu =1/2`$ composite fermion metals. However, when the inter-layer distance $`d`$ is smaller than about twice the magnetic length $`\mathrm{}`$, the system spontaneously develops interlayer phase coherence and forms a $`\nu =1`$ quantum Hall state. This broken symmetry state may be described as a Bose condensate in a bosonic Chern-Simons field theory, as an easy-plane ferromagnet in a theory based on a pseudospin representation for the layer degree of freedom, or as a superfluid excitonic condensate in a theory based on a single-layer particle-hole transformation. We use the pseudospin language below.
In a recent experiment, Spielman et al. observed a qualitative change in the voltage dependence of the interlayer tunneling current $`I(V)`$ upon entering the ordered state. For large $`d/\mathrm{}`$, bilayer $`\nu =1`$ systems exhibit a pseudogap behavior: the tunneling current is extremely small at low bias voltages. This suppression of tunneling is attributed to the slow relaxation of charge characteristic of the $`\nu =1/2`$ state in each layer. In the ordered state, Spielman et al. discovered a strong and sharp zero bias peak in the differential conductance $`dI/dV`$. It appears plausible that this peak is related to the Josephson effect predicted by Wen and Zee and by Ezawa and Iwazaki. In contrast to the conventional Josephson effect, however, no zero-bias supercurrent (infinite tunneling conductance) was observed. The peak conductance, though enormously enhanced, did not exceed $`10^2e^2/h`$.
In this Letter we analyze tunneling in the bi-layer quantum Hall $`\nu =1`$ state. We explain how long range density inhomogeneities introduce topological defects (merons) into the $`SU(2)`$ pseudospin order parameter. These defects carry both charge and vorticity and constitute a dissipative environement which turns the Josephson effect into a finite tunneling peak whose height and width is a measure of the dynamics of the topological defects. We predict dependences of the tunneling current on in-plane magnetic field strength $`B_{||}`$, bias voltage frequency, and on the homogeneity of the 2D layers. In particular, we show that a measurement of $`I(V,B_{||})`$ would test the main premise of our theory, the existence of one low energy excitation mode, and would map the dispersion relation of that mode. Finally, we analyze the current distribution for a perfectly homogeneous sample.
The order parameter field of the quantum Hall ferromagnet is a pseudo-spin unit vector $`\stackrel{}{m}`$. When fluctuations out of the easy plane are small, it can be parametrized by an angle $`\phi `$ and the conjugate ‘charge’ $`m_z`$: $`\stackrel{}{m}=(\mathrm{cos}\phi ,\mathrm{sin}\phi ,m_z)`$. In the absence of tunneling, disorder and topological defects, the long wavelength Hamiltonian density of the $`\nu =1`$bi-layer state is
$$H=\frac{1}{2}\rho _s(\phi )^2+\frac{\left(en_0m_z/2\right)^2}{2\mathrm{\Gamma }},$$
(1)
where $`n_0=\frac{1}{2\pi \mathrm{}^2}`$ is the average density. In Hartree-Fock theory $`\rho _s0.4`$K and the capacitance $`\mathrm{\Gamma }`$ is increased from its electrostatic value. Since the momentum density conjugate to $`\phi `$ is $`p_\phi =\mathrm{}n_0m_z/2`$, the Hamiltonian (1) has a single linearly dispersing collective mode with velocity $`u=\sqrt{\rho _s/\mathrm{\Gamma }}`$. This Goldstone mode signals superfluidity for in-plane currents which are antisymmetric in the layer index. Taking proper account of the significant exchange enhancement of $`\mathrm{\Gamma }`$ yields $`u0.1e^2/\mathrm{}ϵ`$. Eq. (1) is valid only in the limit $`q0`$. Away from that limit this collective mode has a more complicated dispersion, denoted by $`\omega _q,`$ which shows roton effects. It is this dispersion curve that may be extracted from a measurement of $`I(V,B_{||})`$.
The inter-layer tunneling operators are
$$T_\pm =d^2r\lambda (\stackrel{}{r})e^{\pm i\phi (\stackrel{}{r})}e^{\pm iQ_\mathrm{B}x},$$
(2)
where the $`\pm `$ sign refers to the direction of tunneling, $`Q_\mathrm{B}=\frac{edB_{}}{\mathrm{}c}`$ is a characteristic wave vector introduced by the magnetic field $`B_{||}`$ (we choose the gauge $`\stackrel{}{A}_{}=xB_{}\widehat{z}`$). The quantity $`\lambda =\frac{1}{8\pi \mathrm{}^2}\mathrm{\Delta }_{\mathrm{SAS}}`$ is proportional to the tunneling amplitude and may vary with position due to disorder in the tunnel barrier. Here we do not discuss this source of disorder since, on its own, it can not destroy the Josephson effect. As in a Josephson junction, the tunneling term in the Hamiltonian is $`T_++T_{},`$ while the tunneling current operator is $`ie(T_+T_{})/\mathrm{}`$. The striking similarity of the expressions above to their counterparts in superconducting Josephson junctions make it clear that a calculation of the tunneling conductance under Eqs. (1) and (2) leads to a Josephson effect, in contrast to the experiment. We now explain the way that disorder destroys the Josephson effect in the present system.
In the $`\nu =1`$ bi-layer system, a deviation of the total density from $`\nu =1`$ introduces topological defects (merons) into the order parameter vector $`\stackrel{}{m}`$. In terms of bosonic Chern-Simons theory, this statement is a consequence of the residual magnetic field left, away from $`\nu =1,`$ after the external and Hartree-Chern-Simons magnetic fields almost cancel one another. This residual field introduces vortices into the bosonic order parameters of the two layers. In the language of a quantum Hall ferromagnet, this observation is a consequence of the coupling of the symmetric density to the order parameter $`\stackrel{}{m}`$. The symmetric part of the density is constrained to satisfy $`n(\stackrel{}{r})n_0=\frac{1}{8\pi }ϵ_{ab}ϵ_{\mu \nu \kappa }m_\mu _am_\nu _bm_\kappa =\left(\frac{m_z}{8\pi }\widehat{z}\times \phi \right)\frac{m_z}{8\pi }\times \phi .`$ The deviation from $`n_{0\text{ }}`$ is then composed of a charge density carried by an electric dipole field $`\frac{m_z}{8\pi }\widehat{z}\times \phi ,`$ and by a charge density attached to topological defects in $`\stackrel{}{m}`$. The latter are merons of four types, carrying a charge of $`\pm \frac{e}{2}`$, and characterized by their vorticity (the sign of $`\times \phi `$ at the core) and the layer in which their charge resides ($`m_z`$ at the core). Merons interact coulombically due to their charge, and by a logarithmic interaction due to their vorticity. Below the Kosterlitz-Thouless temperature $`T_{\mathrm{KT}}\rho _s`$, merons are bound in pairs of opposite vorticity to avoid the logarithmically diverging energy penalty.
In realistic samples there are long range density fluctuations whose relative magnitude is estimated to be 4%. Thus, the typical distance between the disorder-induced meron pairs is $`12\mathrm{}`$. The separation between the two merons that constitute a pair is estimated to be $`6\mathrm{},`$ comparable to the spacing among different pairs. This estimate is obtained by balancing the Coulomb repulsion and logarithmic attraction. Thus, the $`\nu =1`$ bi-layer sample studied in Ref. is analogous to a superconducting junction with random magnetic flux that introduces many vortices in the two superconductors. Meron pairs may carry a charge $`\pm e`$ (distributed between the two layers) or be charge neutral. The charged pairs affect the longitudinal resistivity to the flow of symmetric current. In the sample of Spielman et al. this resistivity is large ($`1\mathrm{k}\mathrm{\Omega }`$), indicating that the charged vortex pairs are highly mobile. Furthermore, the dissipation is not frozen out at the lowest attainable temperatures indicating that these objects are disorder- rather than thermally-induced. Tunneling in this system is then strongly influenced by these merons, in a way discussed below. The meron pairs do not however destroy the antisymmetric superfluid mode unless they become unbound.
Appealing to the experimental observation that there is no DC Josephson effect (i.e., current linear in the tunneling amplitude) we may use Fermi’s Golden Rule to calculate the tunneling current perturbatively. For a sample of size $`L^2`$
$$I(V)=\frac{2\pi e\lambda ^2L^2}{\mathrm{}}[S(Q_B,eV)S(Q_B,eV)],$$
(3)
where $`S(q,\mathrm{}\omega ),`$ the spectral density for the fluctuations of the operator $`e^{i\phi \text{ }}`$at wavevector $`q`$ and frequency $`\omega ,`$ is proportional to the Fourier transform of $`e^{i\phi (r,t)\text{ }}e^{i\phi (0,0)\text{ }}`$ (where angular brackets denote thermal average). Our prediction regarding the dependence of the tunneling current on $`B_{||}`$ can now be easily understood. For weak disorder, the spectral density $`S(Q_B,eV)`$ is sharply peaked at
$$eV=\mathrm{}\omega _{Q_\mathrm{B}}$$
(4)
Thus, as the parallel field is varied, the peak in the tunneling conductance is shifted in a way that reflects the dispersion of the low energy excitation mode. This is precisely analogous to the Carlson-Goldman experiment measuring the collective oscillations of the pair field in a superconductor. An observation of this dispersing peak would also confirm an essential ingredient of the picture we use, namely the existence of a single branch of low energy excitations. The parallel field allows only tunneling between states that differ by a momentum $`Q_\mathrm{B}`$. Energy conservation requires the energy of these states to differ by $`eV`$. When there is just one low energy excitation branch, there is only one value of the voltage where both these conditions are fulfilled. This is not the case for a fermi liquid (for $`Q_B0`$).
To begin our analysis of the effect of merons on the tunneling conductance we make the simplifying ansatz that the order parameter phase can be separated into the sum a vortex part $`\phi _m`$ and an independent spinwave part $`\phi `$ with the former obeying $`G_m(r,t)e^{i\phi _m(\stackrel{}{r},t)}e^{i\phi _m(\stackrel{}{0},0)}=\mathrm{exp}(\frac{r^2}{2\xi ^2}\frac{t}{\tau _\phi })`$. The gaussian form for the spatial dependence is chosen for algebraic convenience. Applying this ansatz to Eqs.(1-3) yields
$`I(V,B_{||})`$ $`=`$ $`{\displaystyle \frac{4e\lambda ^2L^2}{\mathrm{}^2}}{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle d^2rG_m(r,t)e^{\frac{1}{2}D(r,t)}}`$ (6)
$`\mathrm{sin}{\displaystyle \frac{C(r,t)}{2}}\mathrm{cos}Q_\mathrm{B}x\mathrm{sin}{\displaystyle \frac{eVt}{\mathrm{}}}`$
with a ‘Debye-Waller factor’ $`\mathrm{exp}(D(r,t)/2)`$, where
$`D(r,t)[\phi (\stackrel{}{r},t)\phi (\stackrel{}{0},0)]^2`$ (7)
$`=`$ $`{\displaystyle \underset{q}{}}{\displaystyle \frac{\mathrm{}u}{L^2\rho _sq}}\left[1\mathrm{cos}(\stackrel{}{q}\stackrel{}{r})\mathrm{cos}(uqt)\right]\mathrm{coth}{\displaystyle \frac{\mathrm{}uq}{2T}}`$ (8)
(we set $`k_\mathrm{B}=1`$ throughout) and a commutator term
$`C(r,t)`$ $``$ $`i[\phi (\stackrel{}{r},t),\phi (\stackrel{}{0},0)]`$ (9)
$``$ $`{\displaystyle \frac{\mathrm{}}{2\pi \rho _s}}\theta (utr)\left[t^2\left({\displaystyle \frac{r}{u}}\right)^2\right]^{1/2}`$ (10)
which is independent of temperature, and limits the $`r`$integral in (6) to a “light-cone” of $`r<ut`$. All correlators are evaluated in the absence of tunneling. We rely on the global $`U(1)`$ symmetry and the freedom to renormalize $`\xi `$ and $`\tau _\phi `$ to partially justify the simplifying approximation of neglecting all disorder in the spinwave hamiltonian.
As long as $`2\pi \rho _s\mathrm{}/\tau _\phi T`$ we can expand Eq.(6) to first order in $`C`$ and approximate $`D`$ by its zero temperature value. In this limit the current becomes
$`I(V,B_{||})={\displaystyle \frac{e}{h}}{\displaystyle \frac{\xi ^2\lambda ^2L^2}{4\mathrm{\Gamma }}}e^{\frac{D}{2}}{\displaystyle d^2pe^{|\stackrel{}{p}\stackrel{}{Q}_B|^2\xi ^2/2}\frac{\mathrm{}}{\omega _p}}`$ (11)
$`\left\{{\displaystyle \frac{\delta _\phi }{(eV\mathrm{}\omega _p)^2+(\delta _\phi )^2}}{\displaystyle \frac{\delta _\phi }{(eV+\mathrm{}\omega _p)^2+(\delta _\phi )^2}}\right\}`$ (12)
where $`\delta _\phi \mathrm{}/\tau _\phi .`$ For large $`\tau _\phi ,\xi `$ Eq. (12) shows a peak in the current at the voltage corresponding to the Goldstone mode energy in accordance with Eq.(4). The effect of $`\tau _\phi ,\xi `$ is to smear this peak over a range of $`\mathrm{}/\xi `$ in momentum and $`\mathrm{}/\tau _\phi `$ in voltage. As long as $`Q_B\xi 1`$ and $`uQ_B\tau _\phi 1,`$ this smearing is insignificant.
The expression for the differential conductance simplifies considerably in the limit $`Q_B=0`$, $`\xi u\tau _\phi `$ and $`eV\frac{\mathrm{}u}{\xi }:`$
$$\frac{dI}{dV}=\frac{1}{8}\frac{e^2}{h}\frac{\xi ^2}{\mathrm{}^2}\frac{n_0L^2\mathrm{\Delta }_{SAS}^2}{\rho _s}e^{\frac{D}{2}}\frac{\delta _\phi }{(eV)^2+(\delta _\phi )^2}.$$
(13)
Interestingly, we see that when $`\tau _\phi =\mathrm{},`$ i.e., when the merons provide a random static background phase field, a Josephson-like singularity of $`\frac{dI}{dV}`$ is still present (as is the antisymmetric superfluid property). As shown below, the singularity is present also at finite temperature $`T\rho _s`$. Static topological defects break translational invariance and thus open more phase space for excitation of spin waves in the tunneling process. However, they do not expand the degrees of freedom involved beyond the single spin wave mode, and thus do not dephase the process enough to destroy the zero-voltage singularity.
The temperature dependence of (6) originates from the temperature dependence of $`D`$ and the temperature dependence of $`\rho _s`$ and $`\tau _\phi .`$ Here we calculate the temperature dependence of $`D.`$ At zero temperature it gives the space and time independent result $`D_0_{q\mathrm{}<\sqrt{2}}d^2q\frac{\mathrm{}}{\mathrm{\Gamma }uq}4.8`$. At finite temperature we approximate $`\mathrm{coth}x1+\frac{1}{x}e^x`$, define dimensionless length and time variables, $`\stackrel{~}{r}\frac{rT}{\mathrm{}u}`$ and $`\stackrel{~}{t}\frac{tT}{\mathrm{}}`$, and obtain (for large $`r,t,`$ and $`r<ut`$),
$`D(\stackrel{~}{r},\stackrel{~}{t})D_0+`$ (14)
$`{\displaystyle \frac{T}{2\pi \rho _s}}\mathrm{log}\left|\left(\stackrel{~}{t}+i/2\right)+\sqrt{\left(\stackrel{~}{t}+i/2\right)^2\stackrel{~}{r}^2}\right|^2`$ (15)
The temperature dependence of $`D`$ affects $`I(V)`$ then only at high temperature $`(TeV)`$, where we can approximate $`\stackrel{~}{t}+i/2\stackrel{~}{t}`$. For $`u\tau _\phi \xi `$ and $`B_{}=0`$, Eq.(6) reduces to
$`I\left(V\right)`$ $`{\displaystyle \frac{e\lambda ^2L^2}{\pi \rho _s\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle _{r<ut}}d^2r\mathrm{exp}({\displaystyle \frac{1}{2}}\left({\displaystyle \frac{r}{\xi }}\right)^2{\displaystyle \frac{t}{\tau _\phi }})`$ (17)
$`{\displaystyle \frac{\left|\frac{tT}{\mathrm{}}+\sqrt{\left(\frac{tT}{\mathrm{}}\right)^2\left(\frac{rT}{\mathrm{}u}\right)^2}\right|^{\frac{T}{2\pi \rho _s}}}{\sqrt{t^2(r/u)^2}}}\mathrm{sin}\left({\displaystyle \frac{eVt}{\mathrm{}}}\right)`$
Most of the contribution is then from long times, while $`r`$ is limited to be smaller than $`\xi `$. For a static meron background ($`\tau _\phi =\mathrm{}`$) we find for $`eVT`$
$$\frac{dI}{dV}\frac{\lambda ^2\xi ^2L^2}{\rho _s^2}\left(\frac{T}{V}\right)^{1\frac{T}{2\pi \rho _s}}$$
(18)
which is consistent with the more complete scaling form which can be derived in the classical limit from the expression of Nelson and Fisher for the dynamical structure factor of the XY model. In the presence of a finite $`\tau _\phi `$, the temperature dependence of $`D`$ affects the tunneling current significantly only in the window $`\frac{\mathrm{}}{\tau _\phi },eV<T<T_{\mathrm{KT}}`$. In the experiment of Spielman et al. the peak width is much larger than the temperature. Thus, the observed temperature dependence probably results from the temperature dependence of $`\rho _s`$ and $`\tau _\phi `$ rather than $`D`$.
Using Eq. (12), we can fit the width of the conductance peak in the experiment with a phenomenological value $`\delta _\phi 0.75`$K. This value gives $`u\tau _\phi 11\mathrm{}`$, which is remarkably close to our estimate of $`\xi `$ based on the meron pair spacing. (Sufficiently close that the Lorentzian approximation in Eq. (13) for the peak width will be somewhat inaccurate.)
Our naive estimate for the peak height in the experiment is too large by some two orders of magnitude, but is highly uncertain due to the exponential sensitivity to the ultraviolet cutoff and the acoustic approximation used in computing the Debye-Waller factor. In addition, the estimate $`\mathrm{\Delta }_{SAS}90\mu `$K is exponentially sensitive to the parameters in the modeling of the barrier potential (in particular the poorly understood effective mass appropriate for the high Al concentration in the barrier) and so might be off by a significant factor. It might also be possible that the superfluidity and the tunneling occur predominately in isolated regions close to filling factor $`\nu =1`$ containing few vortices. The parasitic series transport resistance $`1/\sigma _{xx}`$ in this Corbino-like geometry could significantly reduce the peak height.
We now consider inter-layer tunneling under the combined effect of a time independent $`dc`$ voltage and a time dependent $`ac`$ electric field $`E\mathrm{sin}\omega t`$, directed perpendicular to the two layers. As long as the system is not heated, this field can be incorporated into our calculation by writing $`T_\pm =𝑑\stackrel{}{r}\lambda e^{\pm i[\phi (\stackrel{}{r})+\frac{eEd}{\mathrm{}\omega }\mathrm{cos}\omega t]}`$. Repeating the calculation carried out above, we find that the tunneling differential conductance $`\frac{dI}{dV}(V)`$ exhibits peaks at $`eV=n\mathrm{}\omega `$, with $`n`$ an integer. This feature is common to all tunneling systems where the $`dc`$ differential conductance is strongly peaked around zero voltage (for example, a bi-layer system at zero magnetic field). We note, however, that the quantum Hall ferromagnet is relatively less prone to heating, due to the small longitudinal conductivity.
Finally, we discuss the current distribution in an idealized zero-disorder and vortex-free system. Eqs. (1) and (2) then do indeed lead to a Sine-Gordon equation for the phase, as in a long Josephson junction. However, due to the two-dimensionality of the problem, the critical current is not proportional to the area of the sample. Consider a setup where the current is fed into one layer from, say, $`x=\mathrm{}`$, and taken out from the other layer at $`x=\mathrm{}`$, and where tunneling is limited to the region $`\frac{L}{2}<x<\frac{L}{2}`$. Since the symmetric part of the current ($`I_{sym}`$) is conserved, the boundary conditions for the Sine-Gordon equation require $`\frac{\phi }{x}|_{x=\frac{L}{2}}=\frac{\phi }{x}|_{x=\frac{L}{2}}=I_{sym}.`$ For $`L\xi _J\sqrt{4\pi \mathrm{}^2\rho _s/\mathrm{\Delta }_{\mathrm{SAS}}}4\mu `$m the time-independent solution to the Sine-Gordon equation in the tunneling region is $`\phi (x)2\mathrm{arccos}\mathrm{tanh}\frac{\frac{L}{2}|x|}{\xi _J}`$, tunneling takes place only within a distance of order $`\xi _J`$ of the $`x=\pm \frac{L}{2}`$ lines, and the maximal current that can tunnel is $`L`$independent, and is given by $`(2e/\mathrm{})\rho _sW/\xi `$ (here $`W`$ is the width of the current contact). For the parameters we use this current is $`4`$nA/$`\mu `$m$`W80`$nA. The experimental measurement current was much smaller than this value. Thus, the absence of a Josephson effect cannot be attributed to a large measurement current. For the sample geometry described above and used in the experiment, the tunnel resistance is effectively in series with the Hall resistance. The observed tunnel resistance is however much larger, $`10^2h/e^2`$, again indicating that there is no Josephson effect.
To conclude, we have attributed the lack of a Josephson effect in tunneling measurements in a bi-layer quantum Hall $`\nu =1`$ state to density inhomogeneities that introduce dynamical topological defects into the order parameter. The observed peak width is quantitatively consistent with this picture. We showed that a measurement of the tunneling $`I(V)`$ dependence on $`B_{}`$ would map the dispersion relation of the low energy mode of the system, and that tunneling in the presence of an $`ac`$ electric field would result in resonances at voltages corresponding to the $`ac`$ frequency. Finally, we showed that even for a perfect sample where the Josephson effect takes place, the critical current would not scale with the size of the sample.
AS is supported by the US-Israel BSF, and the Israeli Academy of Science, Victor Ehrlich chair, and a DIP-BMBF grant. The Indiana group is supported by NSF DMR-9714055. The authors are grateful to J. P. Eisenstein and I. Ussishkin for numerous useful discussions.
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# NIKHEF 2000-010 June 2000 Deep-inelastic structure functions: Reconstruction from Mellin moments
## 1 INTRODUCTION
The precision measurement of structure functions and their scale evolution in deep inelastic scattering is of basic importance for testing perturbative QCD. By now, experimental data allows for an accurate determination of the parton distribution functions, which serve as input for numerous other hard scattering processes. Therefore, in particular with respect to the upcoming LHC experiments, it is highly desirable to further reduce the theoretical uncertainties on these very fundamental quantities. This task requires the calculation of higher order perturbative QCD corrections. In order to perform a consistent next-to-next-to-leading order analysis, which is expected to significantly reduce the theoretical uncertainties , one has to calculate the still unknown anomalous dimensions of deep-inelastic structure functions to three loops.
The method of calculating Mellin moments of deep-inelastic structure functions is directly connected with the early days of QCD . Since then, it has been applied several times to determine higher order perturbative corrections , some of them related to sum rules. The method rests on the operator product expansion (OPE), which allows the calculation of either all even or all odd integer moments $`N`$ of the anomalous dimensions and coefficient functions. This information by itself then uniquely determines the quantity under consideration for general $`N`$ in Mellin space or, equivalently its corresponding expression in $`x`$-space.
In the following, we will summarize the calculation of the anomalous dimensions and coefficient functions in terms of harmonic sums and the subsequent analytical reconstruction based on an inverse Mellin transformation to harmonic polylogarithms in $`x`$-space. This procedure as detailed in has recently been used to check the original calculation of the two-loop coefficient functions , which was performed with conventional techniques. Currently, work is in progress to obtain the anomalous dimensions and coefficient functions at three loops by this method.
The present approach of analytical reconstruction nicely complements the use of Bernstein polynomials for a direct fit of Mellin moments to experimental data. The latter method is particularly useful if only numerical results for a finite number of fixed moments are available, as it is presently the case .
## 2 METHOD
Deep-inelastic structure functions are defined by the commutator of two local electroweak currents $`j(y)`$ and $`j(z)`$, sandwiched between hadronic states and Fourier transformed into momentum space. In the Bjorken limit, for $`Q^2\mathrm{}`$ and $`x`$ fixed, the OPE allows to express this current product in an asymptotic expansion around the lightcone $`(yz)^20`$ into a series of local composite flavour non-singlet quark $`O^\alpha `$, singlet quark $`O^\mathrm{q}`$ and gluon operators $`O^\mathrm{g}`$ of leading twist and spin $`N`$. The same OPE also holds for the forward Compton amplitude $`T`$ of boson-hadron scattering. Standard perturbation theory applies to $`T`$, which can be written as a series in $`1/x^N`$ in the Euclidean region. Then, by means of the optical theorem one obtains a direct relation between the parameters of the OPE and the Mellin moments of the structure functions. For $`F_2`$ we can write
$`F_2^N(Q^2)`$ $`=`$ $`{\displaystyle \underset{0}{\overset{1}{}}}𝑑xx^{N2}F_2(x,Q^2)`$
$`=`$ $`{\displaystyle \underset{j=\alpha ,\mathrm{q},\mathrm{g}}{}}C_{2,j}^N({\displaystyle \frac{Q^2}{\mu ^2}},\alpha _s)A_{\mathrm{P},N}^j\left(\mu ^2\right),`$
and similar relations define $`F_3^N`$ and $`F_L^N`$. Here $`C_{2,j}^N`$ denote the coefficient functions and $`A_{\mathrm{P},N}^j`$ the spin averaged hadronic matrix elements of the operators $`O^\alpha `$, $`O^\mathrm{q}`$ and $`O^\mathrm{g}`$.
In the dispersive approach the derivation of eq.(2) uses symmetry properties of the Compton amplitude $`T`$ under exchange $`xx`$. As a consequence, eq.(2) is restricted to only either the even or the odd Mellin moments of $`F_2`$. By analytic continuation on the other hand, all moments in the complex $`N`$-plane are fixed if the infinite set of either even or odd moments is known. To be precise, for unpolarized electron-proton scattering eq.(2) determines the even moments of $`F_2^{e\mathrm{P}}`$, while for neutrino-proton scattering the odd moments of $`F_2^{\nu \mathrm{P}\overline{\nu }\mathrm{P}}`$ and $`F_3^{\nu \mathrm{P}+\overline{\nu }\mathrm{P}}`$ and the even moments of $`F_2^{\nu \mathrm{P}+\overline{\nu }\mathrm{P}}`$ and $`F_3^{\nu \mathrm{P}\overline{\nu }\mathrm{P}}`$ are fixed.
The coefficient functions and the renormalized operator matrix elements in eq.(2) both satisfy renormalization group equations. Due to current conservation they are governed by the same anomalous dimensions,
$`{\displaystyle \underset{k=\alpha ,\mathrm{q},\mathrm{g}}{}}[\{\mu ^2{\displaystyle \frac{}{\mu ^2}}+\beta (\alpha _s(\mu ^2)){\displaystyle \frac{}{\alpha _s(\mu ^2)}}\}\delta _{jk}`$
$`+\gamma _{jk}(\alpha _s(\mu ^2))\left]A_{\mathrm{P},N}^j\right(\mu ^2)=0,`$
$`{\displaystyle \underset{k=\alpha ,\mathrm{q},\mathrm{g}}{}}[\{\mu ^2{\displaystyle \frac{}{\mu ^2}}+\beta (\alpha _s(\mu ^2)){\displaystyle \frac{}{\alpha _s(\mu ^2)}}\}\delta _{jk}`$
$`\gamma _{jk}(\alpha _s(\mu ^2))\left]C^N_{2,k}\right({\displaystyle \frac{Q^2}{\mu ^2}},\alpha _s(\mu ^2))=0,`$
where $`j=\alpha ,\mathrm{q},\mathrm{g}`$, while $`\beta `$ and $`\gamma _{jk}`$ represent the QCD $`\beta `$-function and the anomalous dimensions. Both are calculable order by order in $`\alpha _s`$ in perturbative QCD as well as the coefficient functions $`C_{2,j}^N`$.
The anomalous dimensions $`\gamma _{jk}`$ in eqs.(2), (2) determine the scale evolution of deep-inelastic structure functions and, dependent on the particular scattering process, one considers operator matrix elements corresponding to the quark flavour non-singlet and singlet distributions $`\mathrm{q}_{\mathrm{ns}}^\pm `$, $`\mathrm{q}_{\mathrm{ns}}^\mathrm{V}`$ and $`\mathrm{q}_\mathrm{s}`$. Their scale evolution is governed by linear combinations,
$`\mathrm{q}_{\mathrm{ns}}^\pm `$ $``$ $`\gamma _{\mathrm{qq}}^\mathrm{V}\pm \gamma _{\mathrm{q}\overline{\mathrm{q}}}^\mathrm{V}`$ (4)
$`\mathrm{q}_{\mathrm{ns}}^\mathrm{V}`$ $``$ $`\gamma _{\mathrm{qq}}^\mathrm{V}\gamma _{\mathrm{q}\overline{\mathrm{q}}}^\mathrm{V}+n_f\left(\gamma _{\mathrm{qq}}^\mathrm{S}\gamma _{\mathrm{q}\overline{\mathrm{q}}}^\mathrm{S}\right),`$ (5)
$`\mathrm{q}_\mathrm{s}`$ $``$ $`\gamma _{\mathrm{qq}}^\mathrm{V}+\gamma _{\mathrm{q}\overline{\mathrm{q}}}^\mathrm{V}+n_f\left(\gamma _{\mathrm{qq}}^\mathrm{S}+\gamma _{\mathrm{q}\overline{\mathrm{q}}}^\mathrm{S}\right),`$ (6)
and the flavour singlet quark distribution $`\mathrm{q}_\mathrm{s}`$ mixes with the gluon distribution.
The actual calculation of the anomalous dimensions $`\gamma _{jk}`$ and the coefficient functions $`C_{2,j}^N`$ is performed in dimensionally regularized perturbation theory with the external hadron states replaced by quark and gluon states of momentum $`p`$. By means of the OPE, the forward Compton amplitude for the scattering off a parton can be expressed in terms of the coefficient functions of eq.(2) and renormalized partonic operator matrix elements $`A_{\mathrm{p},N}^k`$. To be specific, we write for the flavour singlet case,
$`T_{2,\mathrm{p}}(x,Q^2,\alpha _s,ϵ)=`$
$`{\displaystyle \underset{N}{}}{\displaystyle \underset{j,k=\mathrm{q},\mathrm{g}}{}}\left({\displaystyle \frac{1}{2x}}\right)^NC_{2,j}^N({\displaystyle \frac{Q^2}{\mu ^2}},\alpha _s,ϵ)`$
$`\times Z^{jk}(\alpha _s,{\displaystyle \frac{1}{ϵ}})A_{\mathrm{p},N}^k(\alpha _s,\mu ^2,ϵ)+O(p^2),`$
where $`\mathrm{p}=\mathrm{q},\mathrm{g}`$ and $`O(p^2)`$ denotes higher twist contributions. The $`Z^{jk}`$ represent the matrix of renormalization factors which contain all poles in $`1/ϵ`$. They relate bare to renormalized operators according to $`O^{\mathrm{bare}}=ZO^{\mathrm{ren}}`$ and provide the anomalous dimensions via,
$`\gamma `$ $`=`$ $`\left({\displaystyle \frac{d}{d\mathrm{ln}\mu ^2}}Z\right)Z^1.`$ (8)
The gauge invariant operators $`O^\mathrm{q}`$ and $`O^\mathrm{g}`$ in eq.(2) mix under renormalization with unphysical operators , which are BRST variations of some operators or else disappear by use of the equations of motion. However, if one calculates quantities related to physical $`S`$-matrix elements with physical polarization and on-shell momenta such as in eq.(2), these unphysical operators vanish and have therefore been omitted in eq.(2).
The coefficient function $`C_{2,j}^N`$ and the renormalization factors $`Z^{jk}`$ are calculated from eq.(2) using the method of projection . This method relies on dimensional regularization and amounts to the application of the projection operator,
$`𝒫_N\left[{\displaystyle \frac{q^{\{\mu _1}\mathrm{}q^{\mu _N\}}}{N!}}{\displaystyle \frac{^N}{p^{\mu _1}\mathrm{}p^{\mu _N}}}\right]|_{p=0},`$ (9)
where $`q^{\{\mu _1}\mathrm{}q^{\mu _N\}}`$ is symmetrical and traceless part of the tensor $`q^{\mu _1}\mathrm{}q^{\mu _N}`$.
It is obvious, that on the right hand side of eq.(2) the $`N\text{-th}`$ order differentiation in $`𝒫_N`$ projects on the precisely the $`N\text{-th}`$ moment, which is the coefficient of $`1/(2x)^N`$ being $`(pq/Q^2)^N`$. All higher powers of $`pq/Q^2`$ vanish after nullifying the momentum $`p`$. Furthermore, $`𝒫_N`$ does not act on the renormalization constants $`Z^{jk}`$ and the coefficient functions being only functions of $`N`$, $`\alpha _s`$ and $`ϵ`$. However, the nullification of the momentum $`p`$ in $`𝒫_N`$ effects the partonic matrix elements $`A_{\mathrm{p},N}^k`$. There, it eliminates all diagrams containing loops, as they become massless tadpole diagrams, which are put to zero in dimensional regularization. Only the tree level operator matrix elements $`A_{\mathrm{p},N}^{p,\mathrm{tree}}`$ survive.
Additionally in the flavour singlet case this decouples the operator mixing. Thus, eq.(2) provides two independent identities to separately determine $`Z^{\mathrm{qq}}`$, $`Z^{\mathrm{gq}}`$ and $`Z^{\mathrm{qg}}`$, $`Z^{\mathrm{gg}}`$. We should be aware however, that an $`l`$-loop calculation eq.(2) alone does not suffice to give the full information about the renormalization constants $`Z^{\mathrm{gq}}`$ and $`Z^{\mathrm{gg}}`$. They are determined only up to $`(l1)`$-loops. This limitation is due to the gluonic coefficient function $`C_{2,\mathrm{g}}^N`$ being zero at tree level, as photons couple directly to quarks only. To resolve this situation and to extract the anomalous dimension $`\gamma _{\mathrm{gq}}`$ and $`\gamma _{\mathrm{gg}}`$ to $`l`$-loops we also calculate Green’s functions in which the photon is replaced by an external scalar particle $`\varphi `$ that couples directly only to gluons . These Green’s functions are obtained from a gauge invariant interaction term $`\varphi F_{\mu \nu }^aF_a^{\mu \nu }`$ in the QCD Lagrangian, $`F_{\mu \nu }^a`$ being the QCD field strength. They can be expressed in partonic invariants $`T_{\varphi ,\mathrm{p}}`$ and satisfy an OPE similar to eqs.(2),(2) with the same singlet operators $`O^\mathrm{q}`$ and $`O^\mathrm{g}`$ but with different coefficient functions $`C_{\varphi ,\mathrm{p}}^N`$, where the gluonic coefficient function $`C_{\varphi ,\mathrm{g}}^N`$ starts now already tree level. Thus the invariants $`T_{\varphi ,\mathrm{p}}`$ provide us with the necessary renormalization constants $`Z^{\mathrm{gq}}`$ and $`Z^{\mathrm{gg}}`$ of the singlet operators to $`l`$-loops and complete the determination of all anomalous dimensions in the flavour singlet case to the desired order.
Use of the method of projection to calculate all Feynman diagrams contributing to the forward partonic Compton amplitude, that is to say the application of the projection operator $`𝒫_N`$ to the integrands on the left hand side of eq.(2) proceeds as follows. All ultraviolet divergences as well as those of infrared origin as the momentum $`p0`$ are dimensionally regularized. At the same time, the nullification of $`p`$ drastically simplifies the topological structure of the corresponding graphs, which is reduced to the type of self-energy diagrams. Although the latter appear with symbolic powers of scalar products in the numerator and denominator, they are by far easier to calculate.
In practice this is performed by means of integration-by-parts identities and scaling identities for the $`N\text{-th}`$ Mellin moment of the diagram to be calculated. Combined these identities give rise to difference equations, which relate the $`N\text{-th}`$ Mellin moment of the diagram under consideration to lower Mellin moments and diagrams of a simpler topology. Thereby they introduce a natural hierarchy in the set of diagrams. All difference equations can be summed recursively and the answer is expressible in terms of a basis of harmonic sums of a given weight. The whole approach exploits powerful summation algorithms for a large class of single parameter sums which are reducible to the basis of harmonic sums.
As a demonstration of the method discussed thus far and as a check on , we have performed a recalculation of the two-loop coefficient functions. The contributing Feynman diagrams have been generated with QGRAF and all recursion relations for the evaluation of the individual topologies have been programmed in FORM . The nested sums which one encounters in this way are solved with the SUMMER algorithm in terms of the basis of harmonic sums. In addition, there is the possibility to perform checks at all stages of the calculation by means of the standard MINCER routine . This last feature is very important from a practical point of view as the debugging is greatly simplified.
To choose a specific example, we consider the case of initial state quarks in $`D=42ϵ`$ dimensions and calculate with the help of eqs.(2),(9) the projected invariant $`𝒫_NT_{2,\mathrm{q}}`$. After performing the coupling constant renormalization one is only left with collinear singularities associated with the initial state quarks. Expanding $`𝒫_NT_{2,\mathrm{q}}`$ in $`\alpha _s`$ and $`ϵ`$ gives at leading order
$`T_{2,\mathrm{q}}^{(0)}=1.`$ (10)
At first order in $`\alpha _s`$, we have to expand up to order $`ϵ`$ and find,
$`T_{2,\mathrm{q}}^{(1)}=`$
$`{\displaystyle \frac{\alpha _s}{4\pi }}S_ϵ\left({\displaystyle \frac{\mu ^2}{Q^2}}\right)^ϵ\left[{\displaystyle \frac{1}{ϵ}}\gamma _{\mathrm{qq}}^{(0)}+c_{2,\mathrm{q}}^{(1)}+ϵa_{2,\mathrm{q}}^{(1)}\right],`$
where the factor $`S_ϵ=\mathrm{exp}(ϵ\{\mathrm{ln}(4\pi )\gamma _\mathrm{E}\})`$ and $`\gamma _{\mathrm{qq}}^{(n)}`$, $`c_{2,\mathrm{q}}^{(n)}`$ are the coefficients of $`(\alpha _s/4\pi )^n`$ in the expansion of the anomalous dimension $`\gamma _{\mathrm{qq}}`$ and the quark coefficient function $`C_{2,\mathrm{q}}^N`$.
At second order in $`\alpha _s`$, we need to split up the contributions into flavour non-singlet and singlet parts,
$`T_{2,\mathrm{q}}^{(2),\mathrm{s}}=T_{2,\mathrm{q}}^{(2),\mathrm{ns},+}+T_{2,\mathrm{q}}^{(2),\mathrm{ps}}.`$ (12)
Allowing for electroweak interactions, one can consider in the non-singlet case the structure functions of different physical processes, $`F_2^{\nu \mathrm{P}\pm \overline{\nu }\mathrm{P}}`$, which implies a distinction of even and odd moments, $`T_{2,\mathrm{q}}^{(2),\mathrm{ns},+}`$ and $`T_{2,\mathrm{q}}^{(2),\mathrm{ns},}`$ to be precise. We have
$`T_{2,\mathrm{q}}^{(2),\mathrm{ns},\pm }=\left({\displaystyle \frac{\alpha _s}{4\pi }}\right)^2S_ϵ^2\left({\displaystyle \frac{\mu ^2}{Q^2}}\right)^{2ϵ}`$
$`\times [{\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{1}{2}}\{\left(\gamma _{\mathrm{qq}}^{(0)}\right)^2\beta _0\gamma _{\mathrm{qq}}^{(0)}\}+{\displaystyle \frac{1}{ϵ}}\{{\displaystyle \frac{1}{2}}\gamma _{\mathrm{qq}}^{(1),\pm ,\mathrm{V}}`$
$`+\gamma _{\mathrm{qq}}^{(0)}c_{2,\mathrm{q}}^{(1)}\}+c_{2,\mathrm{q}}^{(2),\mathrm{ns},\pm }+\gamma _{\mathrm{qq}}^{(0)}a_{2,\mathrm{q}}^{(1)}],`$
where $`\gamma _{\mathrm{qq}}^{\pm ,\mathrm{V}}`$ denotes the combinations $`\gamma _{\mathrm{qq}}^\mathrm{V}\pm \gamma _{\mathrm{q}\overline{\mathrm{q}}}^\mathrm{V}`$ as given in eqs.(4)–(6).
The pure-singlet contributions on the other hand give at second order in $`\alpha _s`$,
$`T_{2,\mathrm{q}}^{(2),\mathrm{ps}}=n_f\left({\displaystyle \frac{\alpha _s}{4\pi }}\right)^2S_ϵ^2\left({\displaystyle \frac{\mu ^2}{Q^2}}\right)^{2ϵ}`$
$`\times [{\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{1}{2}}\gamma _{\mathrm{qg}}^{(0)}\gamma _{\mathrm{gq}}^{(0)}+{\displaystyle \frac{1}{ϵ}}\{{\displaystyle \frac{1}{2}}\gamma _{\mathrm{qq}}^{(1),+,\mathrm{S}}+\gamma _{\mathrm{gq}}^{(0)}c_{2,\mathrm{g}}^{(1)}\}`$
$`+c_{2,\mathrm{q}}^{(2),\mathrm{ps}}+\gamma _{\mathrm{gq}}^{(0)}a_{2,\mathrm{g}}^{(1)}],`$
and $`\gamma _{\mathrm{qq}}^{\pm ,\mathrm{S}}`$ represents the combinations $`\gamma _{\mathrm{qq}}^\mathrm{S}\pm \gamma _{\mathrm{q}\overline{\mathrm{q}}}^\mathrm{S}`$ of eqs.(5),(6).
With eqs.(10),(2) and (2),(2) at hand, we have at this point calculated all anomalous dimension and coefficient functions up to two loops in Mellin space as analytical functions of the moment $`N`$. Depending on the physical process under consideration, the quantities of interest are defined for either even or odd $`N`$. In a following step, we perform an inverse Mellin transformation to obtain the corresponding expressions in $`x`$-space. As discussed in , we can do this step analytically as there exists a one-to-one mapping of harmonic sums of weight $`w+1`$ to the set of harmonic polylogarithms of weight $`w`$ multiplied by $`1/(1\pm x)`$. The latter one is the natural class of functions for calculations of deep-inelastic structure functions. Since the resulting expressions in $`x`$-space are unique, we may subsequently execute another Mellin transformation to go back to $`N`$-space, which then serves as an analytical continuation to all non-negative integer values in $`N`$.
In particular, we are then ready to determine individually all anomalous dimensions which control the scale evolution of the structure functions in the flavour non-singlet case as entering in eqs.(4)–(6). This is easily done now, as we obtain from the even moments of $`F_2^{e\mathrm{P}}`$ the sums $`\gamma _{\mathrm{qq}}^{+,\mathrm{V}}`$ and $`\gamma _{\mathrm{qq}}^{+,\mathrm{S}}`$. The odd moments of $`F_2^{\nu \mathrm{P}\overline{\nu }\mathrm{P}}`$ on the other hand determine the differences $`\gamma _{\mathrm{qq}}^{,\mathrm{V}}`$ and $`\gamma _{\mathrm{qq}}^{,\mathrm{S}}`$, the latter being zero up to two loops. As emphasized above, the analytical continuation provides us with expressions valid for all non-negative integer values of $`N`$, such that at this point we can simply take the sums and the differences of $`\gamma _{\mathrm{qq}}^{+,\mathrm{V},\mathrm{S}}`$ and $`\gamma _{\mathrm{qq}}^{,\mathrm{V},\mathrm{S}}`$ to obtain all quantities $`\gamma _{\mathrm{qq}}^\mathrm{V},\gamma _{\mathrm{q}\overline{\mathrm{q}}}^\mathrm{V},\gamma _{\mathrm{qq}}^\mathrm{S}`$ and $`\gamma _{\mathrm{q}\overline{\mathrm{q}}}^\mathrm{S}`$ of eqs.(4)–(6).
Similar reconstructions in Mellin space and $`x`$-space may be carried out for the two loop coefficient functions in eqs.(2),(2). To summarize, we find complete agreement with the published results for the anomalous dimensions and for the coefficient functions .
## 3 CONCLUSION
The calculation of higher order pertubative corrections to deep-inelastic structure functions is vitally important to improve quantitative predictions for many hard scattering processes in QCD.
In the past, this has been done either in $`x`$-space or in Mellin space. Using the method of projection, which enables us to obtain anomalous dimensions and coefficient functions at the same time, we perform the calculation in Mellin space. Due to new insight into the mathematical properties of harmonic sums and their interplay with harmonic polylogarithms, we are first of all able to solve all nested sums in Mellin space. Secondly we can reconstruct the complete analytical expressions of the results in $`x`$-space by means of an inverse Mellin transformation. Thereby we exploit all advantages of working in Mellin space, such as the use of difference equations in the Mellin moment $`N`$ to solve the resulting loop integrals and the possibility for independent checking with the MINCER routine.
Finally, we would like to remark, that the method presented allows for a direct application to the calculation of the three loop anomalous dimensions, as well as a generalization to polarized deep inelastic scattering .
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# Chemical abundances and ionizing clusters of HII regions in the LINER galaxy NGC 4258
## 1 Introduction
Theoretically, the evolution of a young stellar population depends on metallicity at least through two very important effects: the increasing opacity of the stellar material and the dependence of mass loss on metal content in high mass stars. As a consequence of the first, the effective temperature of ionizing stars should be lower in regions of higher metallicity (see for example McGaugh 1991) but, on the other hand, as a consequence of the second, that is if the strength of stellar winds increases with metallicity, the loss of the outer envelopes of the most massive stars can increase their surface temperature to very high values. These highly evolved massive O stars are identified with the Wolf Rayet population.
Most theoretical evolutionary models for ionizing star clusters predict the appearance of WR stars – although at slightly different cluster ages depending on the specific details of the models – and in all of them the fraction of WR stars increases with metallicity since the limiting mass for a star to enter the WR phase decreases with increasing metallicity. However, the effect of these stars on their surrounding ionized gas depends on the characteristics of the assumed wind opacity and therefore while some models predict the existence of high metallicity HII regions of high excitation (e.g. García-Vargas & Díaz 1994; Stasińska & Leitherer 1996), others do not (e.g. Cerviño & Mas-Hesse 1994). The finding of this extreme HII region population would constitute a firm evidence of the existence of WR stars with high effective temperatures, thus favoring the first kind of models but, even if these regions are not definitively identified, the study of high metallicity HII regions is of crucial importance to investigate the processes of star formation and evolution in high metallicity environments. This investigation gains even more relevance if we take into account the recent observations that point out to a firm connection between star formation and activity in active galactic nuclei (Heckman et al. 1997, González-Delgado et al. 1998) already predicted on theoretical grounds (Terlevich & Melnick 1985; Terlevich et al. 1992), since these nuclei seem to reach rather high metallicities (Phillips et al. 1984; González-Delgado & Pérez 1996) as deduced from the study of their circumnuclear star forming regions.
Unfortunately, the term “high metallicity HII region” is rather ill defined. Most authors apply the term to HII regions of solar or oversolar metal content, but the difficulty to derive directly the metallicity of these cool HII regions is well known and the empirical calibrations commonly employed provide estimations with an uncertainty larger than a factor of 2 (see e.g. Díaz et al. 1991). On the other hand, the characteristics of WR features are highly dependent on metallicity. Both the luminosity of the WR “bump”, that is the combination of the NIII and the HeII lines at around $`\lambda `$ 4660 Å, and its equivalent width are predicted to be largest at the highest metallicity and almost an order of magnitude difference is found between the computed WR “bump” intensity relative to H$`\beta `$ in clusters of half solar and solar metallicities (Schaerer & Vacca 1998). Therefore, the detection of WR features combined with a detailed analysis of the emission line spectra of suspected high metallicity HII regions, can constitute excellent tools for the simultaneous determination of both the metallicity and the age rather accurately in these regions.
In the frame of a long term programme to study high metallicity HII regions we are presently carrying out a search for the high metallicity and high excitation HII region population by performing spectrophotometric observations at a resolution high enough to detect and measure WR features. For that, we have selected from the literature the giant HII regions with solar or oversolar abundance, as deduced from empirical calibrations based on the optical forbidden lines, showing the highest excitation. This corresponds to a ratio of \[OIII\]/H$`\beta `$ $``$ 1.0, which is actually quite moderate but still high compared to that expected in high metallicity HII regions. Here we present the first of these investigations involving several HII regions in the galaxy NGC 4258.
NGC 4258 (M106) is classified as SAB(s)bc and Sb(s)II types by De Vaucouleurs \[De Vaucouleurs 1976\] and Sandage & Tammann \[Sandage & Tammann 1981\] respectively. Heckman \[Heckman 1980\] and Stauffer \[Stauffer 1982\] classified its nucleus as belonging to a transition type between LINER and H II-region like emission-line galaxy. As it is addressed by Courtés et al. \[Courtès et al. 1993\](hereinafter C93), the H II region distribution is very peculiar in the sense that three general structures are found: the normal spiral arms, with dense and very bright HII regions following pure spiral shapes as a consequence of the barred spiral type of this galaxy; the faint outer arms (H II region poor), and the anomalous spiral arms, discovered by Courtés & Cruvellier \[Courtès & Cruvellier 1961\], without star formation and with no evidence of very blue stars \[Deharveng & Pellet 1970\]. It is assumed that the anomalous arms were formed by powerful jets ejected by the active nucleus of the galaxy. Stüwe, Schulz & Hühnermann \[Stüwe, Schulz & Hühnermann 1992\] showed that the active nucleus must be obscured along the line of sight as is also indicated by its invisibility in IUE spectra. Cecil, Morse & Veilleux \[Cecil, Morse & Veilleux 1995\] found, from the emission-line flux ratios, LINER-like gaseous excitation along the jets. They also point out the possibility that the circumnuclear region has undergone recent formation of massive stars induced by the jets that are flowing into the disk. However, no evidence for this star formation has been found, maybe due to the obscuration by molecular gas along the line of sight. Finally, a subparsec diameter gaseous disk shows keplerian motion around the centre of the galaxy, which indicates the presence of a mass of 3.6 $`\times `$ 10<sup>7</sup> M within a radius of less than 0.13 pc \[Miyoshi et al. 1995\]. This is probably the best existing evidence for the presence of a massive black hole at the centre of a galaxy.
The problem of the distance determination for this galaxy still remains unsolved. The quoted distance values range from 3.4 Mpc \[Roy, Arsenault & Joncas 1986\] to 10.4 Mpc \[Sandage & Tammann 1981\]. The adopted value in the present work is 5.5 Mpc \[Martin et al. 1989\].
We have observed eight HII regions along two different spiral arms. Four of them have been identified as: 74C, 69C, 5N(A) and 5N(B) in C93. Regions 5N(A) and 5N(B) are located in the northern outer arm at a galactocentric distance of 20.3 Kpc. The other six regions are located in the SE inner arm at 8.05 Kpc from the galaxy centre. Region 74C is a discrete radiosource at 4.9 GHz \[Hummel et al. 1989\].
Regions 74C and 5N were previously observed by Oey & Kennicut \[Oey & Kennicutt 1993\] who derived values of 12+log(O/H) of 8.87 and 8.74 for each of them respectively (the solar value is 8.92). With a ratio of \[OIII\] to H$`\beta `$ of 0.98 and 1.44 both regions fit our selection criteria.
## 2 Observations and Reductions
Our spectrophotometric observations were obtained with the 4.2m William Herschel Telescope at the Roque de los Muchachos Observatory, in 1994 March 13, using the ISIS double spectrograph, with the TEK1 and a EEV CCD detectors in the blue and red arm respectively. Two gratings were used, R300B in the blue and R316R in the red arm, covering from \[O II\] $`\lambda `$3727 to \[S III\] $`\lambda `$9532 in three different spectral ranges of $``$1700 Å each. The dispersion of 1.4 Å pixel<sup>-1</sup> with a slit width of 1$`\stackrel{}{.}`$5 gives a spectral resolution of $``$4 Å. The nominal spatial sampling is 0$`\stackrel{}{.}`$7 pixel<sup>-1</sup>. The seeing was $``$1$`\stackrel{}{.}`$5. The slit was centered at the positions of regions 74C and 5N, as given by Oey & Kennicutt (1993). A third position was observed in order to include other HII regions outlining the SE inner arm. In all cases the position angle was chosen as to maximize the number of HII rgions in the slit. A journal of observations is given in Table 1.
The data were reduced using the IRAF (Image Reduction and Analysis Facility) package following standard methods. The two-dimensional wavelength calibration was accurate to 1 Å in all cases. The two-dimensional frames were flux calibrated using two spectroscopic standard stars observed during the same night with a 5$`\stackrel{}{.}`$ width slit. The agreement between the individual calibration curves was better than 8% in all cases. The spectra were previously corrected for atmospheric extinction using a mean extinction curve applicable to La Palma observing site. The removal of the atmospheric water-vapour absorption bands in the near infrared \[Díaz, Pagel & Wilson 1985\], was achieved dividing by the relatively featureless continuum of a subdwarf star observed in the same night as the galaxy.
The most critical step in the reduction process was the background subtraction in the near infrared frames. There are numerous OH emission lines of the night-sky spectrum that contaminate the near infrared spectra. This contamination is particularly important in the range from 9376 to 9600 Å \[Osterbrock, Fulbright & Bida 1997\]. Hence, the \[SIII\] $`\lambda `$9532 Å and Pa 8 lines intensities could be affected even after the background subtraction is made. Fortunately, the \[SIII\] $`\lambda `$9069 Å and Pa 9 lines are unaffected by these night-sky lines. The observed $`\lambda `$9532/$`\lambda `$9069 ratio is above the theoretical value of 2.44 in four of the observed regions. In these cases, the \[SIII\] $`\lambda `$9532 line was scaled by means of the more accurate \[SIII\] $`\lambda `$9069 line intensity measurement.
## 3 Results
Figure 1 shows the spatial distribution of the H$`\alpha `$ flux along the slit for the three different positions observed. Regions 5N A and B are clearly identified on the H$`\alpha `$ profile corresponding to PA=255°. Also regions 74C and 69C, showing a certain degree of substructure, are resolved at position angle PA=133°. Four different regions are resolved at PA=150°that we have named GA1, GA2, GA3 and GA4. A fifth region is clearly seen in the H$`\alpha `$ profile that is not detected in the near IR frames.
The spectra corresponding to each of the identified regions is shown in Figs. 2 to 4. WR features around $`\lambda `$ 4680 Å are seen in the spectrum of region 74C and to a lesser extent in that of region 5N.
### 3.1 Line intensities
Emission line fluxes were measured using the IRAF SPLOT software package, by integrating the line intensity over a local fitted continuum. The errors in the line fluxes have been calculated from the expression $`\sigma _l`$ = $`\sigma _c`$N<sup>1/2</sup>\[1 + EW/(N$`\mathrm{\Delta }`$)\]<sup>1/2</sup>, where $`\sigma _l`$ is the error in the line flux, $`\sigma _c`$ represents the standard deviation in a box near the measured emission line and stands for the error in the continuum placement, N is the number of pixels used in the measurement of the line flux, EW is the line equivalent width, and $`\mathrm{\Delta }`$ is the wavelength dispersion in angstroms per pixel.
The observed line intensities relative to the H$`\beta `$ line were corrected for interstellar reddening according to an average extinction curve \[Osterbrock 1989\] and assuming the Balmer line theoretical values for case B recombination \[Brocklehurst 1971\]. The presence of an underlying stellar population is evident in the blue spectra of regions GA2, GA3 and GA4. The H$`\gamma `$ and H$`\delta `$ Balmer lines are clearly affected by this stellar absorption. An iterative process was applied in order to fit observed and theoretical Balmer line intensities and to obtain the reddening constant c(H$`\beta `$). In all cases, H$`\alpha `$, H$`\gamma `$ and H$`\delta `$ were fitted, except in regions GA3 and GA4 for which H$`\delta `$ seems to be underestimated. Reddening corrected Paschen lines, when measured, are consistent with their theoretical values.
The reddening corrected line intensities together with their corresponding errors are given in Table 2 for regions 74C and 69C and Table 3 for the rest of the regions. Also given in the tables are the extinction corrected H$`\alpha `$ flux and the H$`\beta `$ equivalent width. For all the regions, H$`\alpha `$ luminosities are lower than 10<sup>38</sup> erg s<sup>-1</sup>, with the exception of region 74C, in which a value of 10<sup>39.3</sup> is found. According to Kennicutt \[Kennicutt 1983\], this region can be classified as a supergiant HII region.
### 3.2 Physical conditions of the gas
Electron densities for each observed region were derived from the ratio of the \[SII\] $`\lambda \lambda `$ 6717, 6731 lines following standard methods (see for example Osterbrock 1989). For regions GA2 and 74C densities of 160 cm<sup>-3</sup> and 60 cm<sup>-3</sup> respectively are obtained. For the rest of the regions the ratio of the \[SII\] lines implies n<sub>e</sub> $``$ 40 cm<sup>-3</sup> .
For two of the observed regions, 74C and 69C, it has been possible to measure the weak auroral \[SIII\] $`\lambda `$ 6312 Å line, which, together with the observed intensities of the nebular lines at $`\lambda \lambda `$ 9069, 9532 Å, can be used to obtain a value for the electron temperature. Using the expressions given by Osterbrock (1989) and the effective collision strengths recently calculated by Tayal (1997) we find for t(S<sup>2+</sup>) values of 7500 $`\pm `$ 300 K and 7600 $`\pm `$ 500 K for regions 74C and 69C respectively.
Photoionization models indicate that, for abundances close to solar, the region where sulphur is twice ionized overlaps both the regions of once and twice ionized oxygen and therefore t(S<sup>2+</sup>) is expected to be intermediate between t(O<sup>+</sup>) and t(O<sup>2+</sup>). Garnett (1992) from single star photoionization models finds a linear relation between the temperatures of O<sup>2+</sup> and S<sup>2+</sup>
$$t(O^{2+})=1.20t(S^{2+})0.20$$
(1)
We have used that relation to obtain t($`O^{2+}`$). Finally, Stasińska’s (1980) relation between t(O<sup>2+</sup>) and t(O<sup>+</sup>) has been used to calculate the temperature of the t(O<sup>+</sup>) zone. These different temperatures for regions 74C and 69C are listed in Table 4.
Another independent way of temperature determination is through the measure of the Paschen discontinuity at $`\lambda `$ 8200 Å. The ratio of the Paschen discontinuity to the H$`\beta `$ intensity is, after correcting for the derived extinction value of c(H$`\beta `$) = 0.48,
$$\frac{\mathrm{\Delta }Pa}{I(H\beta )}=(6.9\pm 0.8)\times 10^{15}Hz^1$$
(2)
According to the models computed by González-Delgado et al. \[González-Delgado et al. 1994\], for the helium abundance derived for the region 74C (see next section), this $`\mathrm{\Delta }`$ Pa/I(H$`\beta `$) implies an electron temperature of 8000 $`\pm `$ 700 K, which agrees with that derived from the \[SIII\] lines within the errors. Therefore, no temperature fluctuations are apparent in the region which might have been expected from the presence of WR stars (Pérez 1997).
For the rest of the regions it was not possible to obtain a direct measure of the electron temperature. An average temperature has been adopted from an empirical calibration by means of the sulphur abundance parameter S<sub>23</sub> (S<sub>23</sub> = (\[SII\]6717, 6731 + \[SIII\]9069, 9532)/H$`\beta `$) (Díaz & Pérez-Montero 1999). The values for these temperatures, together with similar ones derived for regions 74C and 69C, are also given in Table 4.
### 3.3 Chemical abundances
Ionic abundances of oxygen, nitrogen and sulphur have been derived following standard methods \[Pagel et al. 1992\] and using the temperatures found above. These abundances are listed in Table 4. We have assumed that most of the oxygen and sulphur are in the first and second ionization stages and therefore O/H = O<sup>+</sup>/H<sup>+</sup> \+ O<sup>++</sup>/H<sup>+</sup> and S/H = S<sup>+</sup>/H<sup>+</sup> \+ S<sup>++</sup>/H<sup>+</sup>. This assumption seems to be justified given the relatively low estimates of the electron temperature.
The helium abundance was determined from the HeI $`\lambda `$ 4471 and 6678 Å lines. Mean values of 0.070 $`\pm `$ 0.005 and 0.060 $`\pm `$ 0.005 respectively are obtained for regions 74C and 69C. We have estimated the contribution of neutral helium from the expression:
$$He^\mathrm{o}+\frac{He^+}{H^+}=\left(10.25\frac{O^+}{O}\right)^1\frac{He^+}{H^+}$$
(3)
(Kunth & Sargent 1983). The correction factors obtained for regions 74C and 69C are 1.17 and 1.20 respectively which yield He abundances y=0.081 for region 74C and y=0.068 for region 69C.
### 3.4 Wolf-Rayet features
A prominent Wolf-Rayet feature has been observed in region 74C (see Figure 6). Assuming a distance to NGC 4258 of 5.5 Mpc \[Martin et al. 1989\] and a constant extinction value through this region of c(H$`\beta `$) = 0.48, as derived from the Balmer and Paschen recombination lines, the total luminosity of the broad HeII feature, without the contribution of the \[FeIII\] $`\lambda `$ 4658 Å line, is (1.1 $`\pm `$ 0.1) $`\times `$ 10<sup>38</sup> erg s<sup>-1</sup>. This value comprises the broad features of both NIII $`\lambda \lambda `$ 4634, 4640 and HeII $`\lambda `$ 4686 Å lines. The contribution of the NIII lines to the WR bump is metallicity-dependent according to Smith \[Smith 1991\]. In our case this contribution represents 0.4 times the total emission; hence we obtain for the HeII line luminosity:
$$L(HeII\lambda 4686)=(6.6\pm 0.6)\times 10^{37}ergs^1$$
(4)
The non-detection of NV $`\lambda \lambda `$ 4604,4620 Å, suggests that the WR stars are of late N or intermediate type. Using the calibration of Vacca8& Conti \[Vacca & Conti 1992\], 38 $`\pm `$ 5 WN stars are found in this region.
Another fainter Wolf-Rayet feature has been observed in region 5NA. The bump luminosity is 2.2 $`\times `$ 10<sup>36</sup> erg s<sup>-1</sup>. The NIII lines contribution represent 0.45 times the total emission, hence the HeII line luminosity is 1.2 $`\times `$ 10<sup>36</sup> erg s<sup>-1</sup>, which is compatible with the presence of 1 WN star.
## 4 Functional parameters of the observed HII regions
Three are the fundamental parameters which control the emission line spectra of HII regions \[Díaz et al. 1991\]: the ionization parameter, the shape of the ionizing continuum, and the metallicity.
The ionization parameter – i. e. the ratio of the ionizing photon density to the particle density – is a measure of the degree of ionization of the nebula and can be deduced from the ratio of two lines of the same element corresponding to two different ionization states, e.g. \[OII\]/\[OIII\] or \[SII\]/\[SIII\]. Alternatively, it can also be determined from \[OII/H$`\beta `$\] or \[SII/H$`\alpha `$\] if the metallicity of the region is known (Díaz 1994). By using a large grid of single star photoionization models the following expressions are found (Díaz et al. 1991, Díaz 1999):
$$logU=1.39log([OII]/H\beta )+0.87log(Z/Z_{})1.68$$
(5)
$$logU=1.40log([SII]/H\beta )+1.10log(Z/Z_{})3.26$$
(6)
$$logU=0.80log([OII]/[OIII])3.02$$
(7)
$$logU=1.68log([SII]/[SIII])2.99$$
(8)
We have derived $`U`$ from the four expressions above. In all cases the value of U derived from the \[OII\]/\[OIII\] ratio is systematically lower than the rest thus implying low effective temperatures for the ionizing stars (see e.g. Díaz 1999). We have therefore discarded this value and computed U as the mean of the other three. These adopted ionization parameters – $`logU`$ – are listed in Table 5 and their uncertainty is estimated to be around $`\pm `$ 0.2 dex.
The shape of the ionizing continuum is directly related to the effective temperature of the stars that dominate the radiation field responsible for the ionization of the nebula. A recent version of the photoionization code CLOUDY \[Ferland 1996\] has been used to estimate the mean effective temperature of these stars. We have used Mihalas NLTE single-star stellar atmosphere models, with a closed geometry and a constant particle density through the nebula. For all the regions, consistency is found for effective temperatures around 35,000 K - 36,000 K. The results from these photoionization models compared to the observations are presented in Table 5.
## 5 Discussion
The H$`\alpha `$ fluxes for regions 74C and 69C, uncorrected for reddening, are 2692 and 287 $`\times `$ 10<sup>-16</sup> erg s<sup>-1</sup> cm<sup>-2</sup> respectively. These fluxes are to be compared with 3492 and 332 $`\times `$ 10<sup>-16</sup> erg s<sup>-1</sup> cm<sup>-2</sup> measured by C93. Considering that our observations have been obtained through a narrow slit and region 74C is rather extended, the agreement can be considered satisfactory. The combined H$`\alpha `$ flux of our regions 5NA and 5NB is 187 $`\times `$ 10<sup>-16</sup> erg s<sup>-1</sup> cm<sup>-2</sup> close to the value of 249 $`\times `$ 10<sup>-16</sup> erg s<sup>-1</sup> cm<sup>-2</sup> given by C93 for region 5N. Regarding the rest of the regions, we tentatively identify regions GA2, GA3 and GA4 with regions 59C, 58C, 54C of C93. Region GA1 probably corresponds to the outer parts of region 69C.
For each of our observed regions we have calculated the H$`\alpha `$ luminosity and the number of hydrogen ionizing photons from the observed H$`\alpha `$ flux corrected for their derived reddening (see Tables 2 and 3). Both quantities depend on the distance D to NGC 4258 according to the expressions:
$$L(H\alpha )=3.62\times 10^{37}\left(\frac{F(H\alpha )}{10^{14}}\right)\left(\frac{D}{5.5}\right)^2ergs^1$$
(9)
$$Q(H)=2.65\times 10^{49}\left(\frac{F(H\alpha )}{10^{14}}\right)\left(\frac{D}{5.5}\right)^2photonss^1$$
(10)
where F(H$`\alpha `$) is expressed in erg s<sup>-1</sup> cm<sup>-2</sup> and D is given in Mpc.
These values are given in Table 6 for the adopted distance of 5.5 Kpc. Only region 74C has an H$`\alpha `$ luminosity greater than 10<sup>39</sup> erg s<sup>-1</sup> and can be classified as a supergiant HII region as defined by Kennicutt (1983). The rest of the regions have H$`\alpha `$ luminosities typical of HII regions in early spiral galaxies, although all of them are greater than 10<sup>37</sup> erg s<sup>-1</sup> (Q(H) $`>`$ 10<sup>49</sup> photons s<sup>-1</sup>) requiring more than a single star for their ionization (Panagia 1973).
Filling factors for each observed region can be determined from the reddening corrected H$`\alpha `$ flux, F(H$`\alpha `$), and the derived ionization parameter, U, according to the expression:
$`ϵ=`$ $`0.34\left({\displaystyle \frac{F(H\alpha )}{10^{14}}}\right)^{1/2}\left({\displaystyle \frac{D}{5.5}}\right)^1\left({\displaystyle \frac{U}{10^3}}\right)^{3/2}`$
$`\left({\displaystyle \frac{\alpha _B(H^o,T)}{10^{13}}}\right)^1\left({\displaystyle \frac{n_e}{100}}\right)^{1/2}`$
where $`\alpha _B(H^o,T)`$ is the recombination coefficient for hydrogen and $`n_e`$ is the electron density of the emitting gas. The fillin0 factors computed in this way are listed in Table 6. We have used a value of $`\alpha _B(H^o,T)`$ = 3.76 $`\times `$ 10<sup>-13</sup> cm<sup>3</sup> s<sup>-1</sup>, corresponding to T= 7000 K and n<sub>e</sub>= 100 cm<sup>-3</sup> (Osterbrock 1989). For regions 74C and GA2 we have used values of n<sub>e</sub>= 60 and 160 cm<sup>-3</sup> respectively, as derived from the ratio between the red \[SII\] lines. For the rest of the regions for which only upper limits to the density could be derived, a value of n<sub>e</sub>=10 cm<sup>-3</sup> has been assumed.
According to C93 the angular sizes of regions 74C, 69C and 5N are 38.6, 16.8 and 20.5 arcsec respectively. The corresponding linear sizes, at the assumed distance to NGC 4258, are 1.0, 0.45 and 0.55 Kpc . These are comparable to the sizes of HII regions of moderate intensity found in other spiral galaxies (e.g. González-Delgado et al. 1997).
H$`\alpha `$ angular effective diameters – that is the diameters containing half the H$`\alpha `$ emission – of regions 74C, 69C and 5N are 4.3, 3.7 and 6.1 arcsec respectively (C93) which translates into 115, 99 and 163 pc. From the definition of U it is possible to obtain the angular size of the emitting region using the reddening corrected H$`\alpha `$ flux and the derived electron density:
$$\left(\frac{\varphi }{10^{\prime \prime }}\right)=0.06\left(\frac{F(H\alpha )}{10^{14}}\right)^{1/2}\left(\frac{U}{10^3}\right)^{1/2}\left(\frac{n_e}{100}\right)^{1/2}$$
(12)
where $`\varphi `$ is the angular diameter in units of 10 arcsec . This angular diameter does not depend on the assumed distance to the galaxy. For region 74C the derived angular diameter is 3.3 <sup>′′</sup> close to the value of 4.3 <sup>′′</sup> given by C93. For region 69C their observed value of the effective diameter of 3.7 <sup>′′</sup> yields a value of the electron density of n<sub>e</sub> = 5 cm<sup>-3</sup> close to the assumed one of 10 cm<sup>-3</sup>. Using this value of n<sub>e</sub> for the rest of the observed regions, except GA2 for which n<sub>e</sub> = 160 cm<sup>-3</sup> we obtain angular diameters between 1 and 2.6 arcsec.
We can also calculate the mass of ionized hydrogen following the expression:
$$M(HII)=\frac{m_pQ(H)}{n_e(1+y^+)\alpha _B(H^o,T)}$$
(13)
(Osterbrock 1989) where m<sub>p</sub> is the mass of the proton and $`y`$ the abundance by number of ionized helium, that we have taken as 0.10. Using observed and derived quantities this expression becomes
$$M(HII)=538\left(\frac{F(H\alpha )}{10^{14}}\right)\left(\frac{n_e}{100}\right)^1\left(\frac{D}{5.5}\right)^2M_{}$$
(14)
The corresponding masses of ionized hydrogen in each observed region are given in Table 6.
From evolutionary models of single ionizing clusters and radiation bounded HII regions, a relation between the number of ionizing Ly$`\alpha `$ photons per second per solar mass and the H$`\beta `$ equivalent width \[Díaz 1999\] can be found, which allows the estimation of the ionizing cluster mass by taking into account the cluster evolution.
$$log(Q(H)/M_{})=0.86log(EW(H\beta ))+44.48$$
(15)
For our observed regions all values for Q(H) are around 10<sup>49</sup> photon s<sup>-1</sup> except for region 74C for which a value of the order of 10<sup>51</sup> photon s<sup>-1</sup> is found. Hence, in the absence of dust, a lower limit for the mass of the ionizing clusters can be estimated by means of the H$`\beta `$ measured equivalent width and the H$`\alpha `$ luminosity for each region. The estimated ionizing cluster masses range from 4 $`\times `$ 10<sup>3</sup> M for region GA4 to 1.12 $`\times `$ 10<sup>5</sup> M for region 74C.
The high signal-to-noise spectrum of this latter region allows a more detailed modelling. The presence of WR features, in principle, provides a means to constrain the age of the ionizing population. The observed WR “bump” luminosity relative to H$`\beta `$ (L(WR)/H$`\beta `$) is 0.15. According to the models by Schaerer & Vacca (1998), at metallicity Z= 0.008 (0.4 solar), these high values are only found at an age of 4.5 Myr. However, this ratio is strongly dependent on metallicity and changes by more than an order of magnitude between Z=0.008 (0.4 solar) and Z= 0.02 (solar). Therefore, the slightly higher metallicity of region 74C (Z $``$ 0.012) might allow earlier ages for the WR stellar population. Clearly, the available grid for stellar evolution models is too coarse to elucidate this matter. The same applies to the equivalent width of the “bump” (EW(WR)) whose observed value is 10.8 Å. Since the NIII $`\lambda `$ 4640 Å line is strongly dependent on metallicity (Smith et al. 1996) it is probably more advisable to use only the flux and equivalent width of the HeII line. In that case, our observed values of L(HeII)/H$`\beta `$ and EW(HeII) are 0.09 and 6.6 Å respectively. Again, consistency is found for a model of Z=0.008 and an age of 4.5 Myr.
Since, according to the recent evolutionary models of Leitherer et al. (1999; STARBURST99), the observed value of the equivalent width of H$`\beta `$ (81 Å) points also to an age of around 4.5 Myr, we have computed the emission line spectrum corresponding to this single ionizing cluster. <sup>5</sup><sup>5</sup>5 All the models are computed assuming a standard Salpeter IMF for the stellar ionizing cluster and using the photoionization code CLOUDY (Ferland 1996). This set of models makes use of the same evolutionary tracks as the models of Schaerer & Vacca (1998) and therefore provides a selfconsistent frame. The results of the computation can be seen in Table 7 together with the observed emission line ratios. It can be seen that the model ionizing spectrum results too hard and does not reproduce the observations. This is not surprising in view of the low effective temperature found from single star photoionisation models (Table 5). It can be better appreciated in Fig 8 where the spectral energy distribution of the cluster can be compared to that of the Mihalas stellar atmosphere of a star with T<sub>eff</sub> = 35300 K which adequately reproduces the data. A high energy tail is clearly present in the ionising cluster as the result of the WR star contribution.
These models assume an enhanced mass loss during the evolution of massive stars, as prescribed by the stellar evolution models of Meynet et al. (1994). More moderate mass loss rates are assumed in the models of Bressan et al. (1993) and Fagotto et al. (1994) and with them a single population of 5.3 Myr can adequately reproduce the emission line spectrum (see Table 7). The corresponding ionizing spectrum is also shown in Fig 8. However, this ionizing cluster provides only a few WR stars (about 3$`\times `$ 10<sup>-5</sup> per solar mass of the ionizing cluster, for a metallicity Z=0.008; García-Vargas, Bressan & Díaz 1995). This is probably not enough to reproduce the observed values of L(WR)/H$`\beta `$.
The WR features observed in region 74C are similar to those found in other extragalactic HII regions. Table 8 shows the values of the WR feature intensities and equivalent widths for 74C and other three well studied extragalactic HII regions: NGC 604 in M33 (Díaz et al 1987; Terlevich et al. 1996), region A in NGC 3310 (Pastoriza et al. 1993) and region A in NGC 7714 (García-Vargas et al. 1997). Also given in the table are the metallicities of the regions in the form 12+log(O/H) and the effective temperature of the ionizing radiation estimated from single star photo-ionization models.
The tabulated values are difficult to explain in the light of the recent models of Schaerer & Vacca (1998). The expected decrease of the intensities of WR features with decreasing metallicity is not clearly observed. However, a decrease of the temperature of the ionizing radiation with increasing metallicity is apparent. In fact, region A in NGC 3310 and region 74C in NGC 4258 show the same L(WR)/H$`\beta `$ ratio despite their different metallicity (by a factor of about 3). This implies that both regions have similar ratio of WR to O stars, while, at the same time, their effective temperatures are considerably different.
The interpretation of the equivalent widths of WR features is even more difficult, since they may be affected by the continuum from any underlying stellar population. Composite populations have been postulated for region A in NGC 3310 (Pastoriza et al 1993) and region A in NGC 7714 (García-Vargas et al. 1997). Under this assumption, the equivalent widths, EW(WR) and EW(HeII), of the cluster with WR stars are increased to 6 and 4.8 Å for region A in NGC 3310 and 3.6 and 2.5 Å for region A in NGC 7714 which might still be compatible with Scharer & Vacca models. In the case of region 74C, combinations of two ionizing clusters can be found which adequately reproduce the emission line spectrum; however, this assumption would produce equivalent widths for the WR cluster higher than the maximum predicted for the models with Z=0.008, although probably compatible with solar metallicity models. Again, a finer metallicity grid for the WR cluster synthesis models would be needed in order to explore this possibility in further detail.
## 6 Summary and conclusions
We have analyzed eight HII regions in the LINER galaxy NGC 4258 using spectrophotometric observations between 3700 and 9650 Å. For two of the regions it has been possible to measure the electron temperature from the \[SIII\] $`\lambda `$ 6312 Å line, which allows the derivation of accurate abundances following standard methods. For the rest of the regions an empirical calibration based on the sulphur emission lines has been used to determine a mean oxygen content. The derived metallicities range from 0.3 to 0.6 Z. In particular, the metallicities found for regions 74C and 5N, previously reported by Oey & Kennicutt (1993) to be close to solar, are found to be lower by a factor of two.
For each observed region, we have also estimated the functional parameters: ionization parameter and the effective temperature of the ionizing cluster. Most regions show ionization parameters of the order of 10<sup>-3</sup> and effective temperatures of around 35300 K, except regions 74C and 69C for which higher ionization parameters are found. We have also derived the physical properties of the regions and their corresponding ionizing clusters: fillin0 factor, mass of ionized gas and mass of ionizing stars. Most of the regions, except 74C, have small ionizing clusters with masses in the range 4000 to 25000 solar masses. These values constitute in fact lower limits since the regions are assumed to be ionization bounded and the presence of dust has not been taken into account.
WR features have been detected in region 74C and, to a lower extent, in region 5NA. A detailed modelling has been carried out for region 74C using different sets of models: those of Schaerer & Vacca (1998) for WR populations and those of Leitherer et al. (1999) and García-Vargas, Bressan & Díaz (1995) for ionizing populations, to try to reproduce simultaneously both the WR features and the emission line spectrum. No consistent solution has been found: the models which better reproduce the WR features (a cluster 4.5 Myr old in the models of Schaerer & Vacc8 and Leitherer et al.) result too hard to reproduce the emission line spectrum. Conversely, a cluster 5.3 Myr old in the models of García-Vargas, Bressan & Díaz (1995) adequately reproduces the emission line spectrum, but produces too few WR stars to explain the observed WR features.
When comparing the WR “bump” and HeII intensities and equivalent widths with other well studied HII regions of different metallicities, the data on region 74C are found to fall into the observed ranges. However, the expected decrease of WR feauture intensities with decreasing metallicity is not clearly observed. On the other hand a decrease in the effective temperature of the ionizing radiation with increasing metallicity is apparent. This seems to imply that, while the number ratio of WR to O stars is similar for regions of metallicities differing by a factor of 3, the WR stars are cooler for the regions with higher metallicity.
Both more observations of confirmed high metallicity regions and a finer metallicity grid for the evolutionary synthesis models are needed in order to understand the ionizing populations of HII regions.
## Acknowledgements
The WHT is operated in the island of La Palma by the Issac Newton Group in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofísica de Canarias. We would like to thank CAT for awarding observing time. We also thank an anonymous referee for helpful suggestions.
E.T. is grateful to an IBERDROLA Visiting Professorship to UAM during which part of this work was completed. This work has been partially supported by DGICYT project PB-96-052.
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# Smoothed Particle Hydrodynamics for Relativistic Heavy Ion Collisions
## I Introduction
Hydrodynamic descriptions of high energy hadronic and nuclear collisions have a rather long history. Although, from theoretical point of view, it is not a trivial matter to justify their validity, they have been successful in reproducing certain features of these processes, such as the energy dependence of the average multiplicity and the transverse-energy distributions. More recently, relativistic fluid dynamics have become an important tool for the analysis of relativistic heavy-ion-collision processes (for example, and references there in). In these processes, the nuclear matter is expected to be compressed and heated up close to those states of the matter realized in the Big Bang Era of the Universe. Some laboratory data of these fascinating processes have already been obtained in the series of CERN experiments, and data at even higher temperatures and densities will soon be available in the forthcoming RHIC experiments.
The relativistic hydrodynamics is a description based on local conservation laws, together with the hypothesis of local thermodynamical equilibrium. The conservation laws are written in terms of the four-divergence of the energy-momentum tensor. The resulting system of equations is highly nonlinear and analytical solutions are only available for some very particular and limited configurations and equations of state . Thus numerical approaches are resorted to but they usually depend on some sophisticated techniques specific to some symmetry involved in the problem. When no symmetry is involved, these methods become computationally very expensive. However this is exactly the case when we are challenged to a realistic three dimensional simulation for nuclear-collision processes. There, we expect no geometrical symmetry so that a full 3D calculation is required.
One basic point in the hydrodynamic approach of relativistic nuclear collisions is that its principal ingredients, i.e., the equation of state of the matter and the initial conditions for the dynamics are not quite well known. On the contrary, we apply the hydrodynamic models to infer these informations on the properties of the matter in such a highly condensed and excited state. Thus we need to perform many hydrodynamic-model calculations for different equations of state and initial conditions to compare with the experimental data. In such a process, we actually don’t need the very precise solution of hydrodynamic equations, but a general flow pattern which characterizes the final configuration of the system as a response to a given set of equation of state and initial conditions. We are not interested in, and probably neither able to analyze at least in the present stage, any very precise local feature (for example sound ripples, small local perturbations, etc.) in these models. Extremely local properties, therefore, should be averaged out without spoiling the general flow pattern. This is particularly the case when one considers the questionable validity of the local thermal equilibrium in the problem of interest here. In short, for the study of hydrodynamic models of relativistic nuclear collisions, we prefer a rather simple scheme of solving hydrodynamic equations, not unnecessarily too precise but robust enough to deal with any kind of geometry. From this point of view, it is stressed in Ref. the advantage of variational approach to the relativistic hydrodynamics.
Among many numerical approaches, the Smoothed Particle Hydrodynamics (SPH) fits perfectly in with the variational formalism. As a consequence, this approach presents many desirable profiles of the variational approach, such as simplicity and robustness with respect to the changes of geometry, as well as the possibility of smoothing out undesirable local degrees of freedom. Furthermore, the SPH parametrizes the matter flow in terms of discrete Lagrangian coordinates (called “particles”) attached to some conserved quantity, such as baryon number. In this aspect, the SPH method is the best to the studies of relativistic nuclear collisions, where a extremely compressed and high-temperature hadronic matter expands into a very large space region. This is one of the principal advantages of the SPH over the other space-fixed grid algorithms when applied to the study of RHIC physics. The SPH algorithm was first introduced for astrophysical applications and used by now in several calculations, such as fragmentation of asteroids, supernova explosions, collision of neutron stars, etc. Several studies of the SPH algorithm for the ultrarelativistic regime of hydrodynamics have been done.
However, some specific aspects of the relativistic heavy ion collision processes deserve attention when applying the SPH. One of them is that, in the ultrarelativistic regime of central collisions, a large fraction of the incident energy is converted into produced particles. In particular, in the mid-rapidity region of central collisions, the most of these energies are in the form of produced pions and only a very small portion is carried by baryons. This is a very unfavorable situation to the conventional formulation of the SPH algorithm, where particles are associated to the matter defined by the conserved quantity, such as the baryon number. A direct application of the SPH based on the conserved baryon number may fail in the null baryon-number, pion-dominated region.
As mentioned above, the basic point of the SPH method is to introduce a set of “SPH - particles” which follow the flow of the fluid. However, the definition of a flow does not necessarily rely on the conserved quantity. For example, according to Landau a flow is defined in terms of the local Lorentz frame in which the energy-momentum tensor becomes diagonal. In this paper, we explore this aspect and formulate the SPH in terms of any extensive quantities defined in the Landau comoving local frame. We derive the relativistic SPH equations using the variational principle, taking the matter flow as the variable. We argue that the quantity we attribute to the SPH particles convenient for relativistic heavy-ion collisions are the entropy of the fluid. In this way, we can follow directly the entropy content and its change due to the dissipation mechanism, for example, the shock wave. This is particularly interesting for the system where the first order phase transition is present .
Another specific aspect of relativistic heavy-ion collision processes is how to set the initial conditions for the hydrodynamic motion. As mentioned before, this is not a solved problem. However, generally speaking, the hydrodynamic regime will be established from its initial excited state of the microscopic QCD degrees of freedom, only after a certain relaxation time $`\tau _{relax}`$. This means that the hypersurface of the onset of the hydrodynamic regime is characterized by a nearly constant local proper time rather than by a constant global coordinate time $`t`$. For ultrarelativistic regime, the difference becomes crucial. In the limit of very large initial energy, the Bjorken scaling solution is expected to be a good first order approximation at least for the longitudinal motion. Therefore, we shall write the hydrodynamic equations in terms of the Bjorken scaling coordinates $`\tau =\sqrt{t^2z^2}`$ and $`\eta =1/2\mathrm{ln}(t+z)/(tz)`$, replacing $`t`$ and $`z`$, as far as the longitudinal motion is concerned. These variables are also suitable for the setup of the initial conditions. Here, we show that the variational approach is also useful for the derivation of the SPH equations for an arbitrary curvilinear coordinate system.
We organize the present paper as follows. In Sec.II, we introduce the relativistic variational formulation with the SPH parametrization of one of the thermodynamical extensive quantities. In particular, we argue that the entropy is a convenient extensive quantity for the application of relativistic heavy-ion collisions. In this case, the change of entropy due to dissipative processes should be taken into account. This is particularly important in the presence of a shock wave or the first order phase transition. For this purpose, we deduce the entropy based SPH equations in the presence of bulk viscosity. In Sec.III, we apply our formulation for several relativistic systems such as one- and three-dimensional Landau model, the one presenting longitudinal Bjorken scaling solution plus transverse expansion, and one-dimensional relativistic fluid with shock waves. These results are compared with known solutions. Sec. IV is dedicated to the conclusions and perspectives of the present work.
## II Variational derivation of the Relativistic SPH equation
### A SPH Representation of Extensive Variables
The non-relativistic SPH can be formulated in terms of the variational principle. We show here that the relativistic SPH can also be formulated in this way. This guarantees that the SPH coordinates $`\left\{\stackrel{}{r}_a\left(t\right)\right\}`$are the optimal dynamical parameters to minimize the SPH-model action. The other advantage of this approach is that it automatically leads to the symmetrized form of the SPH equation to conserve linear and angular momenta of the system. This is the natural consequence of the Lorentz scalar nature of the action.
In the relativistic hydrodynamics, we assume that at any space-time point $`x=\{\stackrel{}{r},t\}`$ there exists a local reference frame where the energy-momentum tensor becomes diagonal and takes the form
$$T^{\mu \nu }(\stackrel{}{r},t)=\left(\begin{array}{cccc}\epsilon & 0& 0& 0\\ 0& P& 0& 0\\ 0& 0& P& 0\\ 0& 0& 0& P\end{array}\right),$$
(1)
where $`\epsilon `$ and $`P`$ are the (proper) energy density and the pressure of the fluid. From the hypothesis of local equilibrium, we assume that the thermodynamical relations are valid in each local frame.
Let $`A`$ be an arbitrary thermodynamical extensive quantity of the fluid such as baryon number, entropy or specific volume. The amount of $`A`$ contained in the infinitesimal volume element $`dV`$ is denoted by $`dA`$ so that the corresponding (proper) density $`a`$ is
$$a=\frac{dA}{dV}.$$
(2)
This is related to the corresponding density $`a^{}`$ measured in the space-fixed (calculational) frame as
$$a^{}(\stackrel{}{r},t)=\gamma a,$$
(3)
where $`\gamma `$ is the local Lorentz factor associated to the flow.
In the SPH representation, we parametrize this density by the following ansatz,
$$a^{}(\stackrel{}{r},t)=\underset{i}{}\nu _iW(\stackrel{}{r}\stackrel{}{r}_i\left(t\right);h),$$
(4)
where $`W(\stackrel{}{r}\stackrel{}{r}^{};h)`$ is a positive definite kernel function peaked at $`\stackrel{}{r}=\stackrel{}{r}^{}`$ with the normalization
$$d^3\stackrel{}{r}W(\stackrel{}{r}\stackrel{}{r}^{};h)=1.$$
(5)
The parameter $`h`$ represents the width of the kernel. In the limit, $`h0`$, we have
$`\underset{h0}{lim}W(\stackrel{}{r}\stackrel{}{r}^{};h)=\delta ^3\left(\stackrel{}{r}\stackrel{}{r}^{}\right).`$
As will be seen later, it is convenient to choose an even function for $`W.`$ It can be a Gaussian function in $`x=\left|\stackrel{}{r}\stackrel{}{r}^{}\right|/h`$, but in practice we often use the B-spline functions. The role of $`W`$ with a finite value of $`h`$ is to introduce a sort of short-wave length cut filter in the Fourier representation of the density $`a^{}`$. The total amount of $`A`$ of the system is obtained by integrating Eq.(4) over the whole space,
$$A_{tot}=d^3\stackrel{}{r}a^{}(\stackrel{}{r},t)=\underset{i}{}\nu _i$$
(6)
Physically speaking, it is clear from the above expression that we are replacing the system of continuous fluid by a collection of “SPH particles” each of which carries $`\nu _i`$ portion of the extensive quantity $`A`$.
The velocity of these particles are identified as the velocity of the fluid at their position $`\stackrel{}{r}_i(t),`$
$$\stackrel{}{v}_i=\frac{d\stackrel{}{r}_i}{dt},$$
(7)
so that the Lorentz factor of the $`i`$-th particle is given by $`\gamma _i=1/\sqrt{1\stackrel{}{v}_i^2}.`$ If $`A`$ is a conserved quantity, then $`\left\{\nu _i\right\}`$ should be constant in time. Then
$`{\displaystyle \frac{a^{}}{t}}`$ $`={\displaystyle \underset{i}{}}\nu _i\stackrel{}{v}_iW(\stackrel{}{r}\stackrel{}{r}_i\left(t\right);h)`$
$`={\displaystyle \underset{i}{}}\nu _i\stackrel{}{v}_iW(\stackrel{}{r}\stackrel{}{r}_i\left(t\right);h)`$
$`=\stackrel{}{j}_A,`$
where
$`\stackrel{}{j}_A={\displaystyle \underset{i}{}}\nu _i\stackrel{}{v}_iW(\stackrel{}{r}\stackrel{}{r}_i\left(t\right);h)`$
is the current density of $`A`$. Thus, the continuity equation is satisfied by the ansatz, (4), together with Eq.(7). If $`A`$ is not conserved, then $`\nu _i^{}s`$ are not constant in time. In this case, the continuity equation has the contribution from the time derivative of $`\nu _i^{}s`$ as
$$\frac{a^{}}{t}+\stackrel{}{j}_A=\underset{i}{}\dot{\nu }_iW(\stackrel{}{r}\stackrel{}{r}_i\left(t\right);h).$$
(8)
We consider the set of time dependent variables $`\left\{\stackrel{}{r}_i,i=1,\mathrm{},n\right\}`$ as the variational dynamical variables and their equations of motion are determined by minimizing the action for the hydrodynamic system. Here, $`\left\{\nu _i\right\}`$ are not dynamical variables and are determined by the initital conditions together with the constraints for the variational procedure (see the later discussion in 2.3.).
It may seem that the spatial extension of the particle $`i`$ is identified to that of the kernel $`W(\stackrel{}{r}\stackrel{}{r}_i;h)`$. However, the density of $`A`$ at the position of the particle $`i`$ is given by
$$a_i^{}=\underset{j}{}\nu _jW\left(\stackrel{}{r}_i\stackrel{}{r}_j\right),$$
(9)
from Eq.(4). Therefore, we may define the “specific volume” $`V_i`$ of the extensive quantity $`A`$ associated to the particle $`i`$ as
$$V_i\frac{\nu _i}{a_i}=\frac{\gamma _i\nu _i}{_j\nu _jW\left(\stackrel{}{r}_i\stackrel{}{r}_j\right)}.$$
(10)
Any other extensive quantities carried by the particle $`i`$ can be calculated easily. Let $`o^{}`$ be the density of another extensive quantity, say $`O,`$ measured in the calculational frame. Then the amount of $`O`$carried by the particle $`i`$ is
$$\nu _i\left(O/A\right)_i=\nu _i\left(o/a\right)_i=\nu _i\left(o/a\right)_i^{},$$
(11)
so that the density distribution of $`O`$ in the calculational frame is expressed as
$$o^{}(\stackrel{}{r},t)\underset{i}{}\nu _i\left(o/a\right)_i^{}W\left(\stackrel{}{r}\stackrel{}{r}_i\left(t\right)\right).$$
(12)
### B <br>SPH Action
The relativistic hydrodynamic equations can be obtained by the variational principle for the action
$$I=d^4x\epsilon ,$$
(13)
with respect to the matter density distribution, subjected to the constraint expressed by the continuity equation. Here, for the sake of simplicity, we discuss first the case of the Minkowsky metric. The argument can readily be generalized for more general coordinate systems (see Sec.III). The Lagrangian of the system is
$$L=d^3\stackrel{}{r}\epsilon .$$
(14)
We may consider $`\epsilon `$ as the Lagrangian density in the space fixed frame. Therefore, the corresponding Lagrangian for the system of SPH particles can be taken to be
$`L_{SPH}\left(\{\stackrel{}{r}_i,{\displaystyle \frac{d\stackrel{}{r}_i}{dt}}\}\right)`$ $`={\displaystyle \underset{i}{}}\nu _i\left(\epsilon /a^{}\right)_i`$ (15)
$`={\displaystyle \underset{i}{}}\left({\displaystyle \frac{E}{\gamma }}\right)_i,`$ (16)
where $`\left\{\stackrel{}{r}_i\left(t\right)\right\}`$ are the dynamical variables and $`E_i=\nu _i\left(\epsilon /a\right)_i`$ is the “rest energy” of the particle $`i.`$ The SPH model action is then
$$I_{SPH}=𝑑t\underset{i}{}\left(\frac{E}{\gamma }\right)_i.$$
(17)
### C Variational Procedure
The Lagrangian, given by Eq.(16), seems to show that the system is as if it is just a sum of free particles of “rest mass” $`E_i`$. However, there exists a basic difference. Namely, $`E_i^\text{ }`$’s are not constant with respect to the variation of positions. This is because the energy contained in each particle is determined by the configuration of other particles through thermodynamical relations. More explicitly, variations in the particles positions $`\left\{\stackrel{}{r}_i,i=1,\mathrm{},n\right\}`$ cause change of the volume occupied by each particle, which in turn modifies its energy (rest energy).
For our variational procedure, we consider changes of quantities associated only kinematically with virtual variations of the configuration. That is, except for the energy and the volume, all extensive quantities such as the entropy and the particle number should be kept constant while variations of $`\left\{\stackrel{}{r}_i\right\}`$ are taken. This implies that for variations of $`\stackrel{}{r}_i`$, we keep $`\nu _i`$ constant if $`A`$ is not the energy or the volume. Now, if we take $`A`$ as the volume, then its changes associated to variations of $`\stackrel{}{r}_i`$ are described by Eq.(4) where $`\nu _i`$ should be understood as the initial volume of the particle $`i`$. Therefore, even for the case $`A`$ is the volume, $`\nu _i`$ should be kept constant for variations of $`\stackrel{}{r}_i`$. In short, except for the case of the energy, the parameters $`\nu _i`$ should be kept constant while variations of $`\left\{\stackrel{}{r}_i\right\}`$ are taken. If we take $`A`$ as the energy, then the change of $`\nu _i`$ should be calculated as a function of the variation in $`\left\{\stackrel{}{r}_i\right\}`$. This introduces an additional complication without any practical merit. Therefore, for the sake of simplicity, we restrict ourselves to the case where $`A`$ is not the energy but any other extensive quantity.
We can write the change of energy associated with a virtual change of volume $`\delta V`$ as
$`\delta E`$ $`=P\delta V+\delta W`$ (18)
$`=P_{eff}\delta V,`$ (19)
where $`\delta W`$ is an additional work to change irreversibly the volume and $`P_{eff}`$ is the effective pressure. If the change of volume is performed in a quasi-static way and there exists no dissipative force, then this effective pressure coincides with the usual pressure $`P`$. However, if there exists some irreversible process associated with the volume change, then $`P_{eff}P`$ and we may write
$$P_{eff}=P+Q,$$
(20)
where $`Q\delta V`$ is the energy change due to some irreversible process. For a quasi-static adiabatic process, $`Q=0.`$ The introduction of $`Q`$ is important when we deal with shock phenomena (see Sec.III).
The variation of volume $`\delta V`$ for constant $`\nu _i`$’s is calculated from Eq.(10) as
$`\delta V_i`$ $`=V_i\left({\displaystyle \frac{\delta \gamma _i}{\gamma _i}}{\displaystyle \frac{\delta a_i^{}}{a_i^{}}}\right)`$ (21)
$`={\displaystyle \frac{\nu _i}{a_i}}\left(\gamma _i^2\stackrel{}{v}_i\delta \stackrel{}{v}_i+{\displaystyle \frac{1}{a_i^{}}}{\displaystyle \underset{j}{}}\nu _j\left(\delta \stackrel{}{r}_i\delta \stackrel{}{r}_j\right)W_{ij}\right),`$ (22)
where $`W_{ij}W(\stackrel{}{r}_i\stackrel{}{r}_j;h)`$. Using these relations, the variation of the action (17) with respect to $`\left\{\stackrel{}{r}_i,i=1,\mathrm{},n\right\}`$ becomes
$`\delta I_{SPH}`$ $`={\displaystyle 𝑑t\underset{i}{}\frac{1}{\gamma _i}\left(\left(P+Q\right)_i\delta V_iE_i\frac{\delta \gamma _i}{\gamma _i}\right)}`$
$`={\displaystyle }dt{\displaystyle \underset{i}{}}d\stackrel{}{r}_i\{{\displaystyle \frac{d}{dt}}\left[\nu _i\left({\displaystyle \frac{\epsilon +P+Q}{a}}\right)_i\gamma _i\stackrel{}{v}_i\right]`$
$`+\nu _i{\displaystyle \underset{j}{}}\nu _j[{\displaystyle \frac{1}{\gamma _i^2}}\left({\displaystyle \frac{P+Q}{a^2}}\right)_i+{\displaystyle \frac{1}{\gamma _j^2}}\left({\displaystyle \frac{P+Q}{a^2}}\right)_j]_iW_{ij}\}.`$
Here, we used the symmetry of $`W`$ and the property
$$_iW_{ij}=_jW_{ij}.$$
(23)
The requirement $`\delta I_{SPH}=0`$ for $`\delta \stackrel{}{r}_a`$ leads to
$$\frac{d}{dt}\left[\nu _i\left(\frac{P+Q+\epsilon }{a}\right)_i\gamma _i\stackrel{}{v}_i\right]=\underset{j}{}\left[\frac{\nu _i\nu _j}{\gamma _i^2}\left(\frac{P+Q}{a^2}\right)_i+\frac{\nu _i\nu _j}{\gamma _j^2}\left(\frac{P+Q}{a^2}\right)_j\right]_iW_{ij}.$$
(24)
The corresponding hydrodynamic equation in the continuum limit is
$$_\mu T^{\mu \nu }=_\mu \mathrm{\Sigma }^{\mu \nu },$$
(25)
where
$$\mathrm{\Sigma }^{\mu \nu }=Q\left[u^\mu u^\nu g^{\mu \nu }\right]$$
(26)
is the stress tensor for the bulk viscosity $`Q`$. This equation is the same as those discussed in Refs.().
Equation (24) should be complemented by an equation which expresses the conservation of energy (and other thermodynamical quantities, if any diffusion process is present). The energy conservation is written as
$$\frac{dE_i}{dt}=\left(P+Q\right)_i\frac{dV_i}{dt}.$$
(27)
On the other hand, from the second law of thermodynamics, we should have
$$\frac{dE_i}{dt}=P\frac{dV_i}{dt}+T_i\frac{dS_i}{dt}+\mu _i\frac{dN_i}{dt}$$
(28)
in equilibrium, where $`T,S,\mu ,N`$ are the temperature, entropy, chemical potential and particle number, respectively. Thus, we get
$$T_i\frac{dS_i}{dt}+\mu _i\frac{dN_i}{dt}=Q_i\frac{dV_i}{dt}.$$
(29)
In the absence of particle diffusion, the chemical equilibrium requires
$$\mu \frac{dN}{dt}=0$$
(30)
and, in this case,
$$T_i\frac{dS_i}{dt}=Q_i\frac{dV_i}{dt}.$$
(31)
### D Entropy Representation of SPH Equations
As we have mentioned in the Introduction, in applications to ultrarelativisitc nuclear collisions the baryon number is not a suitable quantity to represent the hydrodynamic flow, since most of the energy content is in the form of non-baryonic matter. This is particularly so in the central rapidity region. We may consider the energy content itself as the SPH base. However, as mentioned before, this choice introduces an additional constraint between the coordinates $`\left\{\stackrel{}{r}_i\right\}`$ and the extensive parameters $`\left\{\nu _i\right\}`$ of SPH particles, due to the energy conservation, and not desirable from the practical point of view. Therefore, we propose to take the entropy as the suitable extensive quantity for the SPH representation.
Let $`s^{}(t,\stackrel{}{r})`$ be the entropy density in the space-fixed (calculational) frame. From Eq.(9), we have
$`s_i^{}`$ $`=s^{}(t,\stackrel{}{r}_i\left(t\right))`$
$`={\displaystyle \underset{j}{}}\nu _jW\left(\stackrel{}{r}_i\stackrel{}{r}_j\right),`$
and the equations of motion for $`\left\{\stackrel{}{r}_i\right\}`$ are given by Eq.(24), which we write in the form,
$`{\displaystyle \frac{d\stackrel{}{r}_i}{dt}}`$ $`=\stackrel{}{v}_{i,}`$ (32)
$`{\displaystyle \frac{d\stackrel{}{\pi }_i}{dt}}`$ $`={\displaystyle \underset{j}{}}\left[{\displaystyle \frac{\nu _i\nu _j}{s_i^2}}\left(P+Q\right)_i+{\displaystyle \frac{\nu _i\nu _j}{s_j^2}}\left(P+Q\right)_j\right]_iW_{ij},`$ (33)
where
$$\stackrel{}{\pi }_i=\nu _i\left(\frac{P+Q+\epsilon }{s}\right)_i\gamma _i\stackrel{}{v}_i$$
(34)
is the momentum density associated with the $`ith`$ particle.
For the variational procedure, $`\nu _i`$’s are kept constant but it does not mean that they are constant in time. When there exists some non-adiabatic process, then $`Q0`$ and we have to use the energy conservation equation, Eq.(31) to determine the change of these parameters $`\left\{\nu _i\right\}`$. We have
$$\frac{1}{\nu _i}\frac{d\nu _i}{dt}=\frac{Q_i}{Ts_i^{}}\theta _i,$$
(35)
where
$`\theta _i`$ $`={\displaystyle \frac{1}{V_i}}{\displaystyle \frac{dV_i}{dt}}`$ (36)
$`=_\mu u^\mu .`$ (37)
Eqs. (32,33,35) constitute a system of first-order differential equations for $`\stackrel{}{r}_i,\stackrel{}{\pi }_i`$ and $`\nu _i`$, where the velocity $`\stackrel{}{v}_i`$ should be determined algebraically from Eq.(34). Further, we have to specify the equation of state, for example,
$`E=E(N,S),`$
and the dissipation pressure,
$`Q=Q(N,S,\theta ).`$
We will discuss these points later in practical examples.
### E SPH Equation for Generalized Coordinate System
The variational procedure can readily be extended to coordinate systems with non-Cartesian metric. The use of the generalized coordinate system is particularly important when we consider realistic initial conditions for simulations of RHIC processes. As we know, in a relativistic heavy-ion collisional process, the initial state is a cold, quantum nuclear matters flying separately. Just after the collision, the hadronic matter stays at a highly excited state and the materialization occurs only after $`\mathrm{~}1fm/c`$ in the proper time. Therefore, the local thermodynamical state would emerge for some local proper time and not for the global space-fixed time $`t`$. Thus, it is important to choose a convenient coordinate system for the description of relativistic heavy-ion collisions. For example, it is often used the hyperbolic time and longitudinal coordinates as is described later.
Let us consider a more general coordinate system,
$$ds^2=g_{\mu \nu }dx^\mu dx^\nu .$$
(38)
However, in order to unambiguously define the conserved quantities, we consider only the case when the time-like coordinate is orthogonal to the space-like coordinates,
$$g_{\mu 0}=0.$$
(39)
The action principle for the relativistic fluid motion can be written as
$$\delta I=\delta d^4x\sqrt{g}\epsilon =0,$$
(40)
together with the constraint for the conserved entropy current,
$$\left(su^\mu \right)_{;\mu }=\frac{1}{\sqrt{g}}_\mu \left(\sqrt{g}su^\mu \right)=0$$
(41)
or
$$\frac{1}{\sqrt{g}}_\tau \left(\sqrt{g}s\gamma \right)+\frac{1}{\sqrt{g}}\underset{i}{}_i\left(\sqrt{g}s\gamma v^i\right)=0,$$
(42)
where
$$v^i=\frac{u^i}{u^0}$$
(43)
and we use the notation,
$`\tau =x^0,\gamma =u^0.`$
The generalized gamma factor $`\gamma `$ is related to the velocity $`\stackrel{}{v}_a`$ through $`u_\mu u^\mu =1,`$ so that
$$\gamma =\frac{1}{\sqrt{g_{00}\stackrel{}{v}^T𝐠\stackrel{}{v}}},$$
(44)
where $`𝐠`$ is the $`3\times 3`$ space part of the metric tensor. That is
$$\left(g_{\mu \nu }\right)=\left(\begin{array}{cc}g_{00}& 0\\ 0& 𝐠\end{array}\right).$$
(45)
Let us now introduce the SPH representation. We may, for example, express the entropy density by the ansatz
$$\sqrt{g}s\gamma =s^{}s_{SPH}^{}=\underset{i}{}\nu _iW\left(\stackrel{}{r}\stackrel{}{r}_i\left(\tau \right)\right),$$
(46)
or by
$$s\gamma =s^{}s_{SPH}^{}=\underset{i}{}\nu _iW\left(\stackrel{}{r}\stackrel{}{r}_i\left(\tau \right)\right),$$
(47)
as well. These two possibilities, besides others, are simply different ways to parametrize a variational ansatz in terms of a linear combination of given functions $`W\left(\stackrel{}{r}\stackrel{}{r}_i\left(\tau \right)\right)`$. The most important property of an ansatz should be that $`W`$ satisfies the normalization condition imposed by the basic conserved quantity. Since the total entropy is expressed as
$$S=d^3\stackrel{}{r}\sqrt{g}s\gamma =\underset{i}{}\nu _i,$$
(48)
the normalization of $`W`$ should be taken to be
$$d^3\stackrel{}{r}W\left(\stackrel{}{r}\stackrel{}{r}^{}\right)=1,$$
(49)
for the parametrization Eq.(46) and
$$d^3\stackrel{}{r}\sqrt{g}W\left(\stackrel{}{r}\stackrel{}{r}^{}\right)=1,$$
(50)
for the parametrization Eq.(47). In the usual SPH calculations, it is not desirable to introduce in $`W`$ the space-time dependence through its normalization condition. In this aspect, the most natural way to introduce the SPH representation is Eq.(46). With this choice, the SPH action is given by
$`I_{SPH}`$ $`={\displaystyle 𝑑\tau d^3\stackrel{}{x}\underset{i}{}\nu _i\left(\frac{\sqrt{g}\epsilon }{\sqrt{g}s\gamma }\right)_iW\left(\stackrel{}{r}\stackrel{}{r}_i\left(\tau \right)\right)}`$ (51)
$`={\displaystyle 𝑑\tau \underset{i}{}\nu _i\left(\frac{\epsilon }{s\gamma }\right)_i}.`$ (52)
The variational principle leads to the following equation of motion,
$`{\displaystyle \frac{d}{d\tau }}\stackrel{}{\pi }_i`$ $`={\displaystyle \underset{j}{}}\nu _i\nu _j\left[{\displaystyle \frac{1}{\sqrt{g_i}\gamma _i^2}}{\displaystyle \frac{P_i+Q_i}{s_i^2}}+{\displaystyle \frac{1}{\sqrt{g_j}\gamma _j^2}}{\displaystyle \frac{P_j+Q_j}{s_j^2}}\right]_iW_{ij}`$ (53)
$`+{\displaystyle \frac{\nu _i}{\gamma _i}}{\displaystyle \frac{P_i+Q_i}{s_i}}\left({\displaystyle \frac{1}{\sqrt{g}}}\sqrt{g}\right)_i`$ (54)
$`+{\displaystyle \frac{\nu _i}{2}}\gamma _i\left({\displaystyle \frac{P+Q+\epsilon }{s}}\right)_i\left(g_{00}\stackrel{}{v}_i^T𝐠\stackrel{}{v}_i\right)_i,`$ (55)
where
$$\stackrel{}{\pi }_i=\gamma _i\nu _i\left(\frac{P+Q+\epsilon }{s}\right)_i𝐠\stackrel{}{v}_i$$
(56)
and the operator $``$ is just the simple derivative operator with respect to the coordinate variable in use.
### F Hyperbolic Coordinates
For ultrarelativistic heavy-ion collisions, a useful set of variables are
$`\tau `$ $`=\sqrt{t^2z^2},`$ (57)
$`\eta `$ $`={\displaystyle \frac{1}{2}}\mathrm{tanh}{\displaystyle \frac{t+z}{tz}},`$ (58)
$`\stackrel{}{r}_T`$ $`=\left(\begin{array}{c}x\\ y\end{array}\right).`$ (61)
As mentioned above, the initial conditions for RHIC processes are specified in terms of proper time rather than of coordinate time $`t`$. The variable $`\tau `$ is not exactly the physical proper time of the matter, but for the initial times it may approximate the proper time.
The metric tensor for this coordinate system is given by
$`g_{00}`$ $`=1,`$
$`𝐠`$ $`=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& \tau ^2\end{array}\right),`$
$`\sqrt{g}`$ $`=\tau .`$
Since the metric is space independent, we can use the parametrization,
$`\tau \gamma _is_i=s_i^{}={\displaystyle \underset{j=1}{\overset{n}{}}}\nu _jW\left(q_{ij}\right),`$
where
$`q_{ij}=\sqrt{\left(x_ix_j\right)^2+\left(y_iy_j\right)^2+\tau ^2\left(\eta _i\eta _j\right)^2}`$
and $`W`$ is normalized as
$`4\pi {\displaystyle _0^{\mathrm{}}}q^2𝑑qW\left(q\right)=1.`$
The SPH equation becomes
$`{\displaystyle \frac{d}{d\tau }}\stackrel{}{\pi }_i={\displaystyle \frac{1}{\tau }}{\displaystyle \underset{j}{}}\nu _i\nu _j\left[{\displaystyle \frac{1}{\gamma _i^2}}{\displaystyle \frac{P_i+Q_i}{s_i^2}}+{\displaystyle \frac{1}{\gamma _j^2}}{\displaystyle \frac{P_j+Q_j}{s_j^2}}\right]_iW_{ij},`$
where the $`\eta `$ component of the momentum is related to the velocity $`d\eta /d\tau `$ as
$`\pi _\eta =\tau ^2\nu \gamma \left({\displaystyle \frac{P+Q+\epsilon }{s}}\right){\displaystyle \frac{d\eta }{d\tau }},`$
whereas in the transverse directions, we have
$`\stackrel{}{\pi }_T=\nu \gamma \left({\displaystyle \frac{P+Q+\epsilon }{s}}\right){\displaystyle \frac{d\stackrel{}{r}_T}{d\tau }}.`$
The Lorentz factor is given by
$`\gamma ={\displaystyle \frac{1}{\sqrt{1\stackrel{}{v}_T^2\tau ^2v_\eta ^2}}}.`$
## III Examples
We have formulated the relativistic SPH method, appropriate to the study of RHIC processes. In order to check its validity and efficiency we apply, in this section, our method to several known problems and compare the results to the analytic or known numerical solutions.
### A Adiabatic Flow of Massless Pion Gas
For adiabatic flows of perfect gas, the entropy is conserved. Thus $`\dot{\nu }_i=0`$ and we have $`Q=0`$. In the following examples, we consider the cases where $`Q=0.`$
#### 1 Landau Model
Let us first investigate a well-known analytically soluble flow. The Landau model is one of the few examples of relativistic fluid flow for which analytical solution is available. Consider a one-dimensional, relativistic, massless, baryon-free gas initially at rest. The equation of state is
$`P={\displaystyle \frac{1}{3}}\epsilon =Cs^{4/3},`$
where
$`C=\left({\displaystyle \frac{15}{128\pi ^2}}\right)^{1/3}.`$
Since we are dealing with a perfect gas, $`Q=0`$. To apply the SPH method, we introduce the discrete one dimensional space variable $`\{\eta _i\left(t\right),i=1,..,n\}`$. The relation between the momentum and velocity is then
$$\pi =C\nu s^{1/3}\gamma ^{2/3}v.$$
(62)
In this case, $`v`$ can be solved analytically in terms of $`\pi `$.
In Fig.1-a and -b, we show the results of our SPH calculations, together with the exact solution .
Fig. 1-a, Fig.1-b
In this example, we took only $`100`$ particles with equally spaced $`\left\{\eta _i\right\}`$. As we see, in spite of a rather small number of particles, the SPH solution is quite satisfactory in this example. In particular, when we use the $`\eta \tau `$ coordinates with an appropriate distribution of $`\nu _i`$’s (Fig. 1-b), an excellent agreement with the analytical solution can be obtained.
#### 2 3D Scaling Solution
A simple analytical solution for a 3 dimensional relativistic pion gas is available. It is just a generalization of the one-dimensional scaling solution and given by
$`s={\displaystyle \frac{s_0}{\sqrt{\tau ^2x^2y^2}}}.`$
To see the efficiency of the SPH approach presented here, we reproduce this solution numerically in the full $`3D`$ numerical code, without making use of the spherical symmetry. In Fig.2, we show the result of such calculations.
Fig.2
As we see, our numerical calculations reproduces the analytical solution fairly well. One of the advantages of the SPH approach is that the coding for $`3D`$ cases is almost the same as for the $`1D`$ case. The only problem is a rapid increases of the number of particles for higher dimensions if we want to keep a high accuracy. A direct coding would require a computational time proportional to $`n^2`$, where $`n`$ is the number of particles. However, in the absence of long range forces such as the gravitational or Coulomb interactions, we may apply techniques such as the linked-list method to reduce the computational time to the order of $`n\mathrm{log}n`$. In the example shown above, we used $`50\times 50\times 50`$ particles, which took less than 2 minutes for one time step in a reasonable PC(Pentium-II). Here, we put a rather large number of particles, since it was a somewhat stringent test due to the divergent nature of the solution at the border.
#### 3 Transverse Expansion on Longitudinal Scaling Expansion
As a further test, closer to a realistic situation than that of Fig. 2, we calculated the transverse expansion of a cylindrically symmetric homogeneous massless pion gas, undergoing a longitudinal scaling expansion, and initially at rest in transverse directions. In Fig. 3, we show an example. Such a problem has been discussed by several authors as a useful base to understand the transverse expansion. Here, we compare our results (again a full $`3D`$ calculation without assuming cylindrical symmetry) with (2+1) numerical results, obtained by the use of the method of characteristics.
Fig. 3
In this example, we used also $`50\times 50\times 50`$ particles. The result is quite satisfactory. If we may decrease the accuracy by $`10\%`$, we can reduce the particle number almost by one order of magnitude.
### B Non Adiabatic Case: Shock Waves and Artificial Viscosity
As seen in the previous examples, our entropy-based relativistic SPH method works quite well for the adiabatic dynamics ($`Q=0`$) of the massless pion gas. However, for the application to realistic problems, it is fundamental to see how this scheme works for non-adiabatic cases ($`Q0`$), too. For this purpose, we study some examples of one dimensional shock problems.
#### 1 Compression Shock (Normal Matter)
Whenever there exists a shock wave, always exists the production of entropy through the shock front. The shock front manifests as a discontinuity in thermodynamical quantities in a hydrodynamic solution. Mathematically speaking, the shock front should be treated as a boundary to connect two distinct hydrodynamic solutions. To reproduce such a discontinuous behavior, the full degrees of freedom of hydrodynamics are required. The smoothed particle ansatz excludes such a possibility from the beginning. Since there do not exist short-wavelength excitation modes in the SPH ansatz, the energy and momentum conservation required by the hydrodynamics results in very rapidly oscillating motion of each particle. Such a situation occurs, for example, when a very high energy density gas is released into a low density region. This kind of shock, for the case of a baryon gas, is discussed in Ref.() and also, in the SPH context, in Ref.(). Here, we applied our entropy-based SPH approach to the massless pion gas.
Fig.4 gives a typical behavior of SPH solution for such a situation, if entropy production is not taken into account. As discussed above, there appear rapid oscillations in thermodynamical quantities just behind the shock front. Actually, this oscillation appears always in numerical approaches if entropy production is not included.
In order to avoid these oscillations, von Neuman and Richtmeyer introduced the concept of pseudoviscosity. The idea is to set the dissipative pressure where the shock wave discontinuity is present. To do this, Neuman and Richtmeyer proposed the ansatz
$`Q=\{\begin{array}{ccc}\left(\alpha \mathrm{\Delta }x\right)^2\rho \left(\dot{\rho }/\rho \right)^2,& & \dot{\rho }>0\\ 0,& & \dot{\rho }<0\end{array}`$
for nonrelativistic one-dimensional hydrodynamics. Here, $`\rho `$ is the mass density, $`\mathrm{\Delta }x`$ is the space grid size and $`\alpha `$ is a constant of the order of unity. In order to generalize the above pseudoviscosity for relativistic SPH case, we replace the quantity $`\dot{\rho }/\rho `$ by $`\theta =_\mu u^\mu `$ and $`\mathrm{\Delta }x`$ by $`h`$, where $`h`$ is as before the width of the smoothing kernel $`W`$. More precisely, we take the following form which is a slightly modified expression suggested by Ref.( ),
$$Q=\{\begin{array}{ccc}P\left[\alpha h\theta +\beta \left(h\theta \right)^2\right],& & \theta <0,\\ 0,& & \theta 0.\end{array}$$
(63)
As mentioned before, $`Q`$ is equivalent to the bulk viscosity and therefore there is no heat flow associated with it. What this artificial viscosity does is to convert the collective flow energy into the microscopic thermal energy. As a consequence, the total energy, that is, the sum of the collective flow energy and the internal thermal energy is still conserved.
Fig.5 is the solution of the same problem of Fig.4, but with the entropy production taken into account. In this calculation, the parameters have been chosen as
$`\alpha =2,\beta =4`$
and $`h=0.5fm`$ for $`1000`$ SPH particles. As we see, the rapid oscillations have been smoothed out (and in turn, the numerical calculation became much more efficient).
It is known that the overall energy- and momentum-flux conservation relates the ratio $`s_2/s_1`$ of entropy densities after and before the shock to the velocity $`v_s`$ of the shock front as (Hugoniot-Rankine relation)
$$\frac{s_2}{s_1}=\frac{2}{3^{3/4}}v_s\frac{\left(9v_s^21\right)^{1/4}}{\left(1v_s^2\right)^{5/4}}.$$
(64)
In Fig. 6, we show the velocity of the shock front obtained in our SPH calculations as function of the entropy ratio(dots). Each point corresponds to the different initial condition. They are compared with the Hugoniot-Rankine relation Eq.(64) (curve). The accordance shows that our SPH calculation reproduces faithfully the conservation of kinetic energy and momentum of the flow through the shock front.
#### 2 Rarefaction Shock
When the fluid presents a first order phase transition, a discontinuity appears in an expansion regime. This kind of shock wave has been discussed in connection to the QGP-hadron phase transition. In the present example, we use a simple MIT bag-model equation of state,
$`P=\{\begin{array}{ccc}\frac{\pi ^2}{30}T^4,& & TT_c,\\ \frac{37}{3}\frac{\pi ^2}{30}T^4B,& & TT_c,\end{array}`$
being $`B`$ the bag constant and $`T_c`$ the critical temperature.
Because decompression occurs with a constant pressure, we should have negative $`Q`$ values for $`\theta >0`$. We can use a similar expression for $`Q`$ as before,
$`Q=\{\begin{array}{ccc}0,& & \theta <0,\\ P_c\mathrm{\Theta }\left[\left(\epsilon \epsilon _h\right)\left(\epsilon _{QGP}\epsilon \right)\right]\left[\alpha h\theta \beta \left(h\theta \right)^2\right],& & \theta 0,\end{array}`$
where $`\mathrm{\Theta }`$ represents the Heaviside step function to let the viscosity be effective only in the transition region. In Fig. 6, we show a result of our calculation and compare it to the analytic solution. The real sharp shock front is a little bit smoothed, but this is due to the small number of SPH particles ($`N=1000`$ in this example) and hence a large $`h`$. Using the pseudoviscosity, the shock width turns out to be a few times $`h`$.
## IV Discussion and Perspectives
In the usual hydrodynamic computations using space grids, the symmetry of the problem is often a crucial factor to perform a calculation of reasonable size. The SPH method cures this aspect and furnishes a robust algorithm particularly appropriate to the description of processes where rapid expansions of the fluid should be treated. In this paper, we formulated an entropy-based SPH description of the relativistic hydrodynamics. We have shown that this approach is very promising for the study of ultrarelativistic nucleus-nucleus collision processes. The equations of motion are derived by a variational procedure from the SPH model action with respect to the Lagrangian comoving coordinates. This guarantees that the method furnishes the maximal efficiency for a given number of degrees of freedom, keeping strictly the energy and momentum conservation. For this reason, solutions can be obtained with a very reasonable precision, with a relatively small number of SPH particles. This is the basic advantage of the present method, when we analyze the event-by-event dynamics of the relativistic heavy-ion collisions.
On the other hand, the precision of this method increases rather slowly with the number of SPH particles. Therefore, a relatively large number of particles is required if one wants a very precise numerical solution. However, for the application to the RHIC physics, we may need rather crude precision especially if we consider the dubious validity of the rigorous hydrodynamics. For a calculation with typically $`10\%`$ errors, the SPH algorithm presented here furnishes a very efficient tool to study the flow phenomena in the RHIC physics.
A fundamental difficulty of the relativistic hydrodynamics for viscous fluid is that the dissipation term causes an intrinsic instability to the system described by Eq.(25). This instability basically comes from the fact that the dissipation term contains $`\theta `$ (see Eqs.(26,35)), so that it necessarily introduces the third time-derivative into the equation. This means that we have to specify, at least, a part of the acceleration as the initial condition. Even we specify the initial acceleration, the requirement of the internal self-consistency among the equations above leads to intrinsically unstable solutions. Israel proposed to cure these difficulties by introducing higher-order thermodynamics with respect to deviations from the equilibrium. In the examples presented in the present paper, we did not address this question and simply estimated the quantity $`\theta `$ from the quantities one time step before. In practice, this will cause no numerical instability and the behavior of the solution is quite satisfactory.
For the future application of the present program, we need to specify more realistic initial conditions and also to relate the final state to the physical observable quantities, such as particle spectra. The first point is now being in progress by introducing the initial energy and momentum distributions of the matter using the NEXUS algorithm. As for the second point, that is, the problem of particle production, the possibility of incorporation of the continuous emission mechanism is being studied.
This work was supported in part by PRONEX (contract no. 41.96.0886.00), FAPESP (contract nos. 98/02249-4 and 98/00317-2), FAPERJ (contract no.E-26/150.942/99) and CNPq-Brasil.
FIGURE CAPTIONS
a) Entropy profiles of Landau Model for different times. The numerical solutions correspond to the equal $`\nu `$ for all the particles and the Cartesian coordinate system was used. b) The same but the hyperbolic coordinate system was used (see text).
Three dimensional scaling solution for the massless pion gas. Cartesian coordinates are used for the transverse direction and the hyperbolic coordinates ($`\eta \tau `$) for the longitudial direction. Despite the spherical symmetry, the SPH calculation has been carried out in the full 3D code.
Transverse temperature profiles of a cilyndrically symmetric flow with longitudinally scaling expansion. The SPH results (circles) are compared with the numerical solution of the space-fixed grid method. The SPH calculation has been done in the full 3D code.
Shock wave in one dimensional pion gas. No viscosity is used.
After the introduction of the $`Q`$ term in the SPH calculation.
Test of Hugoniot-Rankine relation. Black circles are those of the SPH calculations with different initial conditions, and the solid line is Eq.(64).
Comparison of the SPH results with the analytical solution for the rarefaction shock. Here we took $`B=400MeV/fm^3`$.
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# OFF-DIAGONAL GENERALIZED VECTOR-DOMINANCE AND COLOUR-DIPOLES IN LOW-x DISaafootnote aContribution to DIS2000, Liverpool, April 2000
## Acknowledgments
It is a pleasure to thank Gorazd Cvetic, Arif Shoshi and Mikhail Tentioukov for a pleasant collaboration and my friends and colleagues at the MPI in Munich for warm hospitality.
## References
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# OUT OF EQUILIBRIUM FIELDS IN SELFCONSISTENT INFLATIONARY DYNAMICS. DENSITY FLUCTUATIONS
## I Introduction and Motivation
Inflationary cosmology has come of age. From its beginnings as a solution to the horizon, flatness, entropy and monopole problems, it has grown into the main contender for the explanation of the source of primordial fluctuations giving rise to large scale structure. There is evidence from the measurements of temperature anisotropies in the cosmic microwave background radiation (CMBR) that the scale invariant power spectrum predicted by generic inflationary models is consistent with observations and we can expect further and more exacting tests of the predictions of inflation when the MAP and PLANCK missions are flown. In particular, if the fluctuations that are responsible for the temperature anisotropies of the CMB truly originate from quantum fluctuations during inflation, determinations of the spectrum of scalar and tensor perturbations will constrain inflationary models based on particle physics scenarios and probably will validate or rule out specific proposals.
The tasks for inflationary universe researchers are then two-fold. First, models of inflation must be constructed on the basis of a realistic particle physics model. This is in contrast to the current situation where most, if not all acceptable inflationary models are ad-hoc in nature, with fields and potentials put in for the sole purpose of generating an inflationary epoch. Second, and equally important, the quantum dynamics of inflation must be understood. This is extremely important, especially in light of the fact that it is exactly this quantum behavior that is supposed to give rise to the primordial metric perturbations which presumably have imprinted themselves in the CMBR. This latter problem is the focus of this review.
The inflaton must be treated as a non-equilibrium quantum field . The simplest way to see this comes from the requirement of having small enough metric perturbation amplitudes which in turn requires that the quartic self coupling $`\lambda `$ of the inflaton be extremely small, typically of order $`10^{12}`$. Such a small coupling cannot establish local thermodynamic equilibrium (LTE) for all field modes; typically the long wavelength modes will respond too slowly to be able to enter LTE. In fact, the superhorizon sized modes will be out of the region of causal contact and cannot thermalize. We see then that if we want to gain a deeper understanding of inflation, non-equilibrium tools must be developed. Such tools exist and have now been developed to the point that they can give quantitative answers to these questions in cosmology -,. These methods permit us to follow the dynamics of quantum fields in situations where the energy density is non-perturbatively large ($`1/\lambda `$). That is, they allow the computation of the time evolution of non-stationary states and of non-thermal density matrices.
Our programme on non-equilibrium dynamics of quantum field theory, started in 1992, is naturally poised to provide a framework to study these problems. The larger goal of the program is to study the dynamics of non-equilibrium processes from a fundamental field-theoretical description, by solving the dynamical equations of motion of the underlying four dimensional quantum field theory for physically relevant problems: the early universe dynamics, high energy particle collisions, phase transitions out of equilibrium, symmetry breaking and dissipative processes.
The focus of our work is to describe the quantum field dynamics when the energy density is high. That is, a large number of particles per volume $`m^3`$, where $`m`$ is the typical mass scale in the theory. Usual S-matrix calculations apply in the opposite limit of low energy density and since they only provide information on in $``$ out matrix elements, are unsuitable for calculations of expectation values.
In high energy density situations such as in the early universe, the particle propagator (or Green function) depends on the particle distribution in momenta in a nontrivial way. This makes the quantum dynamics intrinsically nonlinear and calls to the use of self-consistent non-perturbative approaches as the large $`N`$ limit, Hartree and self-consistent one-loop approximations.
There are basically three different levels to study the early universe dynamics:
1. To work out the nonlinear dynamics of quantum fields in Minkowski spacetime. By non-linear dynamics we understand to solve the quantum equations of motion including the quantum back-reaction quantitatively -, . This level is in fact appropriate to describe high energy particle collisions .
2. To work out the nonlinear dynamics of quantum fields in fixed cosmological backgrounds. New phenomena arise then compared with 1. showing that a Minkowski analysis is not quantitatively precise for expanding universes.
3. A self-consistent treatment of the quantum fields and the cosmological background. That is, the metric is obtained dynamically from the quantum fields (matter source) propagating in the that metric.
We shall successively present the three levels of study. The first stage was reviewed in the 1996 Chalonge School . The second level is the subject of secs. VI and VII. We study the parametric and spinodal resonances both in FRW and de Sitter backgrounds wide range of initial conditions both in FRW and de Sitter backgrounds . \[Parametric resonance appears in chaotic inflationary scenarios for unbroken symmetry whereas spinodal unstabilities show up in new inflation scenarios with broken symmetry\]. Both types of unstabilities shut-off through the non-linear quantum evolution as described in secs. VI and VII both analytically and numerically. We follow the equation of state of the quantum matter during the evolution and analyze its properties.
The third stage of our approach is to apply non-equilibrium quantum field theory techniques to the situation of a scalar field coupled to semiclassical gravity, where the source of the gravitational field is the expectation value of the stress energy tensor in the relevant, dynamically changing, quantum state. In this way we can go beyond the standard analyses which treat the background as fixed and do not consider the non-linear quantum field dynamics.
In all cases 1. - 3. , the quantum fields energy- momentum tensor is covariantly conserved both at the regularized as well as the renormalized levels \- .
We mainly consider for the stage 3. new inflation scenarios where a scalar field $`\varphi `$ evolves under the action of a typical symmetry breaking potential. The initial conditions will be taken so that the initial value of the order parameter is near the top of the potential (the disordered state) with essentially zero time derivative. What we find is that the existence of spinodal instabilities, i.e. the fact that eventually (in an expanding universe) all modes will act as if they have a negative mass squared, drives the quantum fluctuations to grow non-perturbatively large. We have the picture of an initial wave-function or density matrix peaked near the unstable state and then spreading until it samples the stable vacua. Since these vacua are non-perturbatively far from the initial state (typically $`m/\sqrt{\lambda }`$, where $`m`$ is the mass scale of the field and $`\lambda `$ the quartic self-coupling), the spinodal instabilities will persist until the quantum fluctuations, as encoded in the equal time two-point function $`\mathrm{\Phi }(\stackrel{}{x},t)^2`$, grow to $`𝒪(m^2/\lambda `$).
This growth eventually shuts off the inflationary behavior of the scale factor as well as the growth of the quantum fluctuations (this last also happens in Minkowski spacetime ).
The scenario envisaged here is that of a quenched or super-cooled phase transition where the order parameter is zero or very small. Therefore one is led to ask:
a) What is rolling down?.
b) Since the quantum fluctuations are non-perturbatively large ( $`1/\lambda `$), will not they modify drastically the FRW dynamics?.
c) How can one extract (small?) metric perturbations from non-perturbatively large field fluctuations?
We address the questions a)-c) as well as other issues in sec. IX.
We choose such type of new inflationary scenario because the issue of large quantum fluctuations is particularly dramatic there. However, our methods do apply to any inflationary scenario as chaotic, extended and hybrid inflation.
## II Non-Equilibrium Quantum Field Theory, Semiclassical Gravity and Inflation
We present here the framework of the non-equilibrium closed time path formalism. For a more complete discussion, the reader is referred to -.
The time evolution of a system is determined in the Schrödinger picture by the functional Liouville equation
$$i\frac{\rho (t)}{t}=[H(t),\rho (t)],$$
(1)
where $`\rho `$ is the density matrix and we allow for an explicitly time dependent Hamiltonian as is necessary to treat quantum fields in a time dependent background. Formally, the solutions to this equation for the time evolving density matrix are given by the time evolution operator, $`U(t,t^{^{}})`$, in the form
$$\rho (t)=U(t,t_0)\rho (t_0)U^1(t,t_0).$$
(2)
The quantity $`\rho (t_0)`$ determines the initial condition for the evolution. We choose this initial condition to describe a state of local equilibrium in conformal time, which is also identified with the conformal adiabatic vacuum for short wavelengths.
Given the evolution of the density matrix (2), ensemble averages of operators are given by the expression (again in the Schrödinger picture)
$$𝒪(t)=\frac{Tr[U(t_0,t)𝒪U(t,t^{^{}})U(t^{^{}},t_0)\rho (t_0)]}{Tr\rho (t_0)},$$
(3)
where we have inserted the identity, $`U(t,t^{^{}})U(t^{^{}},t)`$ with $`t^{^{}}`$ an arbitrary time which will be taken to infinity. The state is first evolved forward from the initial time $`t_0`$ to $`t`$ when the operator is inserted. We then evolve this state forward to time $`t^{^{}}`$ and back again to the initial time.
We shall study the inflationary dynamics in a spatially flat Friedmann-Robertson-Walker background with scale factor $`a(t)`$ and line element:
$$ds^2=dt^2a^2(t)d\stackrel{}{x}^2.$$
(4)
Our Lagrangian density has the form
$$=\sqrt{g}\left[\frac{1}{2}_\mu \mathrm{\Phi }^\mu \mathrm{\Phi }V(\mathrm{\Phi })\right].$$
(5)
Our approach can be generalized to open as well as closed cosmologies.
Our program incorporates the non-equilibrium behavior of the quantum fields involved in inflation into a framework where the geometry (gravity) is dynamical and is treated self consistently. We do this via the use of semiclassical gravity where we say that the metric is classical and determined through the Einstein equations using the expectation value of the stress energy tensor $`T_{\mu \nu }`$. Such expectation value is taken in the dynamically determined state described by the density matrix $`\rho (t)`$. This dynamical problem can be described schematically as follows:
1. The dynamics of the scale factor $`a(t)`$ is driven by the semiclassical Einstein equations
$$\frac{1}{8\pi G_R}G_{\mu \nu }+\frac{\mathrm{\Lambda }_R}{8\pi G_R}g_{\mu \nu }+\left(\mathrm{higher}\mathrm{curvature}\right)=T_{\mu \nu }_R.$$
(6)
Here $`G_R,\mathrm{\Lambda }_R`$ are the renormalized values of Newton’s constant and the cosmological constant, respectively and $`G_{\mu \nu }`$ is the Einstein tensor. The higher curvature terms must be included to absorb ultraviolet divergences.
2. On the other hand, the density matrix $`\rho (t)`$ of the matter (that determines $`T_{\mu \nu }_R`$) obeys the Liouville equation (1) with $`H`$ is the evolution Hamiltonian, which is dependent on the scale factor, $`a(t)`$.
It is this set of equations we must try to solve specifying the appropriate initial conditions.
### A On the initial state: dynamics of phase transitions
The situations we consider are
1. the theory admits a symmetry breaking potential and in which the field expectation value starts its evolution near the unstable point.
2. The symmetry is not broken and the field expectation value starts its evolution at a finite distance from the absolute minimum.
There is an issue as to how the field got to have an expectation value near the unstable point (typically at $`\mathrm{\Phi }=0`$) as well as an issue concerning the initial state of the non-zero momentum modes.
Since our background is an FRW spacetime, it is spatially homogeneous and we can choose our state $`\rho (t)`$ to respect this symmetry. Starting from the full quantum field $`\mathrm{\Phi }(\stackrel{}{x},t)`$ we can extract its expectation value $`\varphi (t)`$ by writing:
$$\mathrm{\Phi }(\stackrel{}{x},t)=\varphi (t)+\mathrm{\Psi }(\stackrel{}{x},t),\varphi (t)=\mathrm{Tr}[\rho (t)\mathrm{\Phi }(\stackrel{}{x},t)]\mathrm{\Phi }(\stackrel{}{x},t)).$$
(7)
The quantity $`\mathrm{\Psi }(\stackrel{}{x},t)`$ represents the quantum fluctuations about the zero mode $`\varphi (t)`$ and clearly satisfies $`\mathrm{\Psi }(\stackrel{}{x},t)=0`$.
We need to choose a basis to represent the density matrix. A natural choice consistent with the translational invariance of our quantum state is that given by the Fourier modes, in comoving momentum space, of the quantum fluctuations $`\mathrm{\Psi }(\stackrel{}{x},t)`$:
$$\mathrm{\Psi }(\stackrel{}{x},t)=\frac{d^3k}{(2\pi )^3}\mathrm{exp}(i\stackrel{}{k}\stackrel{}{x})\psi _k(t).$$
(8)
In this language we can state our ansatz for the initial condition of the quantum state as follows. We take the zero mode $`\varphi (t=0)=\varphi _0,\dot{\varphi }(t=0)=0`$, where $`\varphi _0`$ will typically be very near the origin for broken symmetry and at a finite distance from it in the unbroken symmetry case. The initial conditions on the the nonzero modes $`\psi _k(t)`$ will be chosen such that the initial density matrix $`\rho (t=0)`$ describes a vacuum state (i.e. an initial state in local thermal equilibrium at a temperature $`T_i=0`$). There are some subtleties involved in this choice. First, as explained in , in order for the density matrix to commute with the initial Hamiltonian, we must choose the modes to be initially in the conformal adiabatic vacuum (these statements will be made more precise below). This choice has the added benefit of allowing for time independent renormalization counterterms to be used in renormalizing the theory.
We are making the assumption of an initial vacuum state in order to be able to proceed with the calculation. It would be interesting to understand what forms of the density matrix can be used for other more general initial conditions.
The assumptions of an initial equilibrium vacuum state are essentially the same used in refs. , and in the analysis of the quantum mechanics of inflation in a fixed de Sitter background.
As discussed in the introduction, if we start from such an initial state, spinodal or parametric instabilities will drive the growth of non-perturbatively large quantum fluctuations. In order to deal with these, we need to be able to perform calculations that take these large fluctuations into account. Although the quantitative features of the dynamics will depend on the initial state, the qualitative features associated with spinodal or parametric unstabilities are fairly robust for a wide choice of initial states that describe a phase transition.
## III The Inflaton Model and the Equations of Motion
Having recognized the appearance of large quantum fluctuations driven by parametric or spinodal unstabilities, we need to study the dynamics within a non-perturbative framework. That is, a framework allowing calculations for non-perturbatively large energy densities. We require that such a framework be: i) renormalizable, ii) covariant energy conserving, iii) numerically implementable. There are very few schemes that fulfill all of these criteria: the large $`N`$ and the Hartree approximation-. Whereas the Hartree approximation is basically a Gaussian variational approximation that in general cannot be consistently improved upon, the large $`N`$ approximation can be consistently implemented beyond leading order. In addition, the presence of a large number of fields in most of the GUT’s models suggest that the large $`N`$ limit will be actually a realistic one. Moreover, for the case of broken symmetry it has the added bonus of providing many light fields (associated with Goldstone modes) that will permit the study of the effects of other fields which are lighter than the inflaton on the dynamics. Thus we will study the inflationary dynamics within the framework of the large $`N`$ limit of a scalar theory in the vector representation of $`O(N)`$ both for unbroken and broken symmetry. In the second case we will have a quenched phase transition.
We assume that the universe is spatially flat with a metric given by eq.(4). The matter action and Lagrangian density are given by eq.(5),
$$S_m=d^4x_m=d^4xa^3(t)\left[\frac{1}{2}\dot{\stackrel{}{\mathrm{\Phi }}}^2(x)\frac{1}{2}\frac{(\stackrel{}{}\stackrel{}{\mathrm{\Phi }}(x))^2}{a^2(t)}V(\stackrel{}{\mathrm{\Phi }}(x))\right]$$
(9)
$$V(\stackrel{}{\mathrm{\Phi }})=\frac{m^2}{2}\stackrel{}{\mathrm{\Phi }}^2+\frac{\lambda }{8N}\left(\stackrel{}{\mathrm{\Phi }}^2+\frac{2Nm^2}{\lambda }\right)^2+\frac{1}{2}\xi \stackrel{}{\mathrm{\Phi }}^2,$$
(10)
where $`m^2>0`$ for unbroken symmetry and $`m^2<0`$ for broken symmetry. Here $`(t)`$ stands for the scalar curvature
$$(t)=6\left(\frac{\ddot{a}(t)}{a(t)}+\frac{\dot{a}^2(t)}{a^2(t)}\right),$$
(11)
The $`\xi `$-coupling of $`\stackrel{}{\mathrm{\Phi }(x)}^2`$ to the scalar curvature $`(t)`$ has been included since arises anyhow as a consequence of renormalization.
The gravitational sector includes the usual Einstein term in addition to a higher order curvature term and a cosmological constant term which are necessary to renormalize the theory. The action for the gravitational sector is therefore:
$$S_g=d^4x_g=d^4xa^3(t)\left[\frac{(t)}{16\pi G}+\frac{\alpha }{2}^2(t)K\right].$$
(12)
with $`K`$ being the cosmological constant. In principle, we also need to include the terms $`R^{\mu \nu }R_{\mu \nu }`$ and $`R^{\alpha \beta \mu \nu }R_{\alpha \beta \mu \nu }`$ as they are also terms of fourth order in derivatives of the metric (fourth adiabatic order), but the variations resulting from these terms turn out not to be independent of that of $`^2`$ in the flat FRW cosmology we are considering.
The variation of the action $`S=S_g+S_m`$ with respect to the metric $`g_{\mu \nu }`$ gives us Einstein’s equation
$$\frac{G_{\mu \nu }}{8\pi G}+\alpha H_{\mu \nu }+Kg_{\mu \nu }=<T_{\mu \nu }>,$$
(13)
where $`G_{\mu \nu }`$ is the Einstein tensor given by the variation of $`\sqrt{g}`$, $`H_{\mu \nu }`$ is the higher order curvature term given by the variation of $`\sqrt{g}^2`$, and $`T_{\mu \nu }`$ is the contribution from the matter Lagrangian. With the metric (4), the various components of the curvature tensors in terms of the scale factor are:
$`G_0^0`$ $`=`$ $`3\left({\displaystyle \frac{\dot{a}}{a}}\right)^2,G_\mu ^\mu ==6\left({\displaystyle \frac{\ddot{a}}{a}}+{\displaystyle \frac{\dot{a}^2}{a^2}}\right),`$ (14)
$`H_0^0`$ $`=`$ $`6\left({\displaystyle \frac{\dot{a}}{a}}\dot{}+{\displaystyle \frac{\dot{a}^2}{a^2}}{\displaystyle \frac{1}{12}}^2\right),H_\mu ^\mu =6\left(\ddot{}+3{\displaystyle \frac{\dot{a}}{a}}\dot{}\right).`$ (15)
Eventually, when we have fully renormalized the theory, we will set $`\alpha _R=0`$ and keep as our only contribution to $`K_R`$ a piece related to the matter fields which we shall incorporate into $`T_{\mu \nu }`$.
The use of semiclassical gravity (i. e. neglecting graviton loops) is justified since graviton loops are suppressed by factors $`(m/M_{Pl})^2`$ where $`m`$ is at the GUT scale and hence $`(m/M_{Pl})^210^8`$.
## IV The Large N Limit for a scalar field with an arbitrary invariant self-interaction
We present here the systematic derivation of the $`1/N`$ expansion for the scalar model with an arbitrary $`O(N)`$-invariant self-interaction and the scalar field in the vector representation of $`O(N)`$-. We consider Minkowski space-time. The generalization to the cosmological space-time (4) is discussed at the end of the section.
The action and Lagrangian density are given by,
$$S=d^4x,=\frac{1}{2}[\stackrel{}{\mathrm{\Phi }}(x)]^2NV\left(\frac{1}{N}\stackrel{}{\mathrm{\Phi }}^2(x)\right).$$
(16)
This choice of the $`N`$ dependence in the interaction ensures that the large $`N`$ limit exists. In particular, for a quartic potential we choose according to eq.(10):
$$V\left(\frac{1}{N}\stackrel{}{\mathrm{\Phi }}^2(x)\right)=\frac{\lambda }{8N^2}\left(\stackrel{}{\mathrm{\Phi }}^2+\frac{2Nm^2}{\lambda }\right)^2.$$
(17)
The functional integral for the model takes the form
$$𝒵[J(.)]=D\stackrel{}{\mathrm{\Phi }}\mathrm{exp}id^4x[+\stackrel{}{\mathrm{\Phi }}(x).\stackrel{}{J}(x)]$$
(18)
where $`\stackrel{}{J}(x)`$ is an external source introduced to generate the correlation functions. For example, the connected two point correlation function is expressed as
$$<T\mathrm{\Phi }_i(x)\mathrm{\Phi }_k(y)>=i^2\frac{\delta ^2}{\delta J_i(x)\delta J_k(y)}\mathrm{log}𝒵[J(.)]|_{J_i(.)=0}.$$
(19)
In order to compute the large $`N`$ limit it is convenient to replace the interaction term by the following functional integral representation :
$`\mathrm{exp}iN{\displaystyle d^4xV(\frac{1}{N}\stackrel{}{\mathrm{\Phi }}^2(x))}`$ $`=`$ $`{\displaystyle Dz\mathrm{\Pi }_x\delta \left(z(x)\frac{1}{N}\stackrel{}{\mathrm{\Phi }}^2(x)\right)e^{iN{\scriptscriptstyle d^4xV(z(x))}}}`$ (20)
$`=`$ $`{\displaystyle DzD\sigma e^{iN{\scriptscriptstyle d^4x\left[{\scriptscriptstyle \frac{1}{2}}\sigma (x)\left(z(x){\scriptscriptstyle \frac{1}{N}}\stackrel{}{\mathrm{\Phi }}^2(x)\right)V(z(x))\right]}}}`$ (22)
Inserting eq.(20) into eq.(18), the integration over $`\stackrel{}{\mathrm{\Phi }}(x)`$ becomes gaussian and therefore can be exactly computed with the result:
$`𝒵[J(.)]={\displaystyle }{\displaystyle }D\stackrel{}{\mathrm{\Phi }}DzD\sigma e^{i{\scriptscriptstyle d^4x}\left[\frac{1}{2}\left[\stackrel{}{\mathrm{\Phi }}(x)\right]^2+\frac{N}{2}\sigma (x)z(x)\frac{1}{2}\sigma (x)\stackrel{}{\mathrm{\Phi }}^2(x)NV(z(x))+\stackrel{}{\mathrm{\Phi }}(x).\stackrel{}{J}(x)\right]}`$ (23)
(24)
$`={\displaystyle }{\displaystyle }DzD\sigma \{det[^2+\sigma (.)]\}^{\frac{N}{2}}e^{iN{\scriptscriptstyle }d^4x[\frac{1}{2}\sigma (x)z(x)V(z(x))]+\frac{i}{2}{\scriptscriptstyle }d^4xd^4yJ_i(x)G(x,y,\sigma (.))J_i(y)}`$ (25)
Here, $`G(x,y,\sigma (.))`$ is the inverse operator of $`^2+\sigma (.)`$. That is,
$$[^2+\sigma (x)]G(x,y,\sigma (.))=\delta (xy)$$
(26)
It is useful to introduce a source $`NK(x)`$ for the field $`z(x)`$. This will permit to generate the correlation functions of the composite field $`\stackrel{}{\mathrm{\Phi }}^2(x)`$. The generating functional now takes the form
$`𝒵[J(.),K(.)]={\displaystyle }{\displaystyle }DzD\sigma \{det[^2+\sigma (.)]\}^{\frac{N}{2}}e^{iN{\scriptscriptstyle d^4x\left[{\scriptscriptstyle \frac{1}{2}}\sigma (x)z(x)V(z(x))\right]}}`$ (27)
(28)
$`e^{\frac{i}{2}{\scriptscriptstyle }d^4xd^4yJ_i(x)G(x,y,\sigma (.))J_i(y)+iN{\scriptscriptstyle }d^4xK(x)z(x)}={\displaystyle DzD\sigma e^{iN𝒮[z(.),\sigma (.)]}}.`$ (29)
Here,
$`𝒮[z(.),\sigma (.)]`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{log}det[^2+\sigma (.)]+{\displaystyle }d^4x[{\displaystyle \frac{1}{2}}\sigma (x)z(x)V(z(x))+K(x)z(x)]`$ (30)
$`+`$ $`{\displaystyle \frac{1}{2N}}{\displaystyle }d^4xd^4yG(x,y,\sigma (.))J_i(x)J_i(y).`$ (32)
Notice that the source terms in $`𝒮`$ are both of order one since $`K(x)`$ and $`J_i(x)`$ are assumed to be of order one.
The functional derivatives with respect to the source $`K(x)`$ at $`K(x)=0`$ produce the insertions
$$\frac{1}{i}\frac{\delta }{\delta K(x)}Nz(x)=\stackrel{}{\mathrm{\Phi }}^2(x).$$
(33)
as follows from eqs.(20) and (27) That is, $`K(x)`$ is the source of the composite field $`\stackrel{}{\mathrm{\Phi }}^2(x)`$. The correlations of $`\stackrel{}{\mathrm{\Phi }}^2(x)`$ follow as functional derivatives of $`\mathrm{log}𝒵[J(.),K(.)]`$ with respect to $`K(x)`$. For example, we have for the two points function
$$<T\stackrel{}{\mathrm{\Phi }}^2(x)\stackrel{}{\mathrm{\Phi }}^2(y)>=i^2\frac{\delta ^2}{\delta K(x)\delta K(y)}\mathrm{log}𝒵[J(.),K(.)]|_{J(.)=,K(.)=0}.$$
Since in eq.(27) the $`N`$-dependence is explicit in the exponent, we can take the large $`N`$ limit by looking for the stationary points of the action there.
Extremizing $`𝒮[z(.),\sigma (.)]`$ with respect to the field $`\sigma (x)`$ yields:
$$iG(x,x,\sigma (.))+z(x)\frac{1}{N}[d^4yG(x,y,\sigma (.))J_i(y)]^2=0$$
(34)
Where we used that,
$$\frac{\delta G(y,z,\sigma (.))}{\delta \sigma (x)}=G(y,x,\sigma (.))G(x,z,\sigma (.)).$$
(35)
Eq.(35) can be derived as follows. Taking the functional derivative of eq.(26) with respect to $`\sigma (.)`$ yields,
$$[^2+\sigma (x)]\frac{\delta G(x,y,\sigma (.))}{\delta \sigma (z)}=\delta (xz)G(x,y,\sigma (.))$$
This is an equation for $`\frac{\delta G(x,y,\sigma (.))}{\delta \sigma (z)}`$ that can be solved using the inverse operator of $`\left[^2+\sigma (x)\right]`$ as given by eq.(26). This gives eq.(35).
Extremizing $`𝒮[z(.),\sigma (.)]`$ with respect to the field $`z(x)`$ yields:
$$\frac{1}{2}\sigma (x)V^{}\left(z(x)\right)+K(x)=0.$$
(36)
Eqs.(34)-(36) define the saddle point $`\sigma (x),z(x)`$ as a functional of the sources $`J(.),K(.)`$. One has thus to make the following shift of the functional integration variables
$$\sigma (x)=\sigma (x)+\xi (x),z(x)=z(x)+w(x)$$
(37)
where $`\xi (x),w(x)`$ are the new functional integration variables.
Now we have to insert the change (37) into eq.(30) and expand $`𝒮[z(.),\sigma (.)]`$ in powers of $`\xi (x)`$ and $`w(x)`$. The zeroth order, that is $`𝒮[z(.),\sigma (.)]`$ yields the $`N=\mathrm{}`$ limit of the model. The quadratic part in $`\xi (x)`$ and $`w(x)`$ provides the propagators and the higher orders provide the vertices of the $`1/N`$ perturbation theory.
### A The Large $`N`$ Limit for the $`[\stackrel{}{\mathrm{\Phi }}^2]^2`$ theory
For simplicity, we will restrict ourselves to the $`\mathrm{\Phi }^4`$ theory with potential (17). That is,
$$V(z)=\frac{\lambda }{8}\left(z+\frac{2m^2}{\lambda }\right)^2=\frac{\lambda }{8}z^2+\frac{m^2}{2}z+\frac{m^4}{2\lambda }.$$
(38)
In this important case we have $`V^{}(z)=\frac{\lambda z}{4}+\frac{m^2}{2}`$ and we can eliminate $`z(x)`$ using the saddle point equation (36),
$$z(x)=\frac{2}{\lambda }\left[\sigma (x)+2K(x)m^2\right].$$
(39)
Equivalently, in this case we can integrate exactly over $`z(x)`$ in eqs.(20), (23) and (27) since these functional integrals become gaussian in the field $`z`$ when $`V(z)`$ is given by eq.(38).
The saddle point equation (34) for $`\sigma (x)`$ becomes
$$iG(x,x,\sigma (.))+\frac{2}{\lambda }[\sigma (x)+2K(x)m^2]\frac{1}{N}[d^4yG(x,y,\sigma (.))J_i(y)]^2=0$$
(40)
where $`G(x,y,\sigma (.))`$ is defined by eq.(26).
The action $`𝒮`$ at the saddle point takes then the form
$`𝒮[\sigma (.)]`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{log}det[^2+\sigma (.)]+{\displaystyle \frac{1}{2\lambda }}{\displaystyle }d^4x[\sigma (x)+2K(x)m^2]^2`$ (41)
$`+`$ $`{\displaystyle \frac{1}{2N}}{\displaystyle }d^4xd^4yG(x,y,\sigma (.))J_i(x)J_i(y).`$ (43)
where we used eqs.(30), (34), (38) and (39).
Hence, from eqs.(27) and (40) we obtain that in the $`N=\mathrm{}`$ limit
$$\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}𝒵[J(.),K(.)]=i𝒮[\sigma (.)].$$
(44)
Furthermore, we can easily compute the two points function of the scalar field $`\stackrel{}{\mathrm{\Phi }}(x)`$ using eqs.(19), (41) and (44) with the result
$$<T\mathrm{\Phi }_i(x)\mathrm{\Phi }_k(y)>=i\delta _{ik}G(x,y,\sigma (.))+𝒪\left(\frac{1}{N}\right).$$
(45)
and $`G(x,y,\sigma (.))`$ was defined through eq.(26). The $`\stackrel{}{\mathrm{\Phi }}(x)`$ propagator turns to be of order one in the $`N=\mathrm{}`$ limit.
The $`O(N)`$ invariance is here explicit. We find for the expectation value of $`\stackrel{}{\mathrm{\Phi }}(x)`$ (one-point function) in the infinite $`N`$ limit,
$`<\stackrel{}{\mathrm{\Phi }}(x)_i>`$ $`=`$ $`i^1{\displaystyle \frac{\delta }{\delta J_i(x)}}\mathrm{log}𝒵[J(.),K(.)]|_{J_i(.)=K(.)=0}=N{\displaystyle \frac{\delta }{\delta J_i(x)}}𝒮[\sigma (.)]|_{J_i(.)=K(.)=0}`$ (46)
$`=`$ $`{\displaystyle }d^4yG(x,y,\sigma (.))J_i(x)+𝒪\left({\displaystyle \frac{1}{N}}\right).`$ (49)
and for the composite field $`\stackrel{}{\mathrm{\Phi }}^2(x)`$
$`{\displaystyle \frac{1}{N}}<\stackrel{}{\mathrm{\Phi }}^2(x)>`$ $`=`$ $`{\displaystyle \frac{i}{N}}{\displaystyle \frac{\delta }{\delta K(x)}}\mathrm{log}𝒵[J(.),K(.)]|_{J_i(.)=K(.)=0}`$ (50)
$`=`$ $`{\displaystyle \frac{\delta }{\delta K(x)}}𝒮[\sigma (.)]|_{J_i(.)=K(.)=0}={\displaystyle \frac{2}{\lambda }}[\sigma (x)m^2]+𝒪\left({\displaystyle \frac{1}{N}}\right)`$ (53)
where we used eqs.(18), (33) and (44).
Once we set the external sources $`J_i(x)`$ and $`K(x)`$ equal to zero, the expectation value of $`\stackrel{}{\mathrm{\Phi }}(x)`$ vanishes as it must be due to the $`O(N)`$ invariance. On the contrary, $`N^1<\stackrel{}{\mathrm{\Phi }}^2(x)>`$ has in general a non-zero value at zero external sources.
It must be stressed that the present derivation applies to an arbitrary quantum state of the theory. Indeed, the equations simplify for the ground state due to translational invariance. In such case $`\sigma (x)`$ must be a constant $`\sigma _0`$ and the propagator $`G(x,y,\sigma _0)`$ takes the form,
$$G(x,y,\sigma _0)=\frac{d^4k}{(2\pi )^4}\frac{e^{ik.(xy)}}{\sigma _0k^2}.$$
The equal-points propagator $`G(x,x,\sigma _0)`$ needs an UV regulator. We obtain using a momentum cutoff $`\mathrm{\Lambda }`$ and Wick rotating $`k_0ik_0`$,
$`G_\mathrm{\Lambda }(x,x,\sigma _0)={\displaystyle \frac{2\pi ^2i}{(2\pi )^4}}{\displaystyle _0^\mathrm{\Lambda }}{\displaystyle \frac{k^3dk}{\sigma _0+k^2}}={\displaystyle \frac{i}{(4\pi )^2}}\left[\mathrm{\Lambda }^2+\sigma _0\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{\sigma _0^2}}+1\right)\right]`$ (54)
The $`\mathrm{\Lambda }`$ dependence can be absorbed into standard mass and coupling constant renormalization. Notice that $`\mu ^2\sigma _0`$ is the physical (renormalized) mass squared of the fundamental boson in the $`N=\mathrm{}`$ limit as we see from eq.(45). $`m^2m_B^2`$ is the bare boson mass and $`\lambda \lambda _B`$ the bare coupling constant.
We find from eq.(40) the relationship between bare and renormalized parameters,
$$\lambda _B=\frac{\lambda _r}{1\frac{\lambda _r}{16\pi ^2}\mathrm{log}\frac{\mathrm{\Lambda }}{\mu }},m_B^2=\frac{\lambda _B\mathrm{\Lambda }^2}{32\pi ^2}+\frac{\mu ^2}{1\frac{\lambda _r}{16\pi ^2}\mathrm{log}\frac{\mathrm{\Lambda }}{\mu }}$$
(55)
where $`\lambda _r`$ stands for the renormalized coupling constant and we dropped all contributions of order $`\mathrm{\Lambda }^2`$ and higher.
We can now choose for simplicity $`J_i(x)=J(x)`$ for $`i=1,\mathrm{},N`$. We proceed now to make a Legendre transformation such that the field expectation values $`\varphi (x)`$ and $`\xi (x)`$ become the independent functional variables. We define following eqs.(46) and (50),
$`\varphi (x)`$ $``$ $`{\displaystyle \frac{1}{iN}}{\displaystyle \frac{\delta }{\delta J(x)}}\mathrm{log}𝒵[J(.),K(.)]={\displaystyle }d^4yG(x,y,\sigma (.))J(y)+𝒪\left({\displaystyle \frac{1}{N}}\right)`$ (56)
$`\rho (x)`$ $``$ $`{\displaystyle \frac{1}{iN}}{\displaystyle \frac{\delta }{\delta K(x)}}\mathrm{log}𝒵[J(.),K(.)]={\displaystyle \frac{2}{\lambda }}[\sigma (x)m^2+2K(x)]+𝒪\left({\displaystyle \frac{1}{N}}\right)`$ (58)
In order to compute these functional derivatives, we used the chain rule,
$$\frac{\delta }{\delta J(x)}|_K=d^4y\frac{\delta \sigma (y)}{\delta J(x)}|_K\frac{\delta }{\delta \sigma (y)}|_{J,K}+\frac{\delta }{\delta J(x)}|_{\sigma ,K}.$$
The effective action functional is then given by,
$$\mathrm{\Gamma }[\varphi (.),\rho (.)]=\frac{1}{iN}\mathrm{log}𝒵[J(.),K(.)]d^4x[J(x)\varphi (x)+K(x)\rho (x)]$$
(59)
We get for this effective action in the infinite $`N`$ limit
$`\mathrm{\Gamma }[\varphi (.),\rho (.)]`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{log}det[^2+m^2+\xi (.)]+{\displaystyle \frac{1}{2\lambda }}{\displaystyle }d^4x[\xi (x)^24K(x)^2]`$ (60)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle }d^4xd^4yG(x,y,m^2+\xi (.))J(x)J(y)+𝒪\left({\displaystyle \frac{1}{N}}\right).`$ (62)
where we used eqs.(41),(44) and (56) and we introduced the new field variable:
$$\xi (x)\sigma (x)m^2=\frac{1}{2}\lambda \rho (x)2K(x).$$
(63)
Using now eqs.(26) and (56) we can easily express $`J(x)`$ in terms of $`\varphi (x)`$ and $`\xi (x)`$ as follows,
$$J(x)=\left[^2+m^2+\xi (x)\right]\varphi (x)+𝒪\left(\frac{1}{N}\right).$$
The effective action can then be written in terms of $`\varphi (x)`$ and $`\xi (x)`$ with the result,
$`\mathrm{\Gamma }[\varphi (.),\xi (.)]`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{log}det[^2+m^2+\xi (.)]+{\displaystyle \frac{1}{2\lambda }}{\displaystyle }d^4x[\xi (x)^24K(x)^2]`$ (64)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^4x\varphi (x)\left[^2+m^2+\xi (x)\right]\varphi (x)}+𝒪\left({\displaystyle \frac{1}{N}}\right).`$ (66)
The equations of motion on the fields $`\varphi (x)`$ and $`\xi (x)`$ follow extremizing the effective action $`\mathrm{\Gamma }[\varphi (.),\xi (.)]`$. We find from eq.(64) always for infinite $`N`$,
$`\left[^2+m^2+\xi (x)\right]\varphi (x)`$ $`=`$ $`0,`$ (67)
$`{\displaystyle \frac{i}{2}}G(x,x,m^2+\xi (.)){\displaystyle \frac{1}{2}}\varphi (x)^2+{\displaystyle \frac{1}{\lambda }}\xi (x)`$ $`=`$ $`0.`$ (69)
Or, in a more explicit form,
$`[^2+m^2+{\displaystyle \frac{\lambda }{2}}\varphi (x)^2{\displaystyle \frac{i\lambda }{2}}G(x,x,m^2+\xi (.))]\varphi (x)`$ $`=`$ $`0,`$ (70)
$`[^2+m^2+{\displaystyle \frac{\lambda }{2}}\varphi (x)^2{\displaystyle \frac{i\lambda }{2}}G(x,x,m^2+\xi (.))]G(x,y,m^2+\xi (.))`$ $`=`$ $`\delta (xy).`$ (72)
As we see, this is a non-linear and non-local set of partial differential equations.
Recall that $`\xi (x)`$ provides the expectation value of $`<\stackrel{}{\mathrm{\Phi }}^2(x)>`$ as derived in eqs.(50) and (63),
$$\xi (x)=\frac{\lambda }{2N}<\stackrel{}{\mathrm{\Phi }}^2(x)>+𝒪\left(\frac{1}{N}\right).$$
The field $`\varphi (x)`$ is the expectation value of $`<\stackrel{}{\mathrm{\Phi }}(x)_i>`$ for all values of $`i`$ as we see from eqs.(46) and (56) (up to $`𝒪\left(\frac{1}{N}\right)`$ corrections).
The ground state correspond to $`\xi (x)=\varphi (x)=0`$. This is the solution $`\sigma _0`$ discussed above. We are interested on solutions with non-zero $`\xi (x)`$ and $`\varphi (x)`$ describing excited states. States with non-zero $`\varphi (x)`$ will not be invariant under $`O(N)`$ transformations.
### B Invariance under spatial translations
Let us now consider states which are invariant under spatial translations. The fields $`\xi (x)`$ and $`\varphi (x)`$ will thus only be functions of time. This fact considerably simplifies the general equations (70). We can Fourier expand $`G(x,y,m^2+\xi (.))`$ as
$$G(x,y,m^2+\xi (.))=e^{i\stackrel{}{k}.(\stackrel{}{x}\stackrel{}{y})}g_k(t,t^{})\frac{d^3k}{(2\pi )^3},$$
where the times $`t`$ and $`t^{}`$ are associated with $`x`$ and $`y`$ respectively and $`g_k(t,t^{})`$ obeys the equation
$$\left[\frac{d^2}{dt^2}+k^2+m^2+\xi (t)\right]g_k(t,t^{})=\delta (tt^{}).$$
(73)
This one-dimensional Green function can be expressed in terms of solutions of the homogeneous equation
$$\left[\frac{d^2}{dt^2}+k^2+m^2+\xi (t)\right]f_k(t)=0.$$
(74)
One gets,
$$g_k(t,t^{})=\frac{f_k^<(t_<)f_k^>(t_>)}{W_k},$$
where $`t_<=`$ min$`(t,t^{})`$, $`t_>=`$ max$`(t,t^{})`$ and $`f_k^<(t)`$ and $`f_k^>(t)`$ are independent solutions of eq.(74). $`W_k`$ stands for the Wronskian between these solutions:
$$W_kW[f_k^<(t),f_k^>(t)].$$
(75)
For causal bondary conditions one has $`f_k^>(t)=[f_k^<(t)]^{}`$.
Before renormalization, eqs.(70) take the form
$`[{\displaystyle \frac{d^2}{dt^2}}+m^2+{\displaystyle \frac{\lambda }{2}}\varphi (t)^2{\displaystyle \frac{i\lambda }{2}}G(x,x,m^2+\xi (.))]\varphi (t)=0`$ (76)
(77)
$`\text{where}G(x,x,m^2+\xi (.))={\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^3}}{\displaystyle \frac{\left|f_k(t)\right|^2}{W_k}}.`$ (78)
and the mode functions $`f_k(t)`$ obey,
$$[\frac{d^2}{dt^2}+k^2+m^2+\frac{\lambda }{2}\varphi (t)^2\frac{i\lambda }{2}G(x,x,m^2+\xi (.))]f_k(t)=0.$$
(79)
We obtained an infinite set of coupled ordinary differential equations on $`\varphi (t),f_k(t)`$ for $`0k<\mathrm{}`$. Notice that eqs.(76)-(79) are local on time. That is they involve the unknown functions $`\varphi (t),f_k(t)`$ always at time $`t`$.
All physical quantities can be computed in terms of the mode functions $`f_k(t),\mathrm{\hspace{0.33em}0}k<\mathrm{}`$ and the order parameter $`\varphi (t)`$. We have described the large $`N`$ approximation. Other non-perturbative approximations are the Hartree approximations and the self-consistent one-loop approximation. They were considered in ref..
In a cosmological spacetime (4) the large $`N`$ evolution equations take the form
$`\left[{\displaystyle \frac{d^2}{dt^2}}+3H(t){\displaystyle \frac{d}{dt}}+\xi (t)+m^2+{\displaystyle \frac{\lambda }{2}}\varphi (t)^2+{\displaystyle \frac{\lambda }{2}}\pi ^2(t)\right]\varphi (t)`$ $`=`$ $`0`$ (80)
$`\left[{\displaystyle \frac{d^2}{dt^2}}+3H(t){\displaystyle \frac{d}{dt}}+{\displaystyle \frac{k^2}{a^2(t)}}+\xi (t)+m^2+{\displaystyle \frac{\lambda }{2}}\varphi (t)^2+{\displaystyle \frac{\lambda }{2}}\pi ^2(t)\right]f_k(t)`$ $`=`$ $`0,`$ (82)
where
$$H(t)\frac{\dot{a}(t)}{a(t)}$$
and we introduced the notation
$$\pi ^2(t)=iG(x,x,m^2+\xi (.))=\frac{d^3k}{(2\pi )^3}\frac{\left|f_k(t)\right|^2}{iW_k}.$$
(83)
### C The Large $`N`$ Limit in Cosmological Spacetimes
Let us connect the functional formalism in previous sections with the operatorial approach. Let us call ‘$`1`$’ the direction in internal space of the expectation value of $`\mathrm{\Phi }(x)`$. We write,
$$\stackrel{}{\mathrm{\Phi }}(\stackrel{}{x},t)=(\sigma (\stackrel{}{x},t),\stackrel{}{\pi }(\stackrel{}{x},t)),$$
with $`\stackrel{}{\pi }`$ an $`(N1)`$-plet, and
$$\sigma (\stackrel{}{x},t)=\sqrt{N}\varphi (t)+\chi (\stackrel{}{x},t);\sigma (\stackrel{}{x},t)=\sqrt{N}\varphi (t);\chi (\stackrel{}{x},t)=0.$$
(84)
We can now write the operator $`\stackrel{}{\pi }(\stackrel{}{x},t)`$ in terms of creation and annihilation operators and mode functions that obey the Heisenberg equations of motion
$$\stackrel{}{\pi }(\stackrel{}{x},t)=\frac{d^3k}{(2\pi )^3}\left[\stackrel{}{a}_kf_k(t)e^{i\stackrel{}{k}\stackrel{}{x}}+\stackrel{}{a}_k^{}f_k^{}(t)e^{i\stackrel{}{k}\stackrel{}{x}}\right].$$
(85)
where $`\stackrel{}{a}_k`$ and $`\stackrel{}{a}_k^{}`$ obey canonical commutation rules. It is easy to derive the two points function eq.(45) from eqs.(76), (75) and (85).
We see that since there are $`N1`$ ‘pion’ fields, contributions from the field $`\chi `$ can be neglected in the $`N\mathrm{}`$ limit as they are of order $`1/N`$ with respect those of $`\pi `$ and $`\varphi `$.
The equations of motion (80) can be written as
$`\ddot{\varphi }(t)+3H(t)\dot{\varphi }(t)+^2(t)\varphi (t)`$ $`=`$ $`0,`$ (86)
$`\left[{\displaystyle \frac{d^2}{dt^2}}+3H(t){\displaystyle \frac{d}{dt}}+{\displaystyle \frac{k^2}{a^2(t)}}+^2(t)\right]f_k(t)`$ $`=`$ $`0,`$ (88)
where
$$^2(t)=m^2+\xi (t)+\frac{\lambda }{2}\varphi ^2(t)+\frac{\lambda }{2}\pi ^2(t).$$
(89)
plays the role off time-dependent effective mass.
An important point to note in the large $`N`$ equations of motion is that the form of the equation for the zero mode (86) is the same as for the $`k=0`$ mode function (88). It is this property that allows solutions of these equations in a symmetry broken scenario to satisfy Goldstone’s theorem in out of equilibrium situations both in Minkowski and cosmological spacetimes .
In this leading order in $`1/N`$ the theory becomes Gaussian, but with the self-consistency condition (83).
The initial conditions on the modes $`f_k(t)`$ must now be determined. At this stage it proves illuminating to pass to conformal time variables in terms of the conformally rescaled fields (see and section V for a discussion) in which the mode functions obey an equation which is very similar to that of harmonic oscillators with time dependent frequencies in Minkowski space-time. It has been realized that different initial conditions on the mode functions lead to different renormalization counterterms; in particular imposing initial conditions in comoving time leads to counterterms that depend on these initial conditions. Thus we chose to impose initial conditions in conformal time in terms of the conformally rescaled mode functions leading to the following choice:
$$f_k^>(t)=f_k(t),f_k^<(t)=[f_k(t)]^{}$$
with initial conditions in comoving time,
$$f_k(t_0)=\frac{1}{\sqrt{\mathrm{\Omega }_k}},\dot{f}_k(t_0)=\left[\frac{\dot{a}(t_0)}{a(t_0)}i\mathrm{\Omega }_k\right]f_k(t_0),$$
(90)
with
$$\mathrm{\Omega }_k^2k^2+^2(t_0)\frac{(t_0)}{6}.$$
(91)
We thus find from eq.(75) that the Wronskian takes the value $`W_k=2i`$ and the quantum fluctuations of the inflaton (83) take the form
$$\pi ^2(t)=\frac{d^3k}{2(2\pi )^3}\left|f_k(t)\right|^2.$$
(92)
For convenience, we have set $`a(t_0)=1`$ in eq.(91). At this point we recognize that when $`^2(t_0)(t_0)/6<0`$ the above initial condition must be modified to avoid imaginary frequencies, which are the signal of instabilities for long wavelength modes in the broken symmetry case. Thus we define the initial frequencies that determine the initial conditions (90) as
$`\mathrm{\Omega }_k^2`$ $``$ $`k^2+\left|^2(t_0){\displaystyle \frac{(t_0)}{6}}\right|\text{ for }k^2<\left|^2(t_0){\displaystyle \frac{(t_0)}{6}}\right|,`$ (93)
$`\mathrm{\Omega }_k^2`$ $``$ $`k^2+^2(t_0){\displaystyle \frac{(t_0)}{6}}\text{ for }k^2\left|^2(t_0){\displaystyle \frac{(t_0)}{6}}\right|.`$ (94)
In the unbroken symmetry case ($`m^2>0`$ ) we use eq.(94) for all $`k`$.
As an alternative we have also used initial conditions which smoothly interpolate from positive frequencies for the unstable modes to the adiabatic vacuum initial conditions defined by (90)-(91) for the high $`k`$ modes. While the alternative choices of initial conditions result in small quantitative differences in the results (a few percent in quantities which depend strongly on these low-$`k`$ modes), all of the qualitative features we will examine are independent of this choice.
In the large $`N`$ limit we find the energy density and pressure density to be given by
$`{\displaystyle \frac{\epsilon }{N}}={\displaystyle \frac{1}{2}}\dot{\varphi }^2+{\displaystyle \frac{1}{2}}m^2\varphi ^2+{\displaystyle \frac{\lambda }{8}}\varphi ^4+{\displaystyle \frac{m^4}{2\lambda }}\xi G_0^0\varphi ^2+6\xi {\displaystyle \frac{\dot{a}}{a}}\varphi \dot{\varphi }+{\displaystyle \frac{1}{2}}\dot{\pi }^2`$ (95)
$`+{\displaystyle \frac{1}{2a^2}}(\pi )^2+{\displaystyle \frac{1}{2}}m^2\pi ^2+{\displaystyle \frac{\lambda }{8}}[2\varphi ^2\pi ^2+\pi ^2^2]\xi G_0^0\pi ^2+6\xi {\displaystyle \frac{\dot{a}}{a}}\pi \dot{\pi },`$ (96)
$`{\displaystyle \frac{\epsilon 3p}{N}}=\dot{\varphi }^2+2m^2\varphi ^2+{\displaystyle \frac{\lambda }{2}}\varphi ^4+{\displaystyle \frac{2m^4}{\lambda }}\xi G_\mu ^\mu \varphi ^2+6\xi \left(\varphi \ddot{\varphi }+\dot{\varphi }^2+3{\displaystyle \frac{\dot{a}}{a}}\varphi \dot{\varphi }\right)\dot{\pi }^2`$ (97)
$`+{\displaystyle \frac{1}{a^2}}(\pi )^2+2m^2\pi ^2\xi G_\mu ^\mu \pi ^2+{\displaystyle \frac{\lambda }{2}}[2\varphi ^2\pi ^2+\pi ^2^2]+6\xi \left(\pi \ddot{\pi }+\dot{\pi }^2+3{\displaystyle \frac{\dot{a}}{a}}\pi \dot{\pi }\right),`$ (98)
where $`\pi ^2`$ is given by equation (92) and we have defined the following integrals:
$`(\pi )^2={\displaystyle }{\displaystyle \frac{d^3k}{2(2\pi )^3}}k^2|f_k(t)|^2,\dot{\pi }^2={\displaystyle }{\displaystyle \frac{d^3k}{2(2\pi )^3}}|\dot{f}_k(t)|^2.`$ (99)
The composite operators $`\pi \dot{\pi }`$ and $`\pi \ddot{\pi }`$ are symmetrized by removing a normal ordering constant. $`\pi \ddot{\pi }`$ may be rewritten using the equation of motion (88):
$$\pi \ddot{\pi }=3\frac{\dot{a}}{a}\pi \dot{\pi }\frac{(\pi )^2}{a^2}^2(t)\pi ^2.$$
(100)
It is straightforward to show that the bare energy is covariantly conserved by using the equations of motion for the zero mode and the mode functions.
## V Renormalization, Conformal time and Initial Conditions
Renormalization is a very subtle but important issue in gravitational backgrounds. The fluctuation contribution $`\pi ^2(\stackrel{}{x},t)`$, the energy, and the pressure all need to be renormalized. The renormalization aspects in curved space times have been discussed at length in the literature and have been extended to the large $`N`$ self-consistent approximations for the non-equilibrium backreaction problem in. More recently, a consistent and covariant regularization scheme that can be implemented numerically has been proposed.
The ultraviolet divergences can be seen in the present framework in the $`k`$-integrals over the modes in eqs.(92) and (99). To analyze these divergences it is convenient to change variables to conformal time $`𝒯`$ defined as
$$𝒯=𝒯_0+_{t_0}^t\frac{dt^{}}{a(t^{})};𝒯(t=t_0)=𝒯_0,$$
(101)
The metric becomes then
$$ds^2=C^2(𝒯)(d𝒯^2d\stackrel{}{x}^2),$$
(102)
where $`C(𝒯)a(t(𝒯));C(𝒯_0)=1`$ stands for the scale factor in conformal time.
The issue of renormalization and initial conditions is best analyzed in conformal time which is a natural framework for adiabatic renormalization and regularization.
Under a conformal rescaling of the field
$$\mathrm{\Phi }(\stackrel{}{x},t)=\frac{\chi (\stackrel{}{x},𝒯)}{C(𝒯)},$$
(103)
the action for a scalar field (with the obvious generalization to $`N`$ components) becomes, after an integration by parts and dropping a surface term
$$S=d^3x𝑑𝒯\left[\frac{1}{2}(\chi ^{})^2\frac{1}{2}(\stackrel{}{}\chi )^2𝒱(\chi )\right],$$
(104)
with
$$𝒱(\chi )=C^4(𝒯)V\left(\frac{\chi }{C(𝒯)}\right)C^2(𝒯)\frac{}{12}\chi ^2,$$
(105)
where $`=6C^{\prime \prime }(𝒯)/C^3(𝒯)`$ is the Ricci scalar, and primes stand for derivatives with respect to conformal time $`𝒯`$.
The conformal time Hamiltonian operator, which is the generator of translations in $`𝒯`$, is given by
$$H_𝒯=d^3x\left\{\frac{1}{2}\mathrm{\Pi }_\chi ^2+\frac{1}{2}(\stackrel{}{}\chi )^2+𝒱(\chi )\right\},$$
(106)
with $`\mathrm{\Pi }_\chi `$ being the canonical momentum conjugate to $`\chi `$, $`\mathrm{\Pi }_\chi =\chi ^{}`$. Separating the zero mode of the field $`\chi `$
$$\chi (\stackrel{}{x},𝒯)=\chi _0(𝒯)+\overline{\chi }(\stackrel{}{x},𝒯),$$
(107)
and in the large $`N`$ approximation we find that the Hamiltonian becomes linear plus quadratic in the fluctuations, and similar to a Minkowski space-time Hamiltonian with a $`𝒯`$-dependent effective mass term given by
$$B(𝒯)=C^2(𝒯)\left[m^2+(\xi \frac{1}{6})+\frac{\lambda }{2}\chi _0^2(𝒯)+\frac{\lambda }{2}\overline{\chi }^2\right].$$
(108)
Notice that this $`B(t)`$ naturally appears in the WKB expansion of the mode functions both in cosmic and conformal time \[see eq.(127)\].
We can now follow the steps and use the results of reference for the conformal time evolution of the density matrix by setting $`a(t)=1`$ in the proper equations of that reference and replacing the frequencies by
$$\omega _k^2(𝒯)=\stackrel{}{k}^2+B(𝒯),$$
(109)
and the expectation value in eq.(108) is obtained in this $`𝒯`$ evolved density matrix. The time evolution of the kernels in the density matrix (see) is determined by the mode functions that obey
$$\left[\frac{d^2}{d𝒯^2}+k^2+B(𝒯)\right]F_k(𝒯)=0.$$
(110)
The Wronskian of these mode functions
$$𝒲(F,F^{})=F_k^{}F_k^{}F_kF_k^{^{}}$$
(111)
is a constant. It is natural to impose initial conditions such that at the initial $`𝒯`$ the density matrix describes a situation of local thermodynamic equilibrium and therefore commutes with the conformal time Hamiltonian at the initial time. This implies that the initial conditions of the mode functions $`F_k(𝒯)`$ be chosen to be (see)
$$F_k(𝒯_o)=\frac{1}{\sqrt{\omega _k(𝒯_o)}};F_k^{}(𝒯_o)=i\sqrt{\omega _k(𝒯_o)}.$$
(112)
With such initial conditions, the Wronskian (111) takes the value
$$𝒲(F,F^{})=2i.$$
(113)
These initial conditions correspond to the choice of mode functions which coincide with the first order adiabatic modes and those of the Bunch-Davies vacuum for large momentum. To see this clearly, we write the solution of eq.(110) in the form,
$$D_k(𝒯)=e^{_{𝒯_o}^𝒯R_k(𝒯^{})𝑑𝒯^{}},$$
(114)
with the function $`R_k(𝒯)`$ obeying the Riccati equation
$$R_k^{}+R_k^2+k^2+B(𝒯)=0.$$
(115)
This equation posses the solution
$$R_k(𝒯)=ik+R_{0,k}(𝒯)i\frac{R_{1,k}(𝒯)}{k}+\frac{R_{2,k}(𝒯)}{k^2}i\frac{R_{3,k}(𝒯)}{k^3}+\frac{R_{4,k}(𝒯)}{k^4}+𝒪\left(\frac{1}{k^5}\right)$$
(116)
and its complex conjugate. We find for the coefficients:
$`R_{0,k}=0;R_{1,k}={\displaystyle \frac{1}{2}}B(𝒯);R_{2,k}={\displaystyle \frac{1}{2}}R_{1,k}^{}`$ (117)
$`R_{3,k}={\displaystyle \frac{1}{2}}\left(R_{2,k}^{}R_{1,k}^2\right);R_{4,k}={\displaystyle \frac{1}{2}}\left(R_{3,k}^{}+2R_{1,k}R_{2,k}\right).`$ (118)
The solutions $`F_k(𝒯)`$ obeying the boundary conditions (112) are obtained as linear combinations of this WKB solution and its complex conjugate
$$F_k(𝒯)=\frac{1}{2\sqrt{\omega _k(𝒯_o)}}\left[(1+\gamma )D_k(𝒯)+(1\gamma )D_k^{}(𝒯)\right],$$
(119)
where the coefficient $`\gamma `$ is obtained from the initial conditions. It is straightforward to find that the real and imaginary parts are given by
$$\gamma _R=1+𝒪(1/k^4);\gamma _I=𝒪(1/k^3).$$
(120)
Therefore the large-$`k`$ mode functions satisfy the adiabatic vacuum initial conditions. This, in fact, is the rationale for the choice of the initial conditions (112).
Following the analysis presented in we find, in conformal time that
$$\overline{\chi }^2(\stackrel{}{x},𝒯)=\frac{d^3k}{2(2\pi )^3}|F_k(𝒯)|^2.$$
(121)
The Heisenberg field operators $`\overline{\chi }(\stackrel{}{x},𝒯)`$ and their canonical momenta $`\mathrm{\Pi }_\chi (\stackrel{}{x},𝒯)`$ can be expanded as:
$`\overline{\chi }(\stackrel{}{x},𝒯)={\displaystyle \frac{d^3k}{\sqrt{2}(2\pi )^{3/2}}\left[a_kF_k(𝒯)+a_k^{}F_k^{}(𝒯)\right]e^{i\stackrel{}{k}\stackrel{}{x}}},`$ (122)
$`\mathrm{\Pi }_\chi (\stackrel{}{x},𝒯)={\displaystyle \frac{d^3k}{\sqrt{2}(2\pi )^{3/2}}\left[a_kF_k^{}(𝒯)+a_k^{}F_k^{^{}}(𝒯)\right]e^{i\stackrel{}{k}\stackrel{}{x}}},`$ (123)
with the time independent creation and annihilation operators $`a_k`$ and $`a_k^{}`$ obeying canonical commutation relations. Since the fluctuation fields in comoving and conformal time are related by the conformal rescaling (103), it is straightforward to see that the mode functions in comoving time are related to those in conformal time simply as
$$f_k(t)=\frac{F_k(𝒯)}{C(𝒯)}.$$
(124)
Therefore the initial conditions (112) on the conformal time mode functions imply the initial conditions for the mode functions in comoving time are given by eq.(90).
For renormalization purposes we need the large-$`k`$ behavior of $`|f_k(t)|^2,|\dot{f}_k(t)|^2`$, which are determined by the large-$`k`$ behavior of the conformal time mode functions and its derivative. These are given by appropriately adapting the Minkowski formulas
$`|F_k(𝒯)|^2`$ $`=`$ $`{\displaystyle \frac{1}{k}}+{\displaystyle \frac{B(𝒯)}{2k^2}}+{\displaystyle \frac{1}{8k^5}}\left[3B^2(𝒯)+B^{\prime \prime }(𝒯)\right]+𝒪\left({\displaystyle \frac{1}{k^7}}\right),`$ (125)
$`|F_k^{}(𝒯)|^2`$ $`=`$ $`k+{\displaystyle \frac{B(𝒯)}{2k}}{\displaystyle \frac{1}{8k^3}}\left[B^2(𝒯)+B^{\prime \prime }(𝒯)\right]+𝒪\left({\displaystyle \frac{1}{k^5}}\right).`$ (127)
We note that the large $`k`$ behavior of the mode functions to the order needed to renormalize the quadratic and logarithmic divergences is insensitive to the initial conditions. This is not the case when the initial conditions are imposed as described in. Thus the merit in considering the initial conditions in conformal time .
There is an important physical consequence of the choice (112) of initial conditions, which is revealed by analyzing the evolution of the density matrix.
In the large $`N`$ or Hartree (also to one-loop) approximation, the density matrix is Gaussian, and defined by a normalization factor, a complex covariance that determines the diagonal matrix elements and a real covariance that determines the mixing in the Schrödinger representation as discussed in reference (and references therein).
In conformal time quantization and in the Schrödinger representation in which the field $`\chi `$ is diagonal the conformal time evolution of the density matrix is via the conformal time Hamiltonian (106). The evolution equations for the covariances is obtained from those given in reference by setting $`a(t)=1`$ and using the frequencies $`\omega _k^2(𝒯)=k^2+^2(𝒯)`$. In particular, by setting the covariance of the diagonal elements (given by equation (2.20) in; see also equation (2.44) of),
$$𝒜_k(𝒯)=i\frac{F_k^{{}_{}{}^{}}(𝒯)}{F_k^{}(𝒯)},$$
(128)
we find that with the initial conditions (112), the conformal time density matrix is that of local equilibrium at $`𝒯_0`$ in the sense that it commutes with the conformal time Hamiltonian. However, it is straightforward to see, that the comoving time density matrix does not commute with the comoving time Hamiltonian at the initial time $`t_0`$.
An important corollary of this analysis and comparison with other initial conditions used in comoving time is that assuming initial conditions of local equilibrium in comoving time leads to divergences that depend on the initial condition as discussed at length in. This dependence of the renormalization counterterms on the initial condition was also realized in ref. within the context of the CTP formulation. Imposing the initial conditions corresponding to local thermal equilibrium in conformal time, we see that: i) the renormalization counterterms do not depend on the initial conditions and ii) the mode functions are identified with those corresponding to the adiabatic vacuum for large momenta. This is why we prefer the initial conditions (112).
For our main analysis we choose this initial temperature to be zero so that the resulting density matrix describes a pure state, which for the large momentum modes coincides with the conformal adiabatic vacuum. Such zero temperature choice seems appropriate after the exponential inflation of the universe.
Particle Number:
We write the Fourier components of the field $`\chi `$ and its canonical momentum $`\mathrm{\Pi }_\chi `$ given by (122) -(123) as:
$`\overline{\chi }_k(𝒯)={\displaystyle \frac{1}{\sqrt{2}}}\left[a_kF_k(𝒯)+a_k^{}F_k^{}(𝒯)\right],`$ (129)
$`\mathrm{\Pi }_{\chi ,k}(𝒯)={\displaystyle \frac{1}{\sqrt{2}}}\left[a_kF_k^{}(𝒯)+a_k^{}F_k^{^{}}(𝒯)\right].`$ (130)
These (conformal time) Heisenberg operators can be written equivalently in terms of the $`𝒯`$ dependent creation and annihilation operators
$`\overline{\chi }_k(𝒯)={\displaystyle \frac{1}{\sqrt{2\omega _k(𝒯_0)}}}\left[\stackrel{~}{a}_k(𝒯)e^{i\omega _k(𝒯_0)𝒯}+\stackrel{~}{a}_k^{}(𝒯)e^{i\omega _k(𝒯_0)𝒯}\right],`$ (131)
$`\mathrm{\Pi }_{\chi ,k}(𝒯)=i\sqrt{{\displaystyle \frac{\omega _k(𝒯_0)}{2}}}\left[\stackrel{~}{a}_k(𝒯)e^{i\omega _k(𝒯_0)𝒯}\stackrel{~}{a}_k^{}(𝒯)e^{i\omega _k(𝒯_0)𝒯}\right].`$ (132)
The operators $`\stackrel{~}{a}_k(𝒯);a_k(𝒯)`$ are related by a Bogoliubov transformation. The number of particles referred to the initial Fock vacuum of the modes $`F_k`$, is given by
$$N_k(𝒯)=\stackrel{~}{a}_k^{}(𝒯)\stackrel{~}{a}_k(𝒯)=\frac{1}{4}\left[\left|\frac{F_k(𝒯)}{F_k(𝒯_0)}\right|^2+\frac{1}{\omega _k^2(𝒯_0)}\left|\frac{F_k^{}(𝒯)}{F_k(𝒯_0)}\right|^2\right]\frac{1}{2},$$
(133)
or alternatively, in terms of the comoving mode functions $`f_k(t)=F_k(𝒯)/C(𝒯)`$ we find
$$N_k(t)=\frac{a^2(t)}{4}\left[\left|\frac{f_k(t)}{f_k(t_0)}\right|^2+\frac{1}{\omega _k^2(t_0)}\left|\frac{\dot{f}_k(t)+H(t)f_k(t)}{f_k(t_0)}\right|^2\right]\frac{1}{2}.$$
(134)
Using the large $`k`$-expansion of the conformal mode functions given by eqs. (127) we find the large-$`k`$ behavior of the particle number to be $`N_k\stackrel{k\mathrm{}}{=}𝒪(1/k^4)`$, and the total number of particles (with reference to the initial state at $`𝒯_0`$) is therefore finite.
Renormalization.
We make our subtractions using an ultraviolet cutoff, $`\mathrm{\Lambda }a(t)`$, constant in physical coordinates. This guarantees that the counterterms will be time independent. The renormalization then proceeds much in the same manner as in reference; the quadratic divergences renormalize the mass and the logarithmic terms renormalize the quartic coupling and the coupling to the Ricci scalar. In addition, there is a quartic divergence which renormalizes the cosmological constant while the leading renormalizations of Newton’s constant and the higher order curvature coupling are quadratic and logarithmic respectively. The renormalization conditions on the mass, coupling to the Ricci scalar and coupling constant are obtained from the requirement that the frequencies that appear in the mode equations are finite, i.e:
$$m_B^2+\xi _B(t)+\frac{\lambda _B}{2}\varphi ^2(t)+\frac{\lambda _B}{2}\pi ^2(t)_B=m_R^2+\xi _R(t)+\frac{\lambda _R}{2}\varphi ^2(t)+\frac{\lambda _R}{2}\pi ^2(t)_R,$$
(135)
while the renormalizations of Newton’s constant, the higher order curvature coupling, and the cosmological constant are given by the condition of finiteness of the semi-classical Einstein-Friedmann equation:
$$\frac{G_0^0}{8\pi G_B}+\alpha _BH_0^0+K_Bg_0^0+T_0^0_B=\frac{G_0^0}{8\pi G_R}+\alpha _RH_0^0+K_Rg_0^0+T_0^0_R.$$
(136)
Finally, we arrive at the following set of renormalizations:
$`{\displaystyle \frac{1}{8\pi NG_R}}={\displaystyle \frac{1}{8\pi NG_B}}{\displaystyle \frac{2}{(4\pi )^2}}\left(\xi _R{\displaystyle \frac{1}{6}}\right)\left[\mathrm{\Lambda }^2m_R^2\mathrm{ln}(\mathrm{\Lambda }/\kappa )\right],\lambda _R=\lambda _B\lambda _R{\displaystyle \frac{\mathrm{ln}(\mathrm{\Lambda }/\kappa )}{(4\pi )^2}},`$ (137)
(138)
$`{\displaystyle \frac{\alpha _R}{N}}={\displaystyle \frac{\alpha _B}{N}}(\xi _R{\displaystyle \frac{1}{6}})^2{\displaystyle \frac{\mathrm{ln}(\mathrm{\Lambda }/\kappa )}{(4\pi )^2}},{\displaystyle \frac{K_R}{N}}={\displaystyle \frac{K_B}{N}}{\displaystyle \frac{\mathrm{\Lambda }^4}{(4\pi )^2}}m_R^2{\displaystyle \frac{\mathrm{\Lambda }^2}{(4\pi )^2}}+{\displaystyle \frac{m_R^4}{2}}{\displaystyle \frac{\mathrm{ln}(\mathrm{\Lambda }/\kappa )}{(4\pi )^2}},`$ (139)
(140)
$`m_R^2=m_B^2+\lambda _R{\displaystyle \frac{\mathrm{\Lambda }^2}{(4\pi )^2}}\lambda _Rm_R^2{\displaystyle \frac{\mathrm{ln}(\mathrm{\Lambda }/\kappa )}{(4\pi )^2}},\xi _R=\xi _B\lambda _R(\xi _R{\displaystyle \frac{1}{6}}){\displaystyle \frac{\mathrm{ln}(\mathrm{\Lambda }/\kappa )}{(4\pi )^2}},`$ (141)
(142)
$`\pi ^2(t)_R={\displaystyle \frac{d^3k}{2(2\pi )^3}\left\{|f_k(t)|^2\frac{1}{ka^2(t)}+\frac{\mathrm{\Theta }(k\kappa )}{2k^3}\left[^2(t)\frac{(t)}{6}\right]\right\}}.`$ (143)
Here, $`\kappa `$ is the renormalization point. As expected, the logarithmic terms are consistent with the renormalizations found using dimensional regularization. Again, we set $`\alpha _R=0`$ and choose the renormalized cosmological constant such that the vacuum energy is zero in the true vacuum. We emphasize that while the regulator we have chosen does not respect the covariance of the theory, the renormalized energy momentum tensor defined in this way nevertheless retains the property of covariant conservation in the limit when the cutoff is taken to infinity.
The logarithmic subtractions can be neglected because of the coupling $`\lambda 10^{12}`$. Using the Planck scale as the cutoff and the inflaton mass $`m_R`$ as a renormalization point, these terms are of order $`\lambda \mathrm{ln}[M_{pl}/m_R]10^{10}`$, for $`M_{Pl}>m10^9\text{ GeV }`$. An equivalent statement is that for these values of the coupling and inflaton masses, the Landau pole is well beyond the physical cutoff $`M_{pl}`$. Our relative error in the numerical analysis is of order $`10^8`$, therefore our numerical study is insensitive to the logarithmic corrections. Though these corrections are fundamentally important, numerically they can be neglected. Therefore, in the numerical computations that follow, we will neglect logarithmic renormalization and subtract only quartic and quadratic divergences in the energy and pressure, and quadratic divergences in the fluctuation contribution.
### A Renormalized Equations of Motion for Dynamical Evolution in the Large $`N`$ limit
It is convenient to introduce the following dimensionless quantities and definitions,
$$\tau =m_Rt;h(\tau )=\frac{H(t)}{m_R};q=\frac{k}{m_R}\omega _q=\frac{\mathrm{\Omega }_k}{m_R};g=\frac{\lambda _R}{8\pi ^2},$$
(144)
$$\eta ^2(\tau )=\frac{\lambda _R}{2m_R^2}\varphi ^2(t);g\mathrm{\Sigma }(\tau )=\frac{\lambda }{2m_R^2}\pi ^2(t)_R;f_q(\tau )\sqrt{m_R}f_k(t).$$
(145)
Choosing $`\xi _R=0`$ (minimal coupling) and the renormalization point $`\kappa =|m_R|`$ and setting $`a(\tau _0)=1`$, the equations of motion become for unbroken symmetry:
$`\left[{\displaystyle \frac{d^2}{d\tau ^2}}+3h(\tau ){\displaystyle \frac{d}{d\tau }}+1+\eta ^2(\tau )+g\mathrm{\Sigma }(\tau )\right]\eta (\tau )=0,`$ (146)
(147)
$`\left[{\displaystyle \frac{d^2}{d\tau ^2}}+3h(\tau ){\displaystyle \frac{d}{d\tau }}+{\displaystyle \frac{q^2}{a^2(\tau )}}+1+\eta ^2+g\mathrm{\Sigma }(\tau )\right]f_q(\tau )=0,`$ (148)
$`f_q(\tau _0)={\displaystyle \frac{1}{\sqrt{\omega _q}}};\dot{f}_q(\tau _0)=[h(\tau _0)i\omega _q]f_q(\tau _0),\omega _q=[q^2+1+\eta ^2(\tau _0){\displaystyle \frac{(\tau _0)}{6m_R^2}}+g\mathrm{\Sigma }(\tau _0)]^{\frac{1}{2}}.`$ (149)
We find for broken symmetry,
$`\left[{\displaystyle \frac{d^2}{d\tau ^2}}+3h(\tau ){\displaystyle \frac{d}{d\tau }}1+\eta ^2(\tau )+g\mathrm{\Sigma }(\tau )\right]\eta (\tau )=0,`$ (150)
(151)
$`\left[{\displaystyle \frac{d^2}{d\tau ^2}}+3h(\tau ){\displaystyle \frac{d}{d\tau }}+{\displaystyle \frac{q^2}{a^2(\tau )}}1+\eta ^2+g\mathrm{\Sigma }(\tau )\right]f_q(\tau )=0,`$ (152)
$`f_q(\tau _0)={\displaystyle \frac{1}{\sqrt{\omega _q}}};\dot{f}_q(\tau _0)=[h(\tau _0)i\omega _q]f_q(\tau _0),`$ (153)
$`\omega _q=\left[q^21+\eta ^2(\tau _0){\displaystyle \frac{(\tau _0)}{6m_R^2}}+g\mathrm{\Sigma }(\tau _0)\right]^{\frac{1}{2}}\text{ for }q^2>1\eta ^2(\tau _0)+{\displaystyle \frac{(\tau _0)}{6m_R^2}}g\mathrm{\Sigma }(\tau _0),`$ (154)
$`\omega _q=\left[q^2+1\eta ^2(\tau _0)+{\displaystyle \frac{(\tau _0)}{6m_R^2}}g\mathrm{\Sigma }(\tau _0)\right]^{\frac{1}{2}}\text{ for }q^2<1\eta ^2(\tau _0)+{\displaystyle \frac{(\tau _0)}{6m_R^2}}g\mathrm{\Sigma }(\tau _0).`$ (155)
Here,
$$\mathrm{\Sigma }(\tau )=_0^{\mathrm{}}q^2𝑑q\left[|f_q(\tau )|^2\frac{1}{qa(\tau )^2}+\frac{\mathrm{\Theta }(q1)}{2q^3}\left(\frac{^2(\tau )}{m_R^2}\frac{(\tau )}{6m_R^2}\right)\right].$$
(156)
both for unbroken and broken symmetry.
The initial conditions for $`\eta (\tau )`$ will be specified later. An important point to notice is that the equation of motion for the $`q=0`$ mode coincides with that of the zero mode (150). Furthermore, for $`\eta (\tau \mathrm{})0`$, a stationary (equilibrium) solution of the eq.(150) is obtained for broken symmetry when the sum rule
$$1+\eta ^2(\mathrm{})+g\mathrm{\Sigma }(\mathrm{})=0$$
(157)
is fulfilled.
This sum rule is nothing but a proof that Goldstone’s theorem holds here out of thermal equilibrium. In addition, it is a result of the fact that the large $`N`$ approximation satisfies the Ward identities associated with the $`O(N)`$ symmetry, since the term $`1+\eta ^2(\tau )+g\mathrm{\Sigma }(\tau )`$ is seen to be the effective mass of the modes transverse to the symmetry breaking direction, i.e. the Goldstone modes in the broken symmetry phase.
The renormalized dimensionless evolution equations in the Hartree approximation are very similar to eqs.(146)-(155). They can be obtained just dividing by three the $`\eta ^2(\tau )`$ term in the zero mode equation .
In terms of the zero mode $`\eta (\tau )`$ and the quantum mode function given by eq.(150) we find that the Friedmann equation for the dynamics of the scale factor in dimensionless variables is given by
$$h^2(\tau )=4h_0^2ϵ_R(\tau );h_0^2=\frac{4\pi Nm_R^2}{3M_{Pl}^2\lambda _R}.$$
(158)
and the renormalized energy and pressure are given by:
$`ϵ_R(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\eta }^2+{\displaystyle \frac{1}{4}}\left(1+\eta ^2+g\mathrm{\Sigma }\right)^2`$ (159)
$`+`$ $`{\displaystyle \frac{g}{2}}{\displaystyle q^2𝑑q\left[|\dot{f_q}|^2𝒮^{(1)}(q,\tau )+\frac{q^2}{a^2}\left(|f_q|^2\mathrm{\Theta }(q1)𝒮^{(2)}(q,\tau )\right)\right]},`$ (160)
$`(p+\epsilon )_R(\tau )`$ $`=`$ $`\dot{\eta }^2+g{\displaystyle q^2𝑑q\left[|\dot{f_q}|^2𝒮^{(1)}(q,\tau )+\frac{q^2}{3a^2}\left(|f_q|^2\mathrm{\Theta }(q1)𝒮^{(2)}(q,\tau )\right)\right]},`$ (161)
where the subtractions $`𝒮^{(1)}`$ and $`𝒮^{(2)}`$ are given by
$`𝒮^{(1)}`$ $`=`$ $`{\displaystyle \frac{k}{a^4(t)}}+{\displaystyle \frac{1}{2ka^4(t)}}\left[B(t)+2\dot{a}^2\right]`$ (162)
$`+`$ $`{\displaystyle \frac{1}{8k^3a^4(t)}}\left[B(t)^2a(t)^2\ddot{B}(t)+3a(t)\dot{a}(t)\dot{B}(t)4\dot{a}^2(t)B(t)\right]`$ (164)
$`𝒮^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{ka^2(t)}}{\displaystyle \frac{1}{2k^3a^2(t)}}B(t)+{\displaystyle \frac{1}{8k^5a^2(t)}}\left\{3B(t)^2+a(t){\displaystyle \frac{d}{dt}}\left[a(t)\dot{B}(t)\right]\right\}`$ (166)
The renormalized energy and pressure are covariantly conserved:
$$\dot{ϵ}_R(\tau )+3h(\tau )(p+\epsilon )_R(\tau )=0.$$
(167)
From the evolution of the mode functions that determine the quantum fluctuations, we can study the growth of correlated domains with the equal time correlation function,
$`S(\stackrel{}{x},t)`$ $`=`$ $`\pi (\stackrel{}{x},t)\pi (\stackrel{}{0},t)={\displaystyle \frac{d^3k}{2(2\pi )^3}e^{i\stackrel{}{k}\stackrel{}{x}}|f_k(t)|^2},`$ (168)
which can be written in terms of the power spectrum of quantum fluctuations, $`|f_q(\tau )|^2`$. It is convenient to define the dimensionless correlation function,
$$𝒮(\rho ,\tau )=\frac{S(|\stackrel{}{x}|,t)}{m_R^2}=\frac{1}{4\pi ^2\rho }_0^{\mathrm{}}q𝑑q\mathrm{sin}[q\rho ]|f_q(\tau )|^2;\rho =m_R|\stackrel{}{x}|.$$
(169)
We now have all the ingredients to study the particular cases of interest.
## VI Scalar Field Dynamics in a fixed FRW background
We consider in this section the evolution of scalar fields in radiation or matter dominated FRW cosmologies. The case for de Sitter expansion will be discussed in sec. IX .
We write the scale factor as $`a(\tau )=(\tau /\tau _0)^n`$ with $`n=1/2`$ and $`n=2/3`$ corresponding to radiation and matter dominated backgrounds, respectively. The value of $`\tau _0`$ determines the initial Hubble constant since
$$h(\tau _0)=\frac{\dot{a}(\tau _0)}{a(\tau _0)}=\frac{n}{\tau _0}.$$
We now solve the system of equations (80) in the large $`N`$ limit. We begin by presenting an early time analysis of the slow roll scenario. We then undertake a thorough numerical investigation of various cases of interest. For the symmetry broken case, we also provide an investigation of the late time behavior of the zero mode and the quantum fluctuations. We use the dimensionless variables (144)-(145).
We will assume minimal coupling to the curvature, $`\xi _r=0`$. In the cases of interest, $`\mu ^2`$, so that finite $`\xi _r`$ has little effect.
### A Early Time Solutions for Slow Roll
For early times in a slow roll scenario \[$`m^2=\mu ^2`$, $`\eta (\tau _0)1`$\], we can neglect in eqs.(155) both the quadratic and cubic terms in $`\eta (\tau )`$ as well as the quantum fluctuations $`\pi ^2(\tau )_r`$ \[recall that $`\pi ^2(\tau _0)_r=0`$\]. Thus, the differential equations for the zero mode and the mode functions (80) become linear equations. In terms of the scaled variables introduced above, with $`a(\tau )=\tau ^n`$ ($`n=2/3`$ for a matter dominated cosmology while $`n=1/2`$ for a radiation dominated cosmology) we have:
$`\ddot{\eta }(\tau )+{\displaystyle \frac{3n}{\tau }}\dot{\eta }(\tau )\eta (\tau )`$ $`=`$ $`0,`$ (170)
$`\left[{\displaystyle \frac{d^2}{d\tau ^2}}+{\displaystyle \frac{3n}{\tau }}{\displaystyle \frac{d}{d\tau }}+k^2\left({\displaystyle \frac{\tau _0}{\tau }}\right)^{2n}1\right]U_k(\tau )`$ $`=`$ $`0.`$ (172)
The solutions to the zero mode equation (170) are
$$\eta (\tau )=c\tau ^\nu I_\nu (\tau )+d\tau ^\nu K_\nu (\tau ),$$
(173)
where $`\nu (3n1)/2`$, and $`I_\nu (\tau )`$ and $`K_\nu (\tau )`$ are modified Bessel functions. The coefficients, $`c`$ and $`d`$, are determined by the initial conditions on $`\eta `$. For $`\eta (\tau _0)=\eta _0`$ and $`\dot{\eta }(\tau _0)=0`$, we have:
$`c=\eta _0\tau _0^{\nu +1}[\dot{K}_\nu (\tau _0){\displaystyle \frac{\nu }{\tau _0}}K_\nu (\tau _0)],d=\eta _0\tau _0^{\nu +1}[\dot{I}_\nu (\tau _0){\displaystyle \frac{\nu }{\tau _0}}I_\nu (\tau _0)].`$ (174)
Taking the asymptotic forms of the modified Bessel functions, we find that for intermediate times $`\eta (\tau )`$ grows as
$$\eta (\tau )\stackrel{\tau 1}{=}\frac{c}{\sqrt{2\pi }}\tau ^{3n/2}e^\tau \left[1\frac{9n^26n}{8\tau }+𝒪(\frac{1}{\tau ^2})\right].$$
(175)
We see that $`\eta (\tau )`$ grows very quickly in time, and the approximations (170) and (172) will quickly break down. For the case shown in fig.1 (with $`n=2/3`$, $`\eta (\tau _0)=10^7`$, and $`\dot{\eta }(\tau _0)=0`$), we find that this approximation is valid up to $`\tau \tau _010`$.
The equations for the mode functions (172) can be solved in closed form for the modes in the case of a radiation dominated cosmology with $`n=1/2`$. The solutions are
$$U_k(\tau )=c_ke^\tau U(\frac{3}{4}\frac{k^2}{2},\frac{3}{2},2\tau )+d_ke^\tau M(\frac{3}{4}\frac{k^2}{2},\frac{3}{2},2\tau ).$$
(176)
Here, $`U()`$ and $`M()`$ are confluent hypergeometric functions (in another common notation, $`M()_1F_1()`$), and the $`c_k`$ and $`d_k`$ are coefficients determined by the initial conditions (90) on the modes. The solutions can also be written in terms of parabolic cylinder functions.
For large $`\tau `$ we have the asymptotic form
$$U_k(\tau )\stackrel{\tau 1}{=}d_ke^\tau (2\tau )^{(3/4+k^2\tau _0/2)}\frac{\sqrt{\pi }}{2\mathrm{\Gamma }\left(\frac{3}{4}\frac{k^2\tau _0}{2}\right)}\left[1+𝒪(\frac{1}{\tau })\right]+c_ke^\tau (2\tau )^{(3/4+k^2\tau _0/2)}\left[1+𝒪(\frac{1}{\tau })\right].$$
(177)
Again, these expressions only apply for intermediate times before the nonlinearities have grown significantly.
### B Numerical Analysis
We now present the numerical analysis of the dynamical evolution of scalar fields in time dependent, matter and radiation dominated cosmological backgrounds. We use initial values of the Hubble constant such that $`h(\tau _0)0.1`$. For expansion rates much less than this value the evolution will look similar to Minkowski space, which has been studied in great detail elsewhere . As will be seen, the equation of state found numerically is, in the majority of cases, that of cold matter. We therefore use matter dominated expansion for the evolution in much of the analysis that follows. The evolution in radiation dominated universes remains largely unchanged, although there is greater initial growth of quantum fluctuations due to the scale factor growing more slowly in time. Using the large $`N`$ approximation to study theories with continuous and discrete symmetries respectively, we treat three important cases. They are 1) $`m^2<0,\eta (\tau _0)1;\mathrm{\hspace{0.33em}2})m^2<0,\eta (\tau _0)1;\mathrm{\hspace{0.33em}3})m^2>0,\eta (\tau _0)1`$.
In presenting the figures, we have shifted the origin of time such that $`\tau \tau ^{}=\tau \tau _0`$. This places the initial time, $`\tau _0`$, at the origin. In these shifted coordinates, the scale factor is given by
$$a(t)=\left(\frac{\tau +\tau _0}{\tau _0}\right)^n,$$
where, once again, $`n=2/3`$ and $`n=1/2`$ in matter and radiation dominated backgrounds respectively, and the value of $`\tau `$ is determined by the Hubble constant at the initial time:
$$h(\tau _0)=\frac{n}{\tau _0}.$$
Case 1: $`m^2<0`$, $`\eta (\tau _0)1`$. This is the case of an early universe phase transition in which there is little or no biasing in the initial configuration (by biasing we mean that the initial conditions break the $`\eta \eta `$ symmetry). The transition occurs from an initial temperature above the critical temperature, $`T>T_c`$, which is quenched at $`t_0`$ to the temperature $`T_fT_c`$. This change in temperature due to the rapid expansion of the universe is modeled here by an instantaneous change in the mass from an initial value $`m_i^2=T^2/T_c^21`$ to a final value $`m_f^2=1`$. We will use the value $`m_i^2=1`$ in what follows. This quench approximation is necessary since the low momentum frequencies (91) appearing in our initial conditions (90) are complex for negative mass squared and small $`\eta (t_0)`$. An alternative choice is to use initial frequencies given by
$$\omega _k(\tau _0)=\left[k^2+^2(\tau _0)\mathrm{tanh}\left(\frac{k^2+^2(\tau _0)}{|^2(\tau _0)|}\right)\right]^{1/2}.$$
These frequencies have the attractive feature that they match the conformal adiabatic frequencies given by eq.(91) for large values of $`k`$ while remaining positive for small $`k`$. We find that such a choice of initial conditions changes the quantitative value of the particle number by a few percent, but leaves the qualitative results unchanged.
We plot the the zero mode $`\eta (\tau )`$, the equal time correlator $`g\mathrm{\Sigma }(\tau )`$, the total number of produced particles $`gN(\tau )`$ (see sec. VI for a discussion of our definition of particles), the number of particles $`gN_k(\tau )`$ as a function of wavenumber for both intermediate and late times, and the ratio of the pressure and energy densities $`p(\tau )/\epsilon (\tau )`$ (giving the equation of state).
Figs. 1a-e shows these quantities in the large $`N`$ approximation for a matter dominated cosmology with an initial condition on the zero mode given by $`\eta (\tau _0)=10^7,\dot{\eta }(\tau _0)=0`$ and for an initial expansion rate of $`h(\tau _0)=0.1`$. This choice for the initial value of $`\eta `$ stems from the fact that the quantum fluctuations only have time to grow significantly for initial values satisfying $`\eta (\tau _0)\sqrt{g}`$; for values $`\eta (\tau _0)\sqrt{g}`$ the evolution is essentially classical. This result is clear from the intermediate time dependence of the zero mode and the low momentum mode functions given by the expressions (175) and (177) respectively.
After the initial growth of the fluctuation $`g\mathrm{\Sigma }(\tau )`$ (fig.1b) we see that the zero mode (fig.1a) approaches the value given by the minimum of the tree level potential, $`\eta =1`$, while $`g\mathrm{\Sigma }(\tau )`$ decays for late times as
$$g\mathrm{\Sigma }(\tau )\frac{𝒞}{a^2(\tau )}=\frac{𝒞}{\tau ^{4/3}}.$$
For these late times, the Ward identity corresponding to the $`O(N)`$ symmetry of the field theory is satisfied, enforcing the condition
$$1+\eta ^2(\tau )+g\mathrm{\Sigma }(\tau )=0.$$
(178)
Hence, the zero mode approaches the classical minimum as
$$\eta ^2(\tau )1\frac{𝒞}{a^2(\tau )}.$$
Figure 1c depicts the number of particles produced. After an initial burst of particle production, the number of particles settles down to a relatively constant value. Notice that the number of particles produced is approximately of order $`1/g`$. In fig.1d, we show the number of particles as a function of the wavenumber, $`k`$. For intermediate times we see the simple structure depicted by the dashed line in the figure, while for late times this quantity becomes concentrated more at low values of the momentum $`k`$.
Finally, fig.1e shows that the field begins with a de Sitter equation of state $`p=\epsilon `$ but evolves quickly to a state dominated by ordinary matter, with an equation of state (averaged over the oscillation timescale) $`p=0`$. This last result is a bit surprising as one expects from the condition (178) that the particles produced in the final state are massless Goldstone bosons which should have the equation of state of radiation. However, as shown in fig.1d, the produced particles are of low momentum, $`q1`$, and while the effective mass of the particles is zero to very high accuracy when averaged over the oscillation timescale, the effective mass makes small oscillations about zero so that the dispersion relation for these particles differs from that of radiation. In addition, since the produced particles have little energy, the contribution to the energy density from the zero mode, which contributes to a cold matter equation of state, remains significant.
Finally, we show the special case in which there is no initial biasing in the field, $`\eta (\tau _0)=0,\dot{\eta }(\tau _0)=0`$, and $`h(\tau _0)=0.1`$ in figs. 2a-d. The zero mode remains zero for all time, so that the quantity $`g\mathrm{\Sigma }(\tau )`$ (fig.2a) satisfies the sum rule (178) by reaching the value one without decaying for late times. Notice that many more particles are produced in this case (fig 2b); the growth of the particle number for late times is due to the expansion of the universe. The particle distribution (fig.2c) is similar to that of the slow roll case in fig.1. The equation of state (fig.2d) is likewise similar.
In each of these cases of slow roll dynamics, increasing the Hubble constant has the effect of slowing the growth of both $`\eta `$ and $`g\mathrm{\Sigma }(t)`$. The equation of state will be that of a de Sitter universe for a longer period before moving to a matter dominated equation of state. Otherwise, the dynamics is much the same as in figs. 1-3.
Case 2: $`m^2<0`$, $`\eta (\tau _0)1`$. We now examine the case of a chaotic inflationary scenario with a symmetry broken potential. In chaotic inflation, the zero mode begins with a value $`\eta (\tau )1`$. During the de Sitter phase, $`h1`$, and the field initially evolves classically, dominated by the first order derivative term appearing in the zero mode equation \[see eq.(80)\]. Eventually, the zero mode rolls down the potential, ending the de Sitter phase and beginning the FRW phase. We consider the field dynamics in the FRW universe after the end of inflation. We thus take the initial temperature to be zero, $`T=0`$.
Figure 4 shows our results for the quantities, $`\eta (\tau ),g\mathrm{\Sigma }(\tau ),gN(\tau ),gN_k(\tau )`$ and $`p(\tau )/\epsilon (\tau )`$ for the evolution in the large $`N`$ approximation within a radiation dominated gravitational background with $`h(\tau _0)=0.1`$. The initial condition on the zero mode is chosen to have the representative value $`\eta (\tau _0)=4`$ with $`\dot{\eta }(\tau _0)=0`$. Initial values of the zero mode much smaller than this will not produce significant growth of quantum fluctuations; initial values larger than this produces qualitatively similar results, although the resulting number of particles will be greater and the time it takes for the zero mode to settle into its asymptotic state will be longer.
We see from fig.3a that the zero mode oscillates rapidly, while the amplitude of the oscillation decreases due to the expansion of the universe.
This oscillation induces particle production through the process of parametric amplification (fig.3c) and causes the fluctuation $`g\mathrm{\Sigma }(\tau )`$ to grow (fig.3b). Eventually, the zero mode loses enough energy that it is restricted to one of the two minima of the tree level effective potential. The subsequent evolution closely follows that of Case 1 above with $`g\mathrm{\Sigma }(\tau )`$ decaying in time as $`1/a^2(\tau )1/\tau `$ with $`\eta (\tau )`$ given by the sum rule (178). The spectrum (fig.3d) indicates a single unstable band of particle production dominated by the modes $`k=1/2`$ to about $`k=3`$ for late times. The structure within this band becomes more complex with time and shifts somewhat toward lower momentum modes.
Such a shift is also observed in Minkowski spacetimes . Figure 4e shows the equation of state which we see to be somewhere between the relations for matter and radiation for times out as far as $`t=400`$, but slowly moving to a matter equation of state. Since matter redshifts as $`1/a^3(t)`$ while radiation redshifts as $`1/a^4(t)`$, the equation of state should eventually become matter dominated. Given the equation of state indicated by fig.3e, we estimate that this occurs for times of order $`t=10^4`$. The reason the equation of state in this case differs from that of cold matter as was seen in figs. 1-3 is that the particle distribution
produced by parametric amplification is concentrated at higher momenta, $`k1`$.
Figure 4 shows the corresponding case with a matter dominated background. The results are qualitatively very similar to those described for fig.3 above. Due to the faster expansion, the zero mode (fig.4a) finds one of the two wells more quickly and slightly less particles are produced. For late times, the fluctuation $`g\mathrm{\Sigma }(t)`$ (fig.4b) decays as $`1/a^2(t)1/t^{4/3}`$. Again we see an equation of state (figs. 4e) which evolves from a state between that of pure radiation or matter toward one of cold matter.
A larger Hubble constant prevents significant particle production unless the initial amplitude of the zero mode is likewise increased such that the relation $`\eta (t_0)h(t_0)`$ is satisfied. For very large amplitude $`\eta (t_0)1`$, to the extent that the mass term can be neglected and while the quantum fluctuation term has not grown to be large, the equations of motion (86) are scale invariant with the scaling $`\eta \mu \eta `$, $`H\mu H`$, $`tt/\mu `$, and $`k\mu k`$, where $`\mu `$ is an arbitrary scale.
Case 3: $`m^2>0`$, $`\eta (t_0)1`$. The final case we examine is that of a simple chaotic scenario with a positive mass term in the Lagrangian. Again, the FRW stage occurs after the inflationary expansion; this allows us to take zero initial temperature.
Figure 5 shows this situation in the large $`N`$ approximation for a matter dominated cosmology. The zero mode, $`\eta (\tau )`$, oscillates in time while decaying in amplitude from its initial value of $`\eta (\tau _0=0)=5`$, $`\dot{\eta }(\tau _0=0)=0`$ (fig.5a), while the quantum fluctuation, $`g\mathrm{\Sigma }(t)`$, grows rapidly for early times due to parametric resonance (figs. 5b). We choose here an initial condition on the zero mode which differs from that of figs 2-3 above since there is no significant growth of quantum fluctuations for smaller initial values. From fig.5d, we see that there exists a single unstable band at values of roughly $`k=1`$ to $`k=3`$, although careful examination reveals that the unstable band extends all the way to $`k=0`$. The equation of state is depicted by the quantity $`p(\tau )/\epsilon (\tau )`$ in fig.5e. As expected in this massive theory, the equation of state is matter dominated.
First, we note that, for early times when $`g\mathrm{\Sigma }(\tau )1`$, the zero mode is well fit by the function $`\eta (\tau )=\eta _0f(\tau )/a(\tau )`$ where $`f(\tau )`$ is an oscillatory function taking on values from $`1`$ to $`1`$. This is clearly seen from the envelope function $`\eta _0/a(\tau )`$ shown in fig.6a (recall that $`g\mathrm{\Sigma }(\tau )1`$ during the entire evolution in this case). Second, the momentum that appears in the equations for the modes (80) is the physical momentum $`k/a(\tau )`$. We therefore write the approximate expressions for the locations of the forbidden bands in FRW by using the Minkowski results of with the substitutions $`\eta _0^2\gamma \eta _0^2/a^2(\tau )`$ (where the factor of $`\gamma `$ accounts for the difference in the definition of the non-linear coupling between this study and ) and $`q^2k^2/a^2(\tau )`$.
Making these substitutions, we find for the location in comoving momentum $`k`$ of the forbidden band in the large $`N`$ (fig.5-6) case:
$`0`$ $`k^2`$ $`{\displaystyle \frac{\eta _0^2}{2}}.`$ (179)
The important feature to notice is that while the location of the unstable band (to a first approximation) in the case of the continuous $`O(N)`$ theory is the same as in Minkowski and does not change in time. Again, the qualitative dynamics remains largely unchanged from the case of a smaller Hubble constant.
As in the symmetry broken case of figs. 1-4, the equations of motion for large amplitude and relatively early times are approximately scale invariant. In fig.6 we show the case of the large $`N`$ evolution in a radiation dominated universe with initial Hubble constant of $`H(\tau _0)=2`$ with appropriately scaled initial value of the zero mode of $`\eta (\tau _0)=16`$.
### C Late Time Behavior
We see clearly from the numerical evolution that in the case of a symmetry broken potential, the late time large $`N`$ solutions obey the sum rule (157). This sum rule is a consequence of the late time Ward identities which enforce Goldstone’s Theorem. Because of this sum rule, we can write down the analytical expressions for the late time behavior of the fluctuations and the zero mode. Using eq.(157), the mode equation (80) becomes
$$\left[\frac{d^2}{d\tau ^2}+3\frac{\dot{a}(\tau )}{a(\tau )}\frac{d}{d\tau }+\frac{k^2}{a^2(\tau )}\right]U_k(\tau )=0.$$
(180)
This equation can be solved exactly if we assume a power law dependence for the scale factor $`a(\tau )=(\tau /\tau _0)^n`$ with solution
$$U_k(\tau )=c_k\tau ^{(13n)/2}J_{\frac{13n}{22n}}\left(\frac{k\tau _0^n\tau ^{1n}}{n1}\right)+d_k\tau ^{(13n)/2}Y_{\frac{13n}{22n}}\left(\frac{k\tau _0^n\tau ^{1n}}{n1}\right),$$
(181)
where $`J_\nu `$ and $`Y_\nu `$ are Bessel and Neumann functions respectively, and the constants $`c_k`$ and $`d_k`$ carry dependence on the initial conditions and on the dynamics up to the point at which the sum rule is satisfied.
These functions have several important properties. In particular, in radiation or matter dominated universes, $`n<1`$, and for values of wavenumber satisfying $`k\tau ^{(1n)}/\tau _0^n`$, the mode functions decay in time as $`1/a(\tau )\tau ^n`$. Since the sum rule applies for late times, $`\tau \tau _01`$ in dimensionless units, we see that all values of $`k`$ except a very small band about $`k=0`$ redshift as $`1/a(\tau )`$. The $`k=0`$ mode for any scale factor $`a(\tau )`$ takes the form
$$U_0(\tau )=C_1+C_2^\tau \frac{d\tau ^{}}{a^3(\tau ^{})}.$$
where we used eq.(180) and $`C_1`$ and $`C_2`$ are constants depending on the initial conditions.
We see that the $`k=0`$ mode freezes out for late times tending to a constant. This explains the support evidenced in the numerical results for values of small $`k`$ (see figs. 1,3).
These results mean that the quantum fluctuation has a late time dependence of $`\pi ^2(\tau )_r1/a^2(\tau )`$. The late time dependence of the zero mode is given by this expression combined with the sum rule (157). These results are accurately reproduced by our numerical analysis. Note that qualitatively this late time dependence is independent of the choice of initial conditions for the zero mode, except that there is no growth of modes near $`k=0`$ in the case in which particles are produced via parametric amplification (figs. 4,5).
For the radiation $`n=\frac{1}{2}`$ and matter dominated $`n=\frac{2}{3}`$ universes, eq.(181) reduces to elementary functions:
$`a(\tau )U_k(\tau )`$ $`=`$ $`c_ke^{2ik\tau _0^{1/2}\tau ^{1/2}}+d_ke^{2ik\tau _0^{1/2}\tau ^{1/2}}\text{(RD) },`$ (182)
$`a(\tau )U_k(\tau )`$ $`=`$ $`c_ke^{3ik\tau _0^{2/3}\tau ^{1/3}}\left[1+{\displaystyle \frac{i}{3k\tau _0^{2/3}\tau ^{1/3}}}\right]+d_ke^{3ik\tau _0^{2/3}\tau ^{1/3}}\left[1{\displaystyle \frac{i}{3k\tau _0^{2/3}\tau ^{1/3}}}\right]\text{ (MD) }.`$ (184)
It is also of interest to examine the $`n>1`$ case. Here, the modes of interest satisfy the condition $`k\tau ^{n1}/t_0^n`$ for late times. These modes are constant in time and one sees that the modes are frozen. In the case of a de Sitter universe, we can formally take the limit $`n\mathrm{}`$ and we see that all modes become frozen at late times. This case is detailed in sec. VII .
### D Discussion, Conclusions and Further Results for the FRW background
We have shown that there can be significant particle production through quantum fluctuations after inflation. However, this production is somewhat sensitive to the expansion of the universe. From our analysis of the equation of state, we see that the late time dynamics is given by a matter dominated cosmology. We have also shown that the quantum fluctuations of the inflaton decay for late times as $`1/a^2(t)`$, while in the case of a symmetry broken inflationary model, the inflaton field moves to the minimum of its tree level potential. The exception to this behavior is the case when the inflaton begins exactly at the unstable extremum of its potential for which the fluctuations grow out to the minimum of the potential and do not decay. Initial production of particles due to parametric amplification is significantly greater in chaotic scenarios with symmetry broken potentials than in the corresponding theories with positive mass terms in the Lagrangian, given similar initial conditions on the zero mode of the inflaton.
In ref. we further investigate a symmetry breaking phase transition triggered by the lowering of the temperature as $`1/a(t)`$ both in radiation dominated and matter dominated FRW spacetimes. We identify three different time scales: an early regime dominated by linear instabilities and the exponential growth of long-wavelength fluctuations, an intermediate scale when the field fluctuations probe the broken symmetry states and an asymptotic scale wherein a scaling regime emerges for modes of wavelength comparable to or larger than the horizon. The scaling regime is characterized by a dynamical physical correlation length $`\xi _{phys}=d_H(t)`$ with $`d_H(t)`$ the size of the causal horizon, thus there is one correlated region per causal horizon. Inside these correlated regions the field fluctuations sample the broken symmetry states. The amplitude of the long-wavelength fluctuations becomes non-perturbatively large due to the early times instabilities and a semiclassical but stochastic description emerges in the asymptotic regime. In the scaling regime, the power spectrum is peaked at zero momentum revealing the onset of a Bose-Einstein condensate. The scaling solution results in that the equation of state of the scalar fields is the same as that of the background fluid. This implies a Harrison-Zeldovich spectrum of scalar density perturbations for long-wavelengths. We discuss the corrections to scaling as well as the universality of the scaling solution and the differences and similarities with the classical non-linear sigma model.
## VII Scalar Field Dynamics in a fixed Inflationary Background (the de Sitter Universe)
We describe in this section the scalar field evolution in a fixed de Sitter background using the large $`N`$ approximation, postponing the evolution with a dynamical background to sec. VII.
We consider an initial state at a non-zero temperature $`T_i`$. This change on the initial conditions does not affect the initial values of the mode functions in eqs.(146) or (155).
Only the expression for the quantum fluctuations $`\mathrm{\Sigma }(\tau )`$ changes due to the fact that the expectation value of the product of a creation and an annihilation operator is temperature dependent. We now have instead of eq.(156)
$$\mathrm{\Sigma }(\tau )=_0^{\mathrm{}}q^2𝑑q\left[|f_q(\tau )|^2\mathrm{coth}\left(\frac{\omega _q}{2T_i}\right)\frac{1}{qa(\tau )^2}+\frac{\mathrm{\Theta }(q1)}{2q^3}\left(\frac{^2(\tau )}{m_R^2}\frac{(\tau )}{6m_R^2}\right)\right].$$
(185)
We consider the case $`m^2<0`$ with a critical temperature $`T_c`$ such that $`T_iT_c|m|`$. The symmetry is initially unbroken and through the expansion of the universe the effective temperature decreases as $`T_i/a(t)`$. Therefore, there is a symmetry breaking phase transition after a few e-folds of inflation. When the temperature falls below the critical value, the effective mass becomes negative. As will be seen explicitly below, when this occurs, long-wavelength modes become unstable and grow. Local thermodynamic equilibrium will set in again if the contribution from the quantum fluctuations can grow and adjust to compensate for the negative mass terms on the same time scales as that in which the temperature drops. However, as discussed below, for very weak coupling the important time scales for the non-equilibrium fluctuations are of the order of $`[H/m_R^2]\mathrm{ln}[1/\lambda ]`$, which are much longer than the time it takes for the temperature to drop well below the critical value to practically zero. Thus, the non-equilibrium dynamics will proceed as if the phase transition occured via a quench, that is with an effective mass term,
$$m_{eff}^2(t)=m_i^2\theta (t_it)m_R^2\theta (tt_i);m_i^2=m_R^2[\frac{T_i^2}{T_c^2}1]>0.$$
(186)
Therefore, we choose the initial conditions on the mode functions at $`t_i=0`$ to be given in terms of the effective mass,
$$M_0^2=m_R^2[\frac{T_i^2}{T_c^2}1]+\frac{\lambda _R}{2}\varphi ^2(0)de;\frac{T_i}{T_c}>1.$$
(187)
### A Evolution for $`\varphi (0)=\dot{\varphi }(0)=0`$. Analytical Results
We begin by considering the broken symmetry situation in which the expectation value of the inflaton field sits atop the potential hill with zero initial velocity. This situation is expected to arise if the system is initially in local thermodynamic equilibrium an initial temperature larger than the critical temperature and cools down through the critical temperature in the absence of an external field or bias.
The order parameter and its time derivative vanish in the local equilibrium high temperature phase, and this condition is a fixed point of the evolution equation for the zero mode of the inflaton. There is no rolling of the inflaton zero mode in this case, although the fluctuations will grow and will be responsible for the dynamics.
We can understand the early stages of the dynamics analytically as follows. For very weak coupling and early time we can neglect the backreation in the mode equations, which become,
$$\left[\frac{d^2}{d\tau ^2}+3h\frac{d}{d\tau }+\frac{q^2}{a^2(\tau )}1\right]f_q(\tau )=0,$$
(188)
$$f_q(0)=\frac{1}{\sqrt{\omega _q}};\dot{f}_q(\tau )=i\sqrt{\omega _q};\omega _q=\sqrt{q^2+r^21};r=\frac{T_i}{T_c}.$$
(189)
The solutions are of the form,
$$f_q(\tau )=\mathrm{exp}[\frac{3}{2}h\tau ]\left\{a(q)J_\nu (z)+b(q)J_\nu (z)\right\};z=\frac{q}{h}\mathrm{exp}[h\tau ];\nu =\sqrt{\frac{1}{h^2}+\frac{9}{4}},$$
(190)
where the coefficients $`a(q)`$ and $`b(q)`$ are determined by the initial conditions:
$$b(q)=\frac{\pi q}{2h\mathrm{sin}\nu \pi }\left[\frac{i\omega _q\frac{3}{2}h}{q}J_\nu \left(\frac{q}{h}\right)J_\nu ^{}\left(\frac{q}{h}\right)\right],$$
(191)
$$a(q)=\frac{\pi q}{2h\mathrm{sin}\nu \pi }\left[\frac{i\omega _q\frac{3}{2}h}{q}J_\nu \left(\frac{q}{h}\right)J_\nu ^{}\left(\frac{q}{h}\right)\right].$$
(192)
For long times, $`e^{h\tau }q/h`$, these mode functions grow exponentially,
$$f_q(\tau )b(q)J_\nu (z)\frac{b(q)}{\mathrm{\Gamma }(1\nu )}\left(\frac{2h}{q}\right)^\nu e^{(\nu 3/2)h\tau }.$$
(193)
The Bessel functions appearing in the expression for the modes $`f_q(\tau )`$ can be approximated by their series expansion,
$$f_q(\tau )=\frac{1}{2}\left[1+\frac{1}{\nu }\left(\frac{3}{2}\frac{q^2}{4h^2}i\frac{\omega _q}{h}\right)+𝒪\left(\frac{1}{\nu ^2}\right)\right]e^{(\nu 3/2)h\tau }.$$
(194)
This is an expansion in powers of $`q^2/(\nu h^2)`$ and we conclude that $`g\mathrm{\Sigma }(\tau )`$ is dominated by the modes with $`q\sqrt{h}`$.
The integral for $`g\mathrm{\Sigma }(\tau )`$ can be approximated by keeping only the modes $`qf\sqrt{h}`$, where $`f`$ is a number of order one, and by neglecting the subtraction term which will cancel the contributions from high momenta. Numerically, even with the backreaction taken into account, the integral is dominated by modes $`qf1020`$ in all of the cases that we studied (see ref. ).
The contribution to the fluctuations from these unstable modes is:
$$g\mathrm{\Sigma }(\tau )\sqrt{\frac{g}{3}}\frac{f^3h^{3/2}rm_R^2}{2\pi M_0^2}\left(1+\frac{M_0^2}{m_R^2}\right)e^{(2\nu 3)h\tau },$$
(195)
where again, we have taken the high temperature limit, $`T_iT_cm_R`$.
From this equation, we can estimate the value of $`\tau _s`$, the ‘spinodal time’, at which the contribution of the quantum fluctuations becomes comparable to the contribution from the tree level terms in the equations of motion. This time scale is obtained from the condition $`g\mathrm{\Sigma }(\tau _s)=𝒪(1)`$:
$$\tau _s\frac{1}{(2\nu 3)h}\mathrm{ln}\left[\sqrt{\frac{g}{3}}\frac{f^3h^{3/2}}{2\pi m_RM_0^2}\frac{T_i}{T_c}\left(1+\frac{M_0^2}{m_R^2}\right)\right],$$
(196)
which is in good agreement with our numerical results (see ref. ). For values of $`h1`$, which, as argued below, lead to the most interesting case, an estimate for the spinodal time is,
$$\tau _s\frac{3h}{2}\mathrm{ln}[1/\sqrt{g}]+𝒪(1)$$
(197)
which is consistent with our numerical results (see ).
For $`\tau >\tau _s`$, the effects of backreaction become very important, and the contribution from the quantum fluctuations competes with the tree level terms in the equations of motion, shutting-off the instabilities. Beyond $`\tau _s`$, only a full numerical analysis will capture the correct dynamics.
It is worth mentioning that had we chosen zero temperature initial conditions, then the coupling $`\overline{g}g`$ (see fig.8) and the estimate for the spinodal time would have been,
$$\tau _s\frac{3h}{2}\mathrm{ln}[1/g]+𝒪(1),$$
(198)
that is, roughly a factor 2 larger than the estimate for which the de Sitter stage began at a temperature above the critical value. Therefore, eq.(197) represents an underestimate of the spinodal time scale at which fluctuations become comparable to tree level contributions.
The number of e-folds occurring during the stage of growth of spinodal fluctuations is therefore,
$$𝒩_e\frac{3h^2}{2}\mathrm{ln}[1/\sqrt{g}]\text{(high temperature) and }𝒩_e\frac{3h^2}{2}\mathrm{ln}[1/g]\text{(zero temperature)}$$
It is a factor 2 larger for zero temperature. Thus, it becomes clear that with $`g10^{12}`$ and $`h2`$, a required number of e-folds, $`𝒩_e100`$ can easily be accommodated before the fluctuations become large, modifying the dynamics and the equation of state.
The implications of these estimates are important. The first conclusion drawn from these estimates is that a ‘quench’ approximation is well justified (see ). While the temperature drops from an initial value of a few times the critical temperature to below critical in just a few e-folds, the contribution of the quantum fluctuations needs a large number of e-folds to grow to compensate for the tree-level terms and overcome the instabilities. Only for a strongly coupled theory is the time scale for the quantum fluctuations to grow short enough to restore local thermodynamic equilibrium during the transition.
The second conclusion is that most of the growth of spinodal fluctuations occurs during the inflationary stage, and with $`g10^{12}`$ and $`Hm_R`$, the quantum fluctuations become of the order of the tree-level contributions to the equations of motion within the number of e-folds necessary to solve the horizon and flatness problems. Since the fluctuations grow to become of the order of the tree level contributions at times of the order of this time scale, for larger times they will modify the equation of state substantially and will be shown in sec. IX to provide a graceful exit from the inflationary phase within an acceptable number of e-folds.
### B The late time limit
For late the dynamics freezes out. The fluctuation, $`g\mathrm{\Sigma }(\tau )=1`$, and the mode functions effectively describe free, minimally coupled, massless particles. The sum rule,
$$1+g\mathrm{\Sigma }(\mathrm{})=0,$$
(199)
is obeyed exactly in the large $`N`$ limit as in the Minkowski case.
We now show that this value is a self-consistent solution of the equations of motion for the mode functions, and the only stationary solution for asymptotically long times.
In the late time limit, the effective time dependent mass term, $`1+\eta ^2+g\mathrm{\Sigma }`$, in the equation for the mode functions, (150), vanishes (in this case with $`\eta =0`$). Therefore, these mode equations asymptotically become,
$$\left[\frac{d^2}{d\tau ^2}+3h\frac{d}{d\tau }+\frac{q^2}{a^2(\tau )}\right]f_q(\tau )=0.$$
(200)
The general solutions are given by,
$$f_q^{asy}(\tau )=\mathrm{exp}\left[\frac{3}{2}h\tau \right]\left[c_+(q)J_{3/2}\left(\frac{q}{h}e^{h\tau }\right)ic_{}(q)N_{3/2}\left(\frac{q}{h}e^{h\tau }\right)\right],$$
(201)
where $`J_{3/2}(z)`$ and $`N_{3/2}(z)`$ are the Bessel and Neumann functions, respectively. The coefficients, $`c_\pm (q)`$ can be computed for large $`q`$ by matching $`f_q^{asy}(\tau )`$ with the WKB approximation to the exact mode functions $`f_q(\tau )`$ that obey the initial conditions (155). The WKB approximation to $`f_q(\tau )`$ has been computed in ref., and we find for large $`q`$,
$$c_\pm (q)=\sqrt{\frac{\pi q}{2h}}\left[1\frac{i}{q}(h+\mathrm{\Delta })+𝒪(q^2)\right]e^{iq/h}\pm \sqrt{\frac{\pi h}{8q}}\left[1+𝒪\left(\frac{1}{q}\right)\right]e^{iq/h},$$
(202)
where
$$\mathrm{\Delta }_0^{\mathrm{}}𝑑\tau e^{h\tau }M^2(\tau ).$$
(203)
In the $`\tau \mathrm{}`$ limit, we have for fixed $`q`$,
$$f_q^{asy}(\tau )\stackrel{\tau \mathrm{}}{=}i\sqrt{\frac{2}{\pi }}\left(\frac{h}{q}\right)^{3/2}c_{}(q).$$
(204)
which are independent of time asymptotically, and explains why the power spectrum of quantum fluctuations freezes at times larger than the spinodal. This behavior is confirmed numerically \[see ref. \]. Clearly at early times the mode functions grow exponentially, and at times of the order of $`\tau _s`$, when $`g\mathrm{\Sigma }(\tau )1`$ the mode functions freeze-out and become independent of time. Notice that the largest $`q`$ modes have grown the least, explaining why the integral is dominated by $`q1020`$.
For asymptotically large times, $`g\mathrm{\Sigma }(\tau )`$ is given by,
$$g\mathrm{\Sigma }(\mathrm{})=gh^2_0^+\mathrm{}\frac{dq}{q}\mathrm{coth}\left(\frac{\omega _q}{2T}\right)\left[\frac{2h}{\pi }c_{}(q)^2q\right],$$
(205)
where only one term in the UV subtraction survived in the $`\tau =\mathrm{}`$ limit. The factor $`\mathrm{coth}\left(\frac{\omega _q}{2T}\right)`$ in eq.(205) takes into account the nonzero initial temperature $`T`$.
For consistency, this integral must converge and be equal to $`1`$ as given by the sum rule. For this to be the case and to avoid the potential infrared divergence in (205), the coefficients $`c_{}(q)`$ must vanish at $`q=0`$. The mode functions are finite in the $`q0`$ limit provided,
$$c_{}(q)\stackrel{q0}{=}𝒞q^{3/2},$$
(206)
where $`𝒞`$ is a constant.
The numerical analysis clearly shows that the mode functions remain finite as $`q0`$, and the coefficient $`𝒞`$ can be read off from these figures. This is a remarkable result. It is well known that for free massless minimally coupled fields in de Sitter space-time with Bunch-Davies boundary conditions, the fluctuation contribution $`\pi ^2(\stackrel{}{x},t)`$ grows linearly in time as a consequence of the logarithmic divergence in the integrals. However, in our case, although the asymptotic mode functions are free, the coefficients that multiply the Bessel functions of order $`3/2`$ have all the information of the interaction and initial conditions and must lead to the consistency of the sum rule. Clearly the sum rule and the initial conditions for the mode functions prevent the coefficients $`c_\pm (q)`$ from describing the Bunch-Davies vacuum. These coefficients are completely determined by the initial conditions and the dynamics. This is the reason why the fluctuation freezes at long times unlike in the free case in which they grow linearly.
It is easy to see from eqs.(159) and (204) that the energy and pressure vanish for $`\tau \mathrm{}`$.
### C Discussion and Conclusions for the de Sitter background
We have identified analytically and numerically two distinct regimes for the dynamics determined by the initial condition on the expectation value of the zero mode of the inflaton .
1. When $`\eta (0)<<g^{1/4}`$ (or $`g^{1/2}`$ for $`T_i=0`$), the dynamics is driven by quantum (and thermal) fluctuations. Spinodal instabilities grow and eventually compete with tree level terms at a time scale, $`\tau _s3h\mathrm{ln}[g]/2`$. The growth of spinodal fluctuations translates into the growth of spatially correlated domains which attain a maximum correlation length (domain size) of the order of the horizon. For very weak coupling and $`h1`$ this time scale can easily accommodate enough e-folds for inflation to solve the flatness and horizon problems. The quantum fluctuations modify the equation of state dramatically providing a means for a graceful exit to the inflationary stage without slow-roll.
This non-perturbative description of the non-equilibrium effects in this regime in which quantum (and thermal) fluctuations are most important provides a reliable understanding of the relevant non-perturbative, non-equilibrium effects of the fluctuations that have not been revealed before in this setting \- .
These initial conditions are rather natural if the de Sitter era arises during a phase transition from a radiation dominated high temperature phase in local thermodynamic equilibrium, in which the order parameter and its time derivative vanish.
2. When $`\eta (0)>>g^{1/4}`$ (or $`g^{1/2}`$ for $`T_i=0`$), the dynamics is driven solely by the classical evolution of the inflaton zero mode. The quantum and thermal fluctuations are always perturbatively small (after renormalization), and their contribution to the dynamics is negligible for weak couplings. The de Sitter era will end when the kinetic contribution to the energy becomes of the same order as the ‘vacuum’ term. This is the realm of the slow-roll analysis whose characteristics and consequences have been analyzed in the literature at length. These initial conditions, however, necessarily imply some initial state either with a biasing field that favors a non-zero initial expectation value, or that in the radiation dominated stage, prior to the phase transition, the state was strongly out of equilibrium with an expectation value of the zero mode different from zero. Although such a state cannot be ruled out and would naturally arise in chaotic scenarios, the description of the phase transition in this case requires further input on the nature of the state prior to the phase transition.
## VIII Self-consistent Evolution of Matter Fields with a dynamical cosmological background
We present in this section the full self-consistent matter-geometry dynamics. That is, the scale factor $`a(t)`$ is here a dynamical variable determined by the Einstein-Friedman eq.(158)-(159) coupled with the scalar field evolution eqs.(150)-(155).
In order to provide the full solution we now must provide the values of $`\eta (0)`$, $`\dot{\eta }(0)`$, and $`h_0`$. Assuming that the inflationary epoch is associated with a phase transition at the GUT scale, this requires that $`Nm_R^4/g(10^{15}\text{ Gev })^4`$ and assuming the bound on the scalar self-coupling $`g10^{12}10^{14}`$ (this will be seen later to be a compatible requirement), we find that $`h_0N^{1/4}`$ which we will take to be reasonably given by $`h_0110`$ (for example in popular GUT’s $`N20`$ depending on particular representations).
We will begin by studying the case of most interest from the point of view of describing the phase transition: $`\eta (0)=0`$ and $`\dot{\eta }(0)=0`$, which are the initial conditions that led to puzzling questions. With these initial conditions, the evolution equation for the zero mode eq.(150) determines that $`\eta (\tau )=0`$ by symmetry.
### A Early time dynamics:
Before engaging in the numerical study, it proves illuminating to obtain an estimate of the relevant time scales and an intuitive idea of the main features of the dynamics. Because the coupling is so weak ($`g10^{12}1`$) and after renormalization the contribution from the quantum fluctuations to the equations of motion is finite, we can neglect all the terms proportional to $`g`$ in eqs.(159) and (150).
For the case where we choose $`\eta (\tau )=0`$ and the evolution equations for the mode functions are those for an inverted oscillator in De Sitter space-time, which have been studied in sec. VII . One obtains the approximate solutions (190)-(191).
After the physical wavevectors cross the horizon, i.e. when $`qe^{h_0\tau }/h_01`$ we find that the mode functions factorize:
$$f_q(\tau )\stackrel{qe^{h_0\tau }h_0}{=}\frac{B_q}{\mathrm{\Gamma }(1\nu )}\left(\frac{2h_0}{q}\right)^\nu e^{(\nu 3/2)h_0\tau }.$$
(207)
This result reveals a very important feature: because of the negative mass squared term in the matter Lagrangian leading to symmetry breaking (and $`\nu >3/2`$), we see that all of the mode functions grow exponentially after horizon crossing (for positive mass squared $`\nu <3/2`$, and they would decrease exponentially after horizon crossing). This exponential growth is a consequence of the spinodal instabilities which is a hallmark of the process of phase separation that occurs to complete the phase transition. We note, in addition that the time dependence is exactly given by that of the $`q=0`$ mode, i.e. the zero mode, which is a consequence of the redshifting of the wavevectors and the fact that after horizon crossing the contribution of the term $`q^2/a^2(\tau )`$ in the equations of motion become negligible. We clearly see that the quantum fluctuations grow exponentially and they will begin to be of the order of the tree level terms in the equations of motion when $`g\mathrm{\Sigma }(\tau )1`$. At large times substantially before the end of inflation
$$\mathrm{\Sigma }(\tau )^2(h_0)h_0^2e^{(2\nu 3)h_0\tau },$$
with $`(h_0)`$ a finite constant that depends on the initial conditions and is found numerically to be of $`𝒪(1)`$ \[see fig.14\].
In terms of the initial dimensionful variables, the condition $`g\mathrm{\Sigma }(\tau )1`$ translates to $`<\pi ^2(\stackrel{}{x},t)>_R2m_R^2/g`$, i.e. the quantum fluctuations sample the minima of the (renormalized) tree level potential. We find that the time at which the contribution of the quantum fluctuations becomes of the same order as the tree level terms is estimated to be
$$\tau _s\frac{1}{(2\nu 3)h_0}\mathrm{ln}\left[\frac{1}{gh_0^2^2(h_0)}\right]=\frac{3}{2}h_0\mathrm{ln}\left[\frac{1}{gh_0^2^2(h_0)}\right]+𝒪(1/h_0).$$
(208)
At this time, the contribution of the quantum fluctuations makes the back reaction very important and, as will be seen numerically, this translates into the fact that $`\tau _s`$ also determines the end of the De Sitter era and the end of inflation. The total number of e-folds during the stage of exponential expansion of the scale factor (constant $`h_0`$) is given by
$$N_e\frac{1}{2\nu 3}\mathrm{ln}\left[\frac{1}{gh_0^2^2(h_0)}\right]=\frac{3}{2}h_0^2\mathrm{ln}\left[\frac{1}{gh_0^2^2(h_0)}\right]+𝒪(1)$$
(209)
For large $`h_0`$ we see that the number of e-folds scales as $`h_0^2`$ as well as with the logarithm of the inverse coupling. These results (207-209) will be confirmed numerically below and will be of paramount importance for the interpretation of the main consequences of the dynamical evolution.
As discussed in sec. VII.C, the early time dynamics is dominated by classical or quantum effects depending on the ratio between the time scales $`\tau _c`$ and $`\tau _s`$.
If $`\tau _c`$ is much smaller than the spinodal time $`\tau _s`$ given by eq.(208) then the classical evolution of the zero mode will dominate the dynamics and the quantum fluctuations will not become very large, although they will still undergo spinodal growth. On the other hand, if $`\tau _c\tau _s`$ the quantum fluctuations will grow to be very large well before the zero mode reaches the non-linear regime. In this case the dynamics will be determined completely by the quantum fluctuations. Then the criterion for the classical or quantum dynamics is given by
$`\eta (0)`$ $``$ $`\sqrt{g}h_0\text{ classical dynamics}`$ (210)
$`\eta (0)`$ $``$ $`\sqrt{g}h_0\text{ quantum dynamics}`$ (211)
or in terms of dimensionful variables $`\varphi (0)H_0`$ leads to classical dynamics and $`\varphi (0)H_0`$ leads to quantum dynamics.
However, even when the classical evolution of the zero mode dominates the dynamics, the quantum fluctuations grow exponentially after horizon crossing unless the value of $`\varphi (t)`$ is very close to the minimum of the tree level potential. In the large $`N`$ approximation the spinodal line, that is the values of $`\varphi (t)`$ for which there are spinodal instabilities, reaches all the way to the minimum of the tree level potential as can be seen from the equations of motion for the mode functions. Therefore even in the classical case one must understand how to deal with quantum fluctuations that grow after horizon crossing.
### B Numerics
The time evolution is carried out by means of a fourth order Runge-Kutta routine with adaptive step-sizing while the momentum integrals are carried out using an 11-point Newton-Cotes integrator. The relative errors in both the differential equation and the integration are of order $`10^8`$. We find that the energy is covariantly conserved throughout the evolution to better than a part in a thousand. Figs. 89 show $`g\mathrm{\Sigma }(\tau )`$ vs. $`\tau `$, $`h(\tau )`$ vs. $`\tau `$ and $`\mathrm{ln}|f_q(\tau )|^2`$ vs. $`\tau `$ for several values of $`q`$ with larger $`q^{}s`$ corresponding to successively lower curves. Figs. 11,11 show $`p(\tau )/\epsilon (\tau )`$ and the horizon size $`h^1(\tau )`$ for $`g=10^{14};\eta (0)=0;\dot{\eta }(0)=0`$ and we have chosen the representative value $`h_0=2.0`$.
Figs. 8 and 8 show clearly that when the contribution of the quantum fluctuations $`g\mathrm{\Sigma }(\tau )`$ becomes of order 1 inflation ends, and the time scale for $`g\mathrm{\Sigma }(\tau )`$ to reach $`𝒪(1)`$ is very well described by the estimate (208). From fig.8 we see that this happens for $`\tau =\tau _s90`$, leading to a number of e-folds $`N_e180`$ which is correctly estimated by eqs. (208)-(209).
Fig. 9 shows clearly the factorization of the modes after they cross the horizon as described by eq.(207). The slopes of all the curves after they become straight lines in fig.9 is given exactly by $`(2\nu 3)`$, whereas the intercept depends on the initial condition on the mode function and the larger the value of $`q`$ the smaller the intercept because the amplitude of the mode function is smaller initially. Although the intercept depends on the initial conditions on the long-wavelength modes, the slope is independent of the value of $`q`$ and is the same as what would be obtained in the linear approximation for the square of the zero mode at times long enough that the decaying solution can be neglected but short enough that the effect of the non-linearities is very small. Notice from the figure that when inflation ends and the non-linearities become important all of the modes effectively saturate. This is also what one would expect from the solution of the zero mode: exponential growth in early-intermediate times (neglecting the decaying solution), with a growth exponent given by $`(\nu 3/2)`$ and an asymptotic behavior of small oscillations around the equilibrium position, which for the zero mode is $`\eta =1`$, but for the $`q0`$ modes depends on the initial conditions. All of the mode functions have this behavior once they cross the horizon. We have also studied the phases of the mode functions and we found that they freeze after horizon crossing in the sense that they become independent of time. This is natural since both the real and imaginary parts of $`f_q(\tau )`$ obey the same equation but with different boundary conditions. After the physical wavelength crosses the horizon, the dynamics is insensitive to the value of $`q`$ for real and imaginary parts and the phases become independent of time. Again, this is a consequence of the factorization of the modes.
The growth of the quantum fluctuations is sufficient to end inflation at a time given by $`\tau _s`$ in eq.(208). Furthermore fig. 11 shows that during the inflationary epoch $`p(\tau )/\epsilon (\tau )1`$ and the end of inflation is rather sharp at $`\tau _s`$ with $`p(\tau )/\epsilon (\tau )`$ oscillating between $`\pm 1`$ with zero average over the cycles, resulting in matter domination. Fig. 11 shows this feature very clearly; $`h(\tau )`$ is constant during the de Sitter epoch and becomes matter dominated after the end of inflation with $`h^1(\tau )\frac{3}{2}(\tau \tau _s)`$. There are small oscillations around this value because both $`p(\tau )`$ and $`\epsilon (\tau )`$ oscillate. These oscillations are a result of small oscillations of the mode functions after they saturate, and are also a feature of the solution for a zero mode.
All of these features hold for a variety of initial conditions. As an example, we show in ref. the case of an initial Hubble parameter of $`h_0=10`$.
### C Zero Mode Assembly:
This remarkable feature of factorization of the mode functions after horizon crossing can be elegantly summarized as
$$f_k(t)|_{k_{ph}(t)H}=g(q,h)f_0(\tau ),$$
(212)
with $`k_{ph}(t)=ke^{Ht}`$ being the physical momentum, $`g(q,h)`$ a complex constant, and $`f_0(\tau )`$ a real function of time that satisfies the mode equation with $`q=0`$ and real initial conditions which will be inferred later.
Then we consider the contribution of these modes to the renormalized quantum fluctuations a long time after the beginning of inflation (so as to neglect the decaying solutions), we find that
$$g\mathrm{\Sigma }(\tau )𝒞e^{(2\nu 3)h\tau }+\text{ small},$$
where ‘small’ stands for the contribution of mode functions associated with momenta that have not yet crossed the horizon at time $`\tau `$, which give a perturbatively small (of order $`g`$) contribution. We find that several e-folds after the beginning of inflation and till inflation ends, this factorization of superhorizon modes implies the following:
$`g{\displaystyle }q^2dq|f_q^2(\tau )||C_0|^2f_0^2(\tau ),g{\displaystyle }q^2dq|\dot{f}_q^2(\tau )||C_0|^2\dot{f}_0^2(\tau ),`$ (213)
(214)
$`g{\displaystyle \frac{q^4}{a^2(\tau )}𝑑q|f_q^2(\tau )|}{\displaystyle \frac{|C_1|^2}{a^2(\tau )}}f_0^2(\tau ),`$ (215)
where we have neglected the weak time dependence arising from the perturbatively small contributions of the short-wavelength modes that have not yet crossed the horizon, and the integrals above are to be understood as the fully renormalized (subtracted), finite integrals. For $`\eta =0`$, we note that (215) and the fact that $`f_0(\tau )`$ obeys the equation of motion for the mode with $`q=0`$ leads at once to the conclusion that in this regime $`\left[g\mathrm{\Sigma }(\tau )\right]^{\frac{1}{2}}=|C_0|f_0(\tau )`$ obeys the zero mode equation of motion
$$\left[\frac{d^2}{d\tau ^2}+3h\frac{d}{d\tau }1+(|C_0|f_0(\tau ))^2\right]|C_0|f_0(\tau )=0.$$
(216)
It is clear that at a time $`\tau _A`$ several e-folds after the beginning of inflation, we can define an effective zero mode as
$$\eta _{eff}^2(\tau )g\mathrm{\Sigma }(\tau ),\text{ or in dimensionful variables, }\varphi _{eff}(t)\left[\pi ^2(\stackrel{}{x},t)_R\right]^{\frac{1}{2}}$$
(217)
Although this identification seems natural, we emphasize that it is by no means a trivial or ad-hoc statement. There are several important features that allow an unambiguous identification: i) $`\left[\pi ^2(\stackrel{}{x},t)_R\right]`$ is a fully renormalized operator product and hence finite, ii) because of the factorization of the superhorizon modes that enter in the evaluation of $`\left[\pi ^2(\stackrel{}{x},t)_R\right]`$, $`\varphi _{eff}(t)`$ (217) obeys the equation of motion for the zero mode, iii) this identification is valid several e-folds after the beginning of inflation, after the transient decaying solutions have died away and the integral in $`\pi ^2(\stackrel{}{x},t)`$ is dominated by the modes with wavevector $`k`$ that have crossed the horizon at $`t(k)t`$. Numerically we see that this identification holds throughout the dynamics except for a very few e-folds at the beginning of inflation. This factorization determines at once the initial conditions of the effective zero mode that can be extracted numerically: after the first few e-folds and long before the end of inflation we find
$$\varphi _{eff}(t)\varphi _{eff}(t_A)e^{(\nu \frac{3}{2})H(tt_A)},$$
(218)
where we parameterized
$$\varphi _{eff}(t_A)\frac{H}{2\pi }(H/m)$$
to make contact with the literature. As is shown in fig. (14), we find numerically that $`(H/m)𝒪(1)`$ for a large range of $`0.1H/m50`$ and that this quantity depends on the initial conditions of the long wavelength modes.
Therefore, in summary, the effective composite zero mode obeys
$$\left[\frac{d^2}{d\tau ^2}+3h\frac{d}{d\tau }1+\eta _{eff}^2(\tau )\right]\eta _{eff}(\tau )=0;\dot{\eta }_{eff}(\tau _A)=\left(\nu \frac{3}{2}\right)\eta _{eff}(\tau _A),$$
(219)
where $`\tau _A`$ is the time at which the composite mode becomes effective, $`\eta _{eff}(\tau _A)\frac{\sqrt{\lambda _R/2}}{m_R}\varphi _{eff}(t_A)`$ is obtained numerically for a given $`h_0`$ by fitting the intermediate time behavior of $`g\mathrm{\Sigma }(\tau )`$ with the growing zero mode solution. Recall that $`\lambda _R=8\pi ^2g`$.
Furthermore, this analysis shows that in the case $`\eta =0`$, the renormalized energy and pressure in this regime in which the renormalized integrals are dominated by the superhorizon modes are given by
$`\epsilon _R(\tau )={\displaystyle \frac{2Nm_R^4}{\lambda _R}}\{{\displaystyle \frac{1}{2}}\dot{\eta }_{eff}^2+{\displaystyle \frac{1}{4}}(1+\eta _{eff}^2)^2\},(p+\epsilon )_R={\displaystyle \frac{2Nm_R^4}{\lambda _R}}\left\{\dot{\eta }_{eff}^2\right\}`$ (220)
where we have neglected the contribution proportional to $`1/a^2(\tau )`$ because it is effectively red-shifted away after just a few e-folds. We found numerically that this term is negligible after the interval of time necessary for the superhorizon modes to dominate the contribution to the integrals. Then the dynamics of the scale factor is given by
$$h^2(\tau )=4h_0^2\left\{\frac{1}{2}\dot{\eta }_{eff}^2+\frac{1}{4}\left(1+\eta _{eff}^2\right)^2\right\}.$$
(221)
We have numerically evolved the set of effective equations (219)-(221) by extracting the initial condition for the effective zero mode from the intermediate time behavior of $`g\mathrm{\Sigma }(\tau )`$. We found a remarkable agreement between the evolution of $`\eta _{eff}^2`$ and $`g\mathrm{\Sigma }(\tau )`$ and between the dynamics of the scale factor in terms of the evolution of $`\eta _{eff}(\tau )`$, and the full dynamics of the scale factor and quantum fluctuations within our numerical accuracy. Figs. 14 and 14 show the evolution of $`\eta _{eff}^2(\tau )`$ and $`h(\tau )`$ respectively from the classical evolution eqs. (219) and (221) using the initial condition $`\eta _{eff}(\tau _A)`$ extracted from the exponential fit of $`g\mathrm{\Sigma }(\tau )`$ in the intermediate regime. These figures should be compared to figs. 8 and 8. We have also numerically compared $`p/\epsilon `$ given solely by the dynamics of the effective zero mode and it is again numerically indistinguishable from that obtained with the full evolution of the mode functions.
This is one of the main results of our work. In summary: the modes that become superhorizon sized and grow through the spinodal instabilities assemble themselves into an effective composite zero mode a few e-folds after the beginning of inflation. This effective zero mode drives the dynamics of the FRW scale factor, terminating inflation when the non-linearities become important. In terms of the underlying fluctuations, the spinodal growth of superhorizon modes gives a non-perturbatively large contribution to the energy momentum tensor that drives the dynamics of the scale factor. Inflation terminates when the mean square root fluctuation probes the equilibrium minima of the tree level potential.
The extension of this analysis to the case for which $`\eta (0)0`$ is straightforward. Since both $`\eta (\tau )`$ and $`\sqrt{g\mathrm{\Sigma }(\tau )}=|C_0|f_0(\tau )`$ obey the equation for the zero mode, eq.(150), it is clear that we can generalize our definition of the effective zero mode to be
$$\eta _{eff}(\tau )\sqrt{\eta ^2(\tau )+g\mathrm{\Sigma }(\tau )}.$$
(222)
which obeys the equation of motion of a classical zero mode:
$$\left[\frac{d^2}{d\tau ^2}+3h(\tau )\frac{d}{d\tau }1+\eta _{eff}(\tau )^2\right]\eta _{eff}(\tau )=0.$$
(223)
If this effective zero mode is to drive the FRW expansion, then the additional condition
$$\dot{\eta }_{eff}^2f_0^22\eta _{eff}\dot{\eta }_{eff}f_0\dot{f_0}+\eta _{eff}^2\dot{f_0}^2=\left[\dot{\eta }_{eff}f_0\eta _{eff}\dot{f_0}\right]^2=0,$$
(224)
must also be satisfied. One can easily show that this relation is indeed satisfied if the mode functions factorize as in (212) and if the integrals in eqs.(215) are dominated by the contributions of the superhorizon mode functions. This leads to the conclusion that the gravitational dynamics is given by eqs. (220) – (221) with $`\eta _{eff}(\tau )`$ defined by (222). Eq.(224) is just the vanishing of the wronskian of $`\eta _{eff}(\tau )`$ and $`f_0(\tau )`$. Namely, $`\eta _{eff}(\tau )`$ and $`f_0(\tau )`$ just differ in a constant factor.
We see that in all cases, the full large $`N`$ quantum dynamics in these models of inflationary phase transitions is well approximated by the equivalent dynamics of a homogeneous, classical scalar field with initial conditions on the effective field $`\eta _{eff}(\tau _A)=\sqrt{g}h_0(h_0)`$. We have verified these results numerically for the field and scale factor dynamics, finding that the effective classical dynamics reproduces the results of the full dynamics to within our numerical accuracy. We have also checked numerically that the estimate for the classical to quantum crossover given by eq.(211) is quantitatively correct. Thus in the classical case in which $`\eta (0)\sqrt{g}h_0`$ we find that $`\eta _{eff}(\tau )=\eta (\tau )`$, whereas in the opposite, quantum case $`\eta _{eff}(\tau )=\sqrt{g\mathrm{\Sigma }(\tau )}`$.
This remarkable feature of zero mode assembly of long-wavelength, spinodally unstable modes is a consequence of the presence of the horizon. It also explains why, despite the fact that asymptotically the fluctuations sample the broken symmetry state, the equation of state is that of matter. Since the excitations in the broken symmetry state are massless Goldstone bosons one would expect radiation domination. However, the assembly phenomenon, i.e. the redshifting of the wave vectors, makes these modes behave exactly like zero momentum modes that give an equation of state of matter (upon averaging over the small oscillations around the minimum).
Subhorizon modes at the end of inflation with $`q>h_0e^{h_0\tau _s}`$ do not participate in the zero mode assembly. The behavior of such modes do depend on $`q`$ after the end of inflation. Notice that these modes have extremely large comoving $`q`$ since $`h_0e^{h_0\tau _s}10^{26}`$. As discussed in sec. VI such modes decrease with time after inflation as $`1/a(\tau )`$.
### D Making sense of small fluctuations:
Having recognized the effective classical variable that can be interpreted as the component of the field that drives the FRW background and rolls down the classical potential hill, we want to recognize unambiguously the small fluctuations. We have argued above that after horizon crossing, all of the mode functions evolve proportionally to the zero mode, and the question arises: which modes are assembled into the effective zero mode whose dynamics drives the evolution of the FRW scale factor and which modes are treated as perturbations? In principle every $`k0`$ mode provides some spatial inhomogeneity, and assembling these into an effective homogeneous zero mode seems in principle to do away with the very inhomogeneities that one wants to study. However, scales of cosmological importance today first crossed the horizon during the last 60 or so e-folds of inflation. Recently Grishchuk has argued that the sensitivity of the measurements of $`\mathrm{\Delta }T/T`$ probe inhomogeneities on scales $`500`$ times the size of the present horizon. Therefore scales that are larger than these and that have first crossed the horizon much earlier than the last 60 e-folds of inflation are unobservable today and can be treated as an effective homogeneous component, whereas the scales that can be probed experimentally via the CMB inhomogeneities today must be treated separately as part of the inhomogeneous perturbations of the CMB.
Thus a consistent description of the dynamics in terms of an effective zero mode plus ‘small’ quantum fluctuations can be given provided the following requirements are met: a) the total number of e-folds $`N_e60`$, b) all the modes that have crossed the horizon before the last 60-65 e-folds are assembled into an effective classical zero mode via $`\varphi _{eff}(t)=\left[\varphi _0^2(t)+\pi ^2(\stackrel{}{x},t)_R\right]^{\frac{1}{2}}`$, c) the modes that cross the horizon during the last 60–65 e-folds are accounted as ‘small’ perturbations. The reason for the requirement a) is that in the separation $`\varphi (\stackrel{}{x},t)=\varphi _{eff}(t)+\delta \varphi (\stackrel{}{x},t)`$ one requires that $`\delta \varphi (\stackrel{}{x},t)/\varphi _{eff}(t)1`$. As argued above, after the modes cross the horizon, the ratio of amplitudes of the mode functions remains constant and given by $`e^{(\nu \frac{3}{2})\mathrm{\Delta }N}`$ with $`\mathrm{\Delta }N`$ being the number of e-folds between the crossing of the smaller $`k`$ and the crossing of the larger $`k`$. Then for $`\delta \varphi (\stackrel{}{x},t)`$ to be much smaller than the effective zero mode, it must be that the Fourier components of $`\delta \varphi `$ correspond to very large $`k`$’s at the beginning of inflation, so that the effective zero mode can grow for a long time before the components of $`\delta \varphi `$ begin to grow under the spinodal instabilities. In fact requirement a) is not very severe; in the figs.8-11 we have taken $`h_0=2.0`$ which is a very moderate value and yet for $`g=10^{12}`$ the inflationary stage lasts for well over 100 e-folds, and as argued above, the larger $`h_0`$ for fixed $`g`$, the longer is the inflationary stage. Therefore under this set of conditions, the classical dynamics of the effective zero mode $`\varphi _{eff}(t)`$ drives the FRW background, whereas the inhomogeneous fluctuations $`\delta \varphi (\stackrel{}{x},t)`$, which are made up of Fourier components with wavelengths that are much smaller than the horizon at the beginning of inflation and that cross the horizon during the last 60 e-folds, provide the inhomogeneities that seed density perturbations.
### E Scalar Metric Perturbations:
Having identified the effective zero mode and the ‘small perturbations’, we are now in position to provide an estimate for the amplitude and spectrum of scalar metric perturbations. We use the clear formulation in ref. in terms of gauge invariant variables. In particular we focus on the dynamics of the Bardeen potential, which in longitudinal gauge is identified with the Newtonian potential. The equation of motion for the Fourier components (in terms of comoving wavevectors) for this variable in terms of the effective zero mode is \[eq.(4.48) in \]
$$\ddot{\mathrm{\Phi }}_k+\left[H(t)2\frac{\ddot{\varphi }_{eff}(t)}{\dot{\varphi }_{eff}(t)}\right]\dot{\mathrm{\Phi }}_k+\left[\frac{k^2}{a^2(t)}+2\left(\dot{H}(t)H(t)\frac{\ddot{\varphi }_{eff}(t)}{\dot{\varphi }_{eff}(t)}\right)\right]\mathrm{\Phi }_k=0.$$
(225)
We are interested in determining the dynamics of $`\mathrm{\Phi }_k`$ for those wavevectors that cross the horizon during the last 60 e-folds before the end of inflation. During the inflationary stage and substantially before its end the numerical analysis yields to a very good approximation
$$H(t)=H_0;\varphi _{eff}(t)=\varphi _{eff}(t_A)e^{(\nu \frac{3}{2})H_0(tt_A)},$$
(226)
where $`H_0`$ is the value of the Hubble constant during inflation, leading to
$$\mathrm{\Phi }_k(t)=e^{(\nu 2)H_0(tt_A)}\left[a_kH_{\nu 1}^{(1)}\left(\frac{ke^{H_0(tt_A)}}{H_0}\right)+b_kH_{\nu 1}^{(2)}\left(\frac{ke^{H_0(tt_A)}}{H_0}\right)\right].$$
(227)
The coefficients $`a_k,b_k`$ are determined by the initial conditions and $`\nu `$ is given by eq.(190).
Since we are interested in the wavevectors that cross the horizon during the last 60 e-folds, the consistency for the zero mode assembly and the interpretation of ‘small perturbations’ requires that there must be many e-folds before the last 60. We are then considering wavevectors that were deep inside the horizon at the onset of inflation. $`\mathrm{\Phi }_k(t)`$ is related to the canonical ‘velocity field’ that determines scalar perturbations of the metric and which is quantized with Bunch-Davies initial conditions for the large $`k`$-mode functions. The relation between $`\mathrm{\Phi }_k`$ and $`v`$ and the initial conditions on $`v`$ lead at once to a determination of the coefficients $`a_k`$ and $`b_k`$ for $`k>>H_0`$ \[see eqs.(13.5) and (13.9) in ref.\],
$$a_k=\frac{\dot{\varphi }_{eff}(t_A)}{k\sqrt{2H_0}M_{Pl}^2}e^{\frac{i\pi }{2}\left(\nu \frac{3}{2}\right)}.$$
(228)
Thus we find that the amplitude of scalar metric perturbations after horizon crossing is given by
$$|\delta _k(t)|=k^{\frac{3}{2}}|\mathrm{\Phi }_k(t)|=\frac{\dot{\varphi }_{eff}(t_A)}{\pi M_{Pl}^2}\mathrm{\Gamma }(\nu 1)\left(\frac{2H_0}{k}\right)^{\nu \frac{3}{2}}e^{(2\nu 3)H_0(tt_A)}.$$
(229)
The power spectrum per logarithmic $`k`$ interval is given by $`|\delta _k(t)|^2`$. The time dependence of $`|\delta _k(t)|`$ displays the unstable growth associated with the spinodal instabilities of super-horizon modes and is a hallmark of the phase transition. This time dependence can be also understood from the constraint equation that relates the Bardeen potential to the gauge invariant field fluctuations, which in longitudinal gauge are identified with $`\delta \varphi (\stackrel{}{x},t)`$. The constraint equation and the evolution equations for the gauge invariant scalar field fluctuations are
$$\frac{d}{dt}(a\mathrm{\Phi }_k)=\frac{4\pi }{M_{Pl}^2}a\delta \varphi _k^{gi}\dot{\varphi }_{eff},$$
(230)
$$\left[\frac{d^2}{dt^2}+3H\frac{d}{dt}+\frac{k^2}{a^2}+^2\right]\delta \varphi _k^{gi}4\dot{\varphi }_{eff}\dot{\mathrm{\Phi }}_k+2V^{}(\varphi _{eff})\mathrm{\Phi }_k=0.$$
(231)
Since the right hand side of (230) is proportional to $`\dot{\varphi }_{eff}/M_{Pl}^21`$ during the inflationary epoch in this model, we can neglect the terms proportional to $`\dot{\mathrm{\Phi }}_k`$ and $`\mathrm{\Phi }_k`$ on the left hand side of (231), in which case the equation for the gauge invariant scalar field fluctuation is the same as for the mode functions. In fact, since $`\delta \varphi _k^{gi}`$ is gauge invariant we can evaluate it in the longitudinal gauge wherein it is identified with the mode functions $`f_k(t)`$. Then absorbing a constant of integration in the initial conditions for the Bardeen variable, we find
$$\mathrm{\Phi }_k(t)=\frac{4\pi }{M_{Pl}^2a(t)}_{t_o}^ta(t^{})\varphi _{eff}(t^{})f_k(t^{})𝑑t^{}+𝒪\left(\frac{1}{M_{Pl}^4}\right),$$
(232)
and using that $`\varphi (t)e^{(\nu 3/2)H_0t}`$ and that after horizon crossing $`f_k(t)e^{(\nu 3/2)H_0t}`$, one obtains at once the time dependence of the Bardeen variable after horizon crossing. In particular the time dependence is found to be $`e^{(2\nu 3)H_0t}`$. It is then clear that the time dependence is a reflection of the spinodal (unstable) growth of the superhorizon field fluctuations.
To obtain the amplitude and spectrum of density perturbations at second horizon crossing we use the conservation law associated with the gauge invariant variable
$$\xi _k=\frac{2}{3}\frac{\frac{\dot{\mathrm{\Phi }}_k}{H}+\mathrm{\Phi }_k}{1+p/\epsilon }+\mathrm{\Phi }_k;\dot{\xi }_k=0,$$
(233)
which is valid after horizon crossing of the mode with wavevector $`k`$. Although this conservation law is an exact statement of superhorizon mode solutions of eq.(225), we have obtained solutions assuming that during the inflationary stage $`H`$ is constant and have neglected the $`\dot{H}`$ term in Eq. (225). Since during the inflationary stage,
$$\dot{H}(t)=\frac{4\pi }{M_{Pl}^2}\dot{\varphi }_{eff}^2(t)H_0^2\left(\frac{d\eta _{eff}(\tau )}{d\tau }\right)^2H_0^2$$
(234)
and $`\ddot{\varphi }/\dot{\varphi }H_0`$, the above approximation is justified. We then see that $`\varphi _{eff}^2(t)e^{(2\nu 3)H_0t}`$ which is the same time dependence as that of $`\mathrm{\Phi }_k(t)`$. Thus the term proportional to $`1/(1+p/\epsilon )`$ in Eq. (233) is indeed constant in time after horizon crossing. On the other hand, the term that does not have this denominator evolves in time but is of order $`(1+p/\epsilon )=2\dot{H}/3H^21`$ with respect to the constant term and therefore can be neglected. Thus, we confirm that the variable $`\xi `$ is conserved up to the small term proportional to $`(1+p/\epsilon )\mathrm{\Phi }_k`$ which is negligible during the inflationary stage. This small time dependence is consistent with the fact that we neglected the $`\dot{H}`$ term in the equation of motion for $`\mathrm{\Phi }_k(t)`$. The validity of the conservation law has been recently studied and confirmed in different contexts. Notice that we do not have to assume that $`\dot{\mathrm{\Phi }}_k`$ vanishes, which in fact does not occur.
However, upon second horizon crossing it is straightforward to see that $`\dot{\mathrm{\Phi }}_k(t_f)0`$. The reason for this assertion can be seen as follows: eq.(231) shows that at long times, when the effective zero mode is oscillating around the minimum of the potential with a very small amplitude and when the time dependence of the fluctuations has saturated (see fig.9), $`\mathrm{\Phi }_k`$ will redshift as $`1/a(t)`$ and its derivative becomes extremely small.
Using this conservation law, assuming matter domination at second horizon crossing, and $`\dot{\mathrm{\Phi }}_k(t_f)0`$, we find
$$|\delta _k(t_f)|=\frac{3}{5\pi }\frac{\mathrm{\Gamma }(\nu )}{(\nu \frac{3}{2})(H_0/m)}\left(\frac{2H_0}{k}\right)^{\nu \frac{3}{2}},$$
(235)
where $`(H_0/m)`$ determines the initial amplitude of the effective zero mode (218). We can now read the power spectrum per logarithmic $`k`$ interval
$$𝒫_s(k)=|\delta _k|^2k^{2(\nu \frac{3}{2})}.$$
(236)
leading to the index for scalar density perturbations
$$n_s=12\left(\nu \frac{3}{2}\right).$$
(237)
For $`H_0/m1`$, we can expand $`\nu 3/2`$ as a series in $`m^2/H_0^2`$ in eq.(235). Given that the comoving wavenumber of the mode which crosses the horizon $`n`$ e-folds before the end of inflation is $`k=H_0e^{(N_en)}`$ where $`N_e`$ is given by (209), we arrive at the following expression for the amplitude of fluctuations on the scale corresponding to $`n`$ in terms of the De Sitter Hubble constant and the coupling $`\lambda =8\pi ^2g`$:
$$|\delta _n(t_f)|\frac{9H^3}{10(2\pi )^{3/2}m^3}\left(2e^n\right)^{m^2/3H_0^2}\sqrt{\lambda }\left[1+\frac{2m^2}{3H_0^2}\left(\frac{7}{6}\mathrm{ln}2\frac{\gamma }{2}\right)+𝒪\left(\frac{m^4}{H_0^4}\right)\right].$$
(238)
Here, $`\gamma `$ is Euler’s constant. Note the explicit dependence of the amplitude of density perturbations on $`\sqrt{g}`$. For $`n60`$, the factor $`\mathrm{exp}(nm^2/3H_0^2)`$ is $`𝒪(100)`$ for $`H_0/m=2`$, while it is $`𝒪(1)`$ for $`H_0/m4`$. Notice that for $`H_0/m`$ large, the amplitude increases approximately as $`(H_0/m)^3`$, which will place strong restrictions on $`g`$ in such models.
We remark that we have not included the small corrections to the dynamics of the effective zero mode and the scale factor arising from the non-linearities. We have found numerically that these nonlinearities are only significant for the modes that cross about 60 e-folds before the end of inflation for values of the Hubble parameter $`H_0/m_R>5`$. The effect of these non-linearities in the large $`N`$ limit is to slow somewhat the exponential growth of these modes, with the result of shifting the power spectrum closer to an exact Harrison-Zeldovich spectrum with $`n_s=1`$. Since for $`H_0/m_R>5`$ the power spectrum given by (237) differs from one by at most a few percent, the effects of the non-linearities are expected to be observationally unimportant. The spectrum given by (235) is similar to that obtained in references although the amplitude differs from that obtained there. In addition, we do not assume slow roll for which $`(\nu \frac{3}{2})1`$, although this would be the case if $`N_e60`$.
We emphasize an important feature of the spectrum: it has more power at long wavelengths because $`\nu 3/2>0`$. This is recognized to be a consequence of the spinodal instabilities that result in the growth of long wavelength modes and therefore in more power for these modes. This seems to be a robust prediction of new inflationary scenarios in which the potential has negative second derivative in the region of field space that produces inflation.
It is at this stage that we recognize the consistency of our approach for separating the composite effective zero mode from the small fluctuations. We have argued above that many more than 60 e-folds are required for consistency, and that the small fluctuations correspond to those modes that cross the horizon during the last 60 e-folds of the inflationary stage. For these modes $`H_0/k=e^{H_0t^{}(k)}`$ where $`t^{}(k)`$ is the time since the beginning of inflation of horizon crossing of the mode with wavevector $`k`$. The scale that corresponds to the Hubble radius today $`\lambda _0=2\pi /k_0`$ is the first to cross during the last 60 or so e-folds before the end of inflation. Smaller scales today will correspond to $`k>k_0`$ at the onset of inflation since they will cross the first horizon later and therefore will reenter earlier. The bound on $`|\delta _{k_0}|\mathrm{\Delta }T/T10^5`$ on these scales provides a lower bound on the number of e-folds required for these type of models to be consistent:
$$N_e>60+\frac{12}{\nu \frac{3}{2}}\frac{\mathrm{ln}(\nu \frac{3}{2})}{\nu \frac{3}{2}},$$
(239)
where we have written the total number of e-folds as $`N_e=H_0t^{}(k_0)+60`$. This in turn can be translated into a bound on the coupling constant using the estimate given by eq.(209).
The four year COBE DMR Sky Map gives $`n1.2\pm 0.3`$ thus providing an upper bound on $`\nu `$
$$0\nu \frac{3}{2}0.05$$
(240)
corresponding to $`h_02.6`$. We then find that these values of $`h_0`$ and $`\lambda 10^{12}10^{14}`$ provide sufficient e-folds to satisfy the constraint for scalar density perturbations.
### F Tensor Metric Perturbations:
The scalar field does not couple to the tensor (gravitational wave) modes directly, and the tensor perturbations are gauge invariant from the beginning. Their dynamical evolution is completely determined by the dynamics of the scale factor. Having established numerically that the inflationary epoch is characterized by $`\dot{H}/H_0^21`$ and that scales of cosmological interest cross the horizon during the stage in which this approximation is excellent, we can just borrow the known result for the power spectrum of gravitational waves produced during inflation extrapolated to the matter era
$$𝒫_T(k)\frac{H_0^2}{M_{Pl}^2}k^0.$$
(241)
Thus the spectrum to this order is scale invariant (Harrison-Zeldovich) with an amplitude of the order $`m^4/\lambda M_{Pl}^4`$. Then, for values of $`m10^{12}10^{14}\text{ Gev }`$ and $`\lambda 10^{12}10^{14}`$ one finds that the amplitude is $`10^{10}`$ which is much smaller than the amplitude of scalar density perturbations. As usual the amplification of scalar perturbations is a consequence of the equation of state during the inflationary epoch.
## IX Conclusions
Since there are a number of articles in the literature treating related problems, it is useful to review the unique features of the present work. First, we have treated the problem dynamically, without using the effective potential (an equilibrium construct) to determine the evolution. Second, we have provided consistent non-perturbative calculations of the non-linear quantum field evolution to bring out some of the most relevant aspects of the late time behavior. In particular, we found that the quantum backreaction naturally inhibits catastrophic growth of fluctuations and provides a smooth transition to the late time regime in which the quantum fluctuations decay as the zero mode approaches its asymptotic state. Third, the dynamics studied obeys the constraint of covariant conservation of the energy momentum tensor.
It can be argued that the inflationary paradigm as currently understood is one of the greatest applications of quantum field theory. The imprint of quantum mechanics is everywhere, from the dynamics of the inflaton, to the generation of metric perturbations, through to the reheating of the universe. It is clear then that we need to understand the quantum mechanics of inflation in as deep a manner as possible so as to be able to understand what we are actually testing via the CMBR temperature anisotropies, say.
What we have found in our work is that the quantum mechanics of inflation is extremely subtle. We now understand that it involves both non-equilibrium as well as non-perturbative dynamics and that what you start from may not be what you wind up with at the end!
In particular, we see now that the correct interpretation of the non-perturbative growth of quantum fluctuations via spinodal decomposition is that the background zero mode must be redefined through the process of zero mode reassembly that we have discovered. When this is done (and only when!) we can interpret inflation in terms of the usual slow-roll approach with the now small quantum fluctuations around the redefined zero mode driving the generation of metric perturbations.
We have studied the non-equilibrium dynamics of a ‘new inflation’ scenario in a self-consistent, non-perturbative framework based on a large $`N`$ expansion, including the dynamics of the scale factor and backreaction of quantum fluctuations. Quantum fluctuations associated with superhorizon modes grow exponentially as a result of the spinodal instabilities and contribute to the energy momentum tensor in such a way as to end inflation consistently.
Analytical and numerical estimates have been provided that establish the regime of validity of the classical approach. We find that these superhorizon modes re-assemble into an effective zero mode and unambiguously identify the composite field that can be used as an effective expectation value of the inflaton field whose classical dynamics drives the evolution of the scale factor. This identification also provides the initial condition for this effective zero mode.
A consistent criterion is provided to extract small fluctuations that will contribute to cosmological perturbations from large non-perturbative spinodal fluctuations. This is an important ingredient for a consistent calculation and interpretation of cosmological perturbations. This criterion requires that the model must provide many more than 60 e-folds to identify the ‘small perturbations’ that give rise to scalar metric (curvature) perturbations. We then use this criterion combined with the gauge invariant approach to obtain the dynamics of the Bardeen variable and the spectrum for scalar perturbations.
We find that during the inflationary epoch, superhorizon modes of the Bardeen potential grow exponentially in time reflecting the spinodal instabilities. These long-wavelength instabilities are manifest in the spectrum of scalar density perturbations and result in an index that is less than one, i.e. a ‘red’ power spectrum, providing more power at long wavelength. We argue that this red spectrum is a robust feature of potentials that lead to spinodal instabilities in the region in field space associated with inflation and can be interpreted as an imprint of the phase transition on the cosmological background. Tensor perturbations on the other hand, are not modified by these features, they have much smaller amplitude and a Harrison-Zeldovich spectrum.
## X Acknowledgements:
We thank F. Cao, D. Cormier, R. Holman, S. P. Kumar, J. Salgado, A. Singh, M. Srednicki all of whom collaborated at different stages on the works reviewed here. We thank J. Baacke, C. Destri, A. Dolgov, E. Kolb, E. Weinberg for conversations and discussions. D. B. thanks the N.S.F for partial support through the grant awards: PHY-9605186 and INT-9815064 and LPTHE for warm hospitality. We thank the CNRS-NSF cooperation programme for partial support.
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# X-rays from the barred galaxy NGC 4303
## 1 Introduction
Barred galaxies constitute a major fraction of all disc galaxies classified in the optical, more than 50% including strong bars and intermediate morphologies (Sellwood & Wilkinson sel93 (1993)). This fraction increases when also near-infrared images are used for classification, thus underlining the importance for the general understanding of the evolution of galaxies. The non-axisymmetric potential has a strong impact on the gas dynamics and the star formation in barred systems. Observations reveal a correlation between the radial abundance gradient and the strength of the bar (Martin & Roy mar94 (1994); Friedli et al. fri94 (1994); Martinet & Friedli mar97 (1997)). This is interpreted as the result of two effects caused by the bar: a stronger radial gas flow and hence a stronger radial mixing of metals and the efficiency of star formation. The radial mass transfer concentrates gas near the galactic center and at the ends of the bar at corotation. Enhanced star formation is the consequence of gas accumulation. The rotating bar potential also heats up the outer disk parts which leads to larger stellar velocity dispersions and a radial diffusion of stars. (Sellwood & Wilkinson sel93 (1993)).
Galactic bars have also been considered to support the central infall of gas to feed a central ”monster”(e.g. Beck et al. bec99 (1999)). Several authors have claimed that active galactic nuclei (AGN) are more likely in barred galaxies than in non-barred ones (e.g. Simkin et al. sim80 (1980); Arsenault ars89 (1989)). Hummel et al. (hum90 (1990)) note that the fraction of central radio sources in barred spirals is by a factor of 5 higher than in non-barred spirals. Other authors doubt that there is a significantly higher number of bars in galaxies harboring an AGN (e.g. Balick & Heckman bal82 (1982); Ho et al. ho97 (1997)). It appears that the concentration of gas on a scale of $``$1 kpc at the galactic center required to enhance the central star formation can easily be achieved by a bar potential. It seems much more difficult, however, to accumulate enough gas on a scale of a few pc to tens of pc in order to produce an AGN . Other effects depending on the environment of the galaxies (interaction: Elmegreen et al. elm90 (1990); H i contents: Cayatte et al. cay90 (1990)) play an important role in mass distribution, gas flow, and therefore in the formation and evolution of bars and the star formation history in these systems.
One of the most famous, closest and most widely studied barred galaxies is NGC 4303 (M61), member of the Virgo Cluster, which is observed at an inclination of 27° (Guhathakurta et al. guh88 (1988)). Optical spectra of this galaxy indicate that it consists of a nuclear starburst and a LINER or Seyfert 2 nucleus (Filippenko & Sargent fil86 (1986); Kennicutt et al. ken89 (1989); Colina et al. col97 (1997); Colina & Arribas 1999, hereafter CA (99)). Indications for a high star formation rate (SFR) in NGC 4303 are given by the numerous H ii regions (Hodge & Kennicutt hod83 (1983); Martin & Roy 1992, hereafter MR (92)) and three observed supernovae (van Dyk van92 (1992)). It also shows strong radio emission distributed over the entire disk (Condon con83 (1983)). Colina et al. (col97 (1997)) and CA (99) discussed the question of a starburst–AGN connection in this barred galaxy, using optical spectroscopy and HST UV images. The data range from a nuclear spiral structure of massive star-forming regions with an outer radius of 225 pc down to the unresolved core of a size $``$ 8 pc. From the UV data it is not clear if the core is a massive stellar cluster or a pure AGN.
VLA observations (Cayatte et al. cay90 (1990)) show that NGC 4303 is not highly H i deficient, which can be explained by only slight environmental influences in the outermost region of the Virgo Cluster. The projected distance to M87 is 8$`\stackrel{}{.}`$2 (Warmels war88 (1988)). No significant difference of the abundance gradient in the disk of NGC 4303 compared to non-barred spiral galaxies has been observed (MR (92)). Martinet & Friedli (mar97 (1997)) discussed the abundance gradient slope in terms of bar age. According to them, a steep gradient in the bar and a flat one in the outer disk are typical for a young bar while a single flat gradient in bar and disk characterizes an old bar. MR (92) did not determine the gradient at large radii because of a small number of H ii regions. Martinet & Friedli (mar97 (1997)) also claimed that bars in late-type spirals with enhanced star formation like NGC 4303 are expected to be young.
Probable interaction companions are the nearby galaxies NGC 4303 A (Condon con83 (1983)) and NGC 4292 (Cayatte et al. cay90 (1990)), at distances of 7$`\stackrel{}{.}`$5 northwest and 10′ northeast, respectively.
## 2 Observations and data reduction
In this paper we present data from the High Resolution Imager (HRI) and the Position Sensitive Proportional Counter (PSPC) on board of the X-ray satellite ROSAT. This X-ray telescope was operating in the energy range of 0.1–2.4 keV. For details concerning ROSAT and its detectors see the ROSAT User’s Handbook (Briel et al. bri96 (1996)).
### 2.1 HRI
We proposed HRI pointed observations of NGC 4303 (ROSAT sequence number 600854, PI: N. Junkes), which were carried out in the time periods July 2–8, 1996, June 9–23, 1997, and January 6–7, 1998 with integration times of 16055 sec, 27538 sec, and 5600 sec, respectively (see Table LABEL:tabobs). All three observations are centered on RA(2000) = 12<sup>h</sup>21<sup>m</sup>55$`\stackrel{s}{.}`$2; Dec(2000) = +04°28′11$`\stackrel{}{.}`$5.
The three photon lists were combined to create one single image. The data were analyzed using the commercial interactive data language software package IDL (Interactive Data Language). All presented images contain pixel values with units of counts per pixel and second. The images were corrected for vignetting but not for background. To quantify any source counts three circular source-free fields with radii of 49″, 34″, and 39″, respectively, were selected to determine the background level (Table LABEL:tabobs). The averaged background flux amounts to 3.52$`\times `$10<sup>-7</sup> cts s<sup>-1</sup> arcsec<sup>-2</sup>.
Fig. 1 shows the HRI field of view (FOV) of the three combined data sets, convolved with a Gaussian of 5″ FWHM. For reasons of identification the 5$`\sigma `$ sources detected with a maximum likelihood method are numbered in Fig. 1 and listed in Table LABEL:tabdetsub. Several single source detections (no. 9–13) coincide with optical peaks in the galaxy. Source no. 18 can be identified as the QSO 1219+044, which served as the central source of the PSPC pointed observation used for spectral analysis in Sect. 2.2. Source no. 7 coincides spatially with the QSO 1219+047 (Bowen et al. bow96 (1996)). It turned out that the source is variable the X-rays. Its flux increased by a factor of 3 between the second and third HRI observation (June 1997 and January 1998).
### 2.2 PSPC
Since there exist no particular PSPC observations pointed on NGC 4303, we use ROSAT observations from the archive (ROSAT sequence number 701095, PI: R. Staubert). NGC 4303 lies in the FOV at an off-axis distance of 17′ which is still within the inner part of the supporting ring structure of the telescope. The position of the X-ray counterpart of NGC 4303 is given in Fig. 2. The central source is the QSO 1219+044 (no. 18 in Fig. 1). The image has been convolved with a Gaussian of 25″ FWHM. Due to strong asymmetry and broadening of the point spread function (PSF) with radial distance from the optical axis ($``$25″ FWHM for an on-axis point source at 1 keV; $``$67″ FWHM for a 17′ off-axis point source at 1 keV; Hasinger et al. has94 (1994)) we cannot get any useful spatial information from the PSPC data.
This PSPC exposure was carried out between December 24 and 26, 1992, with an integration time of 8135 sec. The background flux was determined from three circular source-free areas in the field with radii of 83″, 88″, and 75″ and amounts to 5.14$`\times `$10<sup>-7</sup> cts s<sup>-1</sup> arcsec<sup>-2</sup>. To analyze the PSPC data we used IDL and the software package XSPEC (Arnaud arn96 (1996)) for interactively fitting X-ray spectra.
For NGC 4303, we adopt the same distance as of M100, the brightest spiral in the Virgo Cluster (16.1 Mpc; Ferrarese et al. fer96 (1996)). Then the resolution of the ROSAT detectors of $``$5″ full width at half maximum (FWHM) for the HRI and $``$67″ FWHM for the PSPC at an off-axis distance of 17′ scale to 390 pc and 5.3 kpc, respectively.
## 3 Results
### 3.1 Spatial analysis
As already mentioned, no useful spatial information on the X-ray structure of NGC 4303 can be obtained from the PSPC data due to the spatial resolution of the detector, and moreover to the 17′ off-axis position of the source in the FOV.
In contrast, the more detailed HRI image reveals a number of X-ray sources distributed over the galactic disk (Fig. 4). In comparison to the numbered sources in Fig. 1 the closer view allows to distinguish more details. For example, source no. 9 in Fig. 1 splits into three X-ray spots (labeled A–C in Fig. 3). The most luminous source coincides with the center of NGC 4303 and dominates in the soft X-rays. The count rates and fluxes derived for sources from the HRI are listed in Table LABEL:tabhrisources. The corresponding areas are plotted in Fig. 3. To determine the fluxes we used the energy conversion factor (ECF) from the ROSAT Call For Proposals documentation. The ECF determines the ratio between count rates and unabsorbed source flux in the ROSAT band for given spectral parameters. For the disk sources A–F we assume a 0.3 keV Raymond-Smith model (Raymond & Smith ray77 (1977)) with an absorption column density of 3$`\times `$10<sup>20</sup> cm<sup>-2</sup>. For the nucleus a power law with $`\mathrm{\Gamma }`$=2.6 and a column density of 3$`\times `$10<sup>20</sup> cm<sup>-2</sup> (see Sect. 3.2 and Table LABEL:fittab) is applied as spectral model.
The contours of sources B, C, D, and F in Fig. 3 are all located within the optical arm structure and coincide with bright H$`\alpha `$ emission regions within the spiral arms (Fig. 4). In addition, source E is embedded in the faint outer part of the southwestern spiral arm. MR (92) distinguished 79 H ii regions in NGC 4303, mainly in the spiral arms. The X-ray contour overlay over the H$`\alpha `$ image in Fig. 4 reveals that the X-ray sources B, C, D, and F coincide with H ii regions, while sources A and E are located near such regions.
Gas dynamical models of barred galaxies (Englmaier & Gerhard eng97 (1997)) show strong gas accumulation at the tips of the bars due to corotation of the bar structure with the disk what should lead to enhanced star formation. H i observations as well as the existence of prominent H$`\alpha `$ features strikingly support the outcome of these models. The X-ray contours B and F seem to arise from these regions. The X-ray maximum D is connected with another interesting feature of NGC 4303: in the eastern part the galactic arm seems to be deformed to a boomerang-like bow where source D lies at the bend but without any significant brightening in H$`\alpha `$.
The lower X-ray contours of the nucleus indicate a possible extended source. Recent high-resolution UV observations of the central region with the Hubble Space Telescope reveal a spiral-shaped structure of massive young (2–3 Myr) star-forming regions with an outer radius of 225 pc (CA (99)). This structure cannot be resolved by HRI. Due to the low age of the star clusters almost no thermal X-ray emission is expected at the galactic nucleus (see Sect. 4). The extended X-ray contours may originate from additional sources at distances of about 1 kpc around the nucleus.
No X-ray emission has been detected from the possible interaction companions NGC 4303 A and NGC 4292.
### 3.2 Properties of the X-ray spectrum
Since the HRI maxima are separated by only 50″, and because of the reasons mentioned in Sect. 2.2, the PSPC observations do not allow to study the spectra of the X-ray components of NGC 4303 individually. ROSAT PSPC detected 505$`\pm `$24 backgroung-subtracted source counts from NGC 4303 in a total integration time of 8135 sec. The spectra of the source and the background are shown in Fig. 5.
We fitted the spectrum with several single-component models, as Bremsstrahlung (BS), Raymond-Smith model (RS), and a power law (PO), and a combined RS-PO model. The results are listed in Table LABEL:fittab.
A single power-law model implies the assumption, that the active nucleus of NGC 4303 dominates the X-ray emission. Furthermore, the sources detected by the HRI in the galactic disk would also have to be described with the same power law. The photon index in this model is $`\mathrm{\Gamma }`$=3.2$`\pm `$0.2. The emission of an AGN in the ROSAT energy band is best described by a power law with a photon index of $`\mathrm{\Gamma }`$2.4; nevertheless some cases have been observed with $`\mathrm{\Gamma }>`$3 (MCG –5-23-16: Mulchaey et al. mul93 (1993); Mkn 335: Turner et al. tur93 (1993)). High-mass X-ray binaries (HMXB) found in young star-forming regions in the spiral arms have a similar spectral shape in the 0.1–2.4 keV energy range with a photon index of $`\mathrm{\Gamma }`$2.7 (Mavromatakis mav93 (1993)). The column density of the absorbing component amounts to 5.7$`\times `$10<sup>20</sup> cm<sup>-2</sup>, which is by a factor of 3 higher than the Galactic foreground H i column density (Dickey & Lockman 1990, DL (90)). Nevertheless, self-absorption within NGC 4303 must be expected, and small-scale deviations from the observed Galactic value by DL (90) cannot be ruled out and may result in a higher absorption from the Milky Way. The resulting 0.1–2.4 keV X-ray luminosity amounts to 1.3$`\times `$10<sup>41</sup> erg s<sup>-1</sup>. The flux portion from the sources outside the nuclear region as observed with the HRI amounts to $``$1.4$`\times `$10<sup>40</sup> erg s<sup>-1</sup> in the case of a single power-law emission model with $`\mathrm{\Gamma }`$=3.2 using the corresponding ECF of 1$`\times `$10<sup>10</sup> cts cm<sup>2</sup> erg<sup>-1</sup>. Assuming a mean X-ray luminosity of 10<sup>37</sup> erg s<sup>-1</sup> for an HMXB, as observed in the Milky Way (Fabbiano et al. fab82 (1982); Watson wat90 (1990)) would require an unlikely high number of 1400 of these systems to produce the observed X-ray flux. The ratio of OB stars to HMXBs is assumed to be $``$500 (Fabbiano et al. fab82 (1982)). This means that a total number of 7$`\times `$10<sup>5</sup> OB stars would be required to account for the HMXB X-ray flux in NGC 4303. Even if we consider to have 10<sup>5</sup> OB stars in NGC 4303, as observed e.g. in Mkn 297 (Benvenuti et al. ben79 (1979)), it is still a factor of 7 higher than expected. Moreover, this is the required number only for the disk sources and would involve almost 1.5$`\times `$10<sup>7</sup> M in massive stars with a Salpeter IMF and, by this, would require a moderately high SFR of about 15 M yr<sup>-1</sup> in the disk. On the other hand, the corresponding supernova type ii (SN ii) rate (0.1 yr<sup>-1</sup>) should contribute to the X-ray emission via hot gas.
The single component models BS and RS show similar results. Consequently, we only achieve an adequate fit of RS with very low metallicity, e.g. the portion of emission lines to the spectrum is very small. In contrast, it is expected that emission lines of highly ionized elements, like Fe and Mg, should play an important role in the X-ray spectrum of SNe ii in starburst regions because of the nucleosynthesis of massive SN ii progenitor stars (Woosley & Weaver woo85 (1985)). In both models the column density is about 3$`\times `$10<sup>20</sup> cm<sup>-2</sup> and the plasma temperature is 0.6 keV. For the BS model we get a total X-ray luminosity of 4$`\times `$10<sup>40</sup> erg s<sup>-1</sup>, for the RS model it is 3.5$`\times `$10<sup>40</sup> erg s<sup>-1</sup>. RS models with different higher metallicities yield unacceptable fits.
The fit of the X-ray spectrum with a two-component model (RS+PO) is only slightly better than the one-component fits. Nevertheless, from the points mentioned above and the physical picture discussed in Sect. 4 this model serves as the best explanation for the observed soft X-ray emission. Hydrogen column density ($`N_\mathrm{H}`$=3.3$`\times `$10<sup>20</sup> cm<sup>-2</sup>) and power-law spectral index ($`\mathrm{\Gamma }`$=2.6) lie within the expected range (as discussed for the single power law above). The plasma temperature of 0.3 keV fits with the observed values of other galaxies (e.g. NGC 253: Forbes et al. for99 (1999); NGC 1808: Junkes et al. jun95 (1995)). The total 0.1–2.4 keV luminosity for this model amounts to 4.7$`\times `$10<sup>40</sup> erg s<sup>-1</sup> with 13% contribution from the RS component. The spectral fit together with the residuals is plotted in Fig. 6.
## 4 Discussion
From the quality of the spectral fits alone there is no significance for favoring a single-component model or a combination of two components. Due to the lack of any spatial information in the PSPC image there is no possibility to distinguish between different spatial and spectral components simultaneously. So only the combined information from the PSPC and HRI data allows a more detailed interpretation of the X-ray results.
Several points speak against the single PO component model, as discussed in Sect. 3.2. A more likely scenario is a composition of several different emission sources, like an active nucleus, HMXBs, and supernova remnants (SNRs). In the following we will therefore discuss a composite emission model and the comparison with the UV and optical observations of the galactic core.
### 4.1 The nucleus of NGC 4303
As can be discerned from the HRI image (Fig. 3 and Table LABEL:tabhrisources), most of the X-ray emission of NGC 4303 (83%) comes from the central region of the galaxy. Three different pictures are imaginable for the nucleus: a central active nucleus, a central or circumnuclear region with enhanced star formation, or a combination of both phenomena. Any of these cases requires a sufficient gas density at the galactic center. This can be achieved by a barred potential which triggers radial gas flow from the outer regions toward the nucleus. On the other hand, from numerical simulations including gas dynamics bar formation has proved to be only a transient phenomenon (Combes com00 (2000)). In this picture, the galactic bar will be destroyed by the gas inflow after only a few cycles. A new bar-phase can follow this gas infall due to a subsequent gravitational instability from the accreted central mass. The problem with this picture is the contradiction of a necessary gas inflow to form and feed any nuclear activity (starburst and/or AGN) and the fact that this gas inflow destroys the bar. It seems that a sufficiently massive black hole can provide for its fuelling (Fukuda et al. fuk98 (1998)). Another efficient way for gas to flow further into the center is a second smaller bar embedded into the first one due to a second inner Lindblad resonance (Friedli & Martinet fri93 (1993)). In some cases, a gaseous circumnuclear ring is formed at the end of the second bar.
The concentrated X-ray emission from the galactic core in NGC 4303 may originate from an AGN and/or a nuclear or circumnuclear starburst. A starburst can contribute in two different ways to the X-ray flux. First, the produced star population contains HMXBs, emitting an X-ray radiation in spectral shape similar to an AGN. HMXBs cannot be distinguished from AGN in the ROSAT data. Additionally, high-mass stars (above $``$8–10 M) evolve to SNe ii at an age of $``$10<sup>7</sup> yr, depending on their initial mass. The SNe ii from one star cluster form a cumulative expanding superbubble filled with hot gas which can be described by a thermal Bremsstrahlung spectrum and additional emission lines and recombination edges of highly ionized heavy elements produced in high-mass stars and released by their SN ii explosions, e.g. O, Ne, Mg, Si, and Fe. Theoretical models for the spectral emission of such a hot diffuse gas are MEKA (Mewe et al. mew85 (1985)) and the model by Raymond & Smith (ray77 (1977)), which we used in our spectral fit.
If we apply a two-component model to describe the X-ray spectrum of NGC 4303, the comparison of the flux ratio from the nucleus (83% in the HRI) and the disk sources (17% in the HRI) suggests that the central source can be described by the power-law component (87% in the spectral fit). The thermal RS emission exclusively originates from the disk sources, indicating ongoing star formation. To consider the possible extension of $``$25″ of the central source, as represented by the lowest contours, it is imaginable that a small fraction of the X-ray flux is emitted by a circumnuclear starburst at a distance of $``$1 kpc around the core. This would add a thermal component to the non-thermal X-ray nucleus. On the other hand, a fraction of the X-ray flux from the disk sources may come from HMXBs within these star forming regions.
Fuelling an AGN on scales of a few parsec at the center of the galaxy leads to the problem of reducing the angular momentum of the central gas by several orders of magnitude, as dynamical simulations show (Barnes & Hernquist bar91 (1991)). Concentration of gas in a ring-like feature around the nucleus with a radius of $``$1 kpc is dynamically much easier to achieve. The HRI image agrees with the picture of an extended X-ray source with a diameter of the order of $``$2 kpc at the galactic center of NGC 4303, explained by a circumnuclear starburst region with an additional possible compact nuclear source.The decovered massive rotating circumnuclear disk in NGC 4303 can provide by its spiral-like structure of massive star forming regions an effective mechanism to channel gas from the circumnuclear regions further down to the nucleus to feed the AGN. But one has to consider that the spiral structure in the UV has a diamater of only 225 pc, while the extension of the central X-ray source is about 2 kpc in diameter. The spiral feature detected in the UV cannot be resolved with the HRI.
From the analysis of UV and optical magnitudes and colors of the central 250 pc Colina & Wada (col00 (2000)) estimated ages of 5–25 Myr for the star-forming regions. Consequently a contribution to the central X-ray flux from SNRs and cumulatively expanding hot gas has to be expected. The question remains whether we observe a pure nucleus of massive star-forming clusters or a composition of these star clusters and a low luminous AGN. If NGC 4303 contains a non-thermal active nucleus, the X-ray luminosity of 4$`\times `$10<sup>40</sup> erg s<sup>-1</sup> points to only a low luminous AGN (LINER). Koratkar et al. (kor95 (1995)) found a correlation between $`L_\mathrm{X}`$ and $`L_{\mathrm{H}\alpha }`$ for low luminous AGNs of $`L_\mathrm{X}`$/$`L_{\mathrm{H}\alpha }`$ 14. Pérez-Olea & Colina (per96 (1996)) investigated the correlation between optical and X-ray luminosities of several AGNs with circumnuclear star-forming rings, pure AGNs, and pure starburst galaxies. The pure starbursts in their galaxy sample show $`L_\mathrm{X}`$/$`L_{\mathrm{H}\alpha }`$ values of 0.03–0.3, 100 times smaller than for pure AGNs. If we take the H$`\alpha `$ luminosity of NGC 4303 derived by Keel (kee83 (1983)) and assume that 10% originate from the nucleus, we get log $`L_{\mathrm{H}\alpha }`$(nucleus) = 39.2 (adopted for a distance of 16.1 Mpc). Therefore the X-ray-to-H$`\alpha `$ ratio amounts to log($`L_\mathrm{X}`$/$`L_{\mathrm{H}\alpha }`$) = 1.4, which agrees with the value found by Koratkar et al. (kor95 (1995)). Even the lower $`L_\mathrm{X}`$ value from a single RS model ($``$2.5$`\times `$10<sup>40</sup> erg s<sup>-1</sup> for the nucleus) results in log($`L_\mathrm{X}`$/$`L_{\mathrm{H}\alpha }`$) = 1.2. Typical pure starburst galaxies show H$`\alpha `$ luminosities of the order of their X-ray luminosities or higher.
### 4.2 The galactic disk
At first glance the optical disk of NGC 4303 seems to have the quite symmetrical morphology of a late-type spiral. A closer look reveals that the eastern spiral arm of the galaxy has a much more prominent form with a boomerang-like shape and a lot more bright emission regions than the western counterpart. The northern disk shows a complex structure with many separate features. This asymmetry is more discernible in the H$`\alpha `$ image. The H ii regions are mainly located in the northern part of the disk at the junction of the bar with the eastern spiral arm and along that arm. A close encounter of one or both of the nearby galaxies NGC 4303 A and NGC 4292 may have caused these features. The interaction within the Virgo Cluster is another possible source. Infall into the intracluster medium could cause ram pressure effects. Nevertheless, NGC 4303 is located at the outer edge of the cluster, which may produce only a moderate disturbance. This agrees with the H i distribution over the whole optical disk. Galaxies lying nearer to the cluster center show H i deficiencies and concentration of the neutral hydrogen in the central regions, indicating past interactions with the H i gas been stripped off from the outer disk regions.
As a striking indication for accumulation of gas in these regions, the X-ray sources A–C and F within the galactic disk of NGC 4303 are located at the ends of the bar. Gas dynamical simulations of barred galaxies have shown this accumulation due to mass flows along the bar to the center and to the ends of the bar, respectively (Noguchi nog88 (1988); Englmaier & Gerhard eng97 (1997)). The increased densities lead to enhanced star formation.
From the low inclination of NGC 4303 no direct information can be obtained whether the disk is warped or not. But it is striking that source D lies exactly at the bend of the eastern boomerang-shaped arm. This may indicate that this X-ray source is caused by the tidal force leading to a local gas concentration. Another indirect hint for a past interaction comes from the spectra of the QSO 1219+047 (source no. 7 in Fig. 1), a QSO whose line-of-sight penetrates the outer H i disk of NGC 4303. Bowen et al. (bow96 (1996)) detected complex Mg ii absorption, spanning a velocity range of $``$300 km s<sup>-1</sup>, despite the low inclined galactic disk. This high velocity is not fully understood. One possible explanation could be the result of interactions between NGC 4303 and the nearby companions.
### 4.3 Star formation in the disk
Table LABEL:tabhrisources lists the observed count rates and derived 0.1–2.4 keV luminosities for the single X-ray sources in the disk of NGC 4303. These include the assumptions of a Raymond-Smith plasma with solar abundances at a temperature of 0.3 keV. A power-law model would increase the values by a factor of 2.5. A rough estimation of the SFR in the disk from the X-ray luminosities is done by calculating the SN ii rate $`\nu _{\mathrm{SN}}`$ using a SNR model by Cioffi (cio90 (1990)), and assuming a Salpeter IMF within a mass interval from 0.1 M to 100 M and with all stars with masses above 8 M evolving to SNe ii. According to Cioffi, a SNR expanding into an ISM with a density of 1 cm<sup>-3</sup> radiates a total energy of 4.7$`\times `$10<sup>49</sup> erg in the soft X-ray regime above 0.1 keV for a time of $``$10<sup>4</sup> yr. Norman & Ikeuchi (nor89 (1989)) investigated the cumulative effect of a number of SNRs. The total SFR in the disk of NGC 4303 from the X-ray luminosity amounts to 0.5 M yr<sup>-1</sup>. This however is just a very simple estimation, containing several simplifications, as e.g. the sum of the single SN model from Cioffi (cio90 (1990)) for several cumulatively expanding SNRs in an evolving OB assoziation, or the derivation of the total disk SFR from single X-ray sources. A more detailed determiation of the SFR, for example using an analytic suberbubble model by Suchkov et a. (suc94 (1994)), would need information about the extensions and expansion time of the superbubbles, in order to determine the mechanical energy release by the SNe and, by this, the SN rate. This can be compared with the observed X-ray luminosity.
Besides the mentioned restrictions the difference between the SFR estimated from the X-ray flux and the SFR derived by the H$`\alpha `$ flux by Kennicutt (ken83 (1983)) (14 M yr<sup>-1</sup>) may be due to several reasons. Kennicutt used a Miller-Scalo IMF which increases the SFR by a factor of about 1.5. The total H$`\alpha `$ luminosity underlies a distance determination with a Hubble constant of 50 km s<sup>-1</sup> Mpc<sup>-1</sup>. Taking a radial velocity of 1569 km s<sup>-1</sup> (de Vaucouleurs et al. dev91 (1991)) leads to a 3.8 times higher luminosity than taking the distance of 16.1 Mpc which we used. Additionally, the fact that not all H ii regions may emit an adequate X-ray flux or are not strong enough to be detected lowers the estimated SFR from the X-ray luminosity which implies the existence of high-mass stars having been evolved to SNe ii. Possible X-ray emission from diffuse hot gas within the disk may lie below the detection limit of 1.1$`\times `$10<sup>-6</sup> cts s<sup>-1</sup> arcsec<sup>-2</sup> (5$`\sigma `$ above background level). A very faint component located at the spiral arms can be seen in outlines in Fig. 3, but is not detected at a 3$`\sigma `$ level. This limit corresponds to an X-ray flux of 4.6$`\times `$10<sup>-17</sup> erg s<sup>-1</sup> cm<sup>-2</sup> arcsec<sup>-2</sup> (ECF for a RS model as in Table LABEL:tabhrisources). Kennicutt (ken83 (1983))also admitted to treat the derived H$`\alpha `$ flux and resulting SFR with extreme caution because of only moderate accuracy due to possibly strong extinction effects. The H$`\alpha `$ flux derived by Keel (kee83 (1983)) is by a factor of 30 lower than the one derived by Kennicutt, after adopting the same distance.
Strikingly, the sources B and F both coincide with some of the most H$`\alpha `$ luminous H ii regions (Sources no.s 27 and 69 with log $`F_{\mathrm{H}\alpha }`$=$``$12.12 and log $`F_{\mathrm{H}\alpha }`$=$``$11.99, respectively, in MR (92)). Depending on the fraction of the central X-ray flux steming from SNRs and superbubbles or from an AGN component, the SFR for the core is of the order of 1 M yr<sup>-1</sup>.
## 5 Conclusions
We have analyzed spatial and spectral data from the barred late-type spiral galaxy NGC 4303 in the soft X-ray regime. Several separate X-ray sources can be observed in the core and disk of the galaxy. The locations of the sources correspond to several H ii regions and indicate a concentration of gas at the center and at the ends of the galactic bar, in agreement with numerical simulations of gas dynamics in a barred potential.
The low spatial resolution of the PSPC observation of NGC 4303 does not allow a distinction of several individual X-ray sources within the object. The best fit of the soft X-ray spectrum taking into account the information from the high resolution HRI observation is a combination of a RS component with a temperature of 0.3 keV and a power-law component with a spectral index of 2.6. The total 0.1–2.4 keV X-ray luminosity amounts to 4.7$`\times `$10<sup>40</sup> erg s<sup>-1</sup>, in agreement with other comparable barred galaxies with a nuclear starburst, like e.g. NGC 4569 (Tschöke et al. in preparation). A pure starburst model for the nucleus of NGC 4303 would require a special explanation for the unusually high $`L_\mathrm{X}`$/$`L_{\mathrm{H}\alpha }`$ ratio.
The combination of the flux fraction of the separate sources, the spectral information, and the comparison with the H$`\alpha `$ luminosity from the core leads to the following picture: the soft X-ray emission originates from a composition of several distinct emission regions. The central source consists of a low luminous AGN and a circumnuclear starburst. The disk sources are dominated by SNRs and superbubbles in star forming regions preferably at the ends of the bar and along the eastern spiral arm. Several HMXBs may contribute to the X-ray flux.
The disk X-ray sources are coincident with some of the most luminous H ii regions in the galaxy. The estimated total SFR from the X-ray flux is 1–2 M yr<sup>-1</sup>. Most H ii regions are not detectable in the X-ray, like most H$`\alpha `$ sources in the eastern boomerang-shaped arm. The enhanced star formation in NGC 4303 may have been caused by some kind of interaction although the H i morphology of the galaxy does not support very strong perturbation. If a dwarf galaxy has fallen in and merged with NGC 4303 in the past, the bar may have been produced with the subsequent triggering of the star formation at the center and in the spiral arms. The accreted dwarf galaxy would be resolved and not directly detectable.
###### Acknowledgements.
The authors are grateful to Dominik Bomans for stimulating discussions, and to Dr. Olga Sil’chenko for her substantial and constructive report. The ROSAT project is supported by the German Bundesministerium für Bildung, Wissenschaft, Forschung und Technologie (BMBF) and the Max-Planck-Society. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, Caltech, under contract with the NASA. Observations made with the NASA/ESA Hubble Space Telescope were used, obtained from data archive at STScI. STScI is operated by the Association of Universities for Research in Astronomy, Inc. (AURA) under the NASA contract NAS 5-26555.
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# 1 Introduction
## 1 Introduction
In the work of Lehmann, Newton, and Wu , the Kobayashi-Maskawa matrix is expressed in terms of the masses of the three generations of quarks:
$$\left(\begin{array}{c}u\\ d\end{array}\right)\left(\begin{array}{c}c\\ s\end{array}\right)\left(\begin{array}{c}t\\ b\end{array}\right).$$
(1.1)
This is accomplished by introducing a new horizontal symmetry. Some of the earlier attempts in this direction are given in .
Recent experiments at Super-Kamiokande \[??\] indicate the presence of neutrino oscillations, which would imply that the neutrinos are not all massless. If it is accepted that the neutrinos are not massless, then it is most natural in the Standard Model to introduce three right-handed neutrinos in addition to the three known left-handed ones. In this way, there are six quarks and six leptons.<sup>1</sup><sup>1</sup>1As shown in , neutrino-oscillation experiments cannot distinguish between massive Majorana and Dirac neutrinos. In this paper, the consequences of a universal quark-lepton mixing are studied. In other words, the method of is used to express the lepton KM matrix and the neutrino mixing matrix in terms of the masses of the three generations of leptons:
$$\left(\begin{array}{c}\nu _e\\ e\end{array}\right)\left(\begin{array}{c}\nu _\mu \\ \mu \end{array}\right)\left(\begin{array}{c}\nu _\tau \\ \tau \end{array}\right).$$
(1.2)
Of course the masses of the three charged leptons are accurately known, leaving as unknown parameters the masses of the three neutrinos. Thus there are three parameters to be determined instead of seven, the three masses plus the four in the lepton KM matrix.
It is the purpose of this paper to use the data from solar neutrinos and atmospheric neutrinos to determine the three neutrino masses separately, not only the differences of their squares.
## 2 The mixing matrix
The quark mixing matrix proposed by Lehmann et al. relates the mixing to the actual quark masses. These mass matrices are each given by four parameters $`a`$, $`b`$, $`c`$, and $`d`$, and, when applied to the leptons without modification, take the form
$`M(\mathrm{})`$ $`=`$ $`\left(\begin{array}{ccc}0& d(\mathrm{})& 0\\ d(\mathrm{})& c(\mathrm{})& b(\mathrm{})\\ 0& b(\mathrm{})& a(\mathrm{})\end{array}\right),`$ (2.1)
$`M(\nu )`$ $`=`$ $`\left(\begin{array}{ccc}0& id(\nu )& 0\\ id(\nu )& c(\nu )& b(\nu )\\ 0& b(\nu )& a(\nu )\end{array}\right),`$ (2.2)
with
$$b^2(\mathrm{})=8c^2(\mathrm{}),b^2(\nu )=8c^2(\nu ).$$
(2.3)
The diagonalization of these mass matrices is achieved by the orthogonal matrices $`R(\mathrm{})`$ and $`R(\nu )`$, where
$`M(\mathrm{})`$ $`=`$ $`R(\mathrm{})M_{\mathrm{diag}}(\mathrm{})R^\mathrm{T}(\mathrm{}),`$ (2.4)
$`M(\nu )`$ $`=`$ $`\mathrm{diag}(i,1,1)R(\nu )M_{\mathrm{diag}}(\nu )R^\mathrm{T}(\nu )\mathrm{diag}(i,1,1).`$ (2.5)
If there is $`CP`$ violation in the lepton sector, then the imaginary entries in (2.2) and (2.5) are required. If by any chance there is no $`CP`$ violation in the lepton sector, then both $`i`$ and $`i`$ there should be replaced by 1.
The diagonal mass matrices have the form
$$M_{\mathrm{diag}}=\left(\begin{array}{ccc}\lambda _1& 0& 0\\ 0& \lambda _2& 0\\ 0& 0& \lambda _3\end{array}\right)=\left(\begin{array}{ccc}m_1& 0& 0\\ 0& m_2& 0\\ 0& 0& m_3\end{array}\right),$$
(2.6)
where $`\lambda _2<0`$ and
$$m_1m_3.$$
(2.7)
We shall here only be concerned with $`M(\nu )`$ and $`R(\nu )`$.
In dealing with quarks , the observed quark masses allow a much stronger form for the inequality (2.7), namely
$$m_1m_2m_3.$$
(2.8)
The lack of direct experimental data on the neutrino masses implies that (2.7) can be used, but not (2.8). The first task is therefore to determine the allowed region in the space $`(m_1,m_2,m_3)`$, which must be between those permitted by (2.7) and (2.8).
The parameters $`a`$, $`b`$, $`c`$, and $`d`$ are related to the masses by the following conditions,
$`a+c`$ $`=S_1`$ $`=m_3m_2+m_1,`$
$`8c^2+d^2ac`$ $`=S_2`$ $`=m_3m_2m_3m_1+m_2m_1,`$
$`ad^2`$ $`=S_3`$ $`=m_1m_2m_3.`$ (2.9)
The cubic equation for the parameter $`a`$ is then
$$9a^317S_1a^2+(8S_1^2+S_2)aS_3=0.$$
(2.10)
Any real cubic equation can have either one or three real solutions. Where there is one real solution, that one is negative, and thus unphysical, as is seen from (2). Where there are three real solutions, one of them is negative, while two are positive. We shall refer to these two positive solutions as Solution 1 (larger $`a`$) and Solution 2 (smaller $`a`$).
These considerations can be used to determine the allowed physical region in the $`(m_1/m_3,m_2/m_3)`$ plane, as shown in Fig. 1. This region is only slightly larger than the triangle given by the inequality (2.8), with two additional regions, one where $`m_2>m_3`$ and the other a very small one with $`m_1>m_2`$.
## 3 The three-family MSW mechanism
The coupled equations satisfied by the three neutrino wave functions are
$$i\frac{\text{d}}{\text{d}r}\left(\begin{array}{c}\varphi _1(r)\\ \varphi _2(r)\\ \varphi _3(r)\end{array}\right)=\left[\left(\begin{array}{ccc}D(r)& 0& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)+\frac{1}{2p}\left(\begin{array}{ccc}M_{11}^2& M_{12}^2& M_{13}^2\\ M_{21}^2& M_{22}^2& M_{23}^2\\ M_{31}^2& M_{32}^2& M_{33}^2\end{array}\right)\right]\left(\begin{array}{c}\varphi _1(r)\\ \varphi _2(r)\\ \varphi _3(r)\end{array}\right),$$
(3.1)
where $`D(r)=\sqrt{2}G_\mathrm{F}N_e(r)`$, with $`G_\mathrm{F}`$ the Fermi weak-interaction constant and $`N_e(r)`$ the solar electron density at a distance $`r`$ from the center of the sun. Furthermore, we denote $`\nu _e=\varphi _1`$, $`\nu _\mu =\varphi _2`$, $`\nu _\tau =\varphi _3`$.
The evolution of the neutrino wave functions is determined by the squared mass matrix,
$$[M(\nu )]^2=\left(\begin{array}{ccc}d^2& cd& bd\\ cd& b^2+c^2+d^2& b(a+c)\\ bd& b(a+c)& a^2+b^2\end{array}\right)\left(\begin{array}{ccc}M_{11}^2& M_{12}^2& M_{13}^2\\ M_{21}^2& M_{22}^2& M_{23}^2\\ M_{31}^2& M_{32}^2& M_{33}^2\end{array}\right),$$
(3.2)
the neutrino momentum, $`p`$, and the solar electron density. Here, $`M_{ij}^2(M^2)_{ij}`$. The eigenvalues of the squared mass matrix (multiplied by $`r_0/2p`$, with $`r_0`$ defined below) are denoted $`\mu _1`$, $`\mu _2`$, and $`\mu _3`$, and ordered such that
$$\mu _1\mu _2\mu _3.$$
(3.3)
It is actually a good approximation to take an exponential electron density, $`N_e(r)=N_e(0)\mathrm{exp}(r/r_0)`$. A fit to the solar density as given by leads to $`r_0=6.983\times 10^4`$ km. For this case of an exponential solar density, the three-component wave equation can be solved in terms of generalized hypergeometric functions, $`{}_{2}{}^{}F_{2}^{}`$ .
We scale and shift the radial variable, $`u=r/r_0+u_0`$, with $`u_0`$ determined such that
$$D(0)r_0e^{u_0}=1.$$
(3.4)
The above equation (3.1) may then be written as
$$i\frac{\text{d}}{\text{d}u}\left(\begin{array}{c}\psi _1\\ \psi _2\\ \psi _3\end{array}\right)=\left(\begin{array}{ccc}\omega _1+e^u& \chi _2& \chi _3\\ \chi _2& \omega _2& 0\\ \chi _3& 0& \omega _3\end{array}\right)\left(\begin{array}{c}\psi _1\\ \psi _2\\ \psi _3\end{array}\right),$$
(3.5)
where $`\omega _2`$ and $`\omega _3`$ denote the eigenvalues of the $`2\times 2`$ mass matrix involving $`\psi _2`$ and $`\psi _3`$.
With $`x=e^u`$, one finds for $`\psi _1`$ a differential equation of the form
$$\left[\underset{j=1}{\overset{3}{}}\left(x\frac{\text{d}}{\text{d}x}i\mu _j\right)ix\underset{j=2}{\overset{3}{}}\left(x\frac{\text{d}}{\text{d}x}i\omega _j+1\right)\right]\psi _1=0$$
(3.6)
the solutions of which,
$$\psi _1(u)=\underset{j=1}{\overset{3}{}}C_j\psi _1^{(j)}(u),$$
(3.7)
can be given in terms of generalized hypergeometric functions $`{}_{2}{}^{}F_{2}^{}`$ as
$$\psi _1^{(j)}(u)=e^{i\mu _ju}{}_{2}{}^{}F_{2}^{}\left[\begin{array}{cc}1i(\omega _2\mu _j),& 1i(\omega _3\mu _j)\\ 1i(\mu _k\mu _j),& 1i(\mu _{\mathrm{}}\mu _j)\end{array}|ie^u\right]$$
(3.8)
with $`k,\mathrm{}=1,2,3`$, $`kj`$, $`\mathrm{}j`$. For the other flavours, $`\psi _2(u)`$ and $`\psi _3(u)`$ are given by similar expressions, with shifted parameters.
There is not much information about these functions $`{}_{2}{}^{}F_{2}^{}`$ in the literature . For the case of two flavours, the products in (3.6) go only up to $`j=2`$, and a familiar confluent hypergeometric function $`{}_{1}{}^{}F_{1}^{}`$ (also denoted Whittaker function or parabolic cylinder function) is obtained .
In order to impose the boundary conditions that only electron neutrinos are produced in the sun, we have to determine these functions at large and negative values of $`u`$. The series expansion is in principle convergent, but it is not practical for large absolute values of both parameters and the argument. Instead, methods have been developed for the evaluation by a combination of relating them to $`{}_{3}{}^{}F_{1}^{}`$ and asymptotic methods .
## 4 Results
Since the only unknown parameters for the present theory of universal quark-lepton mixing are the three neutrino masses, it remains to determine these masses using experimental data on solar and atmospheric neutrinos.
For the atmospheric neutrinos we take the 8 data points for $`\nu _\mu `$ and the 8 data points for $`\nu _e`$, as reported recently . These sixteen data points are treated as separate inputs, but we allow two overall normalization constants for the two sets of data.
For the solar-neutrino data, we use the total rates from the Chlorine experiment , the Gallium experiments (we average the two results) and the Super-Kamiokande experiment . We adopt the neutrino energy spectra and detector efficiencies as given by Bahcall et al., and, for the latter detector, we also include the neutral-current cross section .
The determination of the three neutrino masses is to be carried out by a $`\chi ^2`$ fit to these pieces of data. Since the number of degrees of freedom is 14, a good $`\chi ^2`$ fit would give support to this idea that the mass mixing is universal for quarks and leptons. For this purpose, we have scanned the entire ($`m_1`$, $`m_2`$, $`m_3`$) parameter space. It is the necessity to cover this entire space that makes it essential to develop the theory described in Sec. 3 here.
It is found that there are three minima. They are:
(A) Solution 1: (B) Solution 2: (C) Solution 2:
$`m_1=0.0016\mathrm{eV}`$ $`m_1=0.0011\mathrm{eV}`$ $`m_1=0.0011\mathrm{eV}`$
$`m_2=0.013\mathrm{eV}`$ $`m_2=0.013\mathrm{eV}`$ $`m_2=0.013\mathrm{eV}`$
$`m_3=0.034\mathrm{eV}`$ $`m_3=0.058\mathrm{eV}`$ $`m_3=0.094\mathrm{eV}`$
$`\text{with }\chi ^2=31`$ $`\text{with }\chi ^2=23`$ $`\text{with }\chi ^2=20`$
A few simple conclusions can be drawn from this set of mass values. First, with 14 degrees of freedom, the three values of $`\chi ^2`$ must be considered to be quite good. This is evidence in favour of universal quark-lepton mixing. In our mind the difference in these three values of $`\chi ^2`$ is not significant. Secondly, the three values for $`m_2`$ are the same—a surprise to us. This presumably means that the relative accuracy in the present determination of $`m_2`$ is better than those of $`m_1`$ and $`m_3`$. This can be understood in the following way. The relative accuracy for $`m_1`$ is low because $`m_1`$ itself is quite small, and the value of $`\chi ^2`$ is not very sensitive to such small masses. The average of these three $`m_1`$ is given in the abstract, and very roughly we guess the error to be about a factor of two. On the other hand, the uncertainly in the value of $`m_3`$ is simply due to the three minima giving different values. It is a matter of which one of the three minima is physically the correct one.
For completeness, we give the rotation matrix $`R(\nu )`$. The Euler angles $`(\theta ,\varphi ,\psi )`$ are: ($`22{}_{}{}^{},28{}_{}{}^{},8^{}`$) for (A), ($`35{}_{}{}^{},21{}_{}{}^{},5^{}`$) for (B), and ($`31{}_{}{}^{},19{}_{}{}^{},3^{}`$) for (C).
To demonstrate the scanning of the entire region of allowed values of neutrino masses, we show in Fig. 2 an example of such a scan showing contours of constant $`\chi ^2`$. In this figure, the value of $`m_3`$ used is 0.058 eV, corresponding to the minimum (B) above. It is from such scans that we know there are only three minima. In this figure, some minor irregularities along the edges of the allowed region are artifacts of the finite grid spacing.
It is difficult to compare the present result with the previously given allowed regions in the parameter space. The reason is that the allowed regions have typically been given on the basis of the mixing of two neutrino species, while there is significant mixing among all three neutrinos, or the studies concentrate on either solar or atmospheric neutrinos . The following comparisons are nevertheless of interest.
(a) It is probably correct to compare the $`m_3^2m_2^2`$ here with the $`\delta m^2`$ from atmospheric-neutrino data. In this case, the $`m_3^2m_2^2`$ for the three minima covers roughly the allowed range of $`\delta m^2`$, with minimum (B) near the center, and minima (A) and (C) near the edges of the allowed region.
(b) It is probably not too far wrong to identify $`m_2^2m_1^2`$ here with the $`\delta m^2`$ from solar-neutrino data. If so, the values are reasonably close, except with the so-called “low-mass–low-probability” (LOW) solution .
(c) It is more difficult to discuss the strength of the coupling. It can nevertheless be concluded, in the context of solar neutrinos, that the present solutions are significantly closer to the “large-mixing-angle” (LMA) solution than to the “small-mixing-angle” (SMA) solution.
In connection with the problem of distinguishing between the three minima, more accurate data from Super-Kamiokande and related experiments are needed. Also, a better understanding of the high-energy $`hep`$ neutrino flux could make the electron recoil energy spectrum useful for this purpose. Furthermore, important information may be forthcoming also from the long-baseline experiments and exotic atom experiments .
Acknowledgments
We are greatly indebted to Professor Harry Lehmann and Professor Jack Steinberger for the most helpful discussions. We would also like to thank Dr. Steve Armstrong, Professor Alvaro de Rujula, Dr. Conrad Newton, Professor Gabriele Veneziano, and especially Geir Vigdel, for useful discussions. One of us (TTW) wishes to thank the Theory Division of CERN for its kind hospitality.
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# Deformation of surfaces, integrable systems and Chern-Simons theory
## 1 Introduction
Many authors have extensively studied the deep relations among completely integrable systems and the basic equations of the differential geometry, like the Frenet formulae defining curves embedded in $`𝐑^3`$, or their analogous for the surfaces, the Gauss-Weingarten (GW) equations, and the corresponding integrability conditions, i.e. the Gauss-Mainardi-Codazzi (GMC) equations (see for instance ). In these approaches the main idea is to add to a generic differential geometry setting certain auxiliary conditions, containing from the beginning the properties of the completely integrable systems.
A slightly different situation occurs in the study of the so-called Darboux system , which naturally arises in looking for classes of orthogonal curvilinear coordinates in Euclidean spaces and whose integrability has been detected in Ref. 6. Such a system has been investigated mainly in connection with the Topological Field Theory . On the other hand, some years ago some of us proposed a simple method to obtain completely integrable systems in ($`2+1)`$-dimensions, from classes of non Abelian Chern-Simons (CS) field theories, taking values in Hermitian symmetric spaces . In this context completely integrable systems are seen as particular gauge choices in which the theory is formulated. Moreover, linear spectral problems are naturally related to the geometrical constraints imposed on the target space. From this point of view, integrable systems arise as special reductions, which break the general covariance and the gauge invariance of the original field theory, but preserve a residual symmetry in order to allow the Lax representation and the complete integrability, although the solvability is lost.
In the present work we show that this approach can be naturally implemented by resorting to a geometrical interpretation of the completely integrable systems mentioned above. Precisely, we show as the CS equation of motion can describe time evolving 2-dimensional surfaces in such a way that the deformation is not only locally compatible with the GMC equation, but completely integrable as well. The nature and the properties of such relationships are investigated together with the most important consequences. Furthermore, explicit examples of integrable deformations of surfaces are displayed.
The paper is organized as follows. Section 2 contains some results on the CS theory. In Section 3 the fundamental terminology and notations related to the theory of 2-dimensional surfaces are reviewed. In Section 4 the general formulation of deformation of 2-dimensional surfaces is presented. Section 5 is devoted to the analysis of certain spin models in (2+1)-dimensions. Section 6 is addressed to the bilinear representations of the spin systems fields and of the trihedral moving frame. Sections 7 and 8 concern with the deformations of surfaces from integrable (2+1)-dimensional spin systems and equations of the nonlinear Schrödinger type, respectively. In Sections 9 and 10 some solutions and special surfaces associated with spin system vortices are considered. Finally, in Section 11 some concluding remarks are reported.
## 2 Chern-Simons theory and completely integrable systems
Here we shall review some preliminaries concerning the CS theory and we show how one can connect them to the completely integrable systems in 2+1 dimensions. We are dealing with the field theory defined by the action
$$S\left[J\right]=\frac{k}{4\pi }_MTr\left(JdJ+\frac{2}{3}JJJ\right),$$
(1)
where $`J`$ is a 1-form gauge connection taking values in a simple Lie algebra $`\widehat{g}`$ on an oriented closed three-manifold $`M`$, and $`k`$ is a coupling constant which should be quantized in a quantum theory . The related classical equation of motion is the zero - curvature condition
$$FdJ+JJ=0.$$
(2)
The action (1) is manifestly invariant under general coordinate transformations (preserving orientation and volumes). Moreover, under a generic gauge map $`G:M\widehat{G}`$ the gauge connection transforms as $`JG^1JG+G^1dG`$ . Correspondingly, the action (1) changes as $`S[J]S[J]+2\pi kW(G),`$ where
$$W(J)=\frac{1}{24\pi ^2}_MTr(G^1dGG^1dGG^1dG)$$
(3)
is the winding number of the map $`G`$ and takes integer values, because of the result $`\pi _3(\widehat{G})=𝐙`$ from the homotopy theory . This is a Topological Field Theory in the sense that it possesses quantum observables, which are independent of the metric and are related to the Jones polynomials of the knot theory . From other points of view, the action (1) has been used as an effective interaction for quasi-particles and vortices in two space dimensional systems, of interest in the physics of high temperature superconductivity , and in the context of the low dimensional gravity models (see and reference therein). In the static self-dual reductions the system of equations (2) becomes the two - dimensional Toda field theory , the static reductions of the Ishimori model or of the Davey - Stewartson equation .
In the general action (1) was reduced assuming that the Lie algebra $`\widehat{g}`$ admits a $`𝐙_2`$ graduation, in such a way the form $`J`$ splits in two parts, taking values on an isotropy subalgebra and a complement linear space, respectively. The former component will play the role of a gauge field with the isotropy group as a gauge group, the latter could be considered as a sort of coupled ”matter” field. At any point of the corresponding coset space we can introduce in a natural way a Riemannian torsion free connection . Furthermore, the 3-manifold $`M`$ is trivialized into $`\mathrm{\Sigma }\times 𝐑`$, where $`\mathrm{\Sigma }`$ is a Riemann surface, endowed with a set of local complex coordinates $`z=x^1+ix^2,`$ and $`𝐑`$ is interpreted as the time axis. Thus the connection $`J`$ can be further decomposed into time and space-like components. In the simplest $`SU\left(2\right)/U(1)CP^1S^2`$ case the connection takes the form
$$J=\left(\begin{array}{cc}iv_\mu \hfill & q_\mu ^{}\hfill \\ q_\mu \hfill & iv_\mu \hfill \end{array}\right)dx^\mu ,$$
which can be rewritten in terms of complex forms
$$vdz=\frac{1}{2}\left(v_1iv_2\right)dz,\psi _{}dz=\frac{1}{2}\left(q_1iq_2\right)dz,\psi _+d\overline{z}=\frac{1}{2}\left(q_1+iq_2\right)d\overline{z}.$$
This allows us to write the action (1) as follows
$$S=\frac{k}{\pi }_{\mathrm{\Sigma }\times 𝐑}\{\frac{1}{2}\epsilon ^{\lambda \mu \nu }v_\lambda _\mu v_\nu +\frac{i}{2}(\psi _+^{}D_0\psi _+\psi _+(D_0\psi _+)^{}\psi _{}^{}D_0\psi _{}+\psi _{}(D_0\psi _{})^{})$$
$$iq_0^{}(D\psi _+\stackrel{\mathrm{\_}}{D}\psi _{})+iq_0(D\psi _+\stackrel{\mathrm{\_}}{D}\psi _{})^{}\}dx^0dx^1dx^2,$$
(4)
where $`D_0=_02iv_0`$, $`D=_z2iv`$, $`\overline{D}=\overline{_z}2iv^{}`$ ( denotes the complex conjugation). The first order Lagrangian involved in (4) is constrained by the torsion-free condition
$$D\psi _+\overline{D}\psi _{}=0$$
(5)
and by what we call the Gauss - Chern - Simons (GCS) law
$$_zv^{}_{\overline{z}}v=i\left(\left|\psi _+\right|^2\left|\psi _{}\right|^2\right),$$
(6)
enforced by the Lagrangian multipliers $`q_0`$ and $`v_0`$, respectively. Of course, here we are looking in a different way to a subset of the equations of motion (2), in which the general covariance is broken. Indeed, only the isotropic $`U(1)`$ invariance is left. Furthermore, by exploiting the local isomorphism between $`so(3)`$ and $`su(2)`$ realized by the adjoint representation of the connection
$$J^{(ad)}=\left(\begin{array}{ccc}0\hfill & v_0\hfill & Re(q_0)\hfill \\ v_0\hfill & 0\hfill & Im(q_0)\hfill \\ Re(q_0)\hfill & Im(q_0)\hfill & 0\hfill \end{array}\right)dx^0+$$
$$\left(\begin{array}{ccc}0\hfill & 2Re(v)\hfill & Re(\psi _++\psi _{})\hfill \\ 2Re(v)\hfill & 0\hfill & Im(\psi _++\psi _{})\hfill \\ Re(\psi _++\psi _{})\hfill & Im(\psi _++\psi _{})\hfill & 0\hfill \end{array}\right)dx^1+$$
$$\left(\begin{array}{ccc}0\hfill & 2Im(v)\hfill & Im(\psi _+\psi _{})\hfill \\ 2Im(v)\hfill & 0\hfill & Re(\psi _+\psi _{})\hfill \\ Im(\psi _+\psi _{})\hfill & Re(\psi _+\psi _{})\hfill & 0\hfill \end{array}\right)dx^2,$$
we are able to introduce the so-called moving trihedral frame $`\{𝐞_i\}`$ in $`𝐑^3`$ , which satisfies the orthonormal conditions
$$𝐞_i𝐞_j=\delta _{ij},𝐞_i𝐞_j=\epsilon _{ijk}𝐞_k$$
(7)
and changes accordingly to
$$_\mu 𝐞_i=(J_\mu )_{ik}^{(ad)}𝐞_k,\mu =0,1,2,i=1,2,3.$$
(8)
Its integrability is assured by the zero curvature condition, namely by Eq. (2). For instance, assigning to $`𝐞_3`$ the special role of unimodular normal vector to a given surface $`𝒮`$, whose tangent plane is defined by the vectors $`(𝐞_1,𝐞_2)`$, the equations (8) for $`\mu =1,2`$ and $`i=1,2,3`$ can be seen as the Gauss-Weingarten equations of such a surface. Moreover, the mapping $`𝐞_3:𝒮S^2`$ is the well known Gauss map. Furthermore, the corresponding integrability equations, rewritten as
$$_1J_2^{\left(ad\right)}_2J_1^{\left(ad\right)}+[J_1^{\left(ad\right)},J_2^{\left(ad\right)}]=0,$$
are the Gauss-Codazzi-Mainardi equations for a surface $`S`$ immersed in $`𝐑^3.`$ They are the real form of the equations (5) and (6). The $`U(1)`$ invariance of such equations is readily interpreted as the invariance under local rotations of the tangent plane at the surface $`𝒮`$. This identification is one of the motivation of the present article and it will be fully developed from the geometrical point of view in the next Sections. Here we would like to show how to use effectively the remaining equations in (8) for $`\mu =0`$ and $`i=1,2,3`$ and the corresponding integrability conditions. In particular, we ask if structures related to the integrable systems can be detected in the above general picture. Then, since we show that it is the case, we are allowed to introduce a completely integrable dynamics for the trihedral frame and, by consequence, we can infer an integrable dynamics for the corresponding surfaces.
First, we can rewrite Eqs. (5) and (6) by introducing the quantities
$$V=\left(\begin{array}{cc}v\hfill & \\ & v\hfill \end{array}\right),\mathrm{\Psi }_\pm =\left(\begin{array}{cc}& \psi _\pm \hfill \\ \psi _\pm ^{}\hfill & \end{array}\right),B=\frac{i}{2}\sigma _3\left(\mathrm{\Psi }_{}\mathrm{\Psi }_+\right).$$
Indeed, the torsionless condition (5) and its complex conjugate can be written as
$$\left(\begin{array}{cc}_z\hfill & \\ & _{\overline{z}}\hfill \end{array}\right)B+\frac{i}{2}\left(_{\overline{z}}_z\right)\mathrm{\Psi }_{}+V\mathrm{\Psi }_{}\mathrm{\Psi }_+V=0,$$
while combining the GCS law (6) and its complex conjugate yields
$$\mathrm{tr}\left\{\sigma _3\left[\left(\begin{array}{cc}_z\hfill & \\ & _{\overline{z}}\hfill \end{array}\right)V+B\mathrm{\Psi }_{}\mathrm{\Psi }_+B\right]\right\}=0.$$
These equations involve independent components in the basis of the complex $`2\times 2`$ matrices, and we have no information only about the identity component $`\sigma _0`$. So the last two equations provide
$$L_+\left[\frac{i}{2}\left(_{\overline{z}}_z\right)+V+B\right]\left[\frac{i}{2}\left(_{\overline{z}}_z\right)+V+B\right]L_{}=f\sigma _0,$$
(9)
$$L_\pm =\left(\begin{array}{cc}_z\hfill & \\ & _{\overline{z}}\hfill \end{array}\right)\mathrm{\Psi }_\pm $$
(10)
has the form of the two-dimensional principal Zakharov-Shabat spectral problem of elliptic type. Moreover, putting $`f0`$ in equation (9), we obtain the so called space part of the Bäcklund transformations associated with $`L_\pm `$, defining the first-order Bäcklund-gauge operator
$$\widehat{B}=\frac{i}{2}\left(_{\overline{z}}_z\right)+V+B.$$
This means that $`\widehat{B}`$ maps solutions of the two linear problems
$$L_{}\varphi _{}=0L_+B\varphi _{}=L_+\varphi _+=0.$$
(11)
The previous considerations tell us that in the above formalism the Gauss - Codazzi - Mainardi equations are expressed in terms of products of first order differential operators, which have a precise meaning in the theory of the completely integrable systems. Of course, in this context an essential role is played by the second (evolution) linear operator of the Lax pair, in order to introduce a compatible time evolution. The latter, generally speaking, is non linear, while the corresponding equations from the CS theory still contains the arbitrary functions interpreted as Lagrangian multipliers in the action (4). Indeed, the corresponding equations read
$$D_0\psi _+=\overline{D}q_0,D_0\psi _{}=Dq_0,$$
(12)
$$_0v_zv_0=i\left(q_0\psi _+^{}q_0^{}\psi _{}\right)$$
(13)
and their complex conjugated. However, we can exploit the freedom in the choice of $`q_0`$ and $`v_0`$ in order to fix the evolution of $`\psi _\pm `$ and $`v`$ in the $`x^0`$ variable. In fact, let us take
$$q_0=2i\left[\left(\overline{D}\frac{i}{2}\left(_{\overline{z}}\omega i_{\overline{z}}\chi \right)\right)\psi _++\left(D\frac{i}{2}\left(_z\omega +i_z\chi \right)\right)\psi _{}\right],$$
(14)
where we require that the real functions $`\chi `$ and $`\omega `$ satisfy the supplementary conditions
$$_z_{\overline{z}}\chi =4\left(\left|\psi _+\right|^2\left|\psi _{}\right|^2\right),_z_{\overline{z}}\omega =0.$$
(15)
Furthermore, by introducing an irrotational field $`𝐀=(_z\mathrm{\Lambda },_{\overline{z}}\mathrm{\Lambda },_0\mathrm{\Lambda })`$ ( $`\mathrm{\Lambda }`$ is an arbitrary real function) in such a way that
$$v=\frac{1}{4}_z\mathrm{\Lambda }\frac{i}{8}_z\chi ,v_0=\frac{1}{4}\left(_0\mathrm{\Lambda }+u_0\right),\psi _\pm =\mathrm{\Psi }_\pm \mathrm{exp}\left(\frac{i}{2}\mathrm{\Lambda }\right),$$
the ”time evolutions” (12) become
$$i_0\mathrm{\Psi }_\pm +2\left(_z^2+_{\overline{z}}^2\right)\mathrm{\Psi }_\pm +\frac{1}{2}u_0^\pm \mathrm{\Psi }_\pm i\left(_{\overline{z}}\omega _{\overline{z}}+_z\omega _z\right)\mathrm{\Psi }_\pm =0,$$
(16)
where we have suitably defined in terms of $`u_{0,}\chi `$ and $`\omega `$ the scalar fields $`u_0^\pm `$ , which obey the consistency conditions arising from (13)
$$_z_{\overline{z}}u_0^\pm =8\left(_z^2+_{\overline{z}}^2\right)\left|\mathrm{\Psi }_\pm \right|^2.$$
(17)
Thus, to summarize, the gauge fixing conditions (14) and (15) destroy the arbitraryness contained in the equations (12) in favour of a formally decoupled pair of Davey - Stewartson equations (16) and (17). Actually, between the two pairs of fields $`(\mathrm{\Psi }_\pm ,u_0^\pm )`$ there still exists the coupling provided by the torsionless condition (5), which in the new variables takes the form
$$\left(_{\overline{z}}+\frac{1}{4}_{\overline{z}}\chi \right)\mathrm{\Psi }_{}=\left(_z+\frac{1}{4}_z\chi \right)\mathrm{\Psi }_+.$$
(18)
As we discussed above, Eq. (18) is in essence the space part of the Bäcklund transformations. Starting from a known solution, say $`(\mathrm{\Psi }_{},u_0^{})`$ and fixing $`\omega ,`$ one can reconstruct the function $`\chi `$, and solving (18) for $`\mathrm{\Psi }_+,`$ finally we find $`u_0^+`$ from (17). Furthermore, we observe that the gauge choice (14) - (15) is equivalent to fix the second operator of the Lax pair, denoted here by $`M_\pm =_{x^0}+\mathrm{\Sigma }_{k=0}^2M_\pm ^{\left(k\right)}\left[\frac{i}{2}\left(_{\overline{z}}_z\right)\right]^{2k}`$, where $`M_\pm ^{\left(k\right)}`$ are specific matrices. $`M_\pm `$ provide the system (16) - (17) and (18) by the compatibility relations
$$[L_\pm ,M_\pm ]=0,M_+\widehat{B}\widehat{B}M_{}=0.$$
Moreover, by using a suitable particular eigenfunction $`\varphi _{}^0`$ of the equation (11), it is well known (see ,) that one can construct a new spectral problem of the form
$$L_{}^I=\left(\varphi _{}^0\right)^1L_{}\varphi _{}^0=_{x^1}+iS_{x^2},M_{}^I=\left(\varphi _{}^0\right)^1M_{}\varphi _{}^0,$$
where
$$S=i\left(\varphi _{}^0\right)^1\sigma _3\varphi _{}^0$$
is an element of $`SU\left(2\right)/U\left(1\right)`$ coset space and the corresponding eigenfunction is $`\varphi _{}^I=\left(\varphi _{}^0\right)^1\varphi _{}.`$ The resulting integrable system is known as Ishimori model and it describes the evolution of a classical spin in a background generated by the density of the topological charge. Since this equation will be discussed in the next Sections, we do not give other details about it. But here we want to stress that such a system is an alternative integrable restriction of the possible configurations of the CS field, exactly as the Davey-Stewartson equation does. Furthermore, one can put the question if the spin field $`S`$ has something to do with the trihedral frame introduced above. The consequences arising from the identification of $`S`$ with one of the unimodular vector fields $`𝐞_i`$ is the main subject of the next Sections.
## 3 Surfaces in $`R^3`$
To introduce our terminology and notations and to make the exposition self-contained, we recall some basic facts from the theory of 2-dimensional surfaces. So, we consider a smooth surface in a three dimensional Euclidean space $`𝐑^3`$. Let $`x,y`$ be the local coordinates on the surface. At the same time, the surface can be described by the position vector $`𝐫(x_1,x_2,x_3)=𝐫(x,y)`$, where $`x_i`$ are coordinates of $`𝐑^3`$. The surface is uniquely defined within rigid motions by the two fundamental forms
$$I=Edx^2+2Fdxdy+Gdy^2$$
(19)
and
$$II=Ldx^2+2Mdxdy+Ndy^2,$$
(20)
where $`E,F,G,L,M,N`$ can be defined by
$$E=𝐫_x𝐫_x=g_{11},F=𝐫_x𝐫_y=g_{12}=g_{21},G=𝐫_y𝐫_y=g_{22},$$
(21)
$$L=𝐧𝐫_{xx}=b_{11},M=𝐧𝐫_{xy}=b_{12}=b_{21},N=𝐧𝐫_{yy}=b_{22}.$$
(22)
In Equations (21) and (22)
$$𝐧(x,y)=\frac{𝐫_x𝐫_y}{𝐫_x𝐫_y}=\frac{𝐫_x𝐫_y}{\sqrt{g}}$$
(23)
is introduced, where $`g=det(g_{ij})=EGF^2=𝐫_x𝐫_y^2,`$ is the normal vector field at each point of the surface. Then the triple $`(𝐫_x,𝐫_y,𝐧)`$ represents a local frame of $`𝐑^3,`$ the changes of which are characterized by the GW equations
$$𝐫_{xx}=\mathrm{\Gamma }_{11}^1𝐫_x+\mathrm{\Gamma }_{11}^2𝐫_y+L𝐧,$$
(24)
$$𝐫_{xy}=\mathrm{\Gamma }_{12}^1𝐫_x+\mathrm{\Gamma }_{12}^2𝐫_y+M𝐧,$$
(25)
$$𝐫_{yy}=\mathrm{\Gamma }_{22}^1𝐫_x+\mathrm{\Gamma }_{22}^2𝐫_y+N𝐧,$$
(26)
$$𝐧_x=p_{11}𝐫_x+p_{12}𝐫_y,$$
(27)
$$𝐧_y=p_{21}𝐫_x+p_{22}𝐫_y,$$
(28)
where the Christoffel symbols of the second kind are defined by $`g_{ij}`$ ( $`g^{ij}=(g_{ij})^1)`$ as
$$\mathrm{\Gamma }_{jk}^i=\frac{1}{2}g^{il}(\frac{g_{lk}}{x^j}+\frac{g_{jl}}{x^k}\frac{g_{jk}}{x^l}),$$
(29)
and
$$p_{ij}=b_{ik}g^{kj}.$$
(30)
The principal curvatures $`k_1,k_2`$ are the eigenvalues of the Weingarten operator
$$\mathrm{\Lambda }=\left(\begin{array}{cc}E& F\\ F& G\end{array}\right)^1\left(\begin{array}{cc}L& M\\ M& N\end{array}\right).$$
(31)
which for the mean and the Gassian curvature implies
$$H=\frac{k_1+k_2}{2}=tr(\mathrm{\Lambda })=\frac{EN+LG2MF}{2(EGF^2},$$
(32)
$$K=k_1k_2=det(\mathrm{\Lambda })=\frac{LNM^2}{EGF^2}.$$
(33)
One of the global characteristics of surfaces is the integral curvature
$$\chi =\frac{1}{2\pi }K\sqrt{g}𝑑x𝑑y,$$
(34)
which for compact oriented surfaces is the integer
$$\chi =2(1\mathrm{}),$$
(35)
where $`\mathrm{}`$ is the genus of the surface. The compatibility conditions of the GW equations (24)-(28) furnish the GMC equations
$$R_{ijk}^l=b_{ij}b_k^l,\frac{b^{ij}}{x^k}\frac{b_{ik}}{x^j}=\mathrm{\Gamma }_{ik}^sb_{js}\mathrm{\Gamma }_{ij}^sb_{ks},$$
(36)
where $`b_i^j=g^{jl}b_{il}`$ and the curvature tensor is
$$R_{ijk}^l=\frac{\mathrm{\Gamma }_{ij}^l}{x^k}\frac{\mathrm{\Gamma }_{ik}^l}{x^j}+\mathrm{\Gamma }_{ij}^s\mathrm{\Gamma }_{ks}^l\mathrm{\Gamma }_{ik}^s\mathrm{\Gamma }_{js}^l.$$
(37)
For our purposes it is convenient to employ the triad of orthonormal vectors
$$𝐞_1=\frac{𝐫_x}{\sqrt{E}},𝐞_2=𝐧,𝐞_3=𝐞_1𝐞_2.$$
(38)
In terms of these vectors the GW equations (24)-(28) take the form
$$𝐞_{jx}=𝐗𝐞_j,𝐞_{jy}=𝐘𝐞_j,$$
(39)
where
$$𝐗=\tau 𝐞_1+\sigma 𝐞_2+k𝐞_3,𝐘=m_1𝐞_1+m_2𝐞_2+m_3𝐞_3,$$
(40)
and
$$k=\frac{L}{\sqrt{E}},\sigma =\frac{2EF_xEE_yFE_x}{2E\sqrt{g}},\tau =\frac{MELF}{\sqrt{gE}},$$
(41)
$$m_1=\frac{NEMF}{\sqrt{gE}},m_2=\frac{FE_yEG_x}{2E\sqrt{g}},m_3=\frac{M}{\sqrt{E}}.$$
(42)
Similarly, we can rewrite the GMC equations (32) in the following form
$$A_yB_x+[A,B]=0,$$
(43)
with
$$A=\left(\begin{array}{ccc}0& k& \sigma \\ k& 0& \tau \\ \sigma & \tau & 0\end{array}\right),B=\left(\begin{array}{ccc}0& m_3& m_2\\ m_3& 0& m_1\\ m_2& m_1& 0\end{array}\right).$$
(44)
Then, the GMC equation turns out to be equivalent to the set of equations for the coefficients of the first and second fundamental forms. This system, which is in general non integrable, reduces to integrable partial differential equations for certain particular surfaces .
## 4 Deformations of surfaces in 2+1 dimensions: the general formulation
It is well known that in some cases deformations of surfaces can be associated with integrable equations \[1-3\]. Here we are interested in the deformation of the two-dimensional surfaces discussed in Section 3. In other words, we have to deal with the motion of such surfaces. To this aim, let us introduce the vector field
$$𝐫_t=a_1𝐫_x+a_2𝐫_y+a_3𝐧,$$
(45)
where the $`a_i`$’s are some real functions. It is easy to show that the evolution of the local trihedral frame is given by
$$𝐞_{jt}=𝐓𝐞_j,$$
(46)
$$𝐓=\omega _1𝐞_1+\omega _2𝐞_2+\omega _3𝐞_3,$$
(47)
$`\omega _j`$’s being real functions. Summarizing, the changes of the local frame are provided by
$$𝐞_{jx}=𝐗𝐞_j,𝐞_{jy}=𝐘𝐞_j,𝐞_{jt}=𝐓𝐞_j,$$
(48)
where the vectors X and Y are defined by (40). This system is analogous to the system (8) in Section 2.
The system (52) represents the simplest form of the (2+1)-dimensional GW equations.
By introducing the matrix
$$C=\left(\begin{array}{ccc}0& \omega _3& \omega _2\\ \omega _3& 0& \omega _1\\ \omega _2& \omega _1& 0\end{array}\right)$$
(49)
and using the matrices $`A`$ and $`B`$ (see (44)), the compatibility conditions of Eqs. (48) entail
$$A_yB_x+[A,B]=0,$$
(50)
$$A_tC_x+[A,C]=0,$$
(51)
$$B_tC_y+[B,C]=0.$$
(52)
Of the nine funcions $`k,\sigma ,\tau ,m_i,\omega _i`$ involved in $`A,B`$ and $`C`$ only three are independent. In fact, we can express the functions $`m_i,\omega _i`$ in terms of $`k,\sigma ,\tau `$ and their derivatives. This point will be discussed later.
### 4.1 Some geometrical invariants and integrals of motion as consequence of the geometrical formalism
The formalism developed above yields some important invariants having a pure geometrical nature. Indeed, in terms of the triad vectors these geometrical invariants take the form
$$K_1^{(t)}=𝐞_1(𝐞_{1x}𝐞_{1y})𝑑x𝑑y,K_2^{(t)}=𝐞_2(𝐞_{2x}𝐞_{2y})𝑑x𝑑y,$$
$$K_3^{(t)}=𝐞_3(𝐞_{3x}𝐞_{3y})𝑑x𝑑y.$$
(53)
In a similar way we can write down other two classes of invariants with respect to $`x`$ and $`y`$ directions, respectively. These geometrical invariants can be interpreted as ”topological charges”. However, three of them, namely $`K_i^{(t)}`$ $`(i=1,2,3)`$ behave as integrals of motion of the (2+1)-dimensional geometrical models under consideration. This will be elucidated in the next Sections. These invariants can be related to the topological Chern index of a curvature 2-form on a 2-dimensional space .
## 5 Integrable spin models in (2+1)- dimensions
Now let us dwell upon the problem of finding or building up integrable deformations of (2+1)-dimensional surfaces. Among several possibilities, within the geometrical formalism previously presented we shall consider multidimensional integrable spin (field) systems (MISSs) to recognize integrable deformations of surfaces.
### 5.1 The spin model
A few words on MISSs. At present there exist many integrable spin systems in (2+1)-dimensions (see, for example, Refs. \[20,21-25\]). A well known prototype of these systems is the Ishimori model (IM) . A more general (2+1)-dimensional integrable spin model is described by the pair of equations
$$𝐒_t+𝐒\{(b+1)𝐒_{\xi \xi }b𝐒_{\eta \eta }\}+bu_\eta 𝐒_\eta +(b+1)u_\xi 𝐒_\xi =0,$$
(54)
$$u_{\xi \eta }=𝐒(𝐒_\xi 𝐒_\eta ),$$
(55)
where $`\xi ,\eta `$ are real or complex variables, $`b`$ is a real constant, $`𝐒=(S_1,S_2,S_3)`$ is the spin (field) vector, $`𝐒^2=1`$, and $`u`$ is a scalar function. These equations, which are called M-XX equations (about our conditional notations, see e.g. \[21-25\]), are one of the (2+1)-dimensional integrable generalizations of the isotropic Landau-Lifshitz (LL) equation
$$𝐒_t=𝐒𝐒_{xx}.$$
(56)
In (1+1)-dimensions, Eqs.(54) and (55) reduce to the LL equation. In fact, assuming that the variables $`𝐒,u`$ are, for example, independent of $`\eta ,`$ then Eqs. (54) and (55) reproduce Eq. (56) within a simple scale transformation.
We notice that Eqs. (54) and (55) are not the only integrable generalization of the LL equation in (2+1)-dimensions. Actually, other integrable generalizations exist, such as the IM or the model defined by
$$𝐒_t=(𝐒𝐒_y+u𝐒)_x,$$
(57)
$$u_x=𝐒(𝐒_x𝐒_y).$$
(58)
These equations, which are called M-I equations (see ) are again completely integrable. Some properties of these equations are studied in \[21-23\].
#### 5.1.1 The Lax representation
Equations (54) and (55) can be solved by the IST method. The applicability of the IST method to Eqs. (54) and (55) is based on the equivalence of these equations to the compatibility condition of the following linear equations (the Lax representation (LR))
$$\mathrm{\Phi }_{Z^+}=S\mathrm{\Phi }_Z^{},$$
(59)
$$\mathrm{\Phi }_t=2i[S+(2b+1)I]\mathrm{\Phi }_{Z^{}Z^{}}+W\mathrm{\Phi }_Z^{},$$
(60)
where $`Z^\pm =\xi \pm \eta `$ and
$$W=2i\{(2b+1)(F^++F^{}S)+(F^+S+F^{})+(2b+1)SS_Z^{}+\frac{1}{2}S_Z^{}+\frac{1}{2}SS_{Z^+}\},$$
$$S=\left(\begin{array}{cc}S_3& rS^{}\\ rS^+& S_3\end{array}\right),S^\pm =S_1\pm iS_2S^2=EI,E=\pm 1,r^2=\pm 1,$$
$$F^+=2iu_Z^{},F^{}=2iu_{Z^+}.$$
In fact, from the condition $`\mathrm{\Phi }_{Z^+t}=\mathrm{\Phi }_{tZ^+}`$ we deduce
$$iS_t+\frac{1}{2}[S,(b+1)S_{\xi \xi }bS_{\eta \eta }]+ibu_\eta S_\eta +i(b+1)u_\xi S_\xi =0,$$
(61)
$$u_{\xi \eta }=\frac{1}{4i}tr(S[S_\xi ,S_\eta ]),$$
(62)
which is the matrix form of Eqs. (54) and (55).
#### 5.1.2 Special cases
Equations (54) and (55) contain both well known and less known integrable cases in (2+1) and (1+1)-dimensions. Below we shall report some of them.
$`i)`$ If $`b=0`$, Eqs. (54) and (55) yield
$$𝐒_t+𝐒𝐒_{\xi \xi }+w𝐒_\xi =0,$$
(63)
$$w_\eta 𝐒(𝐒_\xi 𝐒_\eta )=0,$$
(64)
where $`w=u_\xi `$. This system, which is known as the M-VIII model , is one of the simplest spin systems in (2+1)-dimensions integrable by IST. It affords different type of solutions (solitons, vortices, etc.). In particular, vortex solutions of Eqs. (54)-(55) can be derived from vortex solutions of the spin system (54)-(55) discussed in Section 9 (for $`b=0).`$
$`ii)`$ Let us introduce the coordinates $`x=\xi \eta ,y=\alpha (\xi +\eta )`$ and put $`b=\frac{1}{2}`$. Then, the spin system (54)-(55) reduces to the IM
$$𝐒_t+𝐒(𝐒_{xx}+\alpha ^2𝐒_{yy})+u_x𝐒_y+u_y𝐒_x=0,$$
(65)
$$u_{xx}\alpha ^2u_{yy}=2\alpha ^2𝐒(𝐒_x𝐒_y).$$
(66)
The IM is the first integrable spin (field) system in the plane which can be solved by IST method. The IM was studied by many authors from different points of view (e.g. ).
$`iii)`$ By setting $`b=0,\eta =t`$, Eqs. (54) and (55) reduce to the following (1+1)-dimensional spin system:
$$𝐒_t+𝐒𝐒_{\xi \xi }+w𝐒_\xi =0,$$
(67)
$$w_t+\frac{1}{2}(𝐒_\xi ^2)_\xi =0.$$
(68)
This integrable model describes the nonlinear dynamics of compressible magnets . It is the first (and, to the best of our knowledge, at present the unique) example of integrable spin system governing the nonlinear interactions of spin ($`𝐒`$) and lattice ($`u`$) subsystems in (1+1)-dimensions.
## 6 Bilinear representations
One of the powerful tools in the soliton theory is the Hirota method. Now we show how to construct the bilinear representations of the fields of the spin system by using geometry. Let $`e_{ji}`$ are the components of the unit vector $`𝐞_j`$, i.e. $`𝐞_j=(e_{j1},e_{j2},e_{j3})`$.We can take the following representation for the components of the vector $`𝐞_1:`$
$$e_1^+=e_{11}+ie_{12}=\frac{2\overline{f}g}{\overline{f}f+\overline{g}g},e_{13}=\frac{\overline{f}f\overline{g}g}{\overline{f}f+\overline{g}g}.$$
(69)
Then, we get
$$e_2^+=e_{21}+ie_{22}=\frac{\overline{f}^2+g^2}{\overline{f}f+\overline{g}g},e_{23}=i\frac{fg\overline{f}\overline{g}}{\overline{f}f+\overline{g}g},$$
(70)
$$e_3^+=e_{31}+ie_{32}=\frac{\overline{f}^2g^2}{\overline{f}f+\overline{g}g},e_{33}=\frac{fg+\overline{f}\overline{g}}{\overline{f}f+\overline{g}g},$$
(71)
with
$$k=i\frac{D_x(gf\overline{g}\overline{f})}{\overline{f}f+\overline{g}g},m_3=i\frac{D_y(gf\overline{g}\overline{f})}{\overline{f}f+\overline{g}g},$$
(72)
$$\sigma =\frac{D_x(gf+\overline{g}\overline{f})}{\overline{f}f+\overline{g}g},m_2=\frac{D_y(gf+\overline{g}\overline{f})}{\overline{f}f+\overline{g}g},$$
(73)
$$\tau =i\frac{D_x(\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g},m_1=i\frac{D_y(\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g},$$
(74)
$$\omega _3=i\frac{D_t(gf\overline{g}\overline{f})}{\overline{f}f+\overline{g}g},\omega _2=\frac{D_t(gf+\overline{g}\overline{f})}{\overline{f}f+\overline{g}g}$$
(75)
$$\omega _1=i\frac{D_t(\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g}.$$
(76)
The Hirota operators $`D_x,D_y`$ and $`D_t`$ are defined by
$`D_x^lD_y^mD_t^nf(x,y,t)g(x,y,t)`$ $`=`$ $`(_x_x^{})^l(_y_y^{})^m(_t_t^{})^nf(x,y,t)`$
$`g(x^{},y^{},t^{})`$ $``$ $`_{x=x^{},y=y^{},t=t^{}}.`$
Now we write down the bilinear representation for the spin vector and for the derivatives of the potential $`u.`$ Taking into account (69)-(76), we find
$$S^+=S_1+iS_2=\frac{2\overline{f}g}{\overline{f}f+\overline{g}g},S_3=\frac{\overline{f}f\overline{g}g}{\overline{f}f+\overline{g}g}.$$
(77)
This is the general representation for the components of the spin vector for all the spin systems. However, for the potential the bilinear forms for every spin system should be different. In the following we shall consider some examples.
$`i)`$ The Ishimori model. In this case we have
$$\tau =\frac{1}{2}u_y,m_1=\frac{1}{2\alpha ^2}u_x.$$
(78)
Hence, from (74) we get
$$u_y=2i\frac{D_x(\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g},u_x=2i\alpha ^2\frac{D_y(\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g}.$$
(79)
On the other hand, from (74) it follows also
$$\tau _x=\alpha ^2m_{1y},$$
(80)
so that
$$m_1=\alpha ^2_y^1\tau _x.$$
(81)
$`ii)`$ The isotropic M-I equation. Let us take
$$\tau =0,m_1=u.$$
(82)
Then, from (74) and (82) we obtain
$$D_x(\overline{f}f+\overline{g}g)=0,u=i\frac{D_y(\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g}.$$
(83)
$`iii)`$The spin system (54)-(55). Let us start from
$$\tau =\frac{1}{2}u_\xi ,m_1=\frac{1}{2}u_\eta .$$
(84)
Then we have
$$u_\xi =2i\frac{D_\xi (\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g},u_\eta =2i\frac{D_\eta (\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g},$$
and
$$\tau _\eta =m_{1\xi },$$
(85)
$$m_1=_\xi ^1\tau _\eta .$$
(86)
An important consequence of these results is the possibility to determine the time evolution of the potential (and/or its derivatives). For instance, for the IM the time evolution of the derivatives of the potential are given by
$$\frac{1}{2}(u_y)_t\omega _{3x}+\sigma \omega _1\tau \omega _2=0,$$
(87)
$$\frac{1}{2\alpha ^2}(u_x)_t\omega _{1y}+m_3\omega _3m_2\omega _3=0.$$
(88)
## 7 Deformations of surfaces by integrable spin systems in 2+1 dimensions
An interesting example of surface integrable deformation can be found out by identifying the tangent unit vector $`𝐞_1`$ with the spin vector, i.e.
$$𝐞_1𝐒.$$
(89)
In such a way, the spin model (54)-(55) takes the form
$$𝐞_{1t}+𝐞_1\{(b+1)𝐞_{1\xi \xi }b𝐞_{1\eta \eta }\}+bu_\eta 𝐞_{1\eta }+(b+1)u_\xi 𝐞_{1\xi }=0,$$
(90)
$$u_{\xi \eta }=𝐞_1(𝐞_{1\xi }𝐞_{1\eta }).$$
(91)
The functions $`m_i,\omega _i`$ can be expressed in terms of the three independent functions $`k,\tau `$,$`\sigma .`$Using the GW equations (48), Eqs. (90)-(91) can be written as
$$𝐞_{1t}=\omega _3𝐞_2\omega _2𝐞_3,$$
(92)
$$u_{\xi \eta }=\sigma m_3km_2,$$
(93)
where
$$\omega _2=b[m_{3\eta }m_2^2u_\eta m_2](b+1)[k_\xi +\sigma \tau +u_\xi \sigma ],$$
(94)
$$\omega _3=(b+1)[\sigma _\xi k\tau ku_\xi ]b[m_{2\eta }m_1m_3+u_\eta m_3].$$
(95)
Now by choosing $`m_1`$ according to the special reduction (84), the remaining functions $`m_2`$ and $`m_3`$ are given by
$$m_2=\frac{\sigma m_3u_{\xi \eta }}{k},m_3=\frac{\sigma \tau u_{\xi \eta }+(a_2/2)u_\xi a_3}{a_1+\sigma ^2\tau },$$
(96)
where
$$a_1=k^2\tau k\sigma _\xi +\sigma k_\xi ,a_2=k^3\sigma ^2k,$$
$$a_3=k^2\sigma _\eta \sigma kk_\eta +ku_{\xi \xi \eta }k_\xi u_{\xi \eta }.$$
(97)
By virtue of these formulae we derive the function $`\omega _1`$ from (50).
Thus, all the unknown functions $`m_i,\omega _i`$ are defined via the three functions $`k,\tau ,\sigma `$ only and their derivatives. This is the consequence of the identification of the motion of surface with the spin system (90)-(91).This means that the motion of surface is fully determined by these three functions. Since the spin model (90)-(91) is integrable, we can conclude that the deformation of the surface characterized by Eqs. (50)-(52) is integrable.
## 8 Deformations of surfaces related to the (2+1)-dimensional NLS-type equation
One of the most remarkable consequence of the geometrical formalism previously outlined is that it allows to find the equivalent counterpart of the spin system (54)-(55). To show this property, let us introduce two complex functions $`q,p`$ according to the following expressions
$$q=a_1e^{ib_1},p=a_2e^{ib_2},$$
(98)
where $`a_j,b_j`$ are real functions. Now let us choose the functions $`a_j,b_j`$ in such a way that
$$a_1^2=\frac{1}{4}k^2+\frac{|\alpha |^2}{4}(m_3^2+m_2^2)\frac{1}{2}\alpha _Rkm_3\frac{1}{2}\alpha _Ikm_2,$$
(99)
$$b_1=_x^1\{\frac{\gamma _1}{2ia_1^^2}(\overline{A}A+D\overline{D})\},$$
(100)
$$a_2^2=\frac{1}{4}k^2+\frac{|\alpha |^2}{4}(m_3^2+m_2^2)+\frac{1}{2}\alpha _Rkm_3\frac{1}{2}\alpha _Ikm_2,$$
(101)
$$b_2=_x^1\{\frac{\gamma _2}{2ia_2^^2}(A\overline{A}+\overline{D}D)\},$$
(102)
where
$$\gamma _1=i\{\frac{1}{2}k^2\tau +\frac{|\alpha |^2}{2}(m_3km_1+m_2k_y)$$
$$\frac{1}{2}\alpha _R(k^2m_1+m_3k\tau +m_2k_x)+\frac{1}{2}\alpha _I[k(2k_ym_{3x})k_xm_3]\},$$
(103)
$$\gamma _2=i\{\frac{1}{2}k^2\tau +\frac{|\alpha |^2}{2}(m_3km_1+m_2k_y)+$$
$$\frac{1}{2}\alpha _R(k^2m_1+m_3k\tau +m_2k_x)+\frac{1}{2}\alpha _I[k(2k_ym_{3x})k_xm_3]\}.$$
(104)
Here $`\alpha =\alpha _R+i\alpha _I`$. In this case, $`q,p`$ satisfy the (2+1)-dimensional equations of the nonlinear Schrödinger (NLS) type
$$iq_t+(1+b)q_{\xi \xi }bq_{\eta \eta }+vq=0,$$
(105)
$$ip_t(1+b)p_{\xi \xi }+bp_{\eta \eta }vp=0,$$
(106)
$$v_{\xi \eta }=2\{(1+b)(pq)_{\xi \xi }b(pq)_{\eta \eta }\}.$$
(107)
These equations are the geometrical equivalent counterpart of the spin system (54)-(55). Therefore, the spin system and the (2+1)-dimensional NLS equations (105)-(107) turn out to be mutually geometrical equivalent.
### 8.1 Gauge equivalence
Now we prove that the spin system (54)-(55) and equations (105)-(107) are not only equivalent in the geometrical sense, but are also each other gauge equivalent. To this purpose, let us perform the gauge transformation $`\mathrm{\Psi }=g\mathrm{\Phi }`$, where the function $`\mathrm{\Phi }`$ is the solution of equations (59)-(60) and $`g`$ is a 2x2 matrix such that
$$S=g^1\sigma _3g,$$
(108)
and
$$g_{Z^+}g^1\sigma _3g_Z^{}g^1=\left(\begin{array}{cc}0& q\\ p& 0\end{array}\right).$$
(109)
Under this transformation the function $`\mathrm{\Psi }`$ obeys the following set of linear equations
$$\mathrm{\Psi }_{Z^+}=\sigma _3\mathrm{\Psi }_Z^{}+B_0\mathrm{\Psi },$$
(110)
$$\mathrm{\Psi }_t=4iC_2\mathrm{\Psi }_{Z^{}Z^{}}+2C_1\mathrm{\Psi }_Z^{}+C_0\mathrm{\Psi },$$
(111)
where $`B_0,C_j`$ are given by
$$B_0=\left(\begin{array}{cc}0& q\\ p& 0\end{array}\right),C_2=\left(\begin{array}{cc}b+1& 0\\ 0& b\end{array}\right),C_1=\left(\begin{array}{cc}0& iq\\ ip& 0\end{array}\right),C_0=\left(\begin{array}{cc}c_{11}& c_{12}\\ c_{21}& c_{22}\end{array}\right),$$
(112)
and the functions $`c_{ij}`$ $`(i,j=1,2)`$ fulfills the equations
$$c_{12}=i[(4b+3)q_Z^{}+q_{Z^+}],c_{21}=i[(4b+1)p_Z^{}+p_{Z^+}],$$
(113)
$$c_{11Z^{}}c_{11Z^+}=i[(4b+3)(pq)_Z^{}+(pq)_{Z^+}],$$
(114)
$$c_{22Z^{}}+c_{22Z^+}=i[(4b+1)(pq)_Z^{}+(pq)_{Z^+}],$$
(115)
with $`v=i(c_{22}c_{11})`$ (see (105)-(107)).
The compatibility condition of Eqs. (110)-(11) gives the equations (105)-(107). This means that the spin model (54)-(55) and the NLS - type equations (105)-(107) are gauge equivalent each other. Moreover, it is easy to check that if $`g`$ satisfies equations Eq. (109), then $`S`$ given by (108) satisfies Eqs. (54)-(55) with
$$u=2\mathrm{ln}detg.$$
(116)
#### 8.1.1 Reductions
Equations (105)-(107), as its and equivalent spin system (54)-(55), contain several integrable cases, namely
$`i)`$ $`b=0`$. Equations (105)-(107) yield the equations
$$iq_t+q_{\xi \xi }+vq=0,$$
(117)
$$ip_tp_{\xi \xi }vp=0,$$
(118)
$$v_\eta =2(pq)_\xi .$$
(119)
$`ii)`$ $`a=b=\frac{1}{2}`$. Then, we give the Davey-Stewartson (DS) equation
$$iq_t+q_{xx}+\alpha ^2q_{yy}+vq=0,$$
(120)
$$ip_tp_{xx}\alpha ^2p_{yy}vq=0,$$
(121)
$$v_{xx}\alpha ^2v_{yy}=2\{(pq)_{xx}+\alpha ^2(pq)_{yy}\},$$
(122)
where $`x=\xi \eta ,y=\alpha (\xi +\eta )`$.
$`iii)`$ Putting $`b=0,\eta =t`$, Eqs. (105)-(107) reduce to the (1+1)-dimensional Ma or Yajima-Oikawa equations
$$iq_t+q_{\xi \xi }+vq=0,$$
(123)
$$ip_tp_{\xi \xi }vp=0,$$
(124)
$$v_t+2(pq)_\xi =0,$$
(125)
which are known to be integrable.
## 9 Solutions of the spin system
It could be of interest to study Eqs. (54)-(55) by the IST method. However, to look for some special solutions, it is convenient to exploit the Hirota bilinear method. To this aim, let us build up the bilinear form of (54)-(55) for the compact case.
We obtain
$$S^+=S_1+iS_2=\frac{2\overline{f}g}{\overline{f}f+\overline{g}g},S_3=\frac{\overline{f}f\overline{g}g}{\overline{f}f+\overline{g}g},$$
(126)
$$u_\xi =2i\frac{D_\xi (\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g},u_\eta =2i\frac{D_\eta (\overline{f}f+\overline{g}g)}{\overline{f}f+\overline{g}g},$$
(127)
where $`S_j`$ $`(j=1,2,3)`$ are the components of spin vector $`𝐒,`$ $`S^\pm `$ $`=S_1\pm iS_2,`$and $`u`$ is the scalar potential.
Hence, from (116) we get
$$u(\xi ,\eta ,t)=2\mathrm{ln}(f^2+g^2).$$
(128)
Substituting formulae (126) and (127) into the spin system (54)-(55), we obtain the bilinear equations
$$[iD_t(b+1)D_\xi ^2+bD_\eta ^2](\overline{f}g)=0,$$
(129)
$$[iD_t(b+1)D_\xi ^2+bD_\eta ^2](\overline{f}f\overline{g}g)=0,$$
(130)
$$\{D_\xi D_\eta +D_\eta D_\xi \}(\overline{f}f+\overline{g}g)(\overline{f}f+\overline{g}g)=0.$$
(131)
Equation (131) coincides with the compatibility condition $`u_{\xi \eta }=u_{\eta \xi }`$.
Now we can construct some special solutions of Eqs. (54)-(55). In particular, to construct vortex solutions, we start from Eqs. (129)-(130) and assume that
$$f=f(\xi ,t),g=g(\xi ,t).$$
(132)
Then Eq. (131) is satisfied automatically. At the same time, Eqs. (129)-(130) are fulfilled if
$$if_t+(b+1)f_{\xi \xi }=0ig_t+(b+1)g_{\xi \xi }=0.$$
(133)
Consequently, we are led to the following multi-vortex solutions
$$g_N=\underset{j=0}{\overset{N}{}}\underset{m+2n=j}{}\frac{a_j}{m!n!}(\frac{2}{b+1})^{\frac{m}{2}}\xi ^m(2it)^n,$$
(134)
$$f_N=\underset{j=0}{\overset{N1}{}}\underset{m+2n=j}{}\frac{b_j}{m!n!}(\frac{2}{b+1})^{\frac{m}{2}}\xi ^m(2it)^n,$$
(135)
where $`a_j`$ and $`b_j`$ are arbitrary complex constants, and $`m,n`$ are non-negative integer numbers. In particular, the 1-vortex solution can be derived by choosing
$$f=b_0,g=a_1^{}\xi +a_0,$$
(136)
where $`a_1^{}=a_1(\frac{2}{b+1})^{\frac{1}{2}}`$.
So, the 1-vortex solution is static. To find a dynamic solution, we have to consider the $`N`$-vortex solution using the forms
$$f(\xi ,t)=b_0\underset{j=1}{\overset{N}{}}[\xi p_j(t)],$$
(137)
$$g(\xi ,t)=a_0\underset{j=1}{\overset{N}{}}[\xi q_j(t)],$$
(138)
where $`p_j`$ and $`q_j`$ denote the positions of the zeros of $`f`$ and $`g`$, and $`a_0,b_0`$ are constants. The evolution of $`p_j`$ and $`q_j`$ is
$$p_{jt}=i(b+1)\underset{kj}{\overset{N}{}}\frac{1}{p_jp_k},$$
(139)
$$q_{jt}=i(b+1)\underset{kj}{\overset{N}{}}\frac{1}{q_jq_k},$$
(140)
where $`j,k=1,2,\mathrm{},N`$. These equations are related to the Calogero-Moser system.
## 10 Special surfaces corresponding to vortex solutions of the spin system
This Section is devoted to the construction of explicit surfaces. To this aim, let us start from the 1-vortex solution of the spin system (54)-(55). By choosing for simplicity $`E=1`$, Eqs. (41)-(42) become
$$k=L,\sigma =\frac{F_x}{\sqrt{g}},\tau =\frac{MLF}{\sqrt{g}},$$
(141)
$$m_1=\frac{NMF}{\sqrt{g}},m_2=\frac{G_x}{2\sqrt{g}},m_3=M.$$
(142)
On the other hand, from (38) and (89) we get
$$𝐫_x(\xi ,\eta ,t)=𝐒(\xi ,\eta ,t).$$
(143)
Now let us consider by way of example the surface associated with the 1-vortex solution of the Ishimori system, whose components are
$`\begin{array}{cc}S_3=\frac{1b^2\left|\mathrm{\Xi }\right|^2}{1+b^2\left|\mathrm{\Xi }\right|^2},\hfill & S^+=2be^{i\delta }\frac{\mathrm{\Xi }}{1+b^2\left|\mathrm{\Xi }\right|^2},\hfill \end{array}`$
$`u=2\left(\mathrm{ln}\left(b\right)+\mathrm{ln}\left[1+b^2\left|\mathrm{\Xi }\right|^2\right]\right),`$
where $`\mathrm{\Xi }`$ denotes the complex variable $`\mathrm{\Xi }=2a\mathrm{exp}\left(i\gamma \right)\left(x+iy\right)+1`$, $`a,b,`$ $`\gamma `$ and $`\delta `$ being real constants.
Then, by resorting to the formula $`e_1=r_x/\sqrt{E}`$(with $`E1)`$, we can integrate to yield the following components for the position vector
$$\begin{array}{c}r_1=\frac{\widehat{c}\mathrm{ln}(1+b^2\left|\mathrm{\Xi }\right|^2)}{ab}\frac{\sqrt{2}\mathrm{arctan}\left[b\sqrt{\frac{2}{\mathrm{\Omega }}}\left(ax+c\right)\right]\left(\stackrel{~}{c}2a\widehat{s}y\right)}{a\sqrt{\mathrm{\Omega }}},\\ r_2=\frac{\widehat{s}\mathrm{ln}(1+b^2\left|\mathrm{\Xi }\right|^2)}{ab}\frac{\sqrt{2}\mathrm{arctan}\left[b\sqrt{\frac{2}{\mathrm{\Omega }}}\left(ax+c\right)\right]\left(2a\widehat{c}y+\stackrel{~}{s}\right)}{a\sqrt{\mathrm{\Omega }}},\\ r_3=x+\frac{2\sqrt{2}\mathrm{arctan}\left[b\sqrt{\frac{2}{\mathrm{\Omega }}}\left(ax+c\right)\right]}{ab\sqrt{\mathrm{\Omega }}},\end{array}$$
(145)
where for the sake of clarity we have introduced the second degree polynomial $`\mathrm{\Omega }=2+b^2\left[1+2a^2y^2\mathrm{cos}\left(2\gamma \right)+4a\mathrm{sin}\left(\gamma \right)y\right]`$ and the constants $`c=\mathrm{cos}\left(\gamma \right),`$ $`s=\mathrm{sin}\left(\gamma \right),`$ $`\widehat{c}=`$ $`\mathrm{cos}\left(\gamma +\delta \right),`$ $`\widehat{s}=`$ $`\mathrm{sin}\left(\gamma +\delta \right),`$ $`\stackrel{~}{c}=\mathrm{cos}(\delta )+\mathrm{cos}(2\gamma +\delta ),`$and $`\stackrel{~}{s}=\mathrm{sin}(\delta )+\mathrm{sin}(2\gamma +\delta ).`$ Furthermore, we have put identically equal to zero any arbitrary function of integration in $`y`$ only.
From them, with the help of the various formulae given in Section 3, we obtain the coefficients of the I-fundamental form
$`\begin{array}{cc}E=1,\hfill & F=\frac{4b(s+ay)(b(c+ax)\sqrt{\mathrm{\Omega }}+\sqrt{2}(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})At[x,y])}{(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})\mathrm{\Omega }^{\frac{3}{2}}}\hfill \end{array}`$
$`\begin{array}{c}G=\frac{2(4b^2(s+ay)^2\mathrm{\Omega }+8(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})At[x,y]{}_{}{}^{2})}{(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})\mathrm{\Omega }^2}\end{array}`$
where $`At[x,y]=\mathrm{arctan}\left[b\sqrt{\frac{2}{\mathrm{\Omega }}}\left(ax+c\right)\right],`$ and analogously for the II-fundamental form
$$\begin{array}{c}\sqrt{g}L=\frac{8ab[\left(b^3(c+ax)(s+ay)^2\sqrt{\mathrm{\Omega }}\right)+\sqrt{2}(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})At[x,y]]}{(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})^2\mathrm{\Omega }^{\frac{3}{2}}},\sqrt{g}M=\frac{4ab^2\left(s+ay\right)}{(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})^2}\hfill \\ \sqrt{g}N=\frac{8ab(2\sqrt{2}(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})\sqrt{\mathrm{\Omega }}At[x,y]+8b(c+ax)(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})At[x,y]{}_{}{}^{2}+b^3(c+ax)(s+ay)^2\mathrm{\Omega }(2+\mathrm{\Omega }))}{(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})^2\mathrm{\Omega }^3}\hfill \end{array}$$
where $`g`$ in the metric factor $`\sqrt{g}`$ is expressed by
$$\begin{array}{c}g=\frac{8}{(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})^2\mathrm{\Omega }^3}\{b^2(s+ay)^2\mathrm{\Omega }2b^2(c+ax)^2+(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})\mathrm{\Omega })\\ 4\sqrt{2}b^3(c+ax)(s+ay)^2(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})\sqrt{\mathrm{\Omega }}At[x,y]+4(1+b^2|\mathrm{\Xi }|{}_{}{}^{2})^2At[x,y]{}_{}{}^{2}\}.\\ .\end{array}$$
The Gauss curvature $`K`$ and the mean curvature $`H`$ are given by (see (32) and (33)) are given by
$$K=\frac{8a^2b^2}{\left(1+b^2\mathrm{\Xi }^2\right)^4\mathrm{\Omega }^{\frac{9}{2}}}\{b^2(s+ay)^2\mathrm{\Omega }^{\frac{3}{2}}(\mathrm{\Omega }^3+4b^4(c+ax)^2(s+ay)^2(2+\mathrm{\Omega }))$$
$$+4\sqrt{2}b^3(c+ax)(s+ay)^2(1+b^2\mathrm{\Xi }^2)\mathrm{\Omega }^2At[x,y]$$
$$16(1+b^2\mathrm{\Xi }^2)[2b^4\left(c+ax\right)^2\left(s+ay\right)^2(1+b^2\mathrm{\Xi }^2)]\sqrt{\mathrm{\Omega }}At[x,y]^2+$$
$$32\sqrt{2}b(c+ax)(1+b^2\mathrm{\Xi }^2)^2At[x,y]^3\},$$
$$H=\frac{4ab}{\left(1+b^2\left|\mathrm{\Xi }\right|^2\right)^3\mathrm{\Omega }^{\frac{7}{2}}\sqrt{g}}\{b^3(c+ax)(s+ay)^2\mathrm{\Omega }^{\frac{3}{2}}[(1+b^2|\mathrm{\Xi }|^2)(2+\mathrm{\Omega })+$$
$$4(2b^2(s+ay)^2+\mathrm{\Omega })]$$
$$+2\sqrt{2}\left(1+b^2\left|\mathrm{\Xi }\right|^2\right)\mathrm{\Omega }\left(14b^2\left(s+ay\right)^2+b^2\left|\mathrm{\Xi }\right|^2+2b^2\left(s+ay\right)^2\mathrm{\Omega }\right)At[x,y]+$$
$$8b\left(c+ax\right)\left(1+b^2\left|\mathrm{\Xi }\right|^2\right)\left(1+2b^2\left(s+ay\right)^2+b^2\left|\mathrm{\Xi }\right|^2\right)\sqrt{\mathrm{\Omega }}At[x,y]^2$$
$$16\sqrt{2}(1+b^2|\mathrm{\Xi }|^2)^2At[x,y]^3\},$$
respectively.
An example of a surface associated with the 1-vortex solution of the Ishimori system is drawn in Fig 1.
## 11 Conclusions
In this paper we have established some notable connections among the purely topological CS theory, deformations of surfaces, and integrable equations in (2+1)-dimensions. However, many questions remain open and deserve further investigation, such as for example the search for other integrable classes of deformations of surfaces, the determination of the Hamiltonian structure and the possible interpretation of the solutions by a physical point of view. To this regard, in particular we have found exact vortex solutions of the (2+1)-dimensional spin system. Furthermore, we have seen that the dynamics of vortices is governed by a system of the Calogero-Moser type. To conclude, we notice that another approach exists to study integrable (2+1)-dimensional deformations of surfaces, i.e. the method developed mainly by Konopelchenko, Taimanov and coworkers . The essential tool of their procedure is the use of a generalized Weierstrass representation for a conformal immersion of surfaces into $`R^3`$ or $`R^4,`$ together with a linear problem related to this representation. The method devised in allows one to express integrable deformations of surfaces via hierarchies of integrable equations, such as the Nizhnik-Veselov-Novikov and the DS equations, and so on. We think that our approach and that described in should be pursued in parallel, with the purpose to achieve possible complementary results on the link between integrable deformations of surfaces and completely integrable partial differential equations.
## 12 Acknowledgments
The authors are grateful to V.S. Dryuma and B.G. Konopelchenko for very helpful discussions. This work was supported in part by MURST of Italy, INFN-Sezione di Lecce and INTAS (grant 99-1782). One of the authors (R.M.) thanks the Department of Physics of the Lecce University for its warm hospitality.
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# Vertex-Facet Incidences of Unbounded Polyhedra
## 1 Introduction
Every (proper) face of a polytope (i.e., a bounded convex polyhedron) is the convex hull of the vertices it contains, and it is also the intersection of the facets that contain it. Thus, the combinatorial structure of a polytope (i.e., its face lattice) is entirely determined by its (matrix of) vertex-facet incidences. Such a vertex-facet incidence matrix is a useful encoding of the combinatorial structure of a polytope. The software package polymake , for instance, represents this matrix rather compactly, in a section called VERTICES\_IN\_FACETS, while the face lattice of a polytope is not stored, but generated “on demand” only if this is really necessary, because typically the entire face lattice is “much too large.”
But how about not necessarily bounded convex polyhedra? The combinatorics of unbounded polyhedra has received only little attention up to now (for some exceptions see Klee , Billera & Lee , Barnette, Kleinschmidt & Lee , and Lee ). One can, of course, reduce the study of geometrically given unbounded polyhedra to the situation of “a polytope with a distinguished face (at infinity).” But what if only the combinatorics of vertices versus facets is given, and not any data about the situation “at infinity?” In other words, how much can really be said/detected/reconstructed if only a matrix of the vertex-facet incidences is given?
As one observes easily from the example of polyhedral cones, in general the combinatorial structure of an unbounded polyhedron is not determined by its vertex-facet incidences. A $`d`$-dimensional cone may have any possible combinatorial structure of a $`(d1)`$-dimensional polytope (via *homogenization*); but from its vertex-facet incidences one can read off only its number of facets. The point is that, for unbounded polyhedra the combinatorial information is based not only on the vertex-facet incidences, but also on the incidences of extremal rays and facets. For cones, nearly the entire information is contained in the latter incidences. The lattice-theoretic reason for such ambiguities is that the face lattice of an unbounded polyhedron is only co-atomic, but not atomic.
One might, however, suspect that cones are (extreme) examples of rather exotic unbounded polyhedra for which one obviously does not have any chance to reconstruct the combinatorial structure from their vertex-facet incidences, while this might be possible for all “reasonable” polyhedra. For instance, a cone is a quite degenerate polyhedron with respect to several criteria: (i) all its facets have the same set of vertices, (ii) its set of vertices does not have the same dimension as the whole polyhedron, and (iii) it does not have any bounded facet. However, the first main point of this paper (in Section 3) is the construction of more convincing examples of unbounded polyhedra whose face-lattices cannot be reconstructed from their vertex-facet incidence matrices; they have the property that the sets of vertices of facets are distinct, and they even form an anti-chain in the Boolean lattice (a *clutter*); they have bounded facets, and their sets of vertices are full-dimensional.
The second main result (in Section 4) will be that one can, however, detect from the vertex-facet matrix whether the polyhedron under consideration is bounded or not.
Thirdly (in Section 5), we discuss the “unbounded version” of a very basic lemma about polytopes. Indeed, Exercise 0.1 of asks one to prove that any $`d`$-polytope that is both simplicial (every facet has $`d`$ vertices) and simple (every vertex is on $`d`$ facets) must either be a simplex, or a polygon ($`d=2`$). But how about unbounded polyhedra? We prove that a polyhedron that is both simple and simplicial (with the definitions as given here) cannot be unbounded. As a byproduct, we obtain a characterization of those polyhedra that have circulant vertex-facet incidence matrices.
In particular, this paper answers a series of questions that arose in Amaldi, Pfetsch, and Trotter , where the structure of certain independence systems is related to the combinatorics of (possibly unbounded) polyhedra.
## 2 Basic Facts
Let $`P`$ be a $`d`$-polyhedron (i.e., the intersection of a finite number of affine halfspaces with $`dim(P)=d`$) with $`m`$ facets and $`n`$ vertices. We will always assume that $`P`$ is pointed (i.e., it has at least one vertex) and that $`d1`$. In particular, these conditions imply $`n1`$ and $`md1`$. For the basic definitions and facts of polyhedral theory we refer to .
A $`0/1`$-matrix $`A=(a_{fv})\{0,1\}^{m\times n}`$ is a *vertex-facet incidence matrix* of $`P`$ if the vertices and facets of $`P`$ can be numbered by $`\{1,\mathrm{},n\}`$ and $`\{1,\mathrm{},m\}`$, respectively, such that $`a_{fv}=1`$ if and only if the vertex with number $`v`$ is contained in the facet with number $`f`$.
By $`\overline{P}`$ we denote any *polytope* which is projectively equivalent to $`P`$. If $`P`$ is unbounded, then there is a unique maximal element $`F_{\mathrm{}}`$ (the *far face*) among the faces of $`\overline{P}`$ that are not images of faces of $`P`$ under the projective transformation mapping $`P`$ to $`\overline{P}`$. If $`P`$ is bounded, then we define $`F_{\mathrm{}}=\mathrm{}`$. Figure 1 illustrates a three-dimensional example.
We denote by $`\left(P\right)`$ the *face poset* of $`P`$, i.e., the set of non-trivial faces of $`P`$ (excluding $`\mathrm{}`$ and $`P`$ itself), ordered by inclusion. The face poset $`\left(P\right)`$ arises from the face poset $`\left(\overline{P}\right)`$ by removing the far face $`F_{\mathrm{}}`$ (and all its faces). While $`\left(\overline{P}\right)`$ is independent of the actual choice of $`\overline{P}`$, in general it depends on the geometry of $`P`$, not only on its combinatorial structure.
The poset $`𝒱\left(P\right)=\left\{\text{vert}(F)\right|F\text{ non-trivial face of }P\}`$ (where $`\text{vert}(F)`$ is the set of vertices of $`F`$) will play an important role. It can be computed from any vertex-facet incidence matrix $`A\{0,1\}^{m\times n}`$ of $`P`$, since it is the set of all non-empty intersections of subsets of $`\{1,\mathrm{},n\}`$ defined by subsets of the rows of $`A`$. Figure 2 shows the three posets $`\left(P\right)`$, $`\left(\overline{P}\right)`$, and $`𝒱\left(P\right)`$ for the example given in Figure 1.
Let the *graph* $`\mathrm{\Gamma }_P`$ of $`P`$ be the graph on the vertices of $`P`$ defined by the bounded one-dimensional faces of $`P`$ (the *edges*), i.e., $`\mathrm{\Gamma }_P`$ is the subgraph of the graph of $`\overline{P}`$ that is induced by those vertices of $`\overline{P}`$ that are not contained in $`F_{\mathrm{}}`$. Two vertices of $`P`$ are connected by an edge of $`P`$ if and only if there is a face of $`P`$ which contains exactly these two vertices. Moreover, we can compute $`𝒱\left(P\right)`$ from any vertex-facet incidence matrix of $`P`$. In particular, we can find $`\mathrm{\Gamma }_P`$ from the vertex-facet incidences of $`P`$.
We will use the following fact, which is a consequence of the correctness of the Simplex-Algorithm for Linear Programming.
###### Lemma 2.1
For every polyhedron $`P`$, the graph $`\mathrm{\Gamma }_P`$ is connected. Moreover, all faces of $`P`$ induce connected subgraphs of $`\mathrm{\Gamma }_P`$.
Let $`P`$ be a pointed $`d`$-polyhedron ($`d1`$). Then, $`P`$ is called *simple* if every vertex of $`P`$ is contained in precisely $`d`$ facets (or, equivalently, if precisely $`d`$ edges and extremal rays are incident to each vertex), and $`P`$ is called *simplicial* if every facet of $`P`$ has precisely $`d`$ vertices. These notions generalize the well-known notions *simple* and *simplicial* for polytopes. While this generalization is standard for simple polyhedra, it is not common for simplicial polyhedra. Thus, it seems to be worth to mention that simplicial unbounded polyhedra form a non-trivial class of polyhedra. For instance, by a modification of the construction of a prism, one easily sees that every simplicial $`d`$-polytope can occur as the far face of a simplicial unbounded $`(d+1)`$-polyhedron.
## 3 Reconstructing Polyhedra from Vertex-Facet Incidences
In this section, we consider conditions under which it is possible to compute $`\left(P\right)`$ from the vertex-facet incidences of an (unbounded) $`d`$-polyhedron $`P`$. Obviously, given any vertex-facet incidence matrix of a pointed $`d`$-polyhedron $`P`$ it is easy to decide whether $`d\{1,2\}`$. Furthermore, if $`d\{1,2\}`$, one can immediately read off $`\left(P\right)`$ from the vertex-facet incidences. Thus, for the rest of this section we restrict our attention to $`d`$-polyhedra with $`d3`$.
The example of cones shows that reconstructing $`\left(P\right)`$ from the vertex-facet incidences of a $`d`$-polyhedron $`P`$ with $`d4`$ is impossible in general, even if additionally the dimension $`d`$ is specified. Furthermore, the same example demonstrates that it is, in general, impossible to detect the dimension of a $`d`$-polyhedron from its vertex-facet incidences for $`d3`$. However, for $`d=3`$ these dimensional ambiguities occur for cones only.
###### Proposition 3.1
Given a vertex-facet incidence matrix of a $`d`$-polyhedron $`P`$ with $`d3`$, it is possible to decide whether $`d=3`$ or $`d4`$, unless $`P`$ is a cone with more than three facets.
Proof. If $`P`$ is a cone with three facets (i.e., $`n=1`$ and $`m=3`$) then clearly $`d=3`$ holds. If $`P`$ is not a cone, then it must have at least two vertices. Thus (by Lemma 2.1) $`P`$ has at least one edge (which we can tell from the vertex-facet incidences of $`P`$). This edge is contained in precisely two facets of $`P`$ if $`d=3`$; otherwise, it is contained in more than two facets. $`\mathrm{}`$
In dimensions larger than three, cones are not the only polyhedra for which one cannot tell the dimension from the vertex-facet incidences. For instance, let $`Q`$ be some $`d^{}`$-polytope and let $`C`$ be a $`d^{\prime \prime }`$-dimensional polyhedral cone with $`m4`$ facets. Then $`P=Q\times C`$ will be a $`(d^{}+d^{\prime \prime })`$-dimensional polyhedron whose vertex-facet incidences only depend on $`Q`$ and $`m`$, while its dimension can be any number between $`d^{}+3`$ and $`d^{}+m`$. In particular, dimensional ambiguities already occur for $`4`$-polyhedra not being cones.
However, the cartesian products constructed above are also “cone-like” in the sense that they do not have any bounded facet.
###### Proposition 3.2
Given a vertex-facet incidence matrix of a $`d`$-polyhedron $`P`$ that has a bounded facet, one can determine $`d`$. Furthermore, one can decide from the vertex-facet incidences of $`P`$ whether it has a bounded facet or not.
Proof. If $`P`$ has a bounded facet, then the maximum length of a chain in $`𝒱\left(P\right)`$ is $`d1`$, thus one can compute $`d`$ from $`𝒱\left(P\right)`$ in this case. Corollary 4.6 proves the second statement of the proposition. $`\mathrm{}`$
Propositions 3.1 and 3.2 might suggest to ask if the entire combinatorial structure of a $`d`$-polyhedron can be reconstructed from its vertex-facet incidences if $`d=3`$ or if $`P`$ has a bounded facet. However, the example given in Figure 3 shows that both answers are “no”.
The crucial feature of the example is that one can reflect the “lower” parts in the drawings without affecting the vertex-facet incidences while changing the face poset (e.g., in contrast to the left polyhedron the right one has two adjacent unbounded facets that have three vertices each). For three-dimensional polyhedra this is more or less the only kind of ambiguity that can arise.
###### Proposition 3.3
Given the vertex-facet incidences of a $`3`$-polyhedron $`P`$, for which $`\mathrm{\Gamma }_P`$ is $`2`$-connected, one can determine $`\left(P\right)`$.
Proof. One can compute $`\mathrm{\Gamma }_P`$ from the vertex-facet incidences of $`P`$, and thus, one finds the graph of each facet of $`P`$. If all these graphs of facets are cycles then $`P`$ is bounded and the statement is clear. Otherwise, due to the $`2`$-connectedness of $`\mathrm{\Gamma }_P`$, there is a unique (up to reorientation) way to arrange the paths that are the graphs of the unbounded facets of $`P`$ as a cycle. From this cycle, it is easy to determine the incidences of extremal rays and facets of $`P`$, which then allow to reconstruct the entire combinatorial structure of $`P`$. $`\mathrm{}`$
In larger dimensions, however, it is not true that higher connectedness of the graph of a polyhedron is a sufficient condition for the possibility to reconstruct its combinatorial structure from its vertex-facet incidences. Figure 4 shows Schlegel-diagrams of (truncations of) two unbounded $`4`$-polyhedra. These two polyhedra have the same vertex-facet incidences and a $`3`$-connected graph, although their face posets are different (e.g., the right polyhedron has an extremal ray more than the left one)<sup>1</sup><sup>1</sup>1The data of these polyhedra as well as explanations on their construction can be found in the EG-Models archive at: http://www-sfb288.math.tu-berlin.de/eg-models/models/polytopes/2000.05.001/\_preview.html.
The examples illustrated in Figures 3 and 4 show that “cone-like” polyhedra are not the only ones that cannot be reconstructed from their vertex-facet incidences (even not in dimensions three and four). The polyhedra in both examples are quite different from cones; each of them has a full-dimensional vertex set, bounded facets, and the property that no two facets have the same vertex set. Furthermore, in the four-dimensional example, the vertex sets of the facets even form an anti-chain (as promised in the introduction).
Nevertheless, any ambiguities in reconstructing the face poset of an unbounded polyhedron from its vertex-facet incidences arise from some degeneracy of $`P`$.
###### Theorem 3.4
Given the vertex-facet incidences of a simple polyhedron $`P`$, one can determine $`\left(P\right)`$.
Proof. Let $`v`$ be a vertex of a simple $`d`$-polyhedron $`P`$ and let $`F_1,\mathrm{},F_d`$ be the facets of $`P`$ that contain $`v`$. Then the edges and extremal rays containing $`v`$ are precisely
$$\underset{i\{1,\mathrm{},d\}\{i_0\}}{}F_i(i_0=1,\mathrm{},d).$$
Since we can compute the edges of $`P`$ from a vertex-facet incidence matrix, we can thus also deduce (combinatorially) the extremal rays of $`P`$ and the information which ray is contained in which facets. From that, we can deduce the entire face poset of $`P`$. $`\mathrm{}`$
Again, the example of cones shows that without dimension information one can (in general) not decide from the vertex-facet incidences of a polyhedron if it is simple.
All algorithms described in this section can be implemented such that their running time is bounded by a polynomial in $`|𝒱\left(P\right)|`$.
To summarize the results in this section: we presented large classes of (unbounded) polyhedra for which the combinatorial structures can be reconstructed from their vertex-facet incidences as well as several examples of polyhedra, for which this is not possible. Unfortunately, these results do not yield a *characterization* of the class of those polyhedra that allow such reconstructions.
## 4 Detecting Boundedness
In this section, we show that one can decide from the vertex-facet incidences of a pointed polyhedron $`P`$ whether it is bounded or not. It turns out that this only depends on the Euler characteristic of (the order complex of) $`𝒱\left(P\right)`$. Thus, it can be read off from the Möbius function of $`𝒱\left(P\right)`$.
We recall some basic facts from topological combinatorics (see Björner ). Let $`\mathrm{\Pi }`$ be a finite poset. The *order complex* $`\mathrm{\Delta }(\mathrm{\Pi })`$ of $`\mathrm{\Pi }`$ is the finite simplicial complex of all chains in $`\mathrm{\Pi }`$. We will use terminology from topology in the context of finite posets such as $`\mathrm{\Pi }`$. Throughout, this is meant to refer to $`\mathrm{\Delta }(\mathrm{\Pi })`$ (i.e., any geometric realization of $`\mathrm{\Delta }(\mathrm{\Pi })`$, endowed with its standard topology).
It is well-known that the order complex $`\mathrm{\Delta }(\left(P\right))`$ of a bounded $`d`$-polytope $`P`$ is isomorphic (as a simplicial complex) to the barycentric subdivision of the boundary $`P`$ of $`P`$. In particular, the topological type of $`𝒱\left(P\right)`$ is well-known in this case.
###### Lemma 4.1
If $`P`$ is a $`d`$-polytope, then $`\left(P\right)`$ is homeomorphic to the $`(d1)`$-sphere.
If $`P`$ is an unbounded (pointed) polyhedron, then we can consider $`\left(P\right)`$ as the sub-poset of $`\left(\overline{P}\right)`$ that consists of all faces of $`\overline{P}`$ that are not contained in $`F_{\mathrm{}}`$. Thus, we will identify $`\mathrm{\Delta }(\left(P\right))`$ with the sub-complex of $`\mathrm{\Delta }(\left(\overline{P}\right))`$ that is induced by all chains $`\{F\}`$, where $`F`$ is a face of $`\overline{P}`$ with $`FF_{\mathrm{}}`$.
###### Lemma 4.2
If $`P`$ is an unbounded (pointed) polyhedron, then $`\left(P\right)`$ is contractible.
Proof. By Lemma 4.1, $`\mathrm{\Delta }(\left(\overline{P}\right))`$ is homeomorphic to a sphere. The induced subcomplexes $`A=\mathrm{\Delta }(\left(P\right))`$ and $`B=\mathrm{\Delta }(\left(F_{\mathrm{}}\right))`$ cover all vertices (i.e., one-element chains of $`\left(\overline{P}\right)`$) of $`\mathrm{\Delta }(\left(\overline{P}\right))`$. Using barycentric coordinates, it is seen that $`\mathrm{\Delta }(\left(\overline{P}\right))B`$ retracts onto $`A`$. Thus, $`A`$ has the same homotopy-type as $`\mathrm{\Delta }(\left(\overline{P}\right))B`$, where the latter is a simplicial sphere minus an induced ball. Hence, $`\left(P\right)`$ is contractible. $`\mathrm{}`$
The two lemmas allow one to distinguish between the face posets of bounded and unbounded polyhedra. Of course, there are simpler ways to decide whether a face poset belongs to a bounded or to an unbounded polyhedron (e.g., checking if every rank one element is a join). However, in general we cannot reconstruct the face poset of a polyhedron $`P`$ from its vertex-facet incidences (see Section 3). Instead, we need criteria allowing to distinguish between bounded and unbounded polyhedra that can be computed from $`𝒱\left(P\right)`$. It turns out that the topological criteria provided by Lemmas 4.1 and 4.2 can be exploited for this.
Consider the poset maps $`\varphi :\left(P\right)𝒱\left(P\right)`$, mapping a face $`F`$ of a pointed polyhedron $`P`$ to $`\text{vert}(F)`$, and $`\psi :𝒱\left(P\right)\left(P\right)`$, mapping the vertex set $`S`$ of a face to the minimal face containing $`S`$. Both $`\varphi `$ and $`\psi `$ are order preserving. Moreover, $`\varphi (\psi (S))=S`$ and $`\psi (\varphi (F))F`$.
###### Lemma 4.3
Let $`P`$ be a pointed polyhedron. Then the face poset $`\left(P\right)`$ is homotopy-equivalent to the poset $`𝒱\left(P\right)`$.
Proof. Setting $`f(F)=\psi (\varphi (F))`$ defines an order preserving map from $`\left(P\right)`$ into itself such that each face $`F`$ is comparable with its image $`f(F)`$. From the Order Homotopy Theorem \[4, Corollary 10.12\], we infer that $`\left(P\right)`$ is homotopy-equivalent to the image $`f(\left(P\right))`$. In fact, $`f(f(F))=f(F)`$, and hence $`f(\left(P\right))`$ is a strong deformation retract of $`\left(P\right)`$. This proves the lemma, since $`\psi `$ is a poset isomorphism from $`𝒱\left(P\right)`$ onto $`\psi (𝒱\left(P\right))=f(\left(P\right))`$. $`\mathrm{}`$
The reduced Euler characteristic of (the order complex of) a poset $`\mathrm{\Pi }`$ is denoted by $`\stackrel{~}{\chi }(\mathrm{\Pi })`$, i.e.,
$$\stackrel{~}{\chi }(\mathrm{\Pi })=\underset{i=1}{\overset{D}{}}(1)^if_i(\mathrm{\Delta }(\mathrm{\Pi }))$$
(where $`f_i(\mathrm{\Delta }(\mathrm{\Pi }))`$ is the number of $`i`$-faces of $`\mathrm{\Delta }(\mathrm{\Pi })`$, and $`D`$ is the dimension of $`\mathrm{\Delta }(\mathrm{\Pi })`$). The following result in particular shows that a polytope and an unbounded polyhedron cannot have isomorphic vertex-facet incidences.
###### Theorem 4.4
Let $`P`$ be a pointed polyhedron. Then $`P`$ is bounded if and only if $`\stackrel{~}{\chi }(𝒱\left(P\right))0`$.
Proof. The reduced Euler characteristic of a $`(d1)`$-sphere equals $`(1)^{d1}`$, while the reduced Euler characteristic of a contractible space vanishes. Thus the claim follows from Lemma 4.1, Lemma 4.2, and Lemma 4.3. $`\mathrm{}`$
As an example consider the case where the unbounded polyhedron $`P`$ has a face $`F`$ which contains all vertices of $`P`$. Then $`\mathrm{\Delta }(𝒱\left(P\right))`$ is a cone over $`F`$ (in the sense of simplicial topology); in particular, it is contractible and thus $`\stackrel{~}{\chi }(𝒱\left(P\right))=0`$.
The reduced Euler characteristic of the poset $`𝒱\left(P\right)`$ can be computed efficiently as follows. By adjoining an artificial top element $`\widehat{1}`$ and an artificial bottom element $`\widehat{0}`$, the poset $`𝒱\left(P\right)`$ becomes a lattice $`\widehat{𝒱}\left(P\right)`$. Note that we adjoin $`\widehat{1}`$ also in the case where $`𝒱\left(P\right)`$ already has a top element corresponding to a face containing all vertices of $`P`$.
For every element $`S\widehat{𝒱}\left(P\right)`$ we define the *Möbius function*, see Rota and Stanley ,
$$\mu (S)=\{\begin{array}{cc}1\hfill & \text{if }S=\widehat{0},\hfill \\ \underset{S^{}S}{}\mu (S^{})\hfill & \text{otherwise}.\hfill \end{array}$$
The *Möbius number* $`\mu (𝒱\left(P\right))=\mu (\widehat{1})`$ of $`𝒱\left(P\right)`$ can be computed in time bounded polynomially in $`|𝒱\left(P\right)|`$. Since it is well-known (see Stanley \[12, 3.8.6\]) that
$$\mu (𝒱\left(P\right))=\stackrel{~}{\chi }(𝒱\left(P\right)),$$
(1)
this proves the following complexity result.
###### Corollary 4.5
There is an algorithm that decides for every vertex-facet incidence matrix of a polyhedron $`P`$ if $`P`$ is bounded. Its running time is bounded by a polynomial in $`|𝒱\left(P\right)|`$.
Actually, Theorem 4.4 allows to decide even more from the vertex-facet incidences of a polyhedron $`P`$. Once we have computed $`𝒱\left(P\right)`$ we clearly can also determine $`\widehat{𝒱}\left(F\right)`$ for every facet $`F`$ of $`P`$ (since we know $`\text{vert}(F)`$ for every facet $`F`$ of $`P`$). This is the interval between $`\widehat{0}`$ and $`\text{vert}(F)`$ in the lattice $`\widehat{𝒱}\left(P\right)`$, where we have to add an additional top element $`\widehat{1}`$ if there is some other facet $`F^{}`$ of $`P`$ containing $`\text{vert}(F)`$.
###### Corollary 4.6
There is an algorithm that tells from a vertex-facet incidence matrix of a polyhedron $`P`$ which facets of $`P`$ are bounded. Its running time is bounded by a polynomial in $`|𝒱\left(P\right)|`$.
## 5 Simple and Simplicial Polyhedra
It is a well-known fact \[13, Exerc. 0.1\] that a $`d`$-*polytope* which is both simple and simplicial is a simplex or a polygon. Both properties (simplicity as well as simpliciality) can be viewed as properties of vertex-facet incidences (see Section 2). In this section, we generalize the known result on polytopes to not necessarily bounded $`d`$-polyhedra with $`d2`$.
###### Theorem 5.1
For $`d2`$, every simple and simplicial $`d`$-polyhedron is a simplex or a polygon. In other words, unbounded simple and simplicial polyhedra do not exist.
Our proof of Theorem 5.1 is organized into two parts. The first part shows that the graph $`\mathrm{\Gamma }_P`$ of a simple and simplicial polyhedron $`P`$ is either a complete graph or a cycle. In the second part, we further deduce that a simple and simplicial polyhedron has a *circulant* vertex-facet incidence matrix. The proof of Theorem 5.1 is then completed by showing that no unbounded $`d`$-polyhedron (with $`d2`$) can have a circulant vertex-facet incidence matrix. Furthermore, Propositions 5.8 and 5.10 yield characterizations of those polyhedra that have circulant vertex-facet incidence matrices.
### 5.1 Graphs of Simple and Simplicial Polyhedra
Throughout this section, let $`P`$ be a pointed simple and simplicial $`d`$-polyhedron with $`n`$ vertices and $`d2`$. Double counting yields that $`P`$ must also have $`n`$ facets. In particular, we have $`n>d`$ (since otherwise $`P`$ would be a cone, which is simple and simplicial only for $`d=1`$). We denote by $`V_P=\text{vert}(P)`$ the set of vertices of $`P`$. For $`SV_P`$ let $`\mathrm{\Theta }(S)`$ be the set of all facets of $`P`$ that contain $`S`$. Recall that (since $`P`$ is simple) two vertices $`v`$ and $`w`$ of $`P`$ form an edge if and only if $`|\mathrm{\Theta }(\{v,w\})|=d1`$.
###### Lemma 5.2
Two different facets of $`P`$ cannot have the same set of vertices.
Proof. Suppose that there are two facets $`F_1`$ and $`F_2`$ of $`P`$ ($`F_1F_2`$) with $`\text{vert}(F_1)=\text{vert}(F_2)=:S`$. Since $`n>d`$, and since $`\mathrm{\Gamma }_P`$ is connected, there must be a vertex $`vS`$ that is a neighbor of some vertex $`wS`$. Hence, we have $`|\mathrm{\Theta }(\{v,w\})|=d1`$. Because of $`|\mathrm{\Theta }(\{w\})|=d`$ and $`F_1,F_2\mathrm{\Theta }(\{w\})\mathrm{\Theta }(\{v,w\})`$ this implies $`F_1\mathrm{\Theta }(\{v,w\})`$ or $`F_2\mathrm{\Theta }(\{v,w\})`$, which in both cases yields a contradiction to $`vS`$. $`\mathrm{}`$
For $`SV_P`$, define $`\mathrm{\Omega }(S)`$ to be the set of those facets of $`P`$ that have non-empty intersection with $`S`$.
###### Lemma 5.3
Let $`SV_P`$ with $`|S|>0`$. Then $`|\mathrm{\Omega }(S)|\mathrm{min}\{n,d+|S|1\}`$.
Proof. If $`|\mathrm{\Omega }(S)|=n`$, then the claim obviously is correct. Therefore, assume $`|\mathrm{\Omega }(S)|<n`$. Since $`\mathrm{\Gamma }_P`$ is connected, the vertices in $`V_PS=\{z_1,\mathrm{},z_r\}`$ ($`r=n|S|`$) can be ordered such that $`z_{i+1}`$ is adjacent to some vertex of $`S_i=S\{z_1,\mathrm{},z_i\}`$ for each $`i\{0,\mathrm{},r1\}`$ (additionally, define $`S_r=S\{z_1,\mathrm{},z_r\}`$). Clearly $`|\mathrm{\Omega }(S_i)||\mathrm{\Omega }(S_{i1})|+1`$, since vertex $`z_i`$ has $`d1`$ facets in common with some vertex in $`S_{i1}`$.
Define $`l`$ to be the last $`i`$, such that $`|\mathrm{\Omega }(S_i)|=|\mathrm{\Omega }(S_{i1})|+1`$, i.e., $`l`$ is the last index, where we encounter a new facet ($`l`$ is well-defined due to $`|\mathrm{\Omega }(S)|<n`$). Since this facet must contain $`d1`$ vertices from $`V_PS_l`$, we have $`rld1`$, which yields $`nld+|S|1`$.
Furthermore, we have $`|\mathrm{\Omega }(S)|+ln`$, since $`S_l`$ intersects all facets. It follows $`|\mathrm{\Omega }(S)|nld+|S|1`$. $`\mathrm{}`$
For $`SV_P`$ let $`\mathrm{\Gamma }_P(S)`$ be the subgraph of $`\mathrm{\Gamma }_P`$ induced by $`S`$.
###### Lemma 5.4
Let $`SV_P`$ with $`0<|S|d`$, such that $`\mathrm{\Gamma }_P(S)`$ is connected. Then $`|\mathrm{\Theta }(S)|=d|S|+1`$ holds.
Proof. Since $`\mathrm{\Gamma }_P(S)`$ is a connected subgraph of the connected graph $`\mathrm{\Gamma }_P`$ (which has $`n>d`$ vertices), there is a chain $`\mathrm{}S_1S_2\mathrm{}S_d`$ with $`S_{|S|}=S`$, such that $`|S_i|=i`$ and $`\mathrm{\Gamma }_P(S_i)`$ is connected for all $`i`$.
For every $`1<id`$, the vertex $`v`$ with $`S_iS_{i1}=\{v\}`$ is connected to some vertex $`wS_{i1}`$. From $`|\mathrm{\Theta }(\{w\})\mathrm{\Theta }(\{v\})|=1`$ we infer $`|\mathrm{\Theta }(S_{i1})\mathrm{\Theta }(\{v\})|1`$, and thus, $`|\mathrm{\Theta }(S_i)||\mathrm{\Theta }(S_{i1})|1`$. Together with $`|\mathrm{\Theta }(S_1)|=d`$ (since $`P`$ is simple) and $`|\mathrm{\Theta }(S_d)|1`$ (by Lemma 5.2), this implies $`|\mathrm{\Theta }(S_i)|=di+1`$ for all $`1id`$. $`\mathrm{}`$
The next three lemmas show that $`\mathrm{\Gamma }_P`$ has a very special structure.
###### Lemma 5.5
If $`\mathrm{\Gamma }_P`$ contains a cycle $`C`$ of size $`k>d`$, then $`\mathrm{\Gamma }_P`$ is the cycle $`C`$ or a complete graph on $`n=d+1`$ nodes.
Proof. Let $`C=(v_0,\mathrm{},v_{k1},v_0)`$ be a cycle of size $`k>d`$ in $`\mathrm{\Gamma }_P`$. In the following, all indices are taken modulo k. For $`0ik1`$ define the set $`C_i=\{v_i,\mathrm{},v_{i+d1}\}`$ of size $`d`$. Clearly, $`\mathrm{\Gamma }_P(C_i)`$ is connected, and, by Lemma 5.4, there exists exactly one facet $`F_i`$ with $`\mathrm{\Theta }(C_i)=\{F_i\}`$. Due to $`k>d`$, the facets $`F_0,\mathrm{},F_{k1}`$ are pairwise distinct. This means that $`\mathrm{\Theta }(\{v_i\})=\{F_{id+1},\mathrm{},F_i\}`$ (since $`P`$ is simple) and $`\text{vert}(F_i)=C_i`$ (since $`P`$ is simplicial). Hence, every vertex that is adjacent to one of the nodes $`v_0,\mathrm{},v_{k1}`$ must be contained in at least one (more precisely, in $`d1>0`$) of the facets $`F_0,\mathrm{},F_{k1}`$, and thus it lies in $`\{v_0,\mathrm{},v_{k1}\}`$.
Since $`\mathrm{\Gamma }_P`$ is connected, this means that $`n=k`$. For $`n=d+1`$ this immediately yields that $`\mathrm{\Gamma }_P`$ is a complete graph on $`n=d+1`$ nodes, while for $`n>d+1`$ one finds that $`\mathrm{\Gamma }_P`$ is the cycle $`C`$ (since, in this case, $`|\mathrm{\Theta }(\{v_i\})\mathrm{\Theta }(\{v_j\})|=d1`$ if and only if $`ji\pm 1modk`$). $`\mathrm{}`$
###### Lemma 5.6
If $`\mathrm{\Gamma }_P`$ contains a cycle of length $`kd`$, then $`\mathrm{\Gamma }_P`$ is a complete graph on $`n=d+1`$ nodes.
Proof. Let $`\stackrel{~}{C}=(v_0,\mathrm{},v_{k1},v_0)`$ be a cycle in $`\mathrm{\Gamma }_P`$ of size $`kd`$. For each $`i\{0,\mathrm{},k1\}`$ define $`\stackrel{~}{C}_i=\{v_0,\mathrm{},v_i\}`$. Taking all indices modulo $`k`$, we have $`|\mathrm{\Theta }(\{v_i,v_{i+1}\})|=d1`$ for each $`i`$, and hence, there are facets $`F_i`$ and $`G_i`$ with
$$\mathrm{\Theta }(\{v_i\})\mathrm{\Theta }(\{v_{i+1}\})=\{F_i\}\text{and}\mathrm{\Theta }(\{v_{i+1}\})\mathrm{\Theta }(\{v_i\})=\{G_i\}.$$
It follows that
$$\mathrm{\Theta }(\stackrel{~}{C}_{k1})=\mathrm{\Theta }(\stackrel{~}{C}_0)\{F_0,\mathrm{},F_{k1}\}.$$
(2)
If $`\mathrm{\Gamma }_P`$ is not complete, then $`n>d+1`$ holds, and we infer from Lemma 5.3 that $`|\mathrm{\Omega }(\stackrel{~}{C}_2)|d+2`$, which implies $`G_0,G_1\mathrm{\Theta }(\stackrel{~}{C}_0)`$ (with $`G_0G_1`$). Due to $`\{F_0,\mathrm{},F_{k1}\}=\{G_0,\mathrm{},G_{k1}\}`$, Equation (2) implies
$$|\mathrm{\Theta }(\stackrel{~}{C}_{k1})||\mathrm{\Theta }(\stackrel{~}{C}_0)|(k2)=dk+2,$$
contradicting Lemma 5.4. $`\mathrm{}`$
By the above two lemmas, $`\mathrm{\Gamma }_P`$ cannot contain any cycles, unless it is complete or a cycle itself. Thus, we are left with the case of $`\mathrm{\Gamma }_P`$ not containing any cycles at all.
###### Lemma 5.7
$`\mathrm{\Gamma }_P`$ is not a tree.
Proof. Assume $`\mathrm{\Gamma }_P`$ is a tree. Let $`vV_P`$ be a leaf of $`\mathrm{\Gamma }_P`$ with $`u`$ being the unique vertex of $`\mathrm{\Gamma }_P`$ adjacent to $`v`$. Due to $`|\mathrm{\Theta }(\{v\})\mathrm{\Theta }(\{v,u\})|=1`$, there is one facet that induces a subgraph of $`\mathrm{\Gamma }_P`$ in which $`v`$ is isolated. This, however, is a contradiction to Lemma 2.1. $`\mathrm{}`$
Altogether this proves the following.
###### Proposition 5.8
Let $`P`$ be a simple and simplicial $`d`$-polyhedron ($`d2`$) with $`n`$ vertices. Then $`\mathrm{\Gamma }_P`$ is an $`n`$-cycle or a complete graph on $`n=d+1`$ nodes.
It is worth to mention that one can generalize Proposition 5.8 in the following way. Let $`A`$ be a $`0/1`$-matrix of size $`n\times n`$ with row and column sums $`d`$. Define a graph $`\mathrm{\Gamma }_A`$ on the columns of $`A`$, such that two columns are adjacent if and only if they have exactly $`d1`$ ones in common rows. Then, by the same arguments as above, one can show that the connectedness of $`\mathrm{\Gamma }_A`$ already implies that $`\mathrm{\Gamma }_A`$ is a cycle or a complete graph. The only difference in the proof arises in Lemma 5.7. Here one has to prove additionally that, for each row, the subgraph of $`\mathrm{\Gamma }_A`$ that is induced by the ones in that row is connected (if $`\mathrm{\Gamma }_A`$ is connected).
### 5.2 Circulant Matrices
We will now exploit Proposition 5.8 to show that every simple and simplicial polyhedron has a very special vertex-facet incidence matrix.
Let $`n,d`$ be integers satisfying $`1dn`$. The $`(n,d)`$-*circulant* $`M(n,d)`$ is the $`n\times n`$-matrix with $`0/1`$ entries whose coefficients $`m_{ij}`$ ($`i,j\{0,\mathrm{},n1\}`$) are defined as follows:
$$m_{ij}=\{\begin{array}{cc}1\hfill & \text{if }j\{i,i+1modn,\mathrm{},i+d1modn\}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
For $`d1`$, the $`(d+1,d)`$-circulant is an incidence matrix of the $`d`$-simplex, and for $`n3`$, the $`(n,2)`$-circulant is an incidence matrix of the ($`2`$-dimensional) $`n`$-gon.
###### Proposition 5.9
A polyhedron $`P`$ is simple and simplicial if and only if it has a circulant $`M(n,d)`$ as a vertex-facet incidence matrix. In this case, $`dim(P)=d`$.
Proof. For the “if”-direction of the proof, let $`P`$ be a polyhedron with a vertex-facet incidence matrix $`M(n,d)`$ ($`1dn`$). The cases $`d=1`$ (implying $`n\{1,2\}`$) as well as $`d=n`$ (implying $`d=n=1`$) are trivial. Therefore, let $`2d<n`$. Obviously, it suffices to show $`dim(P)=d`$. To each row $`i\{0,\mathrm{},n1\}`$ of $`M(n,d)`$ there corresponds a facet $`F_i`$ of $`P`$. For $`0jd1`$ define $`G_j=F_0\mathrm{}F_j`$. Clearly, $`G_jG_{j+1}`$ holds for $`0j<d1`$. Due to $`\text{vert}(G_j)=\text{vert}(F_0)\mathrm{}\text{vert}(F_j)`$ it follows $`\text{vert}(G_j)\text{vert}(G_{j+1})`$ and therefore $`G_jG_{j+1}`$. Now $`F_0=G_0G_1\mathrm{}G_{d1}`$ is a (decreasing) chain of length $`d1`$ in the face poset of $`P`$. Hence we have $`dimPd`$. Since each vertex must be contained in at least $`dimP`$ facets it follows that $`dimPd`$ (because each vertex of $`P`$ is contained in precisely $`d`$ facets).
Conversely, let $`P`$ be a simple and simplicial $`d`$-polyhedron ($`d1`$) with $`n`$ vertices. The case $`d=1`$ is checked easily. Thus, assume $`d2`$. By Proposition 5.8, $`\mathrm{\Gamma }_P`$ either is a complete graph on $`n=d+1`$ nodes or it is a cycle. In the first case, every vertex-facet incidence matrix of $`P`$ is the complement of a permutation matrix, which can be transformed to $`M(n,d)`$ by a suitable permutation of its rows. In the second case, consider any vertex-facet incidence matrix $`A`$ of $`P`$, where the columns are assumed to be ordered according to the cycle $`\mathrm{\Gamma }_P`$. Call two positions $`(i,j)`$ and $`(i,k)`$ in $`A=(a_{fv})`$ ($`f,v\{0,\mathrm{},n1\}`$) *mates* if $`kj+1(modn)`$ and $`a_{ij}=a_{ik}=1`$. Walking around the cycle $`\mathrm{\Gamma }_P`$, we find that the total number of mates in $`A`$ is precisely $`n(d1)`$ (because every edge is contained in precisely $`d1`$ facets). But then, since every row of $`A`$ has only $`d`$ ones (because $`P`$ is simplicial), it follows that in each row the ones must appear consecutively (modulo $`n`$). Denote by $`s(i)`$ the starting position of the block of ones in row $`i`$. Because there are no equal rows in $`A`$ (by Lemma 5.2) we deduce that $`s`$ defines a permutation of the rows of $`A`$ which tells us how to transform $`A`$ to $`M(n,d)`$. $`\mathrm{}`$
The following result finishes the proof of Theorem 5.1 (via Proposition 5.9).
###### Proposition 5.10
If a polyhedron $`P`$ has $`M(n,d)`$ ($`2d<n`$) as a vertex-facet incidence matrix, then $`n=d+1`$ ($`P`$ is a $`d`$-simplex) or $`d=2`$ ($`P`$ is an $`n`$-gon).
Proof. If $`n=d+1`$, then $`M(n,d)`$ is a vertex-facet incidence matrix of a $`d`$-simplex. Hence, by Theorem 4.4, $`P`$ cannot be unbounded, and thus it must be a $`d`$-simplex as well. Therefore, in the following we will assume $`n>d+1`$.
Let us first treat the case $`d+1<n<2d1`$. Consider the facets $`F`$ and $`F^{}`$ corresponding to rows $`0`$ and $`nd+1`$, respectively. If we identify the vertices of $`P`$ with the column indices $`\{0,\mathrm{},n1\}`$ of $`M(n,d)`$, then the vertex set of the face $`G=FF^{}`$ is $`\{0\}\{nd+1,\mathrm{},d1\}`$, where $`\{nd+1,\mathrm{},d1\}\mathrm{}`$ (due to $`n<2d1`$). By Propositions 5.9 and 5.8, $`\mathrm{\Gamma }_P`$ is an $`n`$-cycle (due to $`n>d+1`$). Since neither vertex $`1`$ nor vertex $`n1`$, which are the only neighbors of $`0`$ in $`\mathrm{\Gamma }_P`$, are contained in $`G`$, we conclude that the subgraph of $`\mathrm{\Gamma }_P`$ induced by $`G`$ is disconnected, which is a contradiction to Lemma 2.1.
Hence, we can assume $`n2d1`$. This implies
$$𝒱\left(P\right)=\left\{\{i,\mathrm{},i+s1\}\right|i\{0,\mathrm{},n1\},s\{1,\mathrm{},d\}\}$$
(where, again, all indices are to be taken modulo n), i.e., $`𝒱\left(P\right)`$ consists of all (cyclic) intervals of $`\{0,\mathrm{},n1\}`$ with at least one and at most $`d`$ elements. We will compute the Möbius function $`\mu `$ (see Section 4) on the lattice $`\widehat{𝒱}\left(P\right)`$ (which arises by adding artificial top and bottom elements $`\widehat{1}`$ and $`\widehat{0}`$ to $`𝒱\left(P\right)`$). For each $`s\{1,\mathrm{},d\}`$ let $`\mu (s)=\mu (\{0,\mathrm{},s1\})`$. Obviously, for every $`F𝒱\left(P\right)`$ with $`|F|=s`$ we have $`\mu (F)=\mu (s)`$. In particular, one readily deduces $`\mu (1)=1`$ and $`\mu (2)=(1+2(1))=1`$. For $`3sd`$ we then infer (by induction) $`\mu (s)=(1+s(1)+(s1)(+1))=0`$. Thus, we finally calculate
$$\mu (𝒱\left(P\right))=\mu (\widehat{1})=(1+n(1)+n(+1))=1,$$
which by (1) and Theorem 4.4 implies that $`P`$ is bounded (and, hence, an $`n`$-gon).
(Alternatively, one could derive from the *Nerve Lemma* \[4, Theorem 10.7\] that $`𝒱\left(P\right)`$ is homotopy-equivalent to a circle for $`n2d1`$, and thus, $`P`$ must be a polygon.) $`\mathrm{}`$
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# 1 Introduction
## 1 Introduction
Two dimensional (2D) nonlinear sigma models and four dimensional non-abelian gauge theories have several similarities. Both of them enjoy the property of the asymptotic freedom. They are both massless in the perturbation theory, whereas they acquire the mass gap or the string tension in the non-perturbative treatment. Although it is difficult to solve QCD in analytical way, 2D nonlinear sigma models can be solved by the large $`N`$ expansion and helps us to understand various non-perturbative phenomena in four dimensional gauge theories.
In the nonperturbative treatment of nonlinear $`\sigma `$ models, auxiliary field method play an important role. As an example, let us consider the supersymmetric 2D nonlinear $`\sigma `$ model with $`O(N)`$ symmetry. The bosonic field $`\stackrel{}{A}(x)`$ takes the value on the real $`N1`$ dimensional sphere $`S^{N1}`$ with a radius $`\frac{\sqrt{N}}{g}`$. The corresponding $`N`$ fermionic field $`\stackrel{}{\psi }(x)`$ is a Majorana (real) spinor which has the four-fermi type interaction with $`O(N)`$ symmetry and is called the Gross-Neveu model. The supersymmetric nonlinear sigma model with $`O(N)`$ symmetry is defined simply as a combination of these two models
$$(x)=\frac{1}{2}\{(\stackrel{}{A})^2+\overline{\stackrel{}{\psi }}i\text{/}\stackrel{}{\psi }\}+\frac{g^2}{8N}(\overline{\stackrel{}{\psi }}\stackrel{}{\psi })^2$$
where $`\stackrel{}{A}^2=\frac{N}{g^2}`$ and $`\stackrel{}{A}\stackrel{}{\psi }=0.`$
This model enjoys three kinds of symmetry: $`O(N)`$ symmetry, discrete chiral symmetry $`\psi \gamma _5\psi `$ and the supersymmetry which mixes bosonic and fermionic fields. We can find the solution in the large $`N`$ limit where $`N\mathrm{}`$ with $`g`$ fixed. In order to study the phase structure, it is crucial to introduce auxiliary superfields defined in the superspace with the bosonic coordinate $`x^\mu `$ and the two component fermionic coordinate $`\theta `$
$$\mathrm{\Phi }_0(x,\theta )=A_0(x)+\overline{\theta }\psi _0(x)+\frac{1}{2}\overline{\theta }\theta F_0(x).$$
$`A_0(x)\stackrel{}{\overline{\psi }}\stackrel{}{\psi }`$ describes the scalar bound state of two fermions as the auxiliary field of the Gross-Neveu model. $`\psi _0(x)\stackrel{}{A}\stackrel{}{\psi }`$ describes the fermionic bound state of the fermion and the boson. $`F_0(x)`$ is the Lagrange multiplier for the constraint: $`\stackrel{}{A}^2=\frac{N}{g^2}`$ of the bosonic nonlinear $`\sigma `$ model. These auxiliary fields also play the role of order parameters of various symmetry. With the auxiliary field, the nonlinear lagrangian reduces to a very simple form
$$(x)=\frac{1}{2}\{(\stackrel{}{A})^2+\overline{\stackrel{}{\psi }}i\text{/}\stackrel{}{\psi }+\stackrel{}{F}^2\}+\left(2A_0\stackrel{}{A}\stackrel{}{F}A_0\overline{\stackrel{}{\psi }}\stackrel{}{\psi }+F_0\stackrel{}{A}^2\overline{\psi _0}\stackrel{}{\psi }\stackrel{}{A}\overline{\stackrel{}{\psi }}\psi _0\stackrel{}{A}\right)\frac{N}{g^2}F_0.$$
In the large $`N`$ field theory, we take into account the dominant vacuum fluctuations of $`O(N)`$ vector fields $`\stackrel{}{A},\stackrel{}{\psi }`$ and $`\stackrel{}{F}`$ which can be integrated out easily since the lagrangian is simply of quadratic form in these variables. As is shown in the table 1, $`O(N)`$ symmetry recovers and $`\stackrel{}{A(x)}`$ acquire a mass proportional to $`A_0`$. Chiral symmetry breaks down and fermions have the same mass with bosons.
In this talk, we generalize the auxiliary field formulation to nonlinear $`\sigma `$ models with $`𝒩=2`$ supersymmetry in two dimensions, which is equivalent to $`𝒩=1`$ supersymmetry in four dimensions.
## 2 Nonlinear $`\sigma `$ Models in Four Dimensions
When a global symmetry group $`G`$ breaks down to its subgroup $`H`$ by a vacuum expectation value $`v=\varphi `$, there appear massless Nambu-Goldstone (NG) bosons corresponding to broken generators in $`G/H`$. At low energies, interactions among these NG bosons are described by nonlinear $`\sigma `$ models . In supersymmetric theories, target manifolds of nonlinear $`\sigma `$ models must be Kähler manifolds . A manifold whose metric is given by a Kähler potential $`K(\overline{A},A)`$
$$g_{\overline{m}n}=\frac{^2K(\overline{A},A)}{\overline{A}^{\overline{m}}A^n},$$
(1)
ia called the Kähler manifold. Since NG bosons must be scalar components of complex chiral superfields
$$\varphi (x,\theta )=A(x)+\theta \psi (x)+\frac{1}{2}\theta ^2F(x),$$
(2)
there are two possibilities. If the coset $`G/H`$ itself is a Kähler manifold, both $`\mathrm{Re}A(x)`$ and $`\mathrm{Im}A(x)`$ must be NG-bosons. The effective lagrangian in this case is uniquely determined by the metric of the coset manifold $`G/H`$ . On the other hand, if the coset $`G/H`$ is a submanifold of a Kähler manifold, there is at least one chiral superfield whose real or imaginary part is not a NG-boson. This additional massless boson is called the quasi-Nambu-Goldstone(QNG) boson . In this case, the Kähler metric in the direction of QNG boson is not determined by the metric of its subspace $`G/H`$, and the effective lagrangian is not unique. In this article, we will confine ourselves to the case of Kähler $`G/H`$. The lagrangian of the nonlinear $`\sigma `$ model on a Kähler manifold is
$`(x)`$ $`=g_{\overline{m}n}_\mu \overline{A}^{\overline{m}}^\mu A^n+{\displaystyle \frac{i}{2}}g_{\overline{m}n}\left(\overline{\psi }^{\overline{m}}\sigma ^\mu (D_\mu \psi )^n+\psi ^n\overline{\sigma }^\mu (D_\mu \overline{\psi })^{\overline{m}}\right)`$ (3)
$`+{\displaystyle \frac{1}{4}}R_{k\overline{m}l\overline{n}}(\psi ^k\psi ^l)(\overline{\psi }^{\overline{m}}\overline{\psi }^{\overline{n}}),`$
where
$$(D_\mu \psi )^n=(\delta _m^n_\mu +\mathrm{\Gamma }_{lm}^n_\mu A^l)\psi ^m,\mathrm{\Gamma }_{lm}^n=g^{n\overline{k}}_lg_{\overline{k}m}.$$
(4)
Once we know the Kähler potential $`K(\overline{A},A)`$, we can calculate the metric by (1), the connection by (4) and the lagrangian by (3). This lagrangian is suitable for perturbative calculations. For the nonperturbative study, however, the auxiliary field formulation is hoped for.
## 3 Auxiliary Field Formulation
Let us start from a known example, the $`𝐂P^{N1}`$ model. Chiral superfield $`\varphi _i(x,\theta )(i=1,2,\mathrm{},N)`$ belongs to a fundamental representation of $`G=SU(N)`$ which is the isometry of $`𝐂P^{N1}`$. We introduce $`U(1)`$ gauge symmetry
$$\varphi (x,\theta )\varphi ^{}(x,\theta )=e^{i\mathrm{\Lambda }(x,\theta )}\varphi (x,\theta )$$
(5)
to require that $`\varphi (x,\theta )`$ and $`\varphi ^{}(x,\theta )`$ are physically indistinguishable. With a complex chiral superfield, $`e^{i\mathrm{\Lambda }(x,\theta )}`$ is an arbitrary complex number. $`U(1)`$ gauge symmetry is thus complexified to $`U(1)^𝐂`$. The identification $`\varphi \varphi ^{}`$ defines the complex projective space $`𝐂P^{N1}`$. In order to impose local $`U(1)`$ gauge symmetry, we have to introduce a $`U(1)`$ gauge field $`V(x,\theta ,\overline{\theta })`$ with the transformation property $`e^Ve^Ve^{i\mathrm{\Lambda }+i\mathrm{\Lambda }^{}}`$. Then the lagrangian with a local $`U(1)`$ gauge symmetry is given by
$$=d^2\theta d^2\overline{\theta }(e^V\stackrel{}{\varphi }^{}\stackrel{}{\varphi }cV)$$
(6)
where the last term $`V`$ is called the Fayet-Iliopoulos D-term. In this model, real scalar superfield $`V(x,\theta ,\overline{\theta })`$ is the auxiliary superfield. The Kähler potential $`K(\varphi ,\varphi ^{})`$ is obtained by eliminating $`V`$ by using the equation of motion for $`V`$
$$=d^2\theta d^2\overline{\theta }K(\varphi ,\varphi ^{})=cd^2\theta d^2\overline{\theta }\mathrm{log}(\stackrel{}{\varphi }^{}\stackrel{}{\varphi }).$$
(7)
This Kähler potential reduces to the standard Fubini-Study metric of $`𝐂P^{N1}`$
$$K(\varphi ,\varphi ^{})=c\mathrm{log}(1+\underset{i=1}{\overset{N1}{}}\varphi _i^{}\varphi _i)$$
(8)
by a choice of gauge fixing
$$\varphi _N(x,\theta )=1.$$
(9)
The lagrangian of the $`𝐂P^{N1}`$ model is obtained by substituting the Kähler potential (8) to eqns. (1), (4) and (3). The global symmetry $`G=SU(N)`$, the isometry of the target space $`𝐂P^{N1}`$, is linearly realized on our $`\varphi _i`$ fields and our lagrangian (6) with auxiliary field $`V`$ is manifestly invariant under $`G`$. The gauge fixing condition (9) is not invariant under $`G=SU(N)`$ and we have to perform an appropriate gauge transformation simultaneously to compensate the change of $`\varphi _N`$ caused by the $`SU(N)`$ transformation. Therefore the global symmetry $`G=SU(N)`$ is nonlinearly realized in the gauge fixed theory. In this sense, our lagrangian (6) use the linear realization of $`G`$ in contrast to the nonlinear lagrangian in terms of the Kähler potential which use the nonlinear realization of $`G`$.
Similarly, we introduce $`M(M<N)`$ copies $`\varphi _j(j=1,2,\mathrm{},M)`$ of fundamental representations of $`SU(N)`$ and impose the local gauge symmetry $`U(M)`$ to identify $`\mathrm{\Phi }U\mathrm{\Phi }(UU(M))`$ where $`\mathrm{\Phi }=(\varphi _1,\varphi _2,\mathrm{},\varphi _M)`$ is an $`N\times M`$ matrix of chiral superfields. Then we obtain the nonlinear $`\sigma `$ model on the Grassmann manifold $`SU(N)/SU(NM)\times U(M)`$. The lagrangian with an auxiliary $`U(M)`$ gauge field $`V`$ reads
$$=d^2\theta d^2\overline{\theta }\left(\mathrm{tr}\left(\mathrm{\Phi }^{}\mathrm{\Phi }e^V\right)c\mathrm{tr}V\right)$$
(10)
As in the $`𝐂P^{N1}`$ model, the auxiliary field $`V`$ can be eliminated by the use of its equation motion. We fix the gauge by choosing
$$\mathrm{\Phi }=\left(\begin{array}{c}\mathrm{𝟏}_M\\ \phi \end{array}\right)$$
where $`\phi `$ is an $`(NM)\times M`$ matrix valued chiral superfield. Then the Kähler potential reads
$$K(\phi ,\phi ^{})=c\mathrm{log}det(\mathrm{𝟏}_M+\phi ^{}\phi ).$$
(11)
Again, the global symmetry $`G=SU(N)`$ is linearly realized on our fields $`\mathrm{\Phi }`$ although the gauge fixed fields $`\phi `$ is no longer a linear realization.
## 4 Nonlinear Sigma Models with F-term Constraints
The superspace in the four dimensional space-time consists of four bosonic coordinates $`x^\mu (\mu =0,1,2,3)`$ and a four components Majorana (real) spinor, which is equivalent to a complex Weyl spinor $`\theta ^\alpha (\alpha =1,2)`$ and its hermitian conjugate $`\overline{\theta }`$. The chiral superfield defined by (2) depends only on $`\theta `$ but not on $`\overline{\theta }`$. On the other hand, $`\varphi ^{}`$ depend on $`\overline{\theta }`$ but not on $`\theta `$. If we can make $`G`$ invariant combinations of chiral superfields, it is possible to introduce in the lagrangian another term called the F-term which can be written as an integral over $`\theta `$.
Let us try to impose an F-term constraint on $`𝐂P^{N1}`$ model. The simplest constraint consistent with the $`U(1)`$ gauge symmetry (5) is the quadratic equation
$$\stackrel{}{\varphi }\stackrel{}{\varphi }=0.$$
(12)
With this constraint, the invariance group of the lagrangian is no longer $`SU(N)`$ symmetry but its subgroup $`O(N)`$. Any nonvanishing value is forbidden on the righthand side because of the $`U(1)`$ symmetry. Solution of this F-term constraint: $`\stackrel{}{\varphi }^{\mathrm{\hspace{0.17em}2}}=\stackrel{}{x}^2+(y+iz)(yiz)=0`$ written in terms of $`\stackrel{}{\varphi }=(x_i,y,z)`$ is given by
$$yiz=\frac{x^2}{y+iz}=\frac{x^2}{\sqrt{2}}$$
(13)
where we have chosen a specific gauge $`y+iz=\sqrt{2}`$. If we have imposed the constraint (12) to $`\stackrel{}{\varphi }`$ without introducing the gauge symmetry, the resulting target space would be noncompact Kähler manifold with a QNG boson. The QNG boson is now gauged away by the gauge symmetry (5) as a gauge degree of freedom and has disappeared from the physical spectrum. On substitution of the solution (13) of the constraint to Eq. (7), we obtain the Kähler potential of this model
$$K(x,x^{})=\mathrm{log}\left(1+\underset{i=1}{\overset{N2}{}}x_i^{}x_i+\frac{1}{4}\underset{i,j=1}{\overset{N2}{}}x_i^2x_j^2\right),$$
(14)
which is known as the Kähler potential of the quadratic surface
$$Q^{N2}(𝐂)=\frac{SO(N)}{SO(N2)\times U(1)}.$$
The lagrangian of the model with auxiliary fields is simply obtained from Eq (6) by imposing the F-term constraint (12) with a lagrange multiplier field $`\varphi _0(x,\theta )`$
$$=d^2\theta d^2\overline{\theta }(e^V\stackrel{}{\varphi }^{}\stackrel{}{\varphi }V)+(d^2\theta \varphi _0\stackrel{}{\varphi }\stackrel{}{\varphi }+\mathrm{h}.\mathrm{c}.)$$
(15)
Let us impose an F-term constraint to the model on the Grassmann manifold $`SU(2N)/SU(N)\times U(N)`$. Our basic field is the $`2N\times N`$ matrix valued chiral superfield $`\mathrm{\Phi }`$ which transforms linearly under the global symmetry $`SU(2N)`$. In order to identify $`\mathrm{\Phi }`$ with $`\mathrm{\Phi }U`$ with $`UU(N)`$, we introduce the $`U(N)`$ gauge field $`V`$. The simplest constraint is again the quadratic constraint $`\mathrm{\Phi }^T\mathrm{\Phi }=0`$. Since this constraint transforms as the symmetric second rank tensor under the gauge group $`U(N)`$, we introduce the chiral superfield $`\mathrm{\Phi }_0`$ which also transforms as a symmetric second rank tensor under the gauge group $`U(N)`$. With this constraint, the global symmetry $`SU(2N)`$ reduces to its subgroup $`SO(2N)`$. Thus we obtain the auxiliary field formulation of the $`SO(2N)/U(N)`$ model
$$=d^4\theta (\mathrm{tr}\left(\mathrm{\Phi }^{}\mathrm{\Phi }e^V\right)c\mathrm{tr}V)+(d^2\theta \mathrm{tr}(\mathrm{\Phi }_0\mathrm{\Phi }^T\mathrm{\Phi }))+\mathrm{h}.\mathrm{c}.).$$
(16)
If we insert the simpletic structure
$$J=\left(\begin{array}{cc}\mathrm{𝟎}& \mathrm{𝟏}_N\\ \mathrm{𝟏}_N& \mathrm{𝟎}\end{array}\right).$$
(17)
between $`\mathrm{\Phi }^T`$ and $`\mathrm{\Phi }`$, the global symmetry reduces to the simpletic group $`Sp(N)`$. Therefore, the $`Sp(N)/U(N)`$ model is defined by
$$=d^4\theta (\mathrm{tr}\left(\mathrm{\Phi }^{}\mathrm{\Phi }e^V\right)c\mathrm{tr}V)+(d^2\theta \mathrm{tr}(\mathrm{\Phi }_0\mathrm{\Phi }^TJ\mathrm{\Phi }))+\mathrm{h}.\mathrm{c}.).$$
(18)
In this case, the chiral auxiliary field $`\mathrm{\Phi }_0`$ transforms as an antisymmetric second rank tensor of the gauge group $`U(N)`$.
Similarly, we can formulate supersymmetric nonlinear sigma models on hermitian symmetric spaces shown in the table 2 by using auxiliary fields . It should be noted that the third order and the fourth order polynomials appear as the F-term constraints in the case of exceptional groups.
## 5 Quantum Legendre Transform
If we impose only the symmetry, the nonlinear lagrangian may depend on arbitrary function. For example, consider the simplest $`𝐂P^{N1}`$ model. The following lagrangian with an arbitrary function $`f`$ is allowed by the global as well as the local symmetry
$$=d^2\theta d^2\overline{\theta }(f(e^V\stackrel{}{\varphi }^{}\stackrel{}{\varphi })cV).$$
(19)
We can prove this arbitrariness disappears and (19) reduces to the simplest lagrangian (6) discussed previously . Namely, we can prove that
$`{\displaystyle [dV]\mathrm{exp}[id^4\theta (f(e^V\varphi ^{}\varphi )cV)]}={\displaystyle [dV]\mathrm{exp}[id^4\theta (e^V\varphi ^{}\varphi cV)]}`$
by using a remarkable property of the quantum Legendre transform in supersymmetric theories which is valid for any vector superfields $`\sigma (x,\theta ,\overline{\theta })`$, $`\mathrm{\Phi }(x,\theta ,\overline{\theta })`$
$`{\displaystyle [d\sigma ]\mathrm{exp}\left[id^4xd^4\theta (\sigma \mathrm{\Phi }W(\sigma ))\right]}=\mathrm{exp}\left[i{\displaystyle d^4xd^4\theta U(\mathrm{\Phi })}\right],`$ (20)
where $`U(\mathrm{\Phi })=\widehat{\sigma }(\mathrm{\Phi })\mathrm{\Phi }W(\widehat{\sigma }(\mathrm{\Phi }))`$ is the Legendre transform of $`W`$ and $`\widehat{\sigma }`$ is the stationary point:
$$\frac{}{\sigma }(\sigma \mathrm{\Phi }W(\sigma ))|_{\sigma =\widehat{\sigma }}=\mathrm{\Phi }W(\widehat{\sigma })=0.$$
## Acknowledgments
The work of M.N. is supported in part by JSPS Research Fellowships.
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# Information Bottlenecks, Causal States, and Statistical Relevance Bases: How to Represent Relevant Information in Memoryless Transduction
## I Introduction
Recently, Tishby, Pereira, and Bialek proposed a new method for finding concise representations of the information one set of variables contains about another . This brief note explains the connections between this “information-bottleneck” method and existing mathematical frameworks and techniques. This comparison should enhance the value of research in this promising direction and clarify the relative uses of these techniques in applications.
In the interest of space, we assume readers are familiar with the notation of both and .
## II The Information-Bottleneck Method
In the authors pose the following problem. Given a joint distribution over two random variables—the “input” $`X`$ and the “output” $`Y`$, find an intermediate or “bottleneck” variable $`\stackrel{~}{X}`$ which is a (possibly stochastic) function of $`X`$ such that $`\stackrel{~}{X}`$ is more compressed than $`X`$, but retains predictive information about $`Y`$. More exactly, they ask for a conditional distribution $`\mathrm{Pr}(\stackrel{~}{x}|x)`$ that minimizes the functional
$$=I[\stackrel{~}{X};X]\beta I[\stackrel{~}{X};Y],$$
(1)
where $`I[W,Z]`$ is the mutual information between random variables $`W`$ and $`Z`$ and $`\beta `$ is a positive real number. Minimizing the first term represents the desire to find a compression of the original input data $`X`$; maximizing the second term represents the desire to retain the ability to predict $`Y`$.<sup>*</sup><sup>*</sup>*Since $`\stackrel{~}{X}=g(X,\mathrm{\Omega })`$ for some auxiliary random variable $`\mathrm{\Omega }`$, a theorem of Shannon’s assures us that $`I[\stackrel{~}{X};Y]I[X;Y]`$ and the transformation from $`X`$ to $`\stackrel{~}{X}`$ cannot increase our ability to predict $`Y`$ \[4, App. 7\]. The coefficient $`\beta `$ governs the trade-off between these two goals: as $`\beta 0`$, we lose interest in prediction in favor of compression; whereas as $`\beta \mathrm{}`$, predictive ability becomes paramount.
Extending classical rate-distortion theory, the authors are not only able to state self-consistent equations that determine which distributions satisfy this variational problem, but give a convergent iterative procedure that finds one of these distributions. They do not address the rate of convergence.
## III Causal States for Transducer-Functionals
In and earlier publications, we defined causal states for stationary stochastic processes as follows. Two histories $`\stackrel{}{s}`$ and $`\stackrel{}{s}^{}`$ belong to the same causal state $`𝒮`$ if and only if
$$\mathrm{Pr}(\stackrel{}{S}^L=s^L|\stackrel{}{S}=\stackrel{}{s})=\mathrm{Pr}(\stackrel{}{S}^L=s^L|\stackrel{}{S}=\stackrel{}{s}^{}),$$
(2)
for all $`s^L`$ and for all $`L`$. That is, two histories belong to the same causal state if and only if they give the same conditional distribution for futures. In we showed that the causal states defined by Eq. (2) possess two kinds of optimality. First, their ability to predict the future $`\stackrel{}{S}`$ is maximal. Second, they are the simplest such set of states.
Thus, we showed that the set $`𝓢`$ of causal states is the solution to the following optimization problem. Given the joint distribution $`\mathrm{Pr}(\stackrel{}{S},\stackrel{}{S})`$ over the past $`\stackrel{}{S}`$ and future $`\stackrel{}{S}`$, find the function (equivalently, partition) $`ϵ`$ of $`\stackrel{}{S}`$ such that (i) the conditional entropy $`H[\stackrel{}{S}^L|ϵ(\stackrel{}{S})]`$ is minimized for all $`L`$, and (ii) the entropy $`H[ϵ(\stackrel{}{S})]`$ is minimized among all functions $`\widehat{\eta }`$ that satisfy condition (i).The states induced by the partition $`\widehat{\eta }`$ are called the prescient rivals of the causal states induced by $`ϵ`$. The entropy $`H[ϵ(\stackrel{}{S})]`$ is called the statistical complexity and measures the “size” of, or amount of information stored in, the causal states. The equivalence classes induced by $`ϵ`$ are the causal states, and they are the uniqueMore precisely, any other function satisfying conditions (i) and (ii) may differ from $`ϵ`$ on at most a set of histories of measure zero. solution to the optimization problem. A moment’s reflection shows that this optimization is equivalent to first maximizing the mutual information between the effective states and the futures and then minimizing the mutual information between the effective states and the histories. The first step maximizes the predictive ability of the effective states and the second selects the most concise set of states. Note that the opposite sequence of optimizations—first minimizing complexity and then maximizing predictability—is trivial, since it produces a single-state model that describes an IID sequence of random variables; e.g., a biased coin or die.
## IV Memoryless Transducers
As has been remarked earlier—e.g., and —the causal-state construction is not intrinsically limited to time series. Of particular interest here is the case of transducer-functionals: when one sequence of variables is a functional of another sequence, possibly a stochastic functional. In this case, one can construct causal states that (i) retain all predictive information about the output series, (ii) are deterministic functions of the prior causal state and the most recent value of the input series, and (iii) minimize the statistical complexity of (information stored in) the causal states. We present the theory of causal states for general transducers elsewhere.
In the case where the output depends on the current input alone—the case of memoryless transduction—the causal states assume a particularly simple form: two inputs $`x`$ and $`x^{}`$ belong to the same causal state if and only if
$$\mathrm{Pr}(Y=y|X=x)=\mathrm{Pr}(Y=y|X=x^{}),$$
(3)
for all $`y`$.<sup>§</sup><sup>§</sup>§A somewhat more complicated construction is necessary when the transducer exhibits memory. In this case, it can be shown that
$$H[Y|ϵ(X)]H[Y|\eta (X)],$$
(4)
for any other partition of the inputs $`\eta `$ and that, among all the rival partitions $`\widehat{\eta }`$ minimizing the conditional entropy of the outputs (the prescient rivals of $`ϵ`$),
$$H[ϵ(X)]H[\widehat{\eta }(X)].$$
(5)
This is to say, the causal states are the most compressed hidden variables. In the sense of , they are optimal bottleneck variables.
One concludes that these are precisely what should be delivered by the information-bottleneck method in the limit where $`\beta \mathrm{}`$. It is not immediately obvious that the iterative procedure of is still valid in this limit. Nonetheless, that $`ϵ`$ is the partition satisfying their original constraints is evident.
We note in passing that, as shown in , prescient rivals—those sets of states that retain all predictive information from the original inputs while compressing them—are sufficient statistics. Conversely, states that, when sufficient statistics exist, then compression-with-prediction is possible.
## V The Statistical Relevance Basis
Before closing, we point out another solution to the problem of discovering concise and predictive hidden variables. In his books and , Salmon put forward a construction, under the name of the “statistical relevance basis”, that is identical in its essentials with that of causal states for memoryless transducers.Salmon’s work only came to our attention in mid-1998, and thus it is not cited in our publications including and prior to . Owing to the rather different aims for which Salmon’s construction was intended—explicating the notion of “causation” in the philosophy of science, no one seems to have proved its information-theoretic optimality properties nor even to have noted its connection to sufficient statistics. (Briefly: if a nontrivial sufficient partition of the input variables exists, then the relevance basis is the minimal sufficient partition. These proofs will appear elsewhere.)
## VI Comparison and Conclusion
Recapitulating, the causal states for memoryless transduction coincide with the cells of the “bottleneck” partition in the limit $`\beta \mathrm{}`$. Moreover, both are identical with the statistical relevance basis of Salmon.
The construction of the causal states does not allow us to discard any predictive information about the output $`Y`$, even if this might allow for a substantial reduction in the statistical complexity. The bottleneck method, by contrast, generally throws away some predictive information. It trades one bit less statistical complexity for $`1/\beta `$ bits less predictive information. Of course, the linear trade-off and the particular value of the coefficient $`\beta `$ that controls it are ad hoc choices. Whether this is acceptable in applications would seem to depend on the goal. For example, if the goal is a practical “lossy” data-compression scheme, the bottleneck method recommends itself. However, if the goal is representing the intrinsic computation or causal structure of some natural process, causal states are better suited to the task.
## Acknowledgments
This work was partially supported by the SFI Computation, Dynamics, and Learning Program, by AFOSR via NSF grant PHY-9970158, and by DARPA under contract F30602-00-2-0583.
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# References
The Scale Dependence of Inclusive ep Scattering in the Resonance Region.
Simonetta Liuti<sup>a,b</sup>
<sup>a</sup> Department of Physics, University of Virginia
McCormick Road. Charlottesville, Virginia 22901, USA.
<sup>b</sup> INFN, Sezione di Roma Tre
Dipartimento di Fisica E. Amaldi
Via Vasca Navale, 84. 00146 Roma, Italy.
Abstract
We examine the scale dependence in the resonance region of inclusive ep scattering. In particular we discuss the invariants other than $`Q^2`$, which have been proposed as a scale for pQCD evolution in a kinematical regime where both infrared singularities and power corrections are expected to be largest. We show that the region where most of the present data are available, can be described using NLO pQCD evolution at fixed invariant mass, $`W^2`$, plus a leading order power correction term. We find that the coefficient of the power correction at $`W^2<4GeV^2`$ is relatively small, i.e. comparable in size to the one obtained in the large $`W^2`$ region.
It has been long known that a fully quantitative description of proton structure in terms of parton distribution functions must account for power corrections to the $`Q^2`$ dependence of the data, in addition to the predicted perturbative-QCD (pQCD) behavior. Power corrections are indeed observed in experiments as discrepancies between fixed order perturbative predictions and the data. Their theoretical interpretation is however a less well defined issue. In Deep Inelastic Scattering (DIS), in particular, the short-distance scattering involving single, non-interacting partons is expected to give way to processes in which two (or more) quarks or gluons participate simultaneously in the scattering. These processes correspond formally to the Higher Twist (HT), or the higher order terms in the “twist” expansion (twist=dimension-spin) . The coefficients of the HTs cannot be evaluated directly within perturbation theory. However, it is also known that power corrections should appear in the coefficient functions for hard processes, generated by the divergence of the perturbative series at large orders (renormalons) . Whether in DIS the two kinds of power corrections can be distinguished from one another and compared to the data is still an open question (see however ).
In recent analyses the HT terms have been extracted from DIS data by applying a cut in the kinematics at $`W^210\mathrm{GeV}^2`$ ($`W^2`$ is the invariant mass of the final hadronic state), that is excluding the kinematical region dominated by nucleon resonances. Following we have shown however that the entire set of inclusive data, including the low $`W^2`$ region can be described by a pQCD based analysis, the contribution from HTs being overall relatively small, i.e. within a factor of two from the one obtained in . The observation of a small power contribution can be otherwise phrased in terms of “approximate Duality”, namely the non-perturbative features of the data appearing as the characteristic peaks describing the nucleon resonances, average out to a curve that can be identified with the DIS one modulo perturbative corrections plus a small size power correction. Based on the results of , we perform here a more accurate analysis in the low $`W^2`$ domain with the aim of understanding the origin of the residual inverse-powerlike $`Q^2`$ dependence of the data.
Our analysis is based on three observations:
(1) The recent Jefferson Lab data on the structure function $`F_2`$ show “scaling” in $`W^2`$, i.e. invariance with $`W^2`$ of the smooth curves which average through the resonances peaks. The smooth fits to the data plotted vs. $`\xi =(2x)/(1++4M^2x^2/Q^2)^{1/2}`$, in order to account for target mass corrections, are shown in Fig.1 (dotted curves), for four different values of $`W^2`$ in the $`W^24\mathrm{GeV}^2`$ range. Fig.1 also shows that as $`W^2`$ increases some scaling violations are present. Eventually in the DIS region, i.e. at $`W^24\mathrm{GeV}^2`$, $`W^2`$-scaling breaks down completely. <sup>1</sup><sup>1</sup>1 This is due to the fact that the contribution of the fastly evolving sea quarks is no longer negligible. We conclude that the data scale in $`W^2`$ so long as one one keeps inside the resonance region.
(2) A pQCD based analysis is in principle possible in the resonance region (see also so long as the values of $`Q^2`$ are larger than $`1\mathrm{GeV}^2`$. Because of the kinematical relation $`W^2=Q^2(1x)/x+M^2,`$ $`M^2`$ being the nucleon mass, this constraint corresponds to large $`x`$ values ($`x0.2`$ at present kinematics). In analyses of DIS $`F_2`$ can be written as:
$$F_2(\xi ,Q^2)=F_2^{NLO}(\xi ,Q^2)\left(1+\frac{C(\xi ,Q^2)}{Q^2}\right),$$
(1)
where $`F_2^{NLO}(\xi ,Q^2)`$ is the NLO pQCD contribution and we disregard $`O(1/Q^4)`$ and higher terms. In order to try to reproduce $`W^2`$ scaling, we consider evolution at fixed $`W^2`$. In other words in Eq.(1) $`Q^2=Q^2(x)(W_R^2M^2)x/(1x)`$ where $`W_R`$ is the fixed invariant mass of a given resonance. In our evolution equations we have used the recent PDFs from , in which $`Q_o^21\mathrm{GeV}^2`$. We limit our analysis to the large $`x`$ region ($`x0.2`$) so that the Singlet contribution that might introduce some ambiguity for the evolution down to a low scale, is negligible. The results of perturbative evolution, determining $`F_2^{NLO}`$ are presented in Fig.1 along with the fits to Jlab data. We note first of all that there is an evident mismatch: the PDF results exceed the experimental values at $`x<0.6`$ and they lie lower at large $`x`$. Secondly pQCD evolution predicts a stronger evolution at fixed $`W^2`$.
Using Eq.(1) we interpret the discrepancies between perturbative evolution and the data as given by the leading non-perturbative contribution to the structure function. This actually enables us to determine the coefficient $`C`$ from the data:
$$C(\xi ,Q^2)=\frac{Q^2\mathrm{\Delta }F}{F_2^{NLO}(x,Q^2)},$$
(2)
where
$$\mathrm{\Delta }F=F_2^{NLO}(\xi ,Q^2(x))F_2^{Exp}(x,Q^2(x))$$
(3)
In Fig.2 we show $`C(\xi ,Q^2)`$ vs. $`\xi `$.
Our extractions are shown at different values of $`Q^2`$ (obtained by transforming back from the fixed $`W^2`$ values). For comparison the results from the analysis of using only $`W^2>10\mathrm{GeV}^2`$ data are also shown (dotted curves). <sup>2</sup><sup>2</sup>2 Note: the appearence of a strong $`Q^2`$ dependence in the coefficient $`C`$ is mainly due to the transformation $`x\xi `$. Our results, which use only data at $`W^24GeV^2`$, are in astonishing agreement with the high mass ones. At larger $`\xi `$ however the value of $`C(\xi ,Q^2)`$ becomes less clearly determined and as $`\xi \xi _{th}`$ (corresponding to $`x0.8`$) it is completely undefined. A possible interpretation of the indetermination at large $`\xi `$ is that standard DGLAP evolution becomes unreliable and that evolution using a z-dependent scale should be performed instead . A quantitative approach to this problem is pursued in . However, we also observe that at large $`x`$ data both in the large $`Q^2`$ region, determining $`F_2^{NLO}`$, and in the low $`W^2`$ region, determining $`F_2^{Exp}(x,Q^2(x))`$ in our analysis, are missing. Our analysis is extended to these regions by extrapolating what available at lower $`x`$ and its results are clearly less reliable here. These regions would be accessible at the $`12`$ GeV program at Jefferson Lab.
(3) Finally, we comment on the interpretation of the power corrections in the resonance region. From a practical point of view, if power corrections are found to mantain the same $`x`$ dependence displayed in Fig.2, namely $`CA/(1x)`$ as $`x1`$, we predict that
$$R=\frac{C(x,Q^2)}{Q^2(x)}\frac{A}{W_R^2M^2},x1$$
(4)
therefore at fixed $`W^2`$ one is approaching the $`x1`$ limit and at the same time having control of the power correction terms. This situation is displayed in Fig.3 where we show R for different values of $`W^2`$. Not being dominated by power corrections, this region is ideal for pursuing further quantitative studies of deviations from NLO DGLAP evolution ( for earlier analyses and e.g. for a more recent review).
What is the physical meaning of a small power correction in the resonance region? Power corrections might originate from the account of multi parton “final state interaction” processes that are more likely to occur at large distances and that correspond to well defined terms in the OPE. It is this type of interactions that are eventually responsible for confinement related features of the cross section such as the production of resonances. It turns out that for some, at present, unknown mechanism, the HTs contributions of increasing order in $`1/Q^2`$ cancel each other in the resonance region, giving origin to the duality phenomenon (this is seen either in the average, Fig.1, or in the moments integrals ). However we find out through an accurate analysis of the data that duality is not exact: a “residual” $`1/Q^2`$ dependence with a coefficient comparable to the large $`W^2`$ analyses is still necessary to interpret the data. This $`Q^2`$ dependence not being ascribed to HT corrections, can be taken as the true contribution from the non-perturbative corrections to the pQCD coefficient functions, namely the renormalon term. In order to confirm this interpretation, further studies addressed at determining the universality of this correction should be performed. These would include studies of different structure functions, such as $`F_L`$ in the resonance region, as well as scattering from different targets, including nuclei and studies of different fragmentation functions, a program accessible at Jefferson Lab at $`12\mathrm{GeV}`$.
On a more speculative level, the interpretation of the physical picture behind this behavior leads to a number of intriguing scenarios: for instance, partons inside hadrons might be arranged in a different way at intermediate/large-distances, e.g. they might be clumped together inside valence quarks and the role of final state interactions might be effectively small down to low $`Q^2`$.
In conclusion, information on the structure of the proton is still abundantly missing. It could be obtained if more data were available in the “transition” regions of $`x`$ and $`Q^2`$ where perturbative QCD (pQCD) evolution regulated by DGLAP equations can no longer be considered to be the main mechanism, and non-perturbative contributions become important.
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# Coherence in Substructural Categories
## 1 Substructural categories
By category with multiplication we mean a category $`𝒜`$ together with a bifunctor $`:𝒜\times 𝒜𝒜`$ and a special object I. Categories with multiplication can be axiomatized by postulating the following equations between arrows
| categorial equations |
| --- |
| $`(cat\text{1})`$ | $`f\text{1}_A=f=\text{1}_Bf`$ for all $`f:AB`$ |
| $`(cat\text{2})`$ | $`h(gf)=(hg)f`$ for all $`f,g,hMor(𝒜)`$ |
| functorial equations |
| --- |
| $`()`$ | $`(g_1f_1)(g_2f_2)=(g_1g_2)(f_1f_2)`$ |
| $`(\text{1})`$ | $`\text{1}_A\text{1}_B=\text{1}_{AB}`$ |
A category with multiplication is monoidal if there are special arrows for all objects $`A`$, $`B`$ and $`C`$
$$𝝈_A:\text{I}AA\text{ }𝜹_A:A\text{I}A$$
$$𝝈_A^i:A\text{I}A\text{ }𝜹_A^i:AA\text{I}$$
$$\stackrel{}{\text{b}}_{A,B,C}:A(BC)(AB)C\text{ }\stackrel{}{\text{b}}_{A,B,C}:(AB)CA(BC)$$
and if it satisfies
| $`𝝈`$$`𝜹`$-equations |
| --- |
| ($`𝝈`$) | For $`f:AB`$, | $`f𝝈_A=𝝈_B(\text{1}_\text{I}f)`$. |
| ($`𝜹`$) | For $`f:AB`$, | $`f𝜹_A=𝜹_B(f\text{1}_\text{I})`$. |
| $`(𝝈𝝈^i)`$ | $`𝝈_A𝝈_A^i=\text{1}_A`$, | $`𝝈_A^i𝝈_A=\text{1}_{\text{I}A}`$ |
| $`(𝜹𝜹^i)`$ | $`𝜹_A𝜹_A^i=\text{1}_A`$, | $`𝜹_A^i𝜹_A=\text{1}_{A\text{I}}`$ |
| $`(𝝈𝜹)`$ | $`𝝈_\text{I}=𝜹_\text{I}`$ |
| b-equations |
| --- |
| (b) | For $`f:AD`$, $`g:BE`$ and $`h:CF`$, | $`((fg)h)\stackrel{}{\text{b}}_{A,B,C}=\stackrel{}{\text{b}}_{D,E,F}(f(gh))`$. |
| (bb) | $`\stackrel{}{\text{b}}_{A,B,C}\stackrel{}{\text{b}}_{A,B,C}=\text{1}_{(AB)C}`$, | $`\stackrel{}{\text{b}}_{A,B,C}\stackrel{}{\text{b}}_{A,B,C}=\text{1}_{A(BC)}`$ |
| ($`𝝈`$$`𝜹`$b) | $`(𝜹_A\text{1}_B)\stackrel{}{\text{b}}_{A,\text{I},B}=\text{1}_A𝝈_B`$ |
| (b5) | $`\stackrel{}{\text{b}}_{AB,C,D}\stackrel{}{\text{b}}_{A,B,CD}=(\stackrel{}{\text{b}}_{A,B,C}\text{1}_D)\stackrel{}{\text{b}}_{A,BC,D}(\text{1}_A\stackrel{}{\text{b}}_{B,C,D})`$ |
A monoidal category is symmetric monoidal if it has the special arrow
$$\text{c}_{A,B}:ABBA$$
for every pair $`(A,B)`$ of its objects, and if the following equations hold
| c-equations |
| --- |
| (c) | For $`f:AC`$ and $`g:BD`$, $`(gf)\text{c}_{A,B}=\text{c}_{C,D}(fg)`$ |
| (cc) | $`\text{c}_{B,A}\text{c}_{A,B}=\text{1}_{AB}`$ |
| ($`𝝈`$$`𝜹`$c) | $`𝝈_A\text{c}_{A,\text{I}}=𝜹_A`$ |
| (bc6) | $`\stackrel{}{\text{b}}_{C,A,B}\text{c}_{AB,C}\stackrel{}{\text{b}}_{A,B,C}=(\text{c}_{A,C}\text{1}_B)\stackrel{}{\text{b}}_{A,C,B}(\text{1}_A\text{c}_{B,C})`$ |
A symmetric monoidal category is relevant if it has the special arrow
$$\text{w}_A:AAA$$
for every object $`A`$, and if the following equations hold
| w-equations |
| --- |
| (w) | For $`f:AB`$, $`(ff)\text{w}_A=\text{w}_Bf`$. |
| ($`𝝈`$$`𝜹`$w) | $`𝝈_\text{I}\text{w}_\text{I}=\text{1}_\text{I}`$ |
| (bw) | $`\stackrel{}{\text{b}}_{A,A,A}(\text{1}_A\text{w}_A)\text{w}_A=(\text{w}_A\text{1}_A)\text{w}_A`$ |
| (cw) | $`\text{c}_{A,A}\text{w}_A=\text{w}_A`$ |
| (bcw8) | $`\text{c}_{A,B,A,B}^m\text{w}_{AB}=\text{w}_A\text{w}_B`$, where |
| | $`\text{c}_{A,B,C,D}^m=^{df}\stackrel{}{\text{b}}_{A,C,BD}(\text{1}_A(\stackrel{}{\text{b}}_{C,B,D}(\text{c}_{B,C}\text{1}_D)\stackrel{}{\text{b}}_{B,C,D}))\stackrel{}{\text{b}}_{A,B,CD}`$ |
A symmetric monoidal category is affine if it has the special arrow
$$\text{k}_A:A\text{I}$$
for every object $`A`$ and if the following equations hold
| k-equations |
| --- |
| (k) | For $`f:AB`$, $`\text{k}_A=\text{k}_Bf`$ |
| (1k) | $`\text{k}_\text{I}=\text{1}_\text{I}`$ |
If a symmetric monoidal category is both relevant and affine (described in the same language) and if its arrows satisfy the following equations
$$(𝝈\text{k}\text{w})\text{ }𝝈_A(\text{k}_A\text{1}_A)\text{w}_A=\text{1}_A,\text{ }(𝜹\text{k}\text{w})\text{ }𝜹_A(\text{1}_A\text{k}_A)\text{w}_A=\text{1}_A$$
then we say it is cartesian.
We call this axiomatization of cartesian categories structural-equational. It differs from the standard equational axiomatization (see ) of these categories. The latter is based on the universality of product and uses as primitive, arrows $`\text{1}_A:AA`$, $`𝝅_{A,B}:ABA`$, $`𝝅_{A,B}^{}:ABB`$ and $`\text{k}_A:A\text{I}`$ for all objects $`A`$ and $`B`$, and a partial binary operation on arrows $`,`$, such that
$$\frac{f:CA\text{ }g:CB}{f,g:CAB}$$
Equations that hold are the categorial equations plus
| (E2) | $`f=\text{k}_A`$, for every $`f:A\text{I}`$. |
| --- | --- |
| (E3a.) | $`𝝅_{A,B}f,g=f`$, for $`f:CA`$ and $`g:CB`$. |
| (E3b.) | $`𝝅_{A,B}^{}f,g=g`$, for $`f:CA`$ and $`g:CB`$. |
| (E3c.) | $`𝝅_{A,B}h,𝝅_{A,B}^{}h=h`$, for $`h:CAB`$. |
To show that these two axiomatizations are extensionally equivalent we have to define
$$𝝅_{A,B}=_{df}𝜹_A(\text{1}_A\text{k}_B),\text{ }𝝅_{A,B}^{}=_{df}𝝈_B(\text{k}_A\text{1}_B),$$
and for $`f:CA`$ and $`g:CB`$
$$f,g=_{df}(fg)\text{w}_C,$$
in structural case. Then it is easy to show that (E2)–(E3c.) hold.
Conversely, if we start with the standard axiomatization, then we can define
$$𝝈_A=_{df}𝝅_{\text{I},A}^{},\text{ }𝝈_A^i=_{df}\text{k}_A,\text{1}_A,$$
$$𝜹_A=_{df}𝝅_{A,\text{I}},\text{ }𝜹_A^i=_{df}\text{1}_A,\text{k}_A,$$
$$\stackrel{}{\text{b}}_{A,B,C}=_{df}𝝅_{A,BC},𝝅_{B,C}𝝅_{A,BC}^{},𝝅_{B,C}^{}𝝅_{A,BC}^{}$$
$$\stackrel{}{\text{b}}_{A,B,C}=_{df}𝝅_{A,B}𝝅_{AB,C},𝝅_{A,B}^{}𝝅_{AB,C},𝝅_{AB,C}^{},$$
$$\text{c}_{A,B}=_{df}𝝅_{A,B}^{},𝝅_{A,B},\text{ }\text{w}_A=_{df}\text{1}_A,\text{1}_A,$$
and for $`f:AC`$ and $`g:BD`$
$$fg=_{df}f𝝅_{A,B},g𝝅_{A,B}^{}.$$
It is straightforward to prove that functorial, $`𝝈`$$`𝜹`$, b, c, w, k\- equations, as well as ($`𝝈`$kw) and ($`𝜹`$kw) hold.
In addition we have to prove that the “double translation” will take us to the same notions, i.e., to show that in the structural axiomatization the following equations hold
$$𝝈_A=𝝈_A(\text{k}_\text{I}\text{1}_A),\text{ }𝝈_A^i=(\text{k}_A\text{1}_A)\text{w}_A$$
$$𝜹_A=𝜹_A(\text{1}_A\text{k}_\text{I}),\text{ }𝜹_A^i=(\text{1}_A\text{k}_A)\text{w}_A$$
$`\stackrel{}{\text{b}}_{A,B,C}`$ $`=`$ $`(((𝜹_A(\text{1}_A\text{k}_{BC}))(𝜹_B(\text{1}_B\text{k}_C)𝝈_{BC}(\text{k}_A\text{1}_{BC})))\text{w}_{A(BC)})`$
$`(𝝈_C(\text{k}_B\text{1}_C)𝝈_{BC}(\text{k}_A\text{1}_{BC})))\text{w}_{A(BC)}`$
$`\stackrel{}{\text{b}}_{A,B,C}`$ $`=`$ $`((𝜹_A(\text{1}_A\text{k}_B)𝜹_{AB}(\text{1}_{AB}\text{k}_C))(((𝝈_B(\text{k}_A\text{1}_B)𝜹_{AB}(\text{1}_{AB}\text{k}_C))`$
$`(𝝈_C(\text{k}_{AB}\text{1}_C)))\text{w}_{(AB)C}))\text{w}_{(AB)C}`$
$`\text{c}_{A,B}`$ $`=`$ $`((𝝈_B(\text{k}_A\text{1}_B))(𝜹_A(\text{1}_A\text{k}_B)))\text{w}_{AB}`$
$`\text{w}_A`$ $`=`$ $`(\text{1}_A\text{1}_A)\text{w}_A`$
$`fg`$ $`=`$ $`((f𝜹_A(\text{1}_A\text{k}_B))(g𝝈_B(\text{k}_A\text{1}_B)))\text{w}_{AB},\text{ }f:AC,\text{ }g:BD`$
and in the standard axiomatization,
$$𝝅_{A,B}=𝝅_{A,\text{I}}\text{1}_A𝝅_{A,B},\text{k}_B𝝅_{A,B}^{},\text{ }𝝅_{A,B}^{}=𝝅_{\text{I},B}\text{k}_A𝝅_{A,B},\text{1}_B𝝅_{A,B}^{}$$
$$f,g=f𝝅_{C,C},g𝝅_{C,C}^{}\text{1}_C,\text{1}_C,\text{ }f:CA,\text{ }g:CB.$$
Some derivations are straightforward and some, like those concerning b-arrows, require some effort to be proven. We leave them as an exercise and suggest to the reader who is familiar with the coherence results in symmetric monoidal categories to use them in these proofs.
## 2 G-natural transformations and transformational graphs
Let $`𝒜`$ be an arbitrary category and $`F:𝒜^m𝒜`$, $`G:𝒜^n𝒜`$, $`m,n0`$, two functors ($`𝒜^0`$ is trivial category). Let $`\mathrm{\Gamma }`$ be a function from $`\{1,\mathrm{},n\}`$ to $`\{1,\mathrm{},m\}`$ called graph (if $`n=0`$, then $`\{1,\mathrm{},n\}`$ is $`\mathrm{}`$). We say that an indexed family of morphisms from $`𝒜`$
$$\alpha =\{\alpha (A_1,\mathrm{},A_m):F(A_1,\mathrm{},A_m)G(A_{\mathrm{\Gamma }(1)},\mathrm{},A_{\mathrm{\Gamma }(n)})A_1,\mathrm{},A_m𝒜\}$$
is a g-natural transformation from $`F`$ to $`G`$ with the graph $`\mathrm{\Gamma }`$, denoted by $`\alpha :F\text{}G`$, if for every $`i`$, $`1im`$, arbitrary $`A_1,\mathrm{},A_i,A_i^{},A_{i+1},\mathrm{},A_m`$ and $`f:A_iA_i^{}`$ from $`𝒜`$, the following diagram commutes
where for $`ji`$, $`h_j\text{1}_{A_j}`$, $`A_j^{}A_j`$ and for $`i`$, $`h_if`$, $`A_i^{}A_i^{}`$.
This definition follows the one given in (, page 94).
Example Let $`𝒜`$ be a relevant category. Denote by w the indexed set
$$\{\text{w}_A:AAAA𝒜\}.$$
Then by the equation (w), w is a g-natural transformation from the identity functor $`1_𝒜:𝒜𝒜`$ to the multiplication functor $`:𝒜^2𝒜`$ with the graph $`\mathrm{\Gamma }:\{1,2\}\{1\}`$, $`\mathrm{\Gamma }(1)=\mathrm{\Gamma }(2)=1`$.
If $`\alpha :F\text{}G`$ for $`F`$, $`G`$ and $`\mathrm{\Gamma }`$ as above, then it is easy to see that $`\alpha `$ is a classical natural transformation between $`F`$ and $`G^{}:𝒜^m𝒜`$ where
$$G^{}(A_1,\mathrm{},A_m)=_{df}G(A_{\mathrm{\Gamma }(1)},\mathrm{},A_{\mathrm{\Gamma }(n)}),$$
$$G^{}(f_1,\mathrm{},f_m)=_{df}G(f_{\mathrm{\Gamma }(1)},\mathrm{},f_{\mathrm{\Gamma }(n)}).$$
As in the classical case, g-natural transformations can be composed in the following way. Let $`F:𝒜^m𝒜`$, $`G:𝒜^n𝒜`$ and $`H:𝒜^l𝒜`$ be functors. Let $`\alpha :F\text{}G`$ and $`\beta :G\text{}H`$ for some graphs $`\mathrm{\Phi }:\{1,\mathrm{},n\}\{1,\mathrm{},m\}`$ and $`\mathrm{\Psi }:\{1,\mathrm{},l\}\{1,\mathrm{},n\}`$. We define its composition as
$$\beta \alpha =_{df}\{\beta \alpha (A_1,\mathrm{},A_m)\beta (A_{\mathrm{\Phi }(1)},\mathrm{},A_{\mathrm{\Phi }(n)})\alpha (A_1,\mathrm{},A_m)(A_1,\mathrm{},A_m)𝒜^m\}.$$
Then it is easy to prove (as in the case of classical natural transformations) that $`\beta \alpha `$ is a g-natural transformation from $`F`$ to $`H`$ with the graph $`\mathrm{\Phi }\mathrm{\Psi }`$ (the usual composition of functions $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$).
Generalization of g-natural transformations to the case of several categories (here we have only $`𝒜`$) is not essential, and serves just to complicate the notation.
## 3 Canonical transformations in substructural categories
Throughout this section, $`𝒜`$ denotes an arbitrary substructural category.
Let $``$ be a set of terms obtained from symbols $`\mathrm{}`$, I and binary operation $``$. Its elements we call shapes.
In a natural way, we define correspondence between shapes and functors of type $`𝒜^n𝒜`$ for some $`n0`$.
1. Functor $`\text{1}_𝒜:𝒜𝒜`$ corresponds to the term $`\mathrm{}`$.
2. Functor $`I:𝒜^0𝒜`$, which maps the unique object from $`𝒜^0`$ to the object I from $`𝒜`$, corresponds to the term I.
3. If $`F:𝒜^m𝒜`$ corresponds to the term $`𝑭`$ and $`G:𝒜^n𝒜`$ corresponds to the term $`𝑮`$, then the functor $`H:𝒜^{m+n}𝒜`$ such that for every $`m+n`$-tuple $`(A_1,\mathrm{},A_{m+n})`$ of $`𝒜`$ objects, $`H(A_1,\mathrm{},A_{m+n})=_{df}F(A_1,\mathrm{},A_m)G(A_{m+1},\mathrm{},A_{m+n})`$, and for every $`m+n`$-tuple $`(f_1,\mathrm{},f_{m+n})`$ of $`𝒜`$ arrows, $`H(f_1,\mathrm{},f_{m+n})=_{df}F(f_1,\mathrm{},f_m)G(f_{m+1},\mathrm{},f_{m+n})`$, corresponds to the term $`𝑭𝑮`$.
However, depending on the category $`𝒜`$, two different shapes may define the same functor. From now on, if we say that a functor $`F:𝒜^n𝒜`$ is from $``$, that means it corresponds to some shape from $``$.
Functors from $``$ will serve as domains and codomains of canonical transformations we are going to introduce below.
Let $`F:𝒜^m𝒜`$, $`G:𝒜^n𝒜`$ and $`H:𝒜^l𝒜`$ be functors from $``$.
a) Denote by $`\text{1}_F`$ the indexed family $`\{\text{1}_{F(A_1,\mathrm{},A_m)}(A_1,\mathrm{},A_m)𝒜^m\}`$ and let $`\mathrm{\Gamma }`$ be the identity function from $`\{1,\mathrm{},m\}`$ to $`\{1,\mathrm{},m\}`$. It is easy to see that $`\text{1}_F`$ is a g-natural transformation from $`F`$ to $`F`$ with the graph $`\mathrm{\Gamma }`$.
If $`𝒜`$ is monoidal
b) Denote by $`𝝈_F`$ the indexed family $`\{𝝈_{F(A_1,\mathrm{},A_m)}(A_1,\mathrm{},A_m)𝒜^m\}`$ and let $`\mathrm{\Gamma }`$ be as above. Then $`𝝈_F:\text{I}F\text{}F`$.
In a similar way we define $`𝝈_F^i`$, $`𝜹_F`$ and $`𝜹_F^i`$.
c) Denote by $`\stackrel{}{\text{b}}_{F,G,H}`$ the indexed family
$`\{\stackrel{}{\text{b}}_{F(A_1,\mathrm{},A_m),G(A_{m+1},\mathrm{},A_{m+n}),H(A_{m+n+1},\mathrm{},A_{m+n+l})}(A_1,\mathrm{},A_{m+n+l})𝒜^{m+n+l}\}`$ and let $`\mathrm{\Gamma }`$ be the identity function from $`\{1,\mathrm{},m+n+l\}`$ to $`\{1,\mathrm{},m+n+l\}`$. Then, $`\stackrel{}{\text{b}}_{F,G,H}:F(GH)\text{}(FG)H`$.
In a similar way we define $`\stackrel{}{\text{b}}_{F,G,H}`$.
If $`𝒜`$ is symmetric monoidal
d) Denote by $`\text{c}_{F,G}`$ the indexed set $`\{\text{c}_{F(A_1,\mathrm{},A_m),G(A_{m+1},\mathrm{},A_{m+n})}(A_1,\mathrm{},A_{m+n})𝒜^{m+n}\}`$ and let $`\mathrm{\Gamma }`$ be the function from $`\{1,\mathrm{},m+n\}`$ to $`\{1,\mathrm{},m+n\}`$ that satisfies $`\mathrm{\Gamma }(m+i)=i`$ for $`1in`$ and $`\mathrm{\Gamma }(j)=n+j`$ for $`1jm`$. Then $`\text{c}_{F,G}:FG\text{}GF`$.
e) If $`𝒜`$ is relevant category, we denote by $`\text{w}_F`$ the indexed family $`\{\text{w}_{F(A_1,\mathrm{},A_m)}(A_1,\mathrm{},A_m)𝒜^m\}`$. Let $`\mathrm{\Gamma }`$ be the function from $`\{1,\mathrm{},2m\}`$ to $`\{1,\mathrm{},m\}`$ defined as $`\mathrm{\Gamma }(i)=\mathrm{\Gamma }(m+i)=i`$ for $`1im`$. Then $`\text{w}_F:F\text{}FF`$.
f) If $`𝒜`$ is affine, we denote by $`\text{k}_F`$ the indexed family $`\{\text{k}_{F(A_1,\mathrm{},A_m)}(A_1,\mathrm{},A_m)𝒜^m\}`$. Let $`\mathrm{\Gamma }`$ be the empty function from $`\mathrm{}`$ to $`\{1,\mathrm{},m\}`$. Then $`\text{k}_F:F\text{}I`$.
The g-natural transformations from above, which exist in the category $`𝒜`$ constitute the class of $`𝒜`$ primitive canonical transformations. If we declare $`𝒜`$ is a monoidal category, its primitive canonical transformations are those from a) to c), though $`𝒜`$ may have the structure of a cartesian category.
Let $`F_1:𝒜^m𝒜`$, $`F_2:𝒜^n𝒜`$, $`G_1:𝒜^k𝒜`$, $`G_2:𝒜^l𝒜`$ be functors from $``$, and $`\alpha :F_1\text{}G_1`$ and $`\beta :F_2\text{}G_2`$. Denote by $`\alpha \beta `$ the family
$$\{\alpha (A_1,\mathrm{},A_m)\beta (A_{m+1},\mathrm{},A_{m+n})A_1,\mathrm{},A_{m+n}𝒜\},$$
and let $`\mathrm{\Gamma }`$ as a function from $`\{1,\mathrm{},k+l\}`$ to $`\{1,\mathrm{},m+n\}`$ satisfy the following:
$`\mathrm{\Gamma }(i)=\mathrm{\Phi }(i)`$ for $`1ik`$ and $`\mathrm{\Gamma }(k+j)=m+\mathrm{\Psi }(j)`$ for $`(1jl)`$. Then it is easy to see that $`\alpha \beta :F_1F_2\text{}G_1G_2`$.
Now, we can define canonical transformations in $`𝒜`$ as follows
1. Primitive canonical transformations from $`𝒜`$ are canonical transformations.
2. If $`\alpha :F_1\text{}G_1`$ and $`\beta :F_2\text{}G_2`$ are canonical transformations in $`𝒜`$, then $`\alpha \beta :F_1F_2\text{}G_1G_2`$ is canonical.
3. If $`\alpha :F\text{}G`$ and $`\beta :G\text{}H`$ are canonical transformations in $`𝒜`$, then $`\beta \alpha :F\text{}H`$ is canonical.
It is easy to verify that symmetric monoidal canonical transformations have bijections as graphs, relevant canonical transformations have onto functions as graphs and affine canonical transformations have one-one functions as graphs. Now we can reformulate MacLane’s results from in the following manner:
If $`𝒜`$ is monoidal or symmetric monoidal category and $`\alpha ,\beta :F\text{}G`$ are two canonical transformations (with the same graph $`\mathrm{\Gamma }`$), then $`\alpha `$ and $`\beta `$ are the same indexed sets (i.e., the same functions from the sequences of objects to the morphisms of $`𝒜`$).
This property of a category we call coherence. It completely follows the notion of coherence given in . We can extend this definition to an arbitrary substructural category $`𝒜`$. Namely, we say that an arbitrary substructural category is coherent if for every pair of canonical transformations $`\alpha `$ and $`\beta `$ of the same type and with the same graph we have that $`\alpha =\beta `$ as indexed sets.
## 4 Categories Mon, SyMon, Rel, Aff and Cart
Let $``$ be the category whose objects are monoidal categories and whose arrows are the monoidal structure preserving functors in the language given above. The equational axiomatization of monoidal categories enables us to distinguish a category from $``$ freely generated by a set of objects. Let $`P`$ be an infinite linearly ordered set of objects, whose elements we call letters. We denote by Mon the free monoidal category generated by $`P`$ whose construction is given below. In the same way we will introduce the categories SyMon, Rel, Aff and Cart, namely the free symmetric monoidal, relevant, affine and cartesian category generated by the same set $`P`$ of objects.
The constructions of these categories are algebraic and the set of objects is always the set $`𝒪`$ of terms freely generated by $`P\{\text{I}\}`$ using the binary operation $``$.
Primitive morphism-terms are in the case of Mon
$$\text{1}_A:AA$$
$$𝝈_A:\text{I}AA\text{ }𝝈_A^i:A\text{I}A$$
$$𝜹_A:A\text{I}A\text{ }𝜹_A^i:AA\text{I}$$
$$\stackrel{}{\text{b}}_{A,B,C}:A(BC)(AB)C\text{ }\stackrel{}{\text{b}}_{A,B,C}:(AB)CA(BC)$$
for all $`A,B,C𝒪`$.
SyMon primitive morphism-terms are those of Mon together with
$$\text{c}_{A,B}:ABBA$$
for every $`A,B𝒪`$.
Rel primitive morphism-terms are those of SyMon together with
$$\text{w}_A:AAA$$
for every $`A𝒪`$.
Aff primitive morphism-terms are those of SyMon together with
$$\text{k}_A:A\text{I}$$
for every $`A𝒪`$.
Rel and Aff primitive morphism-terms make the class of Cart primitive morphism-terms.
Morphism-terms are built from the primitive morphism-terms with the help of the binary operations of composition and multiplication.
Morphisms of the category Mon are equivalence classes of Mon morphism-terms modulo monoidal equations. Analogously, we define morphisms of other free categories mentioned above.
Let $`𝒞`$ be one among Mon, SyMon, Rel, Aff and Cart. We define a correspondence between the morphism-terms and canonical transformations of $`𝒞`$ in the following way
1. If $`f:AB`$ is primitive morphism-term, suppose it is of the form $`\text{1}_{F(p_1,\mathrm{},p_m)}`$ for some $`F:𝒞^m𝒞`$ and some, not necessarily distinct, letters $`p_1,\mathrm{},p_m`$. Then the canonical transformation $`\text{1}_F:F\text{}F`$ where $`\mathrm{\Gamma }`$ is the identity function on $`\{1,\mathrm{},m\}`$ corresponds to $`f`$. We procede similarly in the remaining cases.
2. If $`f`$ is of the form $`f_1f_2`$ or $`f_2f_1`$ and if the canonical transformations $`\alpha _1`$ and $`\alpha _2`$ correspond to the morphism-terms $`f_1`$ and $`f_2`$, then the canonical transformation $`\alpha _1\alpha _2`$ respectively $`\alpha _2\alpha _1`$ corresponds to the morphism-term $`f`$.
The graph of a transformation that corresponds to the morphism-term $`f:AB`$ we call simply a graph of $`f`$. It connects the occurrence of a letter in $`A`$ with a set (maybe empty) of occurrences of the same letter in $`B`$.
Let $`\alpha `$ be a term of a $`𝒞`$ canonical transformation of type $`F\text{}G`$, for $`F:𝒞^m𝒞`$, $`G:𝒞^n𝒞`$ $``$. Let $`p_1,p_2,\mathrm{},p_m`$ be distinct letters from $`P`$. We call the morphism-term
$$\alpha (p_1,\mathrm{},p_m):F(p_1,\mathrm{},p_m)G(p_{\mathrm{\Gamma }(1)},\mathrm{},p_{\mathrm{\Gamma }(n)})$$
the representative of the transformation $`\alpha `$.
Lemma 1 Let $`𝒜`$ be an arbitrary substructural category and $`𝒞`$ one of the free categories mentioned above which is of the same type as $`𝒜`$. Let $`F`$ and $`G`$ be from $``$ and let $`\alpha :F\text{}G`$ and $`\beta :F\text{}G`$ be in $`𝒜`$. Denote by $`𝛂`$ and $`𝛃`$ canonical transformations in $`𝒞`$ defined by the same terms as $`\alpha `$ and $`\beta `$ respectively. Let $`f:AB`$ be the representative of $`𝛂`$ and $`g:AB`$ the representative of $`𝛃`$. If $`f=g`$ in $`𝒞`$, then $`\alpha =\beta `$ in $`𝒜`$.
proof Suppose that $`f𝜶(p_1,\mathrm{},p_m):F(p_1,\mathrm{},p_m)G(p_{\mathrm{\Phi }(1)},\mathrm{},p_{\mathrm{\Phi }(n)})`$ for some distinct letters $`p_1,\mathrm{},p_mP`$. Since $`g`$ is the representative equal to $`f`$, they share domains and codomains; hence $`g𝜷(p_1,\mathrm{},p_m):F(p_1,\mathrm{},p_m)G(p_{\mathrm{\Psi }(1)},\mathrm{},p_{\mathrm{\Psi }(n)})`$, which implies that for every $`1in`$, $`\mathrm{\Phi }(i)=\mathrm{\Psi }(i)`$, and so $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are equal graphs. Suppose that $`f^{}\alpha (A_1,\mathrm{},A_m)`$. By the assumption that $`𝒞`$ is free, there is a functor $`U:𝒞𝒜`$ that preserves the structure of $`𝒞`$ and that extends the mapping of the generators given by $`p_1A_1,p_2A_2,\mathrm{},p_mA_m`$ (other generators are mapped arbitrarily). Then we have
$$\alpha f^{}=\alpha (A_1,\mathrm{},A_k)=U(f)=U(g)=\beta (A_1,\mathrm{},A_k)\beta ,$$
hence $`\alpha \beta `$. In the same way we prove that $`\beta \alpha `$.
q.e.d.
Our goal is to prove that every substructural category is coherent, and the following lemma will serve to reduce the problem to the case of the free category in the type. For an object $`A𝒪`$ we say that it is diversified if no letter occurs twice in it.
Lemma 2 If for every pair of $`𝒞`$ morphism-terms $`f,g:AB`$, such that $`A`$ is diversified, holds that $`f=g`$, then every category of the same type as $`𝒞`$ is coherent.
proof Let $`𝒜`$ be an arbitrary substructural category of the same type as $`𝒞`$. Suppose that $`F`$ and $`G`$ are functors from $``$ and $`\alpha ,\beta :F\text{}G`$ in $`𝒜`$. Let $`f𝜶(p_1,\mathrm{},p_m):AB`$ and $`g𝜷(p_1,\mathrm{},p_m):CD`$ be the representatives of $`𝜶`$ and $`𝜷`$ respectively, where $`𝜶`$ and $`𝜷`$ are $`𝒞`$ canonical transformations defined by the same terms as $`\alpha `$ and $`\beta `$. Since they have the same graph, $`A`$ must be identical to $`C`$ and $`B`$ to $`D`$. By the definition of graph, it follows that $`A`$ is diversified and, by the assumption, $`f`$ is equal to $`g`$. Hence, by Lemma 1, we have $`\alpha =\beta `$.
q.e.d.
By the following series of definitions we introduce some auxiliary notions that will help us in proving our coherence results.
Denote by $`𝒫`$ the set of terms generated by the binary operation $``$ from the elements of $`P\{\mathrm{},\text{I}\}`$ (e.g. $`(\mathrm{}p)((\text{I}\mathrm{})q)`$ is in $`𝒫`$). As in the case of terms from $``$, we can define in the same way the correspondence between terms from $`𝒫`$ and functors (with parameters) of the type $`𝒞^n𝒞`$ for some $`n0`$, where $`𝒞`$ is one of the free categories mentioned above.
A product term of $`𝒞`$ is a morphism-term defined recursively as follows
1. The primitive terms (if they exist in $`𝒞`$)
$$𝝈_Q,𝝈_Q^i,𝜹_Q,𝜹_Q^i,\stackrel{}{\text{b}}_{Q,S,R},\stackrel{}{\text{b}}_{Q,S,R},\text{c}_{Q,S},\text{w}_Q,\text{k}_Q.$$
are product terms, called determining factors.
2. The terms $`\text{1}_Q`$ are product terms.
3. If $`f`$ is a product term, then $`\text{1}_Qf`$ and $`f\text{1}_Q`$ are product terms.
The determining factor of a product term $`f`$, if it exists, is denoted by $`d(f)`$ (we call such a term structural product). A structural product $`f`$ is a b-product iff $`d(f)`$ is a b term, c-product iff $`d(f)`$ is a c term, and similarly for $`𝝈`$, $`𝝈^i`$, $`𝜹`$, $`𝜹^i`$, k and w-products.
For a w-product-term we say that it is atomic if the index of its determining factor is a letter.
We say that an atomic w-product is left if there is not any 1 with the letter $`p`$ in the index, on the left of its determining factor $`\text{w}_p`$ (e.g. $`(\text{1}_{qr}\text{w}_p)\text{1}_p`$ is the left atomic w-product, while $`(\text{1}_{pr}\text{w}_p)\text{1}_p`$ is not left.)
For a c-product we say that it is atomic if the index of its determining factor is a pair of atoms (an atom is a letter or I).
We say that an atomic c-product is diversified if its determining factor is not of the form $`\text{c}_{p,p}`$ for some letter $`p`$.
For a k-product we say that it is atomic if the index of its determining factor is a letter.
We say that a composition of atomic w-products (k-products) is ordered, if a $`\text{w}_p`$-product ($`\text{k}_p`$-product) is to the right of $`\text{w}_q`$-product ($`\text{k}_q`$-product) in this composition iff the letter $`p`$ precedes the letter $`q`$ in the ordering of $`P`$.
## 5 Coherence in Relevant categories
At the beginning of this section, we prove a lemma that states coherence for bw fragments of relevant categories.
Lemma 3 Let $`F`$ be from $``$ and let $`f:pF(p,\mathrm{},p)`$, be a composition of atomic w-products. Then $`f`$ is equal to a term of the form $`hg`$ where $`g`$ is a composition of left atomic w-products and $`h`$ is a composition of b-products.
proof For the sake of clarity we introduce a tree that corresponds to $`f`$, denoted by $`\tau _f`$, in the following way.
If $`f\text{w}_p`$, then $`\tau _f`$ is
If $`f`$ is of the form $`G(\text{w}_p)f_1`$, where $`G`$ is from $`𝒫`$, and if in the shape of $`G`$, $`i1`$ letters $`p`$ precede (from the left) the symbol $`\mathrm{}`$, then $`\tau _f`$ is obtained from $`\tau _{f_1}`$ by forking the $`i`$-th leaf (from the left) and concatenating simple segments to remaining leaves.
For example, if $`f`$ is of the form
$$((\text{1}_p\text{w}_p)\text{1}_{(pp)p})(\text{1}_{pp}(\text{w}_p\text{1}_p))(\text{1}_{pp}\text{w}_p)(\text{w}_p\text{1}_p)\text{w}_p$$
then the corresponding tree is
Denote by $`\mathrm{\Lambda }`$ the set of forking vertices in such a tree. Let $`k_\lambda `$ be the number of right branches of these forkings in the path from the vertex $`\lambda `$ to the root. The complexity of the tree is measured by the number $`n_f`$ defined as
$$n_f=_{df}\underset{\lambda \mathrm{\Lambda }}{}k_\lambda $$
In the example above $`n_f`$ is $`3`$. We prove the lemma by induction on $`n_f`$.
If $`n_f=0`$, then $`f`$ is itself a composition of left atomic w-products.
If $`n_f>0`$ and if there is no subtree of $`\tau _f`$ of the form
then by the functoriality of $``$ we obtain a term $`f^{}`$ equal to $`f`$ such that $`n_f^{}=n_f`$, and there is a subtree of the above form in $`\tau _f^{}`$. In the example, we obtain the term
$$((\text{1}_p\text{w}_p)\text{1}_{(pp)p})(\text{1}_{pp}(\text{w}_p\text{1}_p))(\text{w}_p\text{1}_{pp})(\text{1}_p\text{w}_p)\text{w}_p,$$
whose tree is
By using the (bw) equality and the naturality of b-products we transform this term into the form $`h_1f_1`$ where $`f_1`$ is a composition of atomic w-products with $`n_f^{}=n_f1`$, and $`h_1`$ is a composition of b-products. In our example we obtain the term
$$\stackrel{}{\text{b}}_{p(pp),pp,p}(((\text{1}_p\text{w}_p)\text{1}_{pp})\text{1}_p)((\text{1}_{pp}\text{w}_p)\text{1}_p)((\text{w}_p\text{1}_p)\text{1}_p)(\text{w}_p\text{1}_p)\text{w}_p,$$
and the tree corresponding to its initial part ($`f_1`$) is
By the induction hypothesis the term $`f_1`$ is equal to the term $`h_2g`$ of the desired form, and therefore $`f=h_1h_2g`$ is such.
q.e.d.
In our example, the last term is transformed into
$$\stackrel{}{\text{b}}_{p(pp),pp,p}(((\text{1}_p\text{w}_p)\text{1}_{pp})\text{1}_p)((\text{w}_p\text{1}_{pp})\text{1}_p)((\text{1}_p\text{w}_p)\text{1}_p)(\text{w}_p\text{1}_p)\text{w}_p,$$
and the tree corresponding to its initial part is
(we do this to obtain a subtree of the form )
and then by (bw) and (b) this term is transformed into
$$\stackrel{}{\text{b}}_{p(pp),pp,p}(\stackrel{}{\text{b}}_{p(pp),p,p}\text{1}_p)((((\text{1}_p\text{w}_p)\text{1}_p)\text{1}_p)\text{1}_p)((\text{w}_p\text{1}_p)\text{1}_p)\text{1}_p)((\text{w}_p\text{1}_p)\text{1}_p)(\text{w}_p\text{1}_p)\text{w}_p,$$
to whose initial part corresponds the tree
Then, again by (bw) and (b) the last term is transformed into the term
$`\stackrel{}{\text{b}}_{p(pp),pp,p}(\stackrel{}{\text{b}}_{p(pp),p,p}\text{1}_p)(((\stackrel{}{\text{b}}_{p,p,p}\text{1}_p)\text{1}_p)\text{1}_p)`$
$`((((\text{w}_p\text{1}_p)\text{1}_p)\text{1}_p)\text{1}_p)(((\text{w}_p\text{1}_p)\text{1}_p)\text{1}_p)((\text{w}_p\text{1}_p)\text{1}_p)(\text{w}_p\text{1}_p)\text{w}_p,`$
of the desired form, whose tree is of the form
Corollary Let $`F:\text{Rel}^k\text{Rel}`$ be from $``$ and let $`f:pF(p,\mathrm{},p)`$, be a composition of atomic w-products. Then for every $`i`$, $`1ik1`$ there is a morphism term of the form
$$v((\text{1}_{\underset{i1}{\underset{}{pp\mathrm{}p}}}\text{w}_p)\text{1}_{\underset{ki1}{\underset{}{pp\mathrm{}p}}})u$$
equal to $`f`$ (all products of $`p`$’s are associated to the left, i.e., $`ppp`$ means $`(pp)p`$), where $`u:p(pp\mathrm{}p)(pp\mathrm{}p)`$ is a composition of atomic w-products, and $`v`$ is a composition of b-products.
proof By Lemma 3, there is a term of the form $`h_1g`$ equal to $`f`$, where $`g`$ is a composition of left atomic w-products, and $`h_1`$ is a composition of b-products. Let
$$u:p\underset{i}{\underset{}{(pp\mathrm{}p)}}\underset{ki1}{\underset{}{(pp\mathrm{}p)}}$$
be a composition of atomic w-products (such always exists). Again, by Lemma 3 there is a term of the form $`h_2g`$ equal to the term
$$((\text{1}_{\underset{i1}{\underset{}{pp\mathrm{}p}}}\text{w}_p)\text{1}_{\underset{ki1}{\underset{}{pp\mathrm{}p}}})u,$$
where $`g`$ is as before, and $`h_2`$ is a composition of b-products. Then,
$$f=h_1g=h_1h_2^1h_2g=h_1h_2^1((\text{1}_{(pp)\mathrm{}p}\text{w}_p)\text{1}_{(pp)\mathrm{}p})u,$$
where $`h_2^1`$ denotes a composition of b-products inverse to $`h_2`$ (by (bb) equalities).
q.e.d.
This corollary will be of use in the proof of the main result of this section, which states that
Theorem 1 Every relevant category is coherent.
In the proof of the theorem, the following lemma, that gives the normal form of a morphism from Rel, is crucial.
Lemma 4 Let $`h:AB`$ be a morphism-term from Rel with $`A`$ diversified. Then $`h`$ is equal to the morphism-term of the form $`h^{\prime \prime }h^{}`$ where $`h^{}`$ is an ordered composition of atomic left w-products, and $`h^{\prime \prime }`$ is a composition of products with no w-products and with all c-products atomic diversified.
proof The transformation of the term $`h`$ is made in several steps. For the sake of clarity, we illustrate every step starting with the term
$$(\text{1}_q(\text{w}_{pp}\text{w}_p))\text{c}_{p,q}.$$
$`1^{}`$ The term $`h`$ is equal to a composition of product terms. This follows from the functoriality of multiplication. In our example,
$$h=(\text{1}_q\text{w}_{pp})(\text{1}_q\text{w}_p)\text{c}_{p,q}.$$
$`2^{}`$ By (bcw8) and ($`𝝈`$$`𝜹`$w), $`h`$ is equal to a term with all w-products atomic. In our example
$$h=t(\text{1}_q(\text{w}_p\text{1}_{pp}))(\text{1}_q(\text{1}_p\text{w}_p))(\text{1}_q\text{w}_p)\text{c}_{p,q},$$
where
$$t(\text{1}_q\stackrel{}{\text{b}}_{pp,p,p})(\text{1}_q(\stackrel{}{\text{b}}_{p,p,p}\text{1}_p))(\text{1}_q((\text{1}_p\text{c}_{p,p})\text{1}_p))(\text{1}_q(\stackrel{}{\text{b}}_{p,p,p}\text{1}_p))(\text{1}_q\stackrel{}{\text{b}}pp,p,p)$$
$`3^{}`$ By the multiplication functoriality and the naturality of $`𝝈`$, $`𝜹`$, b, c-products, atomic w-products permute towards the right end, in order to obtain a term whose initial part (from the right) consists of atomic w-products and whose tail is a w-free composition of products. In our example this term is
$$t\text{c}_{(pp)(pp),q}((\text{w}_p\text{1}_{pp})\text{1}_q)((\text{1}_p\text{w}_p)\text{1}_q)(\text{w}_p\text{1}_q)$$
$`4^{}`$ Using (bc6), this term is equal to a term whose c-products are atomized. In the example, this atomization is not essential because its application to the unique nonatomic c-product $`\text{c}_{(pp)(pp),q}`$ produces only diversified atomic c-products (see the following step), and therefore we write it in abbreviated form, which is enough for the further analysis:
$$h=t_2(\text{1}_q((\text{1}_p\text{c}_{p,p})\text{1}_p))t_1((\text{w}_p\text{1}_{pp})\text{1}_q)((\text{1}_p\text{w}_p)\text{1}_q)(\text{w}_p\text{1}_q),$$
$$\text{where}t_2(\text{1}_q\stackrel{}{\text{b}}_{pp,p,p})(\text{1}_q(\stackrel{}{\text{b}}_{p,p,p}\text{1}_p)),$$
$$\text{and}t_1(\text{1}_q(\stackrel{}{\text{b}}_{p,p,p}\text{1}_p))(\text{1}_q\stackrel{}{\text{b}}pp,p,p)(\text{c}_{(pp)(pp),q})^{}$$
($``$ means a developed form with atomic c-products.)
$`5^{}`$ Suppose that the tail (w-free composition of products) from the last term is in the form $`t_2F(\text{c}_{p,p})t_1`$, where $`F:\text{Rel}\text{Rel}`$ is from $`𝒫`$ and all the c-products in $`t_1`$ are atomic diversified. Let $`F(\text{c}_{p,p})t_1`$ be of the type $`G(p,p)F(pp)`$ for some $`G:\text{Rel}^2\text{Rel}`$ from $`𝒫`$, where the left $`p`$ from $`F(p,p)`$ is mapped to the right $`p`$ from $`G(pp)`$ by the graph of $`F(\text{c}_{p,p})t_1`$, and the right $`p`$ from $`F(p,p)`$ is mapped to the left $`p`$ from $`F(pp)`$ by the same graph. By the assumption that $`A`$ is diversified, that all c-products in $`t_1`$ are atomic diversified and the initial part consists of atomic w-products only, we have that two emphasized letters $`p`$ in $`G(p,p)`$ occur consecutively (not necessarily in the form $`(pp)`$).
Using the functoriality of $``$, we can push all $`\text{w}_p`$ products to the (left) end of an initial part consisting of w-products only. Now, we apply the corollary of Lemma 3, to show that such an initial part is equal to a term of the form $`i_2H(\text{w}_p)i_1`$, where $`i_1`$ is a composition of atomic w-products, $`H:\text{Rel}\text{Rel}`$ is from $`𝒫`$, $`i_2`$ is a composition of b-products, the term $`i_2H(\text{w}_p)`$ is of the type $`H(p)G(p,p)`$ and its graph maps both distinguished $`p`$’s from $`G(p,p)`$ to the distinguished $`p`$ in $`H(p)`$. (Assume that distinguished $`(p,p)`$ in $`G(p,p)`$ are $`i`$-th and $`i+1`$-th occurrence of the letter $`p`$ and just apply the corollary of Lemma 3.)
In our example
$$h=t_2(\text{1}_q((\text{1}_p\text{c}_{p,p})\text{1}_p))t_1i_2(((\text{1}_p\text{w}_p)\text{1}_p)\text{1}_q)i_1,$$
where $`i_1((\text{w}_p\text{1}_p)\text{1}_q)(\text{w}_p\text{1}_q)`$, $`i_2(\stackrel{}{\text{b}}_{pp,p,p}\text{1}_q)((\stackrel{}{\text{b}}_{p,p,p}\text{1}_p)\text{1}_q)`$,
and $`F(q(p\mathrm{}))p)`$, $`G((p\mathrm{})(\mathrm{}p))q`$, $`H((p\mathrm{})p)q`$.
The SyMon terms $`F(\text{c}_{p,p})t_1i_2`$ and $`t_1i_2H(\text{c}_{p,p})`$ are of the same type. Denote by $`\alpha `$ and $`\beta `$ canonical transformations in SyMon corresponding to these terms. It is easy to see that $`\alpha `$ and $`\beta `$ have the same graphs (the graph of $`H(\text{c}_{p,p})`$ “commutes” in composition with the graph of $`t_1i_2`$ and is transformed into the graph of $`F(\text{c}_{p,p})`$). By MacLane’s coherence for symmetric monoidal categories, we have that $`\alpha =\beta `$. From the property that every canonical transformation of SyMon contains at most one morphism of a certain type (this is because the set $`𝒪`$ of its objects is freely generated by $`P\{\text{I}\}`$), we conclude that $`F(\text{c}_{p,p})t_1i_2=t_1i_2H(\text{c}_{p,p})`$ holds in SyMon. Now, because all the SyMon-equalities hold in Rel, these terms are equal in Rel too.
Therefore, $`h`$ is equal to a term of the form $`t_2t_1i_2H(\text{c}_{p,p})H(\text{w}_p)i_1`$, which is by (cw), equal to $`t_2t_1i_2H(\text{w}_p)i_1`$. Repeating this procedure we can eliminate non diversified c products occurring in $`t_2`$ obtaining a term equal to $`h`$ whose initial part consists of atomic w-products and whose tail is a w-free composition of products, whose c-products are atomic diversified.
In the example, this procedure includes the following steps
$`t_2(\text{1}_q((\text{1}_p\text{c}_{p,p})\text{1}_p))t_1i_2(((\text{1}_p\text{w}_p)\text{1}_p)\text{1}_q)i_1=`$
$`t_2t_1i_2(((\text{1}_p\text{c}_{p,p})\text{1}_p)\text{1}_q)(((\text{1}_p\text{w}_p)\text{1}_p)\text{1}_q)i_1=`$
$`t_2t_1i_2(((\text{1}_p\text{w}_p)\text{1}_p)\text{1}_q)i_1,`$
where $`t_2(\text{1}_q\stackrel{}{\text{b}}_{pp,p,p})(\text{1}_q(\stackrel{}{\text{b}}_{p,p,p}\text{1}_p))`$.
$`6^{}`$ By Lemma 3 and the functoriality of multiplication, this term is equal to the one whose initial part is an ordered composition of atomic left w-products.
q.e.d.
In the example, we have
$`(((\text{1}_p\text{w}_p)\text{1}_p)\text{1}_q)i_1`$
$`(((\text{1}_p\text{w}_p)\text{1}_p)\text{1}_q)((\text{w}_p\text{1}_p)\text{1}_q)(\text{w}_p\text{1}_q)=`$
$`((\stackrel{}{\text{b}}_{p,p,p}\text{1}_p)\text{1}_q)(((\text{w}_p\text{1}_p)\text{1}_p)\text{1}_q)((\text{w}_p\text{1}_p)\text{1}_q)(\text{w}_p\text{1}_q).`$
Lemma 5 Let $`f,g:AB`$ be two morphism-terms in Rel with $`A`$ diversified. Then $`f=g`$ in Rel.
proof By Lemma 4, $`f=f^{\prime \prime }f^{}`$ and $`g=g^{\prime \prime }g^{}`$, with $`f^{},f^{\prime \prime },g^{},g^{\prime \prime }`$ of the given form. The terms $`f^{}`$ and $`g^{}`$ are completely determined by the codomain $`B`$ (the number of occurrences of each letter in $`B`$ determines $`f^{}`$ and $`g^{}`$) and therefore $`f^{}`$ and $`g^{}`$ are identical. Suppose its type is $`AA^{}`$. Then the terms $`f^{\prime \prime }`$ and $`g^{\prime \prime }`$ are of the same type $`A^{}B`$. Let $`\alpha `$ and $`\beta `$ be the SyMon canonical transformations corresponding to these terms. Taking $`\mathrm{\Gamma }_\alpha `$ and $`\mathrm{\Gamma }_\beta `$ (the graphs of these transformations) as connections that connect an occurrence of a letter in $`B`$ with occurrence of the same letter in $`A^{}`$ and since all c-products in $`f^{\prime \prime }`$ and $`g^{\prime \prime }`$ are atomic diversified, they must connect the first (from the left) occurrence of one letter in $`B`$ with the first (again from the left) occurrence of the same letter in $`A^{}`$, the second with the second etc. This means that $`\mathrm{\Gamma }_\alpha `$ and $`\mathrm{\Gamma }_\beta `$ are the same graphs and by coherence in symmetric monoidal categories, $`\alpha `$ and $`\beta `$ are the same canonical transformations in SyMon. Therefore, $`f^{\prime \prime }`$ and $`g^{\prime \prime }`$ are equal in SyMon, and since all symmetric monoidal equalities hold in Rel, they are equal in Rel too. Hence, $`f=g`$ in Rel.
q.e.d.
Now Theorem 1 follows from lemmata 2 and 5.
## 6 Coherence in Affine Categories
As in the case of relevant categories, first we prove a lemma about representation of Aff morphism terms.
Lemma 6 Every Aff morphism-term is equal to a term of the form $`h_2h_1`$, where $`h_1`$ is an ordered composition of atomic k-products, and no k-product occurs in $`h_2`$.
proof As in the proof of Lemma 4, we transform the Aff morphism-term $`h`$ in several steps.
$`1^{}`$ By the functoriality of multiplication, $`h`$ is equal to a composition of product terms.
$`2^{}`$ By the equalities (1k) and $`\text{k}_{AB}=𝝈_\text{I}(\text{1}_\text{I}\text{k}_B)(\text{k}_A\text{1}_B)`$, which is derivable from (k) and (1k), this term is equal to a composition of products with all k-products atomic.
$`3^{}`$ By the naturality of $`𝝈`$,$`𝜹`$,b,c-products and functoriality of multiplication, atomic k-products permute to the right in order to obtain a term whose initial part consists of all k-products present in the term.
$`4^{}`$ By the functoriality of multiplication, we can order this initial part to obtain a term in the desired form, which is equal to $`h`$.
q.e.d.
Lemma 7 Let $`f,g:AB`$ be two morphism terms from Aff with $`A`$ diversified. Then $`f=g`$ in Aff.
proof By Lemma 6, there are terms $`f^{},f^{\prime \prime },g^{},g^{\prime \prime }`$ such that $`f=f^{\prime \prime }f^{}`$ and $`g=g^{\prime \prime }g^{}`$, where $`f^{}`$ and $`g^{}`$ are ordered compositions of atomic k-products, and $`f^{\prime \prime }`$ and $`g^{\prime \prime }`$ are SyMon terms. The objects $`A`$ and $`B`$ (the letters occurring in $`A`$ not in $`B`$) completely determine terms $`f^{}`$ and $`g^{}`$, hence they are identical morphism terms, suppose of the type $`AA^{}`$. Therefore $`f^{\prime \prime }`$ and $`g^{\prime \prime }`$ are of the type $`A^{}B`$. By the assumption concerning $`A`$ it follows that $`A^{}`$ and $`B`$ are diversified. Let $`\alpha `$ and $`\beta `$ be the canonical transformations in SyMon corresponding to $`f^{\prime \prime }`$ and $`g^{\prime \prime }`$ respectively. Taking $`\mathrm{\Gamma }_\alpha `$ and $`\mathrm{\Gamma }_\beta `$ (theirs graphs) as connections between letter occurrences in $`B`$ with letter occurrences in $`A^{}`$, since they connect a letter occurrence with an occurrence of the same letter, and by the assumption about $`A^{}`$ and $`B`$, we must have that $`\mathrm{\Gamma }_\alpha =\mathrm{\Gamma }_\beta `$. This implies, by the coherence in symmetric monoidal categories, that $`\alpha =\beta `$, which has as consequence that $`f^{\prime \prime }=g^{\prime \prime }`$ in SyMon, hence in Aff. We conclude that $`f=f^{\prime \prime }f^{}=g^{\prime \prime }g^{}=g`$.
q.e.d.
From lemmata 2 and 7 it follows that
Theorem 2 Every affine category is coherent.
## 7 Coherence in Cartesian Categories
The coherence in cartesian categories is not a new result. For the first time it was mentioned in and more recently in and . Since we would like to keep to the definition of coherence given above, we give another proof of this result here.
One could expect that the proof of the coherence in cartesian categories will follow the proofs given in the last two sections. However, this method turns out to be too complicated and we use the standard equational axiomatization of cartesian categories to avoid this.
Denote by $`𝒫`$ the set of translations of morphism terms from Cart into the language of standard axiomatization (see Section 1).
Distributed terms form the smallest class of morphism terms from $`𝒫`$ that satisfies:
1. For all Cart objects $`A,B,C,D,E`$, the term $`\text{1}_A`$ as well as well-founded compositions of $`𝝅_{A,B},𝝅_{C,D}^{},\text{k}_E`$ are in the class and we call them compat.
2. If $`f:CA`$ and $`g:CB`$ are in the class then $`f,g`$ is in the class.
The following corresponds to the notion of the expanded normal form of a natural deduction proof.
A distributed term is atomic if every compat in this term has an atomic codomain (letter or I).
Lemma 8 Every morphism-term from $`𝒫`$ is equal to an atomic distributed term.
proof First, we show by induction on complexity of $`f`$ from $`𝒫`$ that it is equal to a distributed term.
$`1^{}`$ If $`f`$ is $`\text{1}_A`$, $`𝝅_{A,B}`$, $`𝝅_{A,B}^{}`$ or $`\text{k}_A`$, then it is compat and therefore distributed.
$`2^{}`$ a) Suppose that $`f`$ is of the form $`g,h`$. By the induction hypothesis, $`g`$ and $`h`$ are equal to distributed terms $`g^{}`$ and $`h^{}`$; hence $`f`$ is equal to the distributed term $`g^{},h^{}`$.
b) Suppose that $`f`$ is of the form $`hg`$. Then by the induction hypothesis, $`g`$ and $`h`$ are equal to distributed terms $`g_1`$ and $`h_1`$. Suppose that every composition of lower complexity than $`h_1g_1`$, of distributed terms is equal to a distributed term (if $`g_1`$ and $`h_1`$ are primitive, then its composition is compat and therefore distributed). There are three possibilities.
$`i)`$ If $`g_1`$ and $`h_1`$ are compat, then $`h_1g_1`$ is compat too, hence $`f`$ is equal to the distributed term $`h_1g_1`$.
$`ii)`$ If $`h_1`$ is of the form $`j,l`$, for $`j`$ and $`l`$ distributed, then $`f=j,lg_1=jg_1,lg_1`$. The terms $`jg_1`$ and $`lg_1`$ are of lower complexity than $`h_1g_1`$, and by assumption they are equal to some distributed terms; hence $`f`$ is equal to a distributed term.
$`iii)`$ Suppose that $`h_1`$ is compat and $`g_1`$ is of the form $`j,l`$. If $`h_1\text{1}`$, then $`f`$ is equal to the distributed term $`g_1`$. If $`h_1h_2𝝅`$, then $`f=h_2𝝅j,l=h_2j`$, where $`h_2`$ and $`j`$ are distributed and $`h_2j`$ is of lower complexity than $`h_1g_1`$; therefore it is equal to a distributed term. The case when $`h_1h_2𝝅^{}`$ is analogous. If $`h_1h_2\text{k}`$, then $`f=h_2\text{k}j,l=h_2\text{k}`$. The last term is compat, hence distributed.
This is the end of the induction. It follows that every $`f`$ from $`𝒫`$ is equal to a distributed term $`f_1`$. For every nonatomic compat $`h`$ in $`f_1`$, using the equality $`h=𝝅_{A,B}h,𝝅_{A,B}^{}h`$ for $`h:CAB`$, we can find a distributed term equal to $`h`$, such that every compat in it has a codomain of lower complexity than $`h`$. Substituting this term for $`h`$ and repeating the procedure we obtain an atomic distributed term equal to $`f_1`$ and therefore to $`f`$.
q.e.d.
Lemma 9 If $`f,g:AB`$ are two morphism-terms from $`𝒫`$ with $`A`$ diversified, and $`B`$ an atom, then $`f=g`$.
proof If $`B\text{I}`$, then because it is terminal in Cart, we have that $`f=g`$. Suppose that $`Bq`$ for $`q`$ a letter. By the previous lemma $`f`$ and $`g`$ are equal to distributed terms $`f_1`$ and $`g_1`$, which are compat by the assumption that codomain is $`q`$. Keeping in mind that there is no morphism in Cart of the type $`Aq`$ such that $`q`$ doesn’t occur in $`A`$, we prove the lemma by induction on the complexity of the domain $`A`$.
$`1^{}`$ If $`A`$ is an atom then $`f_1g_1\text{1}_q`$.
$`2^{}`$ Suppose that $`AA_1A_2`$. Then, by the assumption, $`q`$ occurs either in $`A_1`$ or in $`A_2`$. Suppose it occurs in $`A_1`$. Then we must have that $`f_1f_2𝝅`$ and $`g_1g_2𝝅`$ for some compat $`f_2,g_2:A_1q`$. By the induction hypothesis, $`f_2=g_2`$ holds, and therefore $`f=f_1=f_2𝝅=g_2𝝅=g_1=g`$. We prove analogously the case when $`q`$ occurs only in $`A_2`$.
q.e.d.
Lemma 10 Let $`f,g:AB`$ be two morphism-terms from $`𝒫`$ with $`A`$ diversified. Then $`f=g`$.
proof By Lemma 8, $`f`$ and $`g`$ are equal to atomic distributed terms $`f_1`$ and $`g_1`$. The proof follows by induction on the complexity of the codomain $`B`$. If $`B`$ is an atom, then by the previous lemma $`f=g`$ holds. If $`B`$ is not an atom, then neither $`f_1`$ nor $`g_1`$ are compat, and therefore $`f_1i,j`$ and $`g_1l,h`$, where $`B=B_1B_2`$. The terms $`i,l:AB_1`$ and $`j,h:AB_2`$ are atomic distributed and $`B_1`$, $`B_2`$ are of lower complexity than $`B`$. Therefore, by the induction hypothesis, $`i=l`$ and $`j=h`$; hence $`f=g`$.
q.e.d.
Corollary Let $`f,g:AB`$ be morphism terms from Cart and let $`A`$ be diversified. Then $`f=g`$ in Cart.
proof Let $`𝒇`$ and $`𝒈`$ from $`𝒫`$ correspond to $`f`$ and $`g`$ respectively. By Lemma 10, $`𝒇`$$`=`$$`𝒈`$ and by the extensional equivalence of these two axiomatizations we have that $`f=g`$ in Cart.
q.e.d.
From this corollary and Lemma 2 it follows that
Theorem 3 Every cartesian category is coherent.
## 8 Some consequences of the coherence
Usually, in the literature, coherence is not related to matters concerning natural transformations, but to conditions that imply equality of morphisms of certain categories. In the previous sections, we have lemmata 3, 5 and 7 as examples of such an opinion. Now we prove some facts that can be of practical interest for substructural categories, especially for the free categories of each type.
In Section 4, a definition of canonical transformation that corresponds to a morphism term $`f`$ of a free category $`𝒞`$ is given. From now on, a graph of this transformation will be called a graph of $`f`$. By a straightforward induction on the complexity of $`f`$ we can prove the following.
Lemma 11 If $`f`$ and $`g`$ are two equal morphism terms in $`𝒞`$, then the graph of $`f`$ is identical to the graph of $`g`$.
Denote by $`\text{Finord}^{op}`$ the dual of the category whose objects are finite ordinals and whose arrows are mappings between them. The coherence of substructural categories together with Lemma 11 is equivalent to the fact that there exist embeddings of SyMon, Rel, Aff and Cart into $`\text{Finord}^{op}`$ given by the “graph” functor $`G`$, such that for every $`A𝒪`$, $`G(A)`$ is the number of occurrences of letters in $`A`$ and $`G(f)`$ is the graph of $`f`$.
In the case of Cart, this embedding is onto on morphisms, and if we restrict ourserlves to one-one functions in $`\text{Finord}^{op}`$, then the embedding of Aff in this category is also onto on morphisms. Similarly we obtain embeddings which are onto on morphisms of Rel and SyMon in $`\text{Finord}^{op}`$ with restrictions to onto functions and bijections respectively. This is an alternative characterization of the coherence of substructural categories.
The results obtained in previous sections imply that the categories Rel, Aff and Cart are trivial in some sense. However, they are not preorders (in a preorder, there is at most one arrow between two objects) as the case is with Mon, though lemmata 5, 7 and 10 come close to preordering. We can use the following consequence of coherence in these categories:
Let $`𝒞`$ be one of the mentioned free categories and let
be one of its diagrams. It commutes iff the graphs of $`f`$ and $`g`$ are identical.
Together with freedom of $`𝒞`$, this consequence can be of practical use because it transforms computations in algebra of morphism terms to the simple calculus of morphism graphs.
Example Suppose we want to simplify the term
$$((𝜹_A(\text{1}_A𝝈_\text{I}))(𝜹_B𝝈_{B\text{I}}))\text{c}_{A,\text{I},\text{I}\text{I},B\text{I}}^m(\text{1}_{A\text{I}}\text{c}_{\text{I},B,\text{I},\text{I}}^m)(𝜹_A^i(𝝈_B^i𝜹_\text{I}^i))$$
representing a morphism of some symmetric monoidal category $`𝒮`$. The type of this term is $`A(B\text{I})AB`$. Consider the SyMon term
$$((𝜹_p(\text{1}_p𝝈_\text{I}))(𝜹_q𝝈_{q\text{I}}))\text{c}_{p,\text{I},\text{I}\text{I},q\text{I}}^m(\text{1}_{p\text{I}}\text{c}_{\text{I},q,\text{I},\text{I}}^m)(𝜹_p^i(𝝈_q^i𝜹_\text{I}^i)):p(q\text{I})pq.$$
Since its domain is diversified, it is enough to find a simple SyMon term of the same type. In this case, the term $`\text{1}_p𝜹_q`$ is imposed. So, these two terms are equal in SyMon, and by the freedom of this category, the initial term is equal to $`\text{1}_A𝜹_B`$ in $`𝒮`$.
Example Prove the equality of the following Cart terms
$$(\text{1}_p\text{c}_{p,qp})\stackrel{}{\text{b}}_{p,p,qp}(\text{1}_{pp}\text{c}_{p,q})(\text{w}_p\text{1}_{pq}),$$
and
$$(\text{1}_p(((𝜹_{qp}\stackrel{}{\text{b}}_{q,p,\text{I}})\text{1}_p)\stackrel{}{\text{b}}_{q,p\text{I},p}))\stackrel{}{\text{b}}_{p,q,(p\text{I})p}((\text{1}_p𝝈_q)((\text{1}_p\text{k}_q)\text{1}_p))((\text{1}_p(\text{k}_p\text{1}_q))\text{c}_{p,pq})\text{w}_{p(pq)}.$$
Their graphs are identical, which can be checked directly from the constructions of these terms, following the linkages between the occurrences of letters in domains and codomains of its primitive components.
Since these terms are of the same type, they are equal in Cart.
We conclude this section with a discussion about the hierarchy (in the sense of embeddability) in the set of free categories mentioned above. All these categories have the same set $`𝒪`$ as the set of objects, and their morphism terms satisfy the following inclusions
where arrows stay for inclusions. The fact that equalities between morphism terms that hold in the “lower” category are also true in the “higher” one was used in this paper several times. To show that a “lower” category is a subcategory of those above, in the diagram, we have to show the following
Lemma 12 Let $`𝒞`$ and $`𝒟`$ be categories from the above diagram such that $`𝒞`$ is below $`𝒟`$. If for $`𝒞`$ morphism terms $`f`$ and $`g`$, the equality $`f=g`$ holds in $`𝒟`$, then they are equal in $`𝒞`$, too.
proof Let $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ be the graphs of $`f`$ and $`g`$ respectively. Since $`f=g`$ in $`𝒟`$, by Lemma 11, the graphs $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are identical. By the coherence of $`𝒞`$ (there is an unique morphism of a certain type in a canonical transformation of the free category $`𝒞`$) it follows that $`f=g`$ in $`𝒞`$.
q.e.d.
So, by the embedding $`E`$ such that $`E(A)=A`$ for every $`A𝒪`$, and $`E([f]_𝒞)=[f]_𝒟`$ where $`[f]_𝒞`$ is the equivalence class of the morphism term $`f`$ in $`𝒞`$, we have the following
Theorem The category Mon is a subcategory of SyMon, Rel, Aff and Cart. The category SyMon is a subcategory of Rel, Aff and Cart. The categories Rel and Aff are subcategories of Cart.
This theorem looks almost trivial, but it is of the same strength as the coherence theorems of substructural categories. An independent proof of this theorem immediately delivers MacLane’s coherence for monoidal and symmetric monoidal categories and the above coherence theorems for relevant and affine categories, from the simplest case of cartesian coherence.
Some other coherence applications in this spirit are given in , and the main result of can be proven more easily by using this apparatus.
Acknowledgements. I would like to thank Professor Kosta Došen for the inspiration and lot of helpful comments on this paper. This work was supported by Grant 0401A of the Science Fund of Serbia.
Zoran Petrić
University of Belgrade
Department of Mining and Geology
Djušina 7
11000 Belgrade, Yugoslavia
e-mail: zpetric@rgf.rgf.bg.ac.yu
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# Self-generated magnetic flux in YBa2Cu3O7-x grain boundaries
## I Introduction
Grain boundaries in high-$`T_c`$ cuprates are interesting and important for both fundamental physics and applications of high-temperature superconductivity. Conventional models of strongly coupled Josephson junctions are applicable to describe electromagnetic properties of grain boundaries in thin films of high-$`T_c`$ superconductors. A remarkable exception of this rule is the -tilt boundary in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-x</sub> films with a misorientation angle close to 45. Indeed, these grain boundaries have an anomalous dependence of the critical current $`I_c`$ on an applied magnetic field $`H_a`$. Contrary to a usual Fraunhofer-type dependence $`I_c(H_a)`$ with a major central peak at $`H_a=0`$ and minor symmetric side-peaks the asymmetric 45 -tilt grain boundaries demonstrate a pattern without a central dominant peak. Instead two symmetric major side-peaks appear at certain applied magnetic fields $`H_a=\pm H_{\mathrm{sp}}0`$.
Several mechanisms have been suggested to explain this phenomena. The anomalous dependence $`I_c(H_a)`$ with symmetric major side-peaks is obviously a result of a specific heterogeneity of electrical properties of the asymmetric 45 -tilt grain boundaries. Two fundamental experimental observations in conjunction explain this heterogeneity in a natural way. First, a fine scale faceting of grain boundaries was discovered in experiments using the transmission electron microscopy (TEM) technique. These facets have a typical length-scale $`l`$ of the order of 10–100 nm and a wide variety of orientations relative to the axis of symmetry of the superconductor. Second, quite a few of recent experiments provide an evidence of a predominant $`d_{x^2y^2}`$ wave symmetry of the order parameter in many of the high-$`T_c`$ cuprates. In some experimental studies the symmetry of the order parameter is more complicated and is shown to be a certain mixture of the $`d_{x^2y^2}`$ and $`s`$ wave components.
These two fundamental experimental observations indicate the existence of two contributions to the phase difference of the order parameter across the grain boundary. Indeed, consider a meandering grain boundary in a film of a superconductor with the $`d_{x^2y^2}`$ wave symmetry of the order parameter and assume that there is a certain magnetic flux inside the grain boundary. In this case there is a phase difference $`\phi `$ caused by the magnetic flux and at the same time there is an additional phase difference $`\alpha `$ caused by the misalignment of the anisotropic $`d_{x^2y^2}`$ wave superconductors. Therefore, the tunneling current density $`j_c`$ is defined by the total phase difference $`\mathrm{\Delta }=\phi \alpha `$. A model describing this Josephson current density $`j(x)`$ results from an assumption that $`j(x)\mathrm{sin}[\phi (x)\alpha (x)]`$, where $`x`$ is along the grain boundary line. The local values of the phase difference $`\alpha (x)`$ depend on the relative orientation of neighboring facets. In the case of an asymmetric 45 grain boundary we have $`\alpha (x)=0`$ or $`\pi `$, and therefore $`j(x)\mathrm{sin}\phi (x)\mathrm{cos}\alpha (x)`$. In other words in the framework of a model relating $`j(x)`$ to the orientation of the facets we arrive to $`j(x)=j_c(x)\mathrm{sin}\phi (x)`$ with an alternating critical current density $`j_c(x)\mathrm{cos}\alpha (x)`$. The dependence $`j_c(x)`$ is imposed by the sequence of facets along a grain boundary line. If this sequence is periodic or almost periodic then the function $`j_c(x)`$ is a periodic or almost periodic alternating function. The typical length-scale for $`j_c(x)`$ is of the same order as the length of the facets $`l`$, i.e., this length-scale is of 10–100 nm.
Variation of orientation of facets along a meandering grain boundary leads to formation of local superconducting current loops even in the absence of an applied magnetic field if the total phase difference $`\mathrm{\Delta }0`$. It was predicted, in particular, that these current loops can generate a certain magnetic flux at a contact of two facets with $`\alpha =0`$ and $`\alpha =\pi `$.
Self-generated randomly distributed magnetic flux was discovered in asymmetric 45 -tilt grain boundaries in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-x</sub> superconducting films in the absence of an applied external magnetic field. This flux $`\varphi _s(x)`$ changes its sign randomly and has an amplitude of variations less than the flux quantum $`\varphi _0`$. The average value of $`\varphi _s(x)`$ along the grain boundary is nearly zero. Noticeable, that this disordered non-quantized magnetic flux was observed only for the samples exhibiting the anomalous dependence of the critical current $`I_c(H_a)`$ on magnetic field with the two symmetric major side-peaks.
It was shown analytically that under certain conditions a stationary state with a self-generated flux exists for a Josephson junction with a periodically alternating critical current density $`j_c(x)`$. The same spatial distributions of $`j_c(x)`$ result in an anomalous dependence of the critical current $`I_c(H_a)`$ on the applied magnetic field. Numerical calculations show that two symmetric major side-peaks appear for a periodically alternating $`j_c(x)`$. The randomness of the critical current density $`j_c(x)`$ smears these peaks but leaves their position in place at weak randomness. We therefore conclude that the experimental observation of the well pronounced major side-peaks on the curve $`I_c(H_a)`$ means that the alternating critical current density is a periodic or almost periodic function of $`x`$. A noticeable randomness of $`j_c(x)`$ would smear out the dominant side-peaks.
In this paper we calculate both analytically and numerically the self-generated flux $`\varphi _s`$ in a Josephson junction with an almost periodically alternating critical current density $`j_c(x)`$. The paper is organized as follows. First, we review briefly the case when $`j_c(x)`$ is a periodic alternating function. We derive the main equations of the two-scale perturbation theory and apply these equations to analyze the non-quantized self-generated flux. This approach forms a basis for the following analytical calculations. Next, we treat the self-generated flux for the case of an almost periodic alternating critical density $`j_c(x)`$ and start with a qualitative approach to the problem. We review then the results of our numerical simulations which verify the qualitative consideration and exhibit a magnetic flux-pattern which is similar to the one that was recently observed experimentally.
## II Main equations
It is convenient for the following analyses to write the function $`j_c(x)`$ in the form
$$j_c=j_c[1+g(x)],$$
(1)
where $`j_c`$ is the average value of the critical current density $`j_c(x)`$ over an interval with a length $`Ll`$
$$j_c=\frac{1}{L}_0^Lj_c(x)𝑑x.$$
(2)
The function $`g(x)`$ introduced in Eq. (1) alternates on a typical length scale of $`l`$. Note that by definition the average value of $`g(x)`$ is zero, i.e., $`g(x)=0`$. The maximum value of $`|g(x)|`$ varies from $`|g(x)|_{\mathrm{max}}1`$ to $`|g(x)|_{\mathrm{max}}1`$. We assume also that $`\lambda l\mathrm{\Lambda }_J`$, where $`\lambda `$ is the London penetration depth and
$$\mathrm{\Lambda }_J^2=\frac{c\varphi _o}{16\pi ^2\lambda j_c}$$
(3)
is an effective Josephson penetration depth. It is worth to mention that in the case of an alternating current density this effective penetration depth is not a local characteristics of a tunnel junction. It is rather an effective quantity defined on the same typical length scale as $`j_c`$.
The phase difference $`\phi (x)`$ satisfies the equation
$$\mathrm{\Lambda }_J^2\phi ^{\prime \prime }[1+g(x)]\mathrm{sin}\phi =0.$$
(4)
In the limiting case $`l\mathrm{\Lambda }_J`$ it is convenient to write a solution of this equation as a sum of a certain smooth function $`\psi (x)`$ with a length scale of order $`\mathrm{\Lambda }_J`$ and a rapidly oscillating function $`\xi (x)`$ with a length scale of order $`ł`$
$$\phi (x)=\psi (x)+\xi (x).$$
(5)
We assume also that $`|\xi (x)||\psi (x)|`$. After substituting Eq. (5) into Eq. (4) and keeping the terms up to the first order in $`\xi (x)`$ we obtain
$$\mathrm{\Lambda }_J^2\psi ^{\prime \prime }+\mathrm{\Lambda }_J^2\xi ^{\prime \prime }[1+g(x)][\mathrm{sin}\psi +\xi \mathrm{cos}\psi ]=0.$$
(6)
Note, that experimentally the self-generated flux was observed by a SQUID pickup loop with a size of several $`\mathrm{\Lambda }_Jl`$. It means that this method is averaging out the fast alternating flux defined by the phase $`\xi (x)`$ and measures the spatially smooth flux defined by $`\psi (x)`$.
Next we consider briefly the case of a periodically alternating critical current density $`j_c(x)`$ which forms the basis of the following analysis of a general case with $`j_c(x)`$ being a randomly alternating function.
## III Periodically alternating critical current density
### A Two-scale perturbation theory
If the critical current density $`j_c(x)`$ is a periodic function, then $`g(x)`$ also is a periodic function. In this case an approximate solution of Eq. (6) can be obtained based on a two-scale perturbation theory. As a first step in order to apply this approach to Eq. (6) we separate the fast alternating terms with a typical length scale $`l`$ and the smooth terms varying with a typical length scale $`\mathrm{\Lambda }_J`$
$$(\mathrm{\Lambda }_J^2\psi ^{\prime \prime }\mathrm{sin}\psi g\xi \mathrm{cos}\psi )+(\mathrm{\Lambda }_J^2\xi ^{\prime \prime }g\mathrm{sin}\psi )=0.$$
(7)
In Eq. (7) we omitted two out of three fast alternating terms of Eq. (6) since they are proportional to $`\xi (x)`$ and therefore are smaller than the term proportional to $`g(x)`$. Next, we note that the terms included into the first pair of brackets in Eq. (7) cancel each other independently on the terms included into the second pair of brackets in Eq. (7) as these two type of terms have very different length scales $`l`$ and $`\mathrm{\Lambda }_J`$ and $`l\mathrm{\Lambda }_J`$. The same reasoning is applicable to the terms included into the second pair of brackets in Eq. (7). As a result we obtain the following two equations for $`\xi (x)`$ and $`\psi (x)`$
$`\mathrm{\Lambda }_J^2\xi ^{\prime \prime }=g(x)\mathrm{sin}\psi ,`$ (8)
$`\mathrm{\Lambda }_J^2\psi ^{\prime \prime }\mathrm{sin}\psi g(x)\xi (x)\mathrm{cos}\psi =0.`$ (9)
It is worth to note that we obtain the two functions $`\psi (x)`$ and $`\xi (x)`$ from one equation (7) as only two type of terms with different typical length scales $`l`$ and $`\mathrm{\Lambda }_J`$ appear in Eq. (6). If $`g(x)`$ would have a wide range of typical length scales the above separation would not be possible.
Introducing the Fourier transform of $`g(x)`$ as
$$g(x)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}g_ke^{ikx}$$
(10)
we find the solution of Eq. (8) in the form
$$\xi (x)=\frac{\mathrm{sin}\psi }{\mathrm{\Lambda }_J^2}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{g_ke^{ikx}}{k^2}=\xi _g(x)\mathrm{sin}\psi ,$$
(11)
where the sums in Eqs. (10) and (11) are over the wave vectors $`k=2\pi n/`$, $``$ is the length of the junction, and $`n`$ is an integer. It is worth mentioning that the function $`\xi _g(x)`$ is defined only by the alternating components of the critical current density $`j_c(x)`$. Also while deriving Eq. (11) we ignored the spatial dependence of $`\mathrm{sin}\psi `$. This can be done since on the length-scale $`l`$ the variation of the smooth function $`\mathrm{sin}\psi (x)`$ is of order $`l/\mathrm{\Lambda }_J1`$. The alternating part of the critical current density has the typical wave numbers $`k1/l`$. Therefore, using Eq. (11) we estimate $`\xi (x)`$ as
$$\xi (x)\mathrm{sin}\psi \frac{l^2}{\mathrm{\Lambda }_J^2}g(x).$$
(12)
It follows from this estimate that the typical values of the phase difference $`\xi (x)`$ are small $`(|\xi (x)|1)`$ if
$$|g(x)|\frac{\mathrm{\Lambda }_J^2}{l^2}.$$
(13)
Next, using Eq. (11) we rewrite Eq. (9) for the smooth phase shift $`\psi (x)`$ in the form
$$\mathrm{\Lambda }_J^2\psi ^{\prime \prime }\mathrm{sin}\psi +\gamma \mathrm{sin}\psi \mathrm{cos}\psi =0,$$
(14)
where
$$\gamma =g(x)\xi _g(x)=\frac{1}{\mathrm{sin}\psi }g(x)\xi (x)$$
(15)
is a constant. A similar derivation for the current density across the tunnel junction results in
$$j(x)=j_c\mathrm{sin}\psi (1\gamma \mathrm{cos}\psi ).$$
(16)
The magnetic field $`B_s(x)`$ generated by the alternating component of the critical current $`j_cg(x)`$ is given by
$$B_s=\frac{4\pi }{c}j_cg(x)𝑑x=\frac{\varphi _0}{4\pi \lambda }\frac{d\xi _g}{dx}$$
(17)
and the averaged field $`B_s(x)=0`$. The alternating magnetic flux $`\varphi _s`$ produced by the field $`B_s`$ is equal to
$$\varphi _s=\frac{\varphi _0}{2\pi }\xi _g.$$
(18)
Combining Eqs. (15) and (17) we find for $`\gamma `$ the formula
$$\gamma =\frac{c\lambda }{\varphi _0j_c}B_s^2=\frac{B_s^2}{B_J^2},$$
(19)
where we introduce a characteristic magnetic field
$$B_J=\frac{4\pi }{c}j_c\mathrm{\Lambda }_J.$$
(20)
It follows from equation (19) that $`\gamma `$ is a positive constant which can be estimated as
$$\gamma \frac{l^2}{\mathrm{\Lambda }_J^2}g^2.$$
(21)
The energy of a Josephson junction $``$ takes the form $`=_0+_\phi `$, where $`_0`$ is independent on $`\phi (x)`$ and
$$_\phi =\frac{j_c}{2e}𝑑x\left\{\frac{1}{2}\mathrm{\Lambda }_J^2\phi ^2[1+g(x)]\mathrm{cos}\phi \right\}.$$
(22)
Using Eqs. (8), (11), and the definition of $`\gamma `$, we obtain the energy $`_\phi `$ in terms of the smooth phase shift $`\psi (x)`$
$$_\phi =\frac{\mathrm{}j_c}{4e}𝑑x\left\{\mathrm{\Lambda }_J^2\psi ^22\mathrm{cos}\psi \gamma \mathrm{sin}^2\psi \right\}.$$
(23)
Note, that solutions $`\psi (x)`$ of Eq. (14) correspond to the minima and to the maxima of the energy functional (23).
### B Non-quantized self-generated flux
Let us apply Eqs. (8), (9), and (15) to a consider the stationary states of a Josephson junction with a certain length $`\mathrm{\Lambda }_J`$ in an absence of applied magnetic field. In this case the average flux inside the junction is zero and thus an alternating self-generated flux $`\varphi _s(x)`$ appears simultaneously with a certain phase $`\psi =\mathrm{const}`$ as it follows from Eqs. (11) and (18).
In the stationary state with $`\psi =\mathrm{const}`$ the values of $`\psi `$ are determined by Eq. (14) which takes the form
$$\mathrm{sin}\psi (1\gamma \mathrm{cos}\psi )=0.$$
(24)
Note, that this equation means also that the current density $`j(x)`$ across the tunnel junction is equal to zero.
In the case $`\gamma <1`$ equation (24) has two solutions, namely, $`\psi =0`$ and $`\psi =\pi `$ and thus, as follows from Eq. (11), there is no self-generated flux. It is also worth mentioning that the energy of a Josephson junction $``$ has a minimum for $`\psi =0`$ and maximum for $`\psi =\pi `$.
In the case $`\gamma 1`$ there are four solutions of equation (24), namely, $`\psi =\psi _\gamma ,\mathrm{\hspace{0.17em}0},\psi _\gamma ,\pi `$, where
$$\psi _\gamma =\mathrm{arccos}(1/\gamma ).$$
(25)
The energy $`[\psi (x)]`$ has a minimum for $`\psi =\pm \psi _\gamma `$ and a maximum for $`\psi =0,\pi `$. The self-generated flux
$$\varphi _s=\frac{\varphi _0\xi }{2\pi }=\frac{\varphi _0\xi _g}{2\pi }\mathrm{sin}\psi _\gamma =\varphi _0\frac{\xi _g}{2\pi }\frac{\sqrt{\gamma ^21}}{\gamma }$$
(26)
arises in the two stationary states with $`\psi =\pm \psi _\gamma `$, each of these states corresponds to a minimum energy $``$. The assumption $`\xi (x)1`$ restricts the value of $`\gamma `$. However, it follows from Eqs. (13) and (21) that $`\xi (x)1`$ and $`\gamma >1`$ hold simultaneously only if
$$\frac{\mathrm{\Lambda }_J}{l}<|g(x)|\frac{\mathrm{\Lambda }_J^2}{l^2}.$$
(27)
Using equations (12) and (26) we estimate $`|\varphi _s(x)|`$ as
$$|\varphi _s(x)|\varphi _0\frac{\sqrt{\gamma ^21}}{\gamma }\frac{l^2}{\mathrm{\Lambda }_J^2}|g(x)|\varphi _0.$$
(28)
The above results hold only for a periodic critical current density $`j_c(x)`$ and the predicted self-generated flux $`\varphi _s(x)`$ has a typical amplitude of variations which is much less than the one observed experimentally.
## IV Non-periodic alternating critical current density
The above analytical approach to the problem of a self-generated flux in a non-uniform Josephson junction is based on an assumption that the critical current density $`j_c(x)`$ is an alternating periodic function. This model allows for analytical calculation and provides a reasonable preliminary insight into the properties of an idealized Josephson junction with an alternating $`j_c(x)`$. At the same time this simple model fails for a quantitative description of any real system with a certain randomness of the spatial distribution of Josephson critical current density $`j_c(x)`$. In this section we generalize the above approach assuming that the alternating critical current density $`j_c(x)`$ is almost periodic, i.e., we assume that there is a typical length of interchange of sign of $`j_c(x)`$ which is distributed randomly near some mean value $`l`$.
In the case of an almost periodic $`j_c(x)`$ we can not apply the two-scale perturbation theory in the same way as we did it in the previous section. Indeed, an arbitrary solution of Eq. (8) takes the form
$$\xi (x)=\frac{\mathrm{sin}\psi }{\mathrm{\Lambda }_J^2}G(x),$$
(29)
where
$$G(x)=_a^x𝑑x^{}_a^{}^x^{}𝑑x^{\prime \prime }g(x^{\prime \prime }),$$
(30)
and the integration constants $`a`$ and $`a^{}`$ are defined by the boundary conditions. The random function $`G(x)`$ increases with an increase of the integration interval. In general, the value of $`|G(x)|`$ can become arbitrarily large if the length of the tunnel junction $``$ becomes large enough. This is in a contradiction with our main assumption that the phase $`\xi (x)`$ is a small and fast varying component of the total phase difference $`\phi (x)`$. To solve this contradiction we write $`\xi (x)`$ as
$$\xi (x)=\frac{\mathrm{sin}\psi }{\mathrm{\Lambda }_J^2}\left[G(x)G_a(x)\right].$$
(31)
The function $`G_a(x)`$ is a smoothing average of $`G(x)`$ over an interval with a certain length $`_a`$, where $`l<_a`$.
The procedure of filtering out the smooth part of $`G(x)`$ is especially evident if we use Fourier series for $`g(x)`$ and $`G(x)`$. Introducing Fourier transform for $`g(x)`$ as
$$g(x)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}g_ke^{ikx}$$
(32)
we find Fourier series for the function $`G(x)`$ in the form
$$G(x)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{g_k}{k^2}e^{ikx},$$
(33)
where the sums in Eqs. (32) and (33) are over the wave vectors $`k=2\pi n/`$ and $`n`$ is an integer. The smooth part of the function $`G(x)`$ can be obtained by extracting the fast Fourier harmonics, i.e., by extracting from the sum (33) terms with wave vectors $`|k|>k_a2\pi /_a`$. As a result we find for $`G_a(x)`$ and $`\xi (x)`$ the series
$`G_a(x)`$ $`=`$ $`{\displaystyle \underset{k_a}{\overset{k_a}{}}}{\displaystyle \frac{g_k}{k^2}}e^{ikx},`$ (34)
$`\xi (x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}\psi }{\mathrm{\Lambda }_J^2}}\left({\displaystyle \underset{\mathrm{}}{\overset{k_a}{}}}+{\displaystyle \underset{k_a}{\overset{\mathrm{}}{}}}\right)\left[{\displaystyle \frac{g_k}{k^2}}e^{ikx}\right].`$ (35)
The small and fast alternating part $`\xi (x)`$ of the phase difference $`\phi (x)`$ is thus defined by Eq. (31). This equation is a straightforward generalization of the two-scale perturbation theory approach to a real case of an almost periodic $`g(x)`$.
Next, we use Eq. (31) to derive an equation for the smooth part $`\psi (x)`$ of the phase $`\phi (x)`$. First, combining Eqs. (31), (7), and (30) we arrive to the relation
$$\mathrm{\Lambda }_J^2\psi ^{\prime \prime }[1+G_a^{\prime \prime }(x)]\mathrm{sin}\psi g(x)\xi (x)\mathrm{sin}\psi =0.$$
(36)
Second, we average Eq. (36) over an interval with a certain length $`L`$ assuming that $`lL_a`$. This averaging results in an equation describing the phase $`\psi (x)`$
$$\mathrm{\Lambda }_J^2\psi ^{\prime \prime }[1+G_a^{\prime \prime }(x)]\mathrm{sin}\psi +\gamma (x)\mathrm{sin}\psi \mathrm{cos}\psi =0,$$
(37)
where $`\gamma (x)`$ is defined by Eq. (15). It is worth mentioning that equation (37) differs from the analogues equation (9) by an additional term $`G_a^{\prime \prime }(x)`$ in the coefficient before $`\mathrm{sin}\psi `$ and by the fact that the parameter $`\gamma =\gamma (x)`$ is a function of the coordinate $`x`$ along the junction.
The coefficient $`1+G_a^{\prime \prime }(x)`$ is defined by magnetic field $`B_c(x)`$ which would be generated by a current with the density $`j_c(x)`$. Indeed, using Eqs. (1), (30), and Maxwell’s equation $`dB_c/dx=4\pi j_c/c`$ we obtain for $`G(x)`$
$`G(x)`$ $`=`$ $`{\displaystyle \frac{c}{4\pi j_c}}{\displaystyle 𝑑x^{}[B_c(x^{})B_a(x^{})]}=`$ (38)
$`=`$ $`{\displaystyle \frac{2\pi \mathrm{\Lambda }_J^2}{\varphi _0}}[\varphi _c(x)\varphi _a(x)],`$ (39)
where the magnetic field $`B_a(x)`$ would be generated by a constant current density $`j_c`$, i.e., $`dB_a/dx=4\pi j_c/c`$, and $`\varphi _c(x)`$ and $`\varphi _a(x)`$ are the fluxes of the fields $`B_c(x)`$ and $`B_a(x)`$. It follows now from Eq. (39) that
$$1+G_a^{\prime \prime }(x)=\frac{2\pi \mathrm{\Lambda }_J^2}{\varphi _0}\varphi _c(x)^{\prime \prime }.$$
(40)
As shown above the value of the parameter $`\gamma `$ is determining the existence or absence of the self-generated flux. Using Eq. (8) and the definition of $`\gamma (x)`$ given by Eq. (15) we obtain for $`\gamma (x)`$ an expression
$$\gamma (x)=\frac{\mathrm{\Lambda }_J^2}{\mathrm{sin}^2\psi }\xi ^{\prime \prime }\xi =\frac{\mathrm{\Lambda }_J^2}{\mathrm{sin}^2\psi }\xi _{}^{}{}_{}{}^{2}>0$$
(41)
demonstrating that the condition $`\gamma (x)>0`$ holds also in the case of an almost periodic critical current density.
### A Self-generated flux in a stationary state
Let us now consider stationary solutions for the smooth part $`\psi (x)`$ of the phase difference $`\phi (x)`$ qualitatively. Assume first that there are sufficiently large intervals with lengths $`L_i\mathrm{\Lambda }_J`$, where the function $`\psi (x)`$ is constant or varies with a typical space-scale of order $`_i\mathrm{\Lambda }_J`$. In this case equation (37) reduces to
$$\left[1+G_a^{\prime \prime }(x)\gamma (x)\mathrm{cos}\psi \right]\mathrm{sin}\psi =0.$$
(42)
This equation is similar to Eq. (24) which we derived for the case of a periodic critical current density $`j_c(x)`$ and has different solutions depending on the value of the parameter
$$\gamma _r(x)=\frac{\gamma (x)}{1+G_a^{\prime \prime }(x)}.$$
(43)
In the regions with $`\gamma _r(x)<1`$ equation (42) has two solutions $`\psi =0`$ and $`\psi =\pi `$ and therefore as it follows from Eq. (29) there is no self-generated flux in these regions.
The energy of a Josephson junction $`_\phi `$ given by Eq. (22) can be written in terms of the smooth part of the phase difference $`\psi (x)`$. In the case of an almost periodic critical current density $`j_c(x)`$ this equation reads
$`_\phi ={\displaystyle \frac{\mathrm{}j_c}{4e}}\times `$ (44)
$`{\displaystyle 𝑑x\left\{\mathrm{\Lambda }_J^2\psi ^22\mathrm{cos}\psi [1+G_a^{\prime \prime }(x)]\gamma (x)\mathrm{sin}^2\psi \right\}}`$ (45)
and if $`\gamma _r(x)<1`$, then the energy $`_\phi `$ has a minimum for $`\psi =0`$ and a maximum for $`\psi =\pi `$.
In the regions where the function $`\gamma _r(x)>1`$ equation (42) has four solutions $`\psi =\psi _\gamma (x),\mathrm{\hspace{0.17em}0},\psi _\gamma (x),\pi `$, with
$$\psi _\gamma (x)=\mathrm{arccos}\left[\frac{1}{\gamma _r(x)}\right]=\mathrm{arccos}\left[\frac{1+G_a^{\prime \prime }(x)}{\gamma (x)}\right].$$
(46)
The energy $`_\phi `$ has a minimum for $`\psi =\pm \psi _\gamma (x)`$ and a maximum for $`\psi =0,\pi `$. The self-generated flux is thus non-zero in the regions with $`\gamma _r>1`$. This flux has a fast and a smoothly varying parts defined by $`\xi (x)`$ and $`\psi _\gamma (x)`$.
The randomness of the function $`g(x)`$ causes variation of $`\gamma _r(x)`$ along the junction. As a result of this variation intervals with $`\gamma _r(x)>1`$ are interlaced with intervals with $`\gamma _r(x)1`$. As it was mentioned above, in the case of $`\gamma _r>1`$ the energy of the Josephson junction $`_\phi `$ has a minimum for $`\psi =\pm \psi _\gamma (x)`$ and a maximum for $`\psi =0,\pi `$. When the value of $`\gamma _r(x)`$ changes from $`\gamma _r(x)>1`$ to $`\gamma _r(x)1`$ the energy $`_\phi `$ still has a maximum if $`\psi =\pi `$, but a state with $`\psi =0`$ becomes a state with a minimum energy.
The above results provide a qualitative description of experimentally observable flux distribution along a Josephson junction with an almost periodic alternating critical current density. This flux distribution spatially averaged by the measurement tools is defined by the function $`\psi (x)`$ (see Fig. 1). Inside the intervals with $`\gamma _r(x)>1`$ the phase $`\psi (x)`$ tends to one of the solutions $`\pm \psi _\gamma (x)`$. The profile of the function $`\psi (x)`$ correlates with the profile of $`\psi _\gamma (x)`$, though does not coincide with it exactly because the solution $`\psi (x)=\psi _\gamma (x)`$ was obtained under assumption $`\psi ^{\prime \prime }=0`$, which does not hold exactly for the intervals $`\gamma _r(x)>1`$. The smooth part of the phase difference inside the intervals with $`\gamma _r(x)1`$ is $`\psi =0`$ which is consistent with the assumption $`\psi ^{\prime \prime }=0`$.
The value of $`\psi _\gamma `$ increases quite fast with an increase of the parameter $`\gamma `$ (see Fig. 2). In particular, for $`\gamma =2`$ the value of $`\psi _\gamma `$ is already about 0.75 of its maximum value $`\pi /2`$. This means that for most of the experimentally observable peaks of the self-generated flux the values of $`\psi `$ will be close to $`\pi /2`$ which corresponds to a magnetic flux $`\varphi _0/4`$. In some places of the junction the phase $`\psi `$ changes from $`\psi _\gamma `$ to $`\psi _\gamma `$. The flux localized in this area of the junction will be close to $`\varphi _0/2`$.
## V Numerical simulations
### A The finite difference scheme
To study the self-generation of magnetic flux in a tunnel junction with an alternating critical current density numerically we introduced time dependence into the main equation (4)
$$\ddot{\phi }+\alpha \dot{\phi }\mathrm{\Lambda }_J^2\psi ^{\prime \prime }+[1+g(x)]\mathrm{sin}\phi =0,$$
(47)
where $`\alpha 1`$ is a decay constant. This approach allows to study both dynamics and statics of the system.
The term $`\alpha \dot{\phi }`$ introduces dissipation. As a result of this dissipation the relaxation of the system ends up in one of the stable stationary states described by a certain solution of the static equation (4). Moreover, for a given distribution of the critical current density $`j_c(x)`$ we obtained different static solutions when we start the numeric simulation from different initial states. We compare and classify these solutions based on the features of the function $`j_c(x)`$. Indeed, this function essentially describes the pinning properties of the junction. Therefore a variety of initial conditions can converge to a similar flux pinning pattern.
To solve Eq. (47) numerically we use the finite difference scheme. We adopted this method to our case and checked stability and convergency of the obtained solutions. As a result we arrived to the following scheme
$`\phi `$ $`{\displaystyle \frac{\phi _{m+1}^n+\phi _{m1}^n}{2}}\stackrel{~}{\phi }_m^n`$ (48)
$`\dot{\phi }`$ $`{\displaystyle \frac{\stackrel{~}{\phi }_m^n\phi _m^{n1}}{\tau }}`$ (49)
$`\ddot{\phi }`$ $`{\displaystyle \frac{\phi _m^{n+1}+\phi _m^{n1}2\stackrel{~}{\phi }_m^n}{\tau ^2}}`$ (50)
$`\phi ^{\prime \prime }`$ $`{\displaystyle \frac{\phi _{m+1}^n+\phi _{m1}^n2\phi _m^n}{h^2}},`$ (51)
where $`f_m^n=f(x_m,t_n)`$, $`\tau `$ and $`h`$ are steps along $`t`$ and $`x`$ correspondingly. Next, we choose units providing $`\mathrm{\Lambda }_J=1`$ and set $`h`$=$`\tau `$. As a result we arrive to the following finite difference scheme
$`\phi _m^{n+1}=`$ $``$ $`(1\alpha \tau )\phi _m^{n1}+(2\alpha \tau )\stackrel{~}{\phi }_m^n`$ (52)
$``$ $`\tau ^2(1g_m)\mathrm{sin}\stackrel{~}{\phi }_m^n.`$ (53)
### B Stationary solutions
Initially a certain random function $`g(x)`$ is generated for an interval with a length $`L`$ with a given values of $`l`$ and $`\delta l`$ (a typical length-scale of the function $`g(x)`$ and its dispersion), $`g`$ and $`\delta g`$ (amplitude of the function $`g(x)`$ and its dispersion). This allows to calculate the function $`\gamma _r(x)`$ for the whole interval. An initial state $`\phi _0(x)`$ is prepared as a random or some specific function. Finally the dynamical rules (53) are applied to the initial state iteratively until a stationary state is established.
In Fig. 3 we show one of the stationary solutions obtained by a numerical simulation and the function $`\gamma _r(x)`$ calculated for this solution. It is clearly seen from Fig. 3 that $`\phi (x)`$ arises at the places where $`\gamma _r(x)`$ exceeds $`1`$. Heights of the peaks are less than $`\pi /2`$, and thus the corresponding magnetic flux amplitudes are less than $`\varphi _0/4`$.
In general a different initial state of the same sample, i.e., for the same function $`\gamma _r(x)`$, generates a different stationary state. Our numerical simulations show that these different states differ only by sign of some peaks of $`\phi (x)`$, but the shapes and locations stay unchanged.
We have compared our results with the experimental data. The typical amplitude of the flux variations measured by a SQUID pickup loop with a size of several $`\mathrm{\Lambda }_J`$ is about $`0.25`$ of $`\varphi _0`$ with rare narrow picks with an amplitude about $`0.5\varphi _0`$ which is in a good agreement with our calculations.
## VI SUMMARY
We treated a Josephson junction with an alternating critical current density $`j_c(x)`$ as a model for considering electromagnetic properties of grain boundaries in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-x</sub> superconducting films. The study is mainly focused on a specific case of an almost periodically alternating function $`j_c(x)`$. We demonstrated both analytically and numerically that under certain conditions a self-generated flux pattern arise for this type of spatial distribution of the critical current density $`j_c(x)`$. The obtained flux pattern with two types of interlacing flux domains is similar to the one which was recently observed experimentally in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-x</sub> superconducting films in the absence of an external magnetic field. The typical amplitude for the magnetic flux peaks is of order $`\varphi _0/4`$ and $`\varphi _0/2`$ and the typical distance between the peaks depends on the spatial distribution of $`j_c(x)`$.
###### Acknowledgements.
One of us (RGM) is grateful to E.H. Brandt, J.R. Clem, H. Hilgenkamp, V.G. Kogan, and J. Mannhart for useful and stimulating discussions. This research was supported in part by grant No. 96-00048 from the United States – Israel Binational Science Foundation (BSF), Jerusalem, Israel.
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# 1 Motivations, and Content Versus Ref.[3]
## 1 Motivations, and Content Versus Ref.
In physics at the highest available energies, it is always important to exploit simple reactions and decays so as to search for new forces, for new dynamics, and for discrete symmetry violations. Because the t-quark weakly decays before hadronization effects are significant, and because of the large t-quark mass, t-quark decay can be an extremely useful tool for such fundamental searches. Initial tests of the Lorentz structure and of symmetry properties of $`tW^+b`$ decay will be carried out at the Tevatron, but the more precise measurements will be possible at the CERN LHC and at a NLC .
It is important to be able to quantitatively assay future measurements of competing observables consistent with the standard model (SM) prediction of only a $`g_{VA}`$ coupling and only its associated discrete symmetry violations. For this purpose, without consideration of possible explicit $`T`$ violation, in Ref. plots were given of the values of the helicity parameters in terms of a “$`(VA)`$ $`+`$ Additional Lorentz Structure” versus effective-mass scales for new physics, $`\mathrm{\Lambda }_i`$, associated with each additional Lorentz structure.
In this contributed paper, to assay future measurements of helicity parameters in regard to $`T`$ violation, the effects of possible explicit $`T`$ violation are briefly reported. A more detailed paper on this latter subject will soon be available. In effective field theory, $`\mathrm{\Lambda }_i`$, is the scale at which new particle thresholds or new dynamics are expected to occur; $`\mathrm{\Lambda }_i`$ can also be interpreted as a measure of a top quark compositeness/condensate scale. In measurement of some of the helicity parameters, the LHC should be sensitive to $`3`$ % and the Tevatron in a Run 3 to perhaps the $`10`$ % level (“ideal statistical error levels”) .
## 2 Consequences of Single Additional Lorentz Structures in Absence of Explicit $`T`$ Violation
In this section, we briefly review the work reported in Ref. . This published paper contains a more detailed discussion, useful simple formulas relating the “ $`\alpha ,\beta ,\gamma `$ ” relative phases of Fig.1a and the helicity parameters of Fig.1b, and plots of the values of the associated helicity parameters in the case of single additional Lorentz structures.
The attached Figs. 1a, 1b provide a good orientation to this topic: a complete measurement of on-shell properties of the $`tW^+b`$ decay mode will have been accomplished when the 4 moduli are determined and any 3 of the relative phases of the helicity amplitudes $`A(\lambda _{W^+},\lambda _b)`$. The helicity parameters appear directly in various polarization and spin-correlation functions such as those obtained in Ref..
The top lines of the first two tables list the standard model(SM)’s numerical values for the quantities shown in Figs. 1a, 1b. In the SM, all the relative phases are either zero or $`\pm \pi `$ so the primed helicity parameters are zero. In Table 1 in the top line are the standard model expectations for the numerical values of the helicity amplitudes $`A(\lambda _{W^+},\lambda _b)`$ for $`tW^+b`$ decay in $`g_L=1`$ units. The input values are $`m_t=175GeV,m_W=80.35GeV,m_b=4.5GeV`$. The $`\lambda _b=1/2`$ b-quark helicity amplitudes would vanish if $`m_b`$ were zero. For this reason, if one is guided by the SM expectations, the most accessible quantities experimentally should be the two moduli and the relative phase shown on the right of Fig. 1a. If the SM is correct, one expects that the $`A(0,1/2)`$ and $`A(1,1/2)`$ moduli and relative phase $`\beta _L`$ will be the first quantities to be determined . The $`\lambda _b=1/2`$ moduli are factors of 30 and 100 smaller in the SM. Interference measurements between the two columns are of order $`𝒪(LR)`$. $`L`$ and $`R`$ denote the $`b`$ quark’s helicity $`\lambda _b=1/2`$. Throughout this moduli-phase analysis of top decays, intrinsic and relative signs of the helicity amplitudes are specified in accordance with the standard Jacob-Wick phase convention.
In Table 2 in the top line are the SM’s numerical values of the associated helicity parameters. Explicit formulas for the standard model helicity amplitudes and for experimental distributions in terms of these helicity parameters are given in Ref..
The layout of the corners in Fig. 1 has been chosen to reflect the layout in the probability plots for $`P(W_L)`$ versus $`P(b_L)`$, see Ref. and Figs.5-6 below. The quantities
$$\begin{array}{c}P(W_L)=\text{Probability}W^+\text{ is longitudinally polarized,}\lambda _{W^+}=0\\ P(b_L)=\text{Probability}b\text{ is left-handed,}\lambda _b=1/2\end{array}$$
In terms of the first two helicity parameters of Table 2, $`P(W_L)=\frac{1+\sigma }{2}=0.705(SM)`$ and $`P(b_L)=\frac{1+\xi }{2}=1.00(SM)`$. So in the standard model, the emitted $`W`$ boson should be $`70\%`$ longitudinally polarized and the emitted b-quark should be almost completely left-handed polarized.
The “arrows” in the upper part of Fig. 1 define the measurable $`\alpha ,\beta ,\gamma `$ relative phases between the four amplitudes. For instance,
$$\alpha _0=\varphi _0^R\varphi _0^L,\beta _L=\varphi _1^L\varphi _0^L,\gamma _+=\varphi _1^R\varphi _0^L$$
(1)
where $`A(\lambda _{W^+},\lambda _b)=|A|\mathrm{exp}(i\varphi _{\lambda _{W^+}}^{L,R})`$. So for a pure $`VA`$ coupling, the $`\beta `$’s vanish and all the $`\alpha `$’s and $`\gamma `$’s equal $`+\pi `$ (or $`\pi `$) to give the intrinsic minus sign of the standard model’s $`b_R`$ amplitudes, see top row of Table 1.
The lower part of Fig. 1 displays the real part and imaginary part (primed) helicity parameters corresponding to interference measurements of the respective relative phases. For instance, c.f. Appendix B of Ref.,
$$\begin{array}{c}\eta _L\frac{1}{\mathrm{\Gamma }}|A(1,\frac{1}{2})||A(0,\frac{1}{2})|\mathrm{cos}\beta _L\\ \eta _L^{}\frac{1}{\mathrm{\Gamma }}|A(1,\frac{1}{2})||A(0,\frac{1}{2})|\mathrm{sin}\beta _L\end{array}$$
(2)
and
$$\eta _{L,R}=\frac{1}{2}(\eta \pm \omega )$$
(3)
Because of the relative magnitudes of the moduli predicted by the SM, in our consideration of information from b-quark polarization measurements, we concentrate on the two b-quark interference parameters $`\kappa _0`$ and $`ϵ_+`$ and on their primed analogues. If surprises are discovered in top quark decay, other phases and/or helicity parameters might be more useful and certainly would be useful as checks and/or constraints. By $`\mathrm{\Lambda }_b`$ polarimetry, or some other $`b`$-polarimetry technique, it would be important to measure the $`\alpha `$ and $`\gamma `$ relative phase. In the standard model, the two helicity parameters between the amplitudes with the largest moduli are
$$\begin{array}{c}\kappa _0\frac{1}{\mathrm{\Gamma }}|A(0,\frac{1}{2})||A(0,\frac{1}{2})|\mathrm{cos}\alpha _0\\ ϵ_+\frac{1}{\mathrm{\Gamma }}|A(1,\frac{1}{2})||A(0,\frac{1}{2})|\mathrm{cos}\gamma _+\end{array}$$
(4)
We refer to $`\kappa _0,ϵ_+`$ as the “$`b`$-polarimetry interference parameters”. For $`\kappa _{0}^{}{}_{}{}^{^{}},ϵ_{+}^{}{}_{}{}^{^{}}`$, the sine function replaces the cosine function in Eqs.(4). Unfortunately from the perspective of a complete measurement of the four helicity amplitudes, the tree-level values of $`\kappa _0,ϵ_+`$ in the SM are only about $`1\%`$. See the top line in both parts of Table 2. Two dimensional plots of the type $`(ϵ_+,\eta _L)`$ and $`(\kappa _0,\eta _L)`$, and of their primed counterparts, have the useful property that the unitarity limit is a circle of radius $`\frac{1}{2}`$ centered on the origin.
In the plots in Ref. and below, the values of the helicity parameters are given in terms of a
$`(VA)`$ \+ Single Additional Lorentz Structure”. Generically, in the case of no explicit $`T`$ violation, we denote these additional couplings by
$$g_{Total}g_L+g_X$$
(5)
$$X=\{\begin{array}{cc}X_c=\text{chiral}=\{V+A,S\pm P,f_M\pm f_E\}\hfill & \\ X_{nc}=\text{non-chiral}=\{V,A,S,P,f_M,f_E\}.\hfill & \end{array}$$
For $`tW^+b`$, the most general Lorentz coupling is $`W_\mu ^{}J_{\overline{b}t}^\mu =W_\mu ^{}\overline{u}_b\left(p\right)\mathrm{\Gamma }^\mu u_t\left(k\right)`$ where $`k_t=q_W+p_b`$, and
$`\mathrm{\Gamma }_V^\mu =g_V\gamma ^\mu +{\displaystyle \frac{f_M}{2\mathrm{\Lambda }}}\iota \sigma ^{\mu \nu }(kp)_\nu +{\displaystyle \frac{g_S^{}}{2\mathrm{\Lambda }}}(kp)^\mu `$
$`+{\displaystyle \frac{g_S}{2\mathrm{\Lambda }}}(k+p)^\mu +{\displaystyle \frac{g_{T^+}}{2\mathrm{\Lambda }}}\iota \sigma ^{\mu \nu }(k+p)_\nu `$ (6)
$`\mathrm{\Gamma }_A^\mu =g_A\gamma ^\mu \gamma _5+{\displaystyle \frac{f_E}{2\mathrm{\Lambda }}}\iota \sigma ^{\mu \nu }(kp)_\nu \gamma _5+{\displaystyle \frac{g_P^{}}{2\mathrm{\Lambda }}}(kp)^\mu \gamma _5`$
$`+{\displaystyle \frac{g_P}{2\mathrm{\Lambda }}}(k+p)^\mu \gamma _5+{\displaystyle \frac{g_{T_5^+}}{2\mathrm{\Lambda }}}\iota \sigma ^{\mu \nu }(k+p)_\nu \gamma _5`$ (7)
For $`g_L=1`$ units with $`g_i=1`$, the nominal size of $`\mathrm{\Lambda }_i`$ is $`\frac{m_t}{2}=88GeV`$, see below.
Lorentz equivalence theorems for these couplings are treated in Appendix A of Ref.. Explicit expressions for the $`A(\lambda _{W^+},\lambda _b)`$ in the case of these additional Lorentz structures are given in Ref. . Other recent general analyses of effects in $`tW^+b`$ decay associated with new physics arising from large effective- mass scales $`\mathrm{\Lambda }_i`$ are in Refs. \[6-12\]. Some work on higher order QCD and EW corrections has been done in .
The partial width $`\mathrm{\Gamma }`$ for $`tW^+b`$ is the remaining and very important moduli parameter for testing for additional Lorentz structures. Since $`\mathrm{\Gamma }`$ sets the overall scale, it cannot be well measured by spin-correlation techniques, which better measure the ratios of moduli and relative phases, so we consider $`\mathrm{\Gamma }`$ separately; see also . From the perspective of possible additional Lorentz structures, measurement of the partial width $`\mathrm{\Gamma }`$ is an important constraint. In particular, this provides a strong constraint on possible $`V+A`$ couplings in contrast to measurement of $`P(W_L)`$ which does not. $`\mathrm{\Gamma }`$ provides a useful constraint for the possibility of additional $`V`$ and $`A`$ couplings which are appealing from the perspective of additional gauge-theoretic structures.
## 3 Moduli Parameters and Phase-Type Ambiguities
Versus predictions based on the SM, two dynamical phase-type ambiguities were found by investigation of the effects of a single additional “chiral” coupling $`g_i`$ on the three moduli parameters $`\sigma =P(W_L)P(W_T),\xi =P(b_L)P(b_R),`$and $`\zeta =\frac{1}{\mathrm{\Gamma }}(\mathrm{\Gamma }_L^{b_Lb_R}\mathrm{\Gamma }_T^{b_Lb_R})`$.
For an additional $`S+P`$ coupling with $`\mathrm{\Lambda }_{S+P}34.5GeV`$ the values of $`(\sigma ,\xi ,\zeta )`$ and also of the partial width $`\mathrm{\Gamma }`$ are about the same as the SM prediction, see Table 2. This is the first dynamical ambiguity. Table 1 shows that this ambiguity will also occur if the sign of the $`A_X(0,\frac{1}{2})`$ amplitude for $`g_L+g_X`$ is taken to be opposite to that of the SM’s amplitude. An additional $`S\pm P`$ only effects the longitudinal $`W^\pm `$ amplitudes and not the transverse $`\lambda _W=1`$ ones. By requiring that
$$\frac{A_X(0,\frac{1}{2})}{A_X(1,\frac{1}{2})}=\frac{A_L(0,\frac{1}{2})}{A_L(1,\frac{1}{2})}$$
(8)
for $`X=S+P`$, we obtain a simple formula
$$\mathrm{\Lambda }_{S+P}=(\frac{g_{S+P}}{g_L})\frac{m_tq_W}{2(E_W+q_W)}(\frac{g_{S+P}}{g_L})\frac{m_t}{4}(1(\frac{m_W}{m_t})^2).$$
(9)
It is important to regard these ambiguities from (i) the signs in their $`b_L`$ amplitudes versus those for the SM and from (ii) the tensorial character and $`\mathrm{\Lambda }`$ value of the associated Lorentz structure.
For an additional $`f_M+f_E`$ coupling with $`\mathrm{\Lambda }_{f_M+f_E}53GeV`$ the values of $`(\sigma ,\xi ,\zeta )`$ are also about the same as the SM prediction, see Table 2. This is the second dynamical ambiguity. In this case, the partial width $`\mathrm{\Gamma }`$ is about half that of the SM due to destructive interference. Table 1 shows that this ambiguity will also occur if the sign of the $`A_X(1,\frac{1}{2})`$ amplitude for $`g_L+g_X`$ is taken to be opposite to that of the SM’s amplitude. Again, from (8) for $`X=f_M+f_E`$, we obtain
$$\mathrm{\Lambda }_{f_M+f_E}=(\frac{g_{f_M+f_E}}{g_L})\frac{m_tE_W}{2(E_W+q_W)}(\frac{g_{f_M+f_E}}{g_L})\frac{m_t}{4}(1+(\frac{m_W}{m_t})^2)$$
(10)
since $`\frac{m_b}{m_t}\frac{\sqrt{E_bq_W}}{\sqrt{E_b+q_W}}10^3`$.
Besides the $`f_M+f_E`$ construction of this second phase- type ambiguity, it should be kept in mind that some other mechanism might produce the relative sign change shown in Table 1, but without also changing the absolute value of the $`b_L`$ amplitudes. In this case the measurement of the partial width $`\mathrm{\Gamma }`$ would not resolve this phase ambiguity.
From consideration of Table 1, a third (phase) ambiguity can be constructed by making an arbitrary sign-flip in the $`b_L`$ amplitudes, so $`A_X(\lambda _{W,}\lambda _b=\frac{1}{2})=A_{VA}(\lambda _{W,}\lambda _b=\frac{1}{2})`$, with no corresponding sign changes in the $`b_R`$ amplitudes.
Resolution of this third ambiguity, as well as determination of two remaining independent relative phases ( e.g. $`\alpha _0`$ and $`\gamma _+`$ ) necessary for a complete amplitude measurement of $`tW^+b`$ decay, will require direct empirical information about the $`b_R`$-amplitudes. One way would be from a $`\mathrm{\Lambda }_b`$ polarimetry measurement of the $`b`$-polarimetry interference parameters $`ϵ_+`$ and $`\kappa _0`$. Even at an NLC, such measurements will be difficult unless certain non-SM couplings occur. In particular, here additional $`S+P`$ and $`f_M+f_E`$ couplings have negligible effects, but non-chiral couplings like $`V`$ or $`A`$, $`f_M`$ or $`f_E`$ (for $`ϵ_+`$), $`S`$ or $`P`$ (for $`\kappa _0`$) can produce large effects.
Since the helicity parameters appear directly in the various polarization and spin-correlation functions, it is clearly more model independent to simply measure them rather than to set limits on an “ ad hoc” set of additional coupling constants. The large $`m_b`$ effects displayed in some of the plots in Ref. explicitly demonstrate this point. In many cases, finite $`m_b`$ effects in both $`b_L`$ and $`b_R`$ amplitudes lead to sizable “ oval shapes” as the effective mass scale $`\mathrm{\Lambda }_i`$ varies. There do not exist “Lorentz equivalence theorems” with-respect-to both $`m_b`$ dependence and a minimal set of couplings when $`m_b`$ is allowed to vary.
In summary, in the absence of explicit $`T`$ violation, three phase-type ambiguities versus the SM prediction exist: two dynamical ones with low effective mass scales, $`g_{VA}+g_{S+P}`$ with $`\mathrm{\Lambda }_{S+P}35GeV`$ and $`g_{VA}+g_{f_M+f_E}`$ with $`\mathrm{\Lambda }_{f_M+f_E}53GeV`$, and a third due to an arbitrary sign-flip in the $`b_L`$-amplitudes $`A_X(\lambda _b=1/2)=A_{VA}(\lambda _b=1/2)`$. The two dynamical ambiguities can be resolved by measurement of the sign of the large interference between the $`W`$ longitudinal/transverse amplitudes. Measurement of the sign of the $`\eta _L`$ helicity parameter will determine the sign of $`cos\beta _L`$ where $`\beta _L`$ is the relative phase of the two $`b_L`$-amplitudes ( $`\eta _L=\pm 0.46`$ where the upper sign is for the SM ). Both from the perspective of carefully testing the SM and that of searching for new physics, we believe that it is very important that experiments measure both this $`W`$ longitudinal/transverse interference parameter and its associated $`T`$ violation parameter $`\eta _{L}^{}{}_{}{}^{^{}}`$. The latter parameter is very important in the following analysis in this paper.
## 4 Consequences of Explicit $`T`$ Violation
To assay future measurements of helicity parameters in regard to $`T`$ violation, the next five sets of figures, Figs. 2-6, are for the case of a single additional pure-imaginary coupling, $`ig_i/2\mathrm{\Lambda }_i`$ or $`ig_i`$, associated with a specific additional Lorentz structure, $`i=S,P,S+P,\mathrm{}`$.
In the $`t`$ rest frame, the matrix element for $`tW^+b`$ is
$$\theta _1^t,\varphi _1^t,\lambda _{W^+},\lambda _b|\frac{1}{2},\lambda _1=D_{\lambda _1,\mu }^{(1/2)}(\varphi _1^t,\theta _1^t,0)A(\lambda _{W^+},\lambda _b)$$
(11)
where $`\mu =\lambda _{W^+}\lambda _b`$ in terms of the $`W^+`$ and $`b`$ helicities. $`\lambda _1`$ gives the $`t`$ quark’s spin component quantized along a $`z_1^t`$ axis, see Fig.1 in 2nd paper in Ref.. So, upon a boost back to the $`(t\overline{t})`$ center-of-mass frame, or to the $`\overline{t}`$ rest frame, $`\lambda _1`$ also specifies the helicity of the t quark. By rotational invariance there are only two amplitudes $`A(0,1/2),A(1,1/2)`$ for $`\lambda _b=1/2`$, and two with $`\lambda _{W^+}=0,1`$ for $`\lambda _b=1/2`$. For the $`CP`$-conjugate process, $`\overline{t}W^{}\overline{b}`$, in the $`\overline{t}`$ rest frame
$$\theta _2^t,\varphi _2^t,\lambda _W^{},\lambda _{\overline{b}}|\frac{1}{2},\lambda _2=D_{\lambda _2,\overline{\mu }}^{(1/2)}(\varphi _2^t,\theta _2^t,0)B(\lambda _W^{},\lambda _{\overline{b}})$$
(12)
with $`\overline{\mu }=\lambda _W^{}\lambda _{\overline{b}}`$.
As shown in Table 3 a specific discrete symmetry implies a specific relation among the associated helicity amplitudes. In the case of $`T`$ invariance, the helicity amplitudes must be purely real. The $`T`$ invariance of Table 3 will be violated if either (i) there is a fundamental violation of canonical $`T`$ invariance, or (ii) there are absorptive final-state interactions. In the SM, there are no such final-state interactions at the level of sensitivities considered in the present analysis. In our earlier papers, we have kept this assumption of “the absence of final-state interactions” manifest by referring to the $`T`$ invariance of Table 3 as “$`\stackrel{~}{T}_{FS}`$ violation”.
Barred parameters $`\overline{\xi },\overline{\zeta },\mathrm{}`$ have the analogous definitions for the $`CP`$ conjugate process, $`\overline{t}W^{}\overline{b}`$. Therefore, any $`\overline{\xi }\xi ,\overline{\zeta }\zeta ,\mathrm{}`$ $``$ CP is violated. That is, “slashed parameters” $`\mathit{\xi ̸}\xi \overline{\xi }`$, …, could be introduced to characterize and quantify the degree of CP violation. This should be regarded as a test for the presence of a non-CKM-type CP violation because, normally, a CKM-phase will contribute equally at tree level to both the $`tW^+b_L`$ decay amplitudes and so a CKM-phase will cancel out in the ratio of their moduli and in their relative phase. There are four tests for non-CKM-type CP violation.
A recent review of $`CP`$-violation in t-quark physics is in .
### 4.1 Additional $`S\pm P,,f_M\pm f_E,S,P,f_M,`$ or $`f_E`$ couplings
The two plots displayed in Fig.2 are for dimensional couplings with chiral $`S\pm P,f_M\pm f_E`$ and non-chiral $`S,P,f_M,f_E`$ Lorentz structures. The upper plot displays the $`\eta _{L}^{}{}_{}{}^{^{}}`$ helicity parameter versus the effective-mass scale $`\mathrm{\Lambda }_i`$ with $`g_i=1`$ in $`g_L=1`$ units. The lower plot displays the induced effect of the additional coupling on the partial width for $`tW^+b`$. The standard model limit is at the “wings” where $`|\mathrm{\Lambda }_i|\mathrm{}`$ for each additional dimensional coupling.
Fig.3 displays plots of the b-polarimetry interference parameters $`ϵ_{+}^{}{}_{}{}^{^{}}`$ and $`\kappa _{0}^{}{}_{}{}^{^{}}`$ versus $`\mathrm{\Lambda }_i`$ for the case of a single additional $`S,P,f_M,f_E`$ and $`S\pm P,f_Mf_E`$ coupling: Curves are omitted in the plots in this paper when the couplings produce approximately zero deviations in the helicity parameter of interest, e.g. this occurs for $`f_M+f_E`$ in both the $`ϵ_{+}^{}{}_{}{}^{^{}}`$ and $`\kappa _{0}^{}{}_{}{}^{^{}}`$ helicity parameters. The unitarity limit for $`ϵ_{+}^{}{}_{}{}^{^{}}`$ and $`\kappa _{0}^{}{}_{}{}^{^{}}`$ is also $`0.5`$ .
### 4.2 Additional $`V+A,V,`$ or $`A`$ couplings
An additional $`VA`$ type coupling with a complex phase versus the SM’s $`g_L`$ is equivalent to an additional overall complex factor in the SM’s helicity amplitudes. This will effect the overall partial width $`\mathrm{\Gamma }`$, but it can’t otherwise be observed by spin-correlation measurements.
For a single additional gauge-type coupling $`V,A,`$ or $`V+A`$, in Fig.4 are plots of the b-polarimetry interference parameters $`ϵ_{+}^{}{}_{}{}^{^{}}`$ and $`\kappa _{0}^{}{}_{}{}^{^{}}`$, and of the partial width for $`tW^+b`$ versus pure-imaginary coupling constant $`ig_i`$. The $`g_i`$ value is in $`g_L=1`$ units. In the cases of the additional dimensionless, gauge-type couplings, the standard model limit is at the origin, $`g_i0`$.
### 4.3 Indirect effects of $`T`$ violation on other helicity parameters
The plots in Fig.5 show the indirect effects of a single additional pure-imaginary chiral coupling, $`ig_i/2\mathrm{\Lambda }_i`$ or $`ig_i`$, on other helicity parameters. For the coupling strength ranges listed in the “middle table”, the upper plot shows the effects on the probability, $`P(W_L)`$, that the emitted $`W^+`$ is “Longitudinally” polarized and the effects on the probability, $`P(b_L)`$, that the emitted b-quark has “Left-handed” helicity. Each curve is parametrized by the magnitude of the associated $`g_i`$ or $`\mathrm{\Lambda }_i`$. On each curve, the central open circle corresponds to the region with a maximum direct $`T`$ violation signature, e.g. for $`f_M+f_E`$ from Fig. 2 this is at $`|\mathrm{\Lambda }_{f_M+f_E}|50GeV`$. The large/small solid circles correspond respectively to the ends of the ranges listed in the middle table where the direct signatures fall to about $`50\%`$ of their maximum values. Similarly the lower plot is for the W-polarimetry interference parameters $`\eta ,\omega `$. Curves are omitted for the remaining moduli parameter $`\zeta `$ since a single additional pure-imaginary coupling in these ranges produces approximately zero deviations from the pure $`VA`$ value of $`\zeta =0.41`$.
The plots in Fig.6 show the indirect effects of a single additional pure-imaginary non-chiral coupling on other helicity parameters. Versus the middle table given here, the curves are labeled as in Fig. 5. The upper plot is for the two probabilities $`P(W_L)`$ and $`P(b_L)`$. The lower plot is for the W-polarimetry interference parameters $`\eta ,\omega `$.
In summary, sizable $`T`$-violation signatures can occur for low-effective mass scales ( $`<320GeV`$) as a consequence of pure-imaginary couplings associated with a specific additional Lorentz structure. However, in most cases, such additional couplings can be more simply excluded by $`10\%`$ precision measurement of the probabilities $`P(W_L)`$ and $`P(b_L)`$. The W-polarimetry interference parameters $`\eta `$ and $`\omega `$ can also be used as indirect tests, or to exclude such additional couplings.
## 5 Tests for $`T`$ Violation Associated with the Dynamical Phase-Type Ambiguities
In Fig.7 are plots of the signatures for a partially-hidden $`T`$ violation associated with a $`S+P`$ phase-type ambiguity: We require Eq.(8) to hold when the additional $`S+P`$ coupling, $`g_{S+P}/2\mathrm{\Lambda }_{S+P}`$ has a complex effective mass scale parameter $`\mathrm{\Lambda }_{S+P}=|\mathrm{\Lambda }_{S+P}|\mathrm{exp}i\theta `$ where $`\theta `$ varies with the mass scale $`|\mathrm{\Lambda }_{S+P}|`$. For $`m_b=0`$, the resulting function $`\theta (|\mathrm{\Lambda }_{S+P}|)`$ is very simple. This construction maintains the standard model values in the massless b-quark limit for the four moduli parameters, $`P(W_L),P(b_L),\zeta ,`$ and $`\mathrm{\Gamma }`$. The function $`\theta (|\mathrm{\Lambda }_{S+P}|)`$ is then used for the $`S+P`$ coupling when $`m_b=4.5GeV`$. The SM values for the moduli parameters are essentially unchanged. There are two cases, $`\mathrm{sin}\theta 0`$ and $`\mathrm{sin}\theta 0`$. The phase choice of $`\varphi _{}^{R}{}_{1}{}^{}=\pm \pi `$, cf. top line in Table 1, has no consequence since it is a $`2\pi `$ phase difference.
For $`\mathrm{sin}\theta 0`$ in Fig.7 is the solid curve for the $`\eta _{L}^{}{}_{}{}^{^{}}`$, the $`T`$ violation W-polarimetry interference parameter, plotted versus $`1/|\mathrm{\Lambda }_{S+P}|`$. The dashed curve is for the W-polarimetry interference parameters $`\eta _L,\eta ,\omega `$ which are degenerate. The dark rectangles show the standard model values at the $`|\mathrm{\Lambda }_{S+P}|\mathrm{}`$ endpoint where $`\theta =\pi /2`$. At the other endpoint $`|\mathrm{\Lambda }_{S+P}|34.5GeV`$, or $`1/|\mathrm{\Lambda }_{S+P}|=0.029GeV^1`$, the coupling is purely real with $`\theta =\pi `$. The unitarity limit for each of these helicity parameters is $`0.5`$.
From the perspectives of (i) measuring the $`W`$ interference parameters and of (ii) excluding this type of $`T`$ violation, it is noteworthy that where $`\eta _{L}^{}{}_{}{}^{^{}}`$ has the maximum deviation, there is a zero in $`\eta _L,\eta ,\omega `$. So if the latter parameters were found to be smaller than expected or with the opposite sign than expected, this would be consistent with this type of $`T`$ violation.
At the maximum of $`\eta _{L}^{}{}_{}{}^{^{}}`$, $`|\mathrm{\Lambda }_{S+P}|49GeV`$ and the other $`T`$ violation parameters are also maximum. The curves for these parameters have the same over all shape as $`\eta _{L}^{}{}_{}{}^{^{}}`$ but their maxima are small, $`ϵ_{+}^{}{}_{}{}^{^{}}0.015`$ and $`\kappa _{0}^{}{}_{}{}^{^{}}0.028`$.
For the other case where $`\mathrm{sin}\theta 0`$, all these $`T`$ violation primed parameters have the opposite overall sign. The signs of other helicity parameters are not changed.
In Fig.8 are plots of the signatures for a partially-hidden $`T`$ violation associated with a $`f_M+f_E`$ phase-type ambiguity: As above for the analogous $`S+P`$ construction, the additional $`f_M+f_E`$ coupling $`g_{f_M+f_E}/2\mathrm{\Lambda }_{f_M+f_E}`$ now has an effective mass scale parameter $`\mathrm{\Lambda }_{f_M+f_E}=|\mathrm{\Lambda }_{f_M+f_E}|\mathrm{exp}i\theta `$ in which $`\theta `$ varies with the mass scale $`|\mathrm{\Lambda }_{f_M+f_E}|`$ to maintain standard model values in the massless b-quark limit for the moduli parameters $`P(W_L),P(b_L),`$ and $`\zeta `$. For the case $`\mathrm{sin}\theta 0`$, in Fig.8 the upper plot shows by the solid curve the $`T`$ violation W-polarimetry interference parameter $`\eta _{L}^{}{}_{}{}^{^{}}`$ versus $`1/|\mathrm{\Lambda }_{f_M+f_E}|`$. By the dashed curve, it shows the W-polarimetry interference parameters $`\eta _L,\eta ,\omega `$ which are degenerate. At the endpoint $`|\mathrm{\Lambda }_{f_M+f_E}|52.9GeV`$, or $`1/|\mathrm{\Lambda }_{f_M+f_E}|=0.0189GeV^1`$, the coupling is purely real with $`\theta =0`$.
Here, as in Fig. 7, where $`\eta _{L}^{}{}_{}{}^{^{}}`$ has the maximum deviation, there is a zero in $`\eta _L,\eta ,\omega `$. The lower plot shows the indirect effect of such a coupling on the partial width $`\mathrm{\Gamma }`$ for $`tW^+b`$.
At the maximum of $`\eta _{L}^{}{}_{}{}^{^{}}`$, $`|\mathrm{\Lambda }_{f_M+f_E}|63GeV`$. The curve for the $`T`$ violation parameter $`\kappa _{0}^{}{}_{}{}^{^{}}`$ has the same shape and is also maxmimum at the same position with a value $`\kappa _{0}^{}{}_{}{}^{^{}}0.005`$. $`ϵ_{+}^{}{}_{}{}^{^{}}`$ remains very small. For the other case where $`\mathrm{sin}\theta 0`$, each of these $`T`$ violation primed parameters has the opposite overall sign.
In summary, sufficiently precise measurement of the W-interference parameter $`\eta _L`$ and of the $`\eta _{L}^{}{}_{}{}^{^{}}`$ parameter can exclude partially-hidden $`T`$ violation associated with either of the two dynamical phase-type ambiguities. However, if $`\eta _L=(\eta +\omega )/2`$ were found to be smaller than expected or with a negative sign, this would be consistent with this type of $`T`$ violation.
Acknowledgments
For computer services, one of us (CAN) thanks John Hagan and Ted Brewster. This work was partially supported by U.S. Dept. of Energy Contract No. DE-FG 02-86ER40291.
Table Captions
Table 1: For the ambiguous moduli points, numerical values of the associated helicity amplitudes $`A(\lambda _{W^+},\lambda _b)`$. The values for the amplitudes are listed first in $`g_L=1`$ units, and second as $`A_{new}=A_{g_L=1}/\sqrt{\mathrm{\Gamma }}`$ which removes the effect of the differing partial width, $`\mathrm{\Gamma }`$ for $`tW^+b`$. \[$`m_t=175GeV,m_W=80.35GeV,m_b=4.5GeV`$ \].
Table 2: For the ambiguous moduli points, numerical values of the associated helicity parameters. Listed first are the four moduli parameters. Listed second are the values of the $`W`$-polarimetry interference parameters which could be used to resolve these dynamical ambiguities.
Table 3: The helicity formalism is based on the assumption of Lorentz invariance but not on any specific discrete symmetry property of the fundamental amplitudes, or couplings. For instance, for $`tW^+b`$ and $`\overline{t}W^{}\overline{b}`$ a specific discrete symmetry implies a definite symmetry relation among the associated helicity amplitudes.
Figure Captions
FIG. 1: For $`tW^+b`$ decay, display of the four helicity amplitudes $`A(\lambda _{W^+},\lambda _b)`$ relative to the $`W^+`$ boson and b-quark helicities. The upper sketch defines the measurable “ $`\alpha ,\beta ,\gamma `$ ” relative phases, c.f. Eqs(1). The lower sketch defines the real part and imaginary part (primed) helicity parameters corresponding to these relative phases. Measurement of a non-zero primed helicity parameter would be a direct signature for $`T`$ violation.
FIG. 2: To assay future measurements of helicity parameters in regard to $`T`$ violation, the next five sets of figures are for the case of a single additional pure-imaginary coupling, $`ig_i/2\mathrm{\Lambda }_i`$ or $`ig_i`$, associated with a specific additional Lorentz structure, $`i=S,P,S+P,\mathrm{}`$ . The two plots displayed here are for dimensional couplings with chiral $`S\pm P,f_M\pm f_E`$ and non-chiral $`S,P,f_M,f_E`$ Lorentz structures. The upper plot displays the $`\eta _{L}^{}{}_{}{}^{^{}}`$ helicity parameter versus the effective-mass scale $`\mathrm{\Lambda }_i`$ with $`g_i=1`$ in $`g_L=1`$ units. The lower plot displays the induced effect of the additional coupling on the partial width for $`tW^+b`$. The standard model limit is at the “wings” where $`|\mathrm{\Lambda }_i|\mathrm{}`$ for each additional dimensional coupling. The unitary limit for $`\eta _{L}^{}{}_{}{}^{^{}}`$ is $`0.5`$.
FIG. 3: Plots of the b-polarimetry interference parameters $`ϵ_{+}^{}{}_{}{}^{^{}}`$ and $`\kappa _{0}^{}{}_{}{}^{^{}}`$ versus $`\mathrm{\Lambda }_i`$ for the case of a single additional $`S,P,f_M,f_E`$ and $`S\pm P,f_Mf_E`$ coupling: Curves are omitted in the plots in this paper when the couplings produce approximately zero deviations in the helicity parameter of interest, e.g. this occurs for $`f_M+f_E`$ in both the $`ϵ_{+}^{}{}_{}{}^{^{}}`$ and $`\kappa _{0}^{}{}_{}{}^{^{}}`$ helicity parameters. The unitarity limit for $`ϵ_{+}^{}{}_{}{}^{^{}}`$ and $`\kappa _{0}^{}{}_{}{}^{^{}}`$ is also $`0.5`$ .
FIG. 4: For a single additional gauge-type coupling $`V,A,`$ or $`V+A`$, plots of the b-polarimetry interference parameters $`ϵ_{+}^{}{}_{}{}^{^{}}`$ and $`\kappa _{0}^{}{}_{}{}^{^{}}`$, and of the partial width for $`tW^+b`$ versus pure-imaginary coupling constant $`ig_i`$. The $`g_i`$ value is in $`g_L=1`$ units. In the cases of the additional dimensionless, gauge-type couplings, the standard model limit is at the origin, $`g_i0`$.
FIG. 5: These plots show the indirect effects of a single additional pure-imaginary chiral coupling, $`ig_i/2\mathrm{\Lambda }_i`$ or $`ig_i`$, on other helicity parameters. For the coupling strength ranges listed in the “middle table”, the upper plot shows the effects on the probability, $`P(W_L)`$, that the emitted $`W^+`$ is “Longitudinally” polarized and the effects on the probability, $`P(b_L)`$, that the emitted b-quark has “Left-handed” helicity. Each curve is parametrized by the magnitude of the associated $`g_i`$ or $`\mathrm{\Lambda }_i`$. On each curve, the central open circle corresponds to the region with a maximum direct $`T`$ violation signature, e.g. for $`f_M+f_E`$ from Fig. 2 this is at $`|\mathrm{\Lambda }_{f_M+f_E}|50GeV`$. The large/small solid circles correspond respectively to the ends of the ranges listed in the middle table where the direct signatures fall to about $`50\%`$ of their maximum values. Similarly the lower plot is for the W-polarimetry interference parameters $`\eta ,\omega `$. Curves are omitted for the remaining moduli parameter $`\zeta `$ since a single additional pure-imaginary coupling in these ranges produces approximately zero deviations from the pure $`VA`$ value of $`\zeta =0.41`$. A dark rectangle denotes the value for the pure $`VA`$ coupling of the standard model.
FIG. 6: These plots show the indirect effects of a single additional pure-imaginary non-chiral coupling on other helicity parameters. Versus the middle table given here, the curves are labeled as in Fig. 5. The upper plot is for the two probabilities $`P(W_L)`$ and $`P(b_L)`$. The lower plot is for the W-polarimetry interference parameters $`\eta ,\omega `$.
FIG. 7: Plots of the signatures for a partially-hidden $`T`$ violation (see text) associated with a $`S+P`$ phase-type ambiguity: In this case, the additional $`S+P`$ coupling, $`g_{S+P}/2\mathrm{\Lambda }_{S+P}`$, has an effective mass scale parameter $`\mathrm{\Lambda }_{S+P}=|\mathrm{\Lambda }_{S+P}|\mathrm{exp}i\theta `$ where $`\theta `$ varies with the mass scale $`|\mathrm{\Lambda }_{S+P}|`$ to maintain standard model values in the massless b-quark limit for the four moduli parameters, $`P(W_L),P(b_L),\zeta ,`$ and $`\mathrm{\Gamma }`$. Plotted versus $`1/|\mathrm{\Lambda }_{S+P}|`$ for the case $`\mathrm{sin}\theta 0`$ is the solid curve for the $`\eta _{L}^{}{}_{}{}^{^{}}`$, the $`T`$ violation W-polarimetry interference parameter and the dashed curve for the W-polarimetry interference parameters $`\eta _L,\eta ,\omega `$ which are degenerate. For $`\mathrm{sin}\theta 0`$, the $`\eta _{L}^{}{}_{}{}^{^{}}`$ sign is opposite. The dark rectangles show the standard model values at the $`|\mathrm{\Lambda }_{S+P}|\mathrm{}`$ endpoint where $`\theta =\pi /2`$. At the other endpoint $`|\mathrm{\Lambda }_{S+P}|34.5GeV`$, or $`1/|\mathrm{\Lambda }_{S+P}|=0.029GeV^1`$, the coupling is purely real with $`\theta =\pi `$. Where $`\eta _{L}^{}{}_{}{}^{^{}}`$ has the maximum deviation, there is a zero in $`\eta _L,\eta ,\omega `$.
FIG. 8: Plots of the signatures for a partially-hidden $`T`$ violation (see text) associated with a $`f_M+f_E`$ phase-type ambiguity: The additional $`f_M+f_E`$ coupling, $`g_{f_M+f_E}/2\mathrm{\Lambda }_{f_M+f_E}`$, has an effective mass scale parameter $`\mathrm{\Lambda }_{f_M+f_E}=|\mathrm{\Lambda }_{f_M+f_E}|\mathrm{exp}i\theta `$ where $`\theta `$ varies with the mass scale $`|\mathrm{\Lambda }_{f_M+f_E}|`$ to maintain standard model values in the massless b-quark limit for the moduli parameters $`P(W_L),P(b_L),`$ and $`\zeta `$. Versus $`1/|\mathrm{\Lambda }_{f_M+f_E}|`$ for $`\mathrm{sin}\theta 0`$, the upper plot shows by the solid curve $`\eta _{L}^{}{}_{}{}^{^{}}`$. By the dashed curve, it shows $`\eta _L,\eta ,\omega `$ which are degenerate. At the endpoint $`|\mathrm{\Lambda }_{f_M+f_E}|52.9GeV`$, or $`1/|\mathrm{\Lambda }_{f_M+f_E}|=0.0189GeV^1`$, the coupling is purely real with $`\theta =0`$. For $`\mathrm{sin}\theta 0`$, the $`\eta _{L}^{}{}_{}{}^{^{}}`$ sign is opposite. The lower plot shows the indirect effect of such a coupling on the partial width $`\mathrm{\Gamma }`$ for $`tW^+b`$.
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# Report of the Beyond the MSSM Subgroup for the Tevatron Run II SUSY/Higgs Workshop
## I Introduction
Most of the studies performed to assess the discovery reach for supersymmetry and most of the current limits on the masses of supersymmetric particles have been obtained assuming R-parity conservation, the minimal matter content of the Minimal Supersymmetric Model (MSSM), and universal boundary conditions at $`M_U`$ for the soft-SUSY-breaking parameters: $`m_0`$ for the scalar masses; $`M_0`$ for the SU(3),SU(2) and U(1) gaugino masses $`M_{3,2,1}`$; and $`A_0`$ for the tri-linear scalar field couplings. Additional parameters of the MSSM include: $`\mathrm{tan}\beta `$, the ratio of Higgs field vacuum expectation values $`H_u/H_d`$; $`\mu `$, the coefficient of the bilinear $`\widehat{H}_u\widehat{H}_d`$ superpotential term; and $`B`$, which specifies the strength of the corresponding $`H_uH_d`$ scalar field mixing term. By requiring correct electroweak symmetry breaking after evolution down to the scale $`m_Z`$, the magnitude of $`\mu `$ is fixed and only its sign remains undetermined. This boundary condition scenario is often referred to as the mSUGRA or CMSSM model. If unification of the $`\tau `$ and $`b`$ Yukawa couplings at $`M_U`$ is required, then correct EWSB strongly constrains $`\mathrm{tan}\beta `$ as well.
While the matter content and boundary conditions of the MSSM have the virtue of simplicity and can be reasonably motivated in the context of several types of gravity-mediated supersymmetry breaking, other possibilities should certainly be considered. For a model with MSSM matter content and R-parity conservation, the most general form of soft-SUSY-breaking allows for a total of 124 parameters (not counting certain additional parameters expected to be suppressed by $`1/M_U`$ factors). This is to be compared to the 19 parameters of the Standard Model. The most general form of R-parity violation increases the parameter count to 315. Some of these parameters are associated with phases and CP violation. A summary appears in Table LABEL:paramcount. A final parameter is the mass of the gravitino, $`\stackrel{~}{G}`$. If the scale of supersymmetry breaking is sufficiently low, $`m_{\stackrel{~}{G}}`$ can be small enough that it is the LSP. In particular, this is the expectation in models where supersymmetry breaking is mediated by gauge interactions.
Although almost all of the most general parameter space is excluded by various phenomenological constraints (no FCNC, proton stability, small EDM’s, etc.), there are sub-spaces that differ drastically from mSUGRA while maintaining consistency with all existing data. Examples include models with R-parity violation and the standard GMSB phenomenology, both of which will be covered in separate reports. We also do not consider non-zero phases. Aside from a discussion of the phenomenology of an extremely light $`\stackrel{~}{G}`$, our focus will be on models with a heavy $`\stackrel{~}{G}`$ and soft-SUSY-breaking parameters that conserve R-parity and CP. Even with these restrictions, there are many well-motivated theories with MSSM matter content that yield vastly different phenomenology than the mSUGRA model. In addition, we will consider a number of models in which the matter/gauge content of the MSSM is extended. We will focus in particular on implications for supersymmetry discovery and study at the upgraded Tevatron.
We now give brief motivation and an introduction to the models considered.
* One of the least satisfactory features of the MSSM is the ad hoc nature of the parameter $`\mu `$, which a priori is most naturally of order $`M_U`$, but which is expected to be $`<\mathrm{\hspace{0.33em}1}\text{ TeV}`$ for natural EWSB. The addition of a singlet superfield $`\widehat{S}`$ provides a very compelling and natural origin of the $`\mu \widehat{H}_u\widehat{H}_d`$ superpotential term. Such a term arises if the scalar component of $`\widehat{S}`$ ($`S`$) acquires a vacuum expectation value. The result is an effective superpotential interaction of the form $`S\widehat{H}_u\widehat{H}_d`$. The $`S`$ quantum degrees of freedom result in 1 CP-even and 1 CP-odd Higgs bosons beyond the 2 CP-even and 1 CP-odd Higgs bosons of the MSSM. The spin-1/2 component of $`\widehat{S}`$ provides an additional neutralino, $`\stackrel{~}{S}`$, that can mix with the usual four neutralinos of the MSSM. It is very natural for the LSP of this model to be the $`\stackrel{~}{S}`$. All supersymmetric particles then cascade decay down to the $`\stackrel{~}{S}`$. The phenomenology of SUSY detection is then significantly altered compared to mSUGRA. A review of this phenomenology is given in Sec. II.
* In mSUGRA, the LSP is essentially always a light bino-like neutralino, $`\stackrel{~}{\chi }_1^0\stackrel{~}{B}`$. However, there is substantial motivation for the possibility that the LSP is a massive gluino. This occurs if $`M_3M_{1,2}`$, as is possible in several well-motivated SUSY-breaking scenarios. Current limits on a heavy gluino LSP are summarized and discovery prospects are discussed in Sec. III.
* Another alternative arrangement of the gaugino masses that arises in string and brane models is $`M_2<M_1<M_3`$. In this case, the LSP $`\stackrel{~}{\chi }_1^0`$ is wino-like and is highly degenerate with the lightest chargino. (This assumes $`|\mu |`$ is large, as usually implied by RGE EWSB.) The resulting phenomenology differs greatly from mSUGRA phenomenology. If $`\mathrm{\Delta }m_{\stackrel{~}{\chi }}m_{\stackrel{~}{\chi }_1^\pm }m_{\stackrel{~}{\chi }_1^0}`$ is not too much larger than $`m_\pi `$, striking background-free signals for $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}+\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0`$ production will be present. For $`\mathrm{\Delta }m_{\stackrel{~}{\chi }}>\mathrm{\hspace{0.33em}300}\text{ MeV}`$, detection of these processes will be very difficult; one will have to hope that other SUSY particles are light. Section IV gives a discussion of the phenomenology and of some very critical detector issues and related discovery strategies.
* Although the scalar masses have the universal value $`m_0`$ at the high scale, significant flavor violation can arise via RGE evolution if this scale is not the same as $`M_U`$. More generally, FCNC will be a problem unless the SUSY breaking mechanism yields either universality for the scalar masses or flavor alignment. An interesting exception to this statement is the possibility that the scalar masses for all the sfermions of the first two generations are simply extremely heavy ($`>10\text{ TeV}`$) and, thus, have greatly suppressed FCNC effects. Very heavy scalars are also helpful for unifying with greater precision the strong coupling constant with the SU(2) and U(1) couplings for the somewhat ‘low’ value, $`\alpha _s(m_Z)0.12`$, preferred by existing data. Of course, to maintain naturalness for the Higgs sector the 3rd generation squarks should be below $`1\text{ TeV}`$. This scenario is sometimes referred to as Superheavy Supersymmetry or More Minimal Supersymmetry. Section V reviews some theoretical issues and constraints on this scenario, including the apparent necessity to have GMSB-like boundary conditions in order to preserve anomaly cancellations.
* As noted earlier, the mass of the $`\stackrel{~}{G}`$ is another crucial parameter of supersymmetry. If $`m_{\stackrel{~}{G}}`$ is very small, the couplings of the $`\stackrel{~}{G}`$ are sufficiently large that processes in which the $`\stackrel{~}{G}`$ is directly produced have observable rates. Since the $`\stackrel{~}{G}`$ is undetectable, the most basic signature is jets plus missing energy. If these processes are detected, they provide a measurement of the scale of supersymmetry breaking, possibly the most important parameter of supersymmetry. The phenomenology and discovery prospects for direct production of a very light $`\stackrel{~}{G}`$ are reviewed in Sec. VI.
* An interesting question is whether superstring theory provides any guidance as regards boundary conditions and matter content for low energy supersymmetry. The detailed predictions of one sample superstring model are outlined in Sec. VII. The model considered has a plethora of additional matter, including exotics, extra Higgs bosons, and extra gauge bosons. This provides further warning against being complacent in our approach to SUSY phenomenology.
* The possibility that left-right symmetry is restored at a high energy scale is very attractive. In LR-symmetric models, proper symmetry breaking requires introduction of triplet Higgs representations that contain a doubly-charged Higgs field. In the supersymmetric context, these doubly-charged scalars have a doubly-charged higgsino partner. Careful investigation reveals that these are likely to be one of the lightest states in the superparticle spectrum. They will appear in cascade decays and can also be directly produced. The phenomenology of the doubly-charged higgsino states is reviewed in Sec. VIII.
* Recently, the possibility that the compactified extra dimensions of the string/brane world are large and that the Kaluza-Klein excited states are within experimental reach has received much attention. Some of the indirect signals for such extra dimension are reviewed in Sec. IX. In addition, Sec. VI considers external KK gravitons production, which provides a signature similar to that of the very light gravitino through jets plus missing energy.
* It is well-known that supersymmetry predicts a rather low mass for the lightest CP-even Higgs boson. Thus, one should ask whether supersymmetry can be rescued if a sufficiently light Higgs boson is not discovered. One means for increasing the upper bound on the light Higgs is to introduce a 4th family. The contraints upon and implications for the Higgs sector of supersymmetry in a 4-family model are discussed in Sec. X.
* Is there room for a 4th family in supersymmetry? If the Yukawa couplings associated with the 4th family are to remain perturbative in evolving from $`m_Z`$ up to some high scale, one finds that the leptons and quarks of the 4th family must be quite light. Experimental constraints are becoming very restrictive. In Sec. XI, the current situation is reviewed with the conclusion that the 4th generation will almost certainly be either discovered or eliminated as a possibility during Run II at the Tevatron.
* Could the gluino be very light? Remarkably, this scenario cannot yet be absolutely excluded. In addition, it might explain some detailed features of Run I jet data at very high $`p_T`$. Section XII presents the case for a very light gluino.
* The DØ detector has been upgraded dramatically for Run II. It is important to understand the extent to which it will be able to probe some of the more exotic supersymmetry scenarios that are discussed here and elsewhere. Particularly interesting are signals associated with long-lived charged particles, photons, vertices etc. The capabilities of the DØ detector as regards such exotic phenomena are discussed in Sec. XIII.
* The other major detector at the Tevatron, CDF, has also been upgraded. Section XIV reviews its capabilities for 4th generation searches via: looking for a long-lived ($`b^{}`$) parent of the $`Z`$; looking for prompt $`b^{}bZ`$ production; and searching for a long-lived heavy quark or similar object.
* New gauges bosons are a common feature of supersymmetric models motivated by string theory. The ability to detect such gauge bosons and to determine their couplings during Run II at the Tevatron is considered in Section XV
* We end with some brief concluding remarks in Section XVI.
## II Cascade decays in the NMSSM
U. Ellwanger and C Hugonie
The NMSSM (Next-to-minimal SSM, or (M+1)SSM) is defined by the addition of a gauge singlet superfield $`S`$ to the MSSM. The superpotential $`W`$ is scale invariant, i.e. there is no $`\mu `$-term. Instead, two Yukawa couplings $`\lambda `$ and $`\kappa `$ appear in $`W`$. Apart from the standard quark and lepton Yukawa couplings, $`W`$ is given by
$`W=\lambda H_1H_2S+{\displaystyle \frac{1}{3}}\kappa S^3+\mathrm{}`$ (1)
and the corresponding trilinear couplings $`A_\lambda `$ and $`A_\kappa `$ are added to the soft susy breaking terms. The vev of $`S`$ generates an effective $`\mu `$-term with $`\mu =\lambda S`$.
The constraint NMSSM (CNMSSM) ell1 is defined by universal soft susy breaking gaugino masses $`M_0`$, scalar masses $`m_0^2`$ and trilinear couplings $`A_0`$ at the GUT scale, and a number of phenomenological constraints:
\- Consistency of the low energy spectrum and couplings with negative Higgs and sparticle searches.
\- In the Higgs sector, the minimum of the effective potential with $`H_1`$ and $`H_20`$ has to be deeper than any minimum with $`H_1`$ and/or $`H_2=0`$. Charge and colour breaking minima induced by trilinear couplings have to be absent. (However, deeper charge and colour breaking minima in ”UFB” directions are allowed, since the decay rate of the physical vacuum into these minima is usually large compared to the age of the universe ell2 .)
Cosmological constraints as the correct amount of dark matter are not imposed at present. (A possible domain wall problem due to the discrete $`Z_3`$ symmetry of the model is assumed to be solved by, e.g., embedding the $`Z_3`$ symmetry into a $`U(1)`$ gauge symmetry at $`M_{GUT}`$, or by adding non-renormalisable interactions which break the $`Z_3`$ symmetry without spoiling the quantum stability ell3 .)
The number of free parameters of the CNMSSM, ($`M_{1/2}`$, $`m_0`$, $`A_0`$, $`\lambda `$, $`\kappa `$ \+ standard Yukawa couplings), is the same as in the CMSSM ($`M_{1/2}`$, $`m_0`$, $`A_0`$, $`\mu `$, $`B+idem`$). The new physical states in the CNMSSM are one additional neutral Higgs scalar and Higgs pseudoscalar, respectively, and one additional neutralino. In general these states mix with the corresponding ones of the MSSM with a mixing angle proportional to the Yukawa coupling $`\lambda `$. However, in the CNMSSM $`\lambda `$ turns out to be quite small, $`\lambda <\mathrm{\hspace{0.33em}0.1}`$ (and $`\lambda 1`$ for most allowed points in the parameter space) ell1 Thus the new physical states are generally almost pure gauge singlets with very small couplings to the standard sector.
The new states in the Higgs sector can be very light, a few GeV or less, depending on $`\lambda `$ ell4 . Due to their small couplings to the $`Z`$ boson they will escape detection at LEP and elsewhere, i.e. the lightest “visible” Higgs boson is possibly the next-to-lightest Higgs of the NMSSM. The upper limits on the mass of this visible Higgs boson (and its couplings) are, on the other hand, very close to the ones of the MSSM, i.e. $`<\mathrm{\hspace{0.33em}140}`$ GeV depending on the stop masses ell4 .
The phenomenology of sparticle production in the CNMSSM can differ considerably from the MSSM, depending on the mass of the additional state $`\stackrel{~}{S}`$ in the neutralino sector: If the $`\stackrel{~}{S}`$ is not the LSP, it will hardly be produced, and all sparticle decays proceed as in the MSSM with a LSP in the final state. If, on the other hand, the $`\stackrel{~}{S}`$ is the LSP, the sparticle decays will proceed differently: First, the sparticles will decay into the NLSP, because the couplings to the $`\stackrel{~}{S}`$ are too small. Only then the NLSP will realize that it is not the true LSP, and decay into the $`\stackrel{~}{S}`$ plus an additional cascade.
The condition for a singlino LSP scenario can be expressed relatively easily in terms of the bare parameters of the CNMSSM: Within the allowed parameter space of the CNMSSM, the lightest non-singlet neutralino is essentially a bino $`\stackrel{~}{B}`$. Since the masses of $`\stackrel{~}{S}`$ and $`\stackrel{~}{B}`$ are proportional to $`A_0`$ and $`M_{1/2}`$, respectively, one finds, to a good approximation, that the $`\stackrel{~}{S}`$ is the true LSP if the bare susy breaking parameters satisfy $`|A_0|<\mathrm{\hspace{0.33em}0.4}M_{1/2}`$. Since $`A_0^2>\mathrm{\hspace{0.33em}9}m_0^2`$ is also a necessary condition within the CNMSSM, the singlino LSP scenario corresponds essentially to the case where the gaugino masses are the dominant soft susy breaking terms.
Note, however, that the $`\stackrel{~}{B}`$ is not necessarily the NLSP in this case: Possibly the lightest stau $`\stackrel{~}{\tau }_1`$ is lighter than the $`\stackrel{~}{B}`$, since the lightest stau can be considerably lighter than the sleptons of the first two generations. Nevertheless, most sparticle decays will proceed via the $`\stackrel{~}{B}\stackrel{~}{S}+\mathrm{}`$ transition, which will give rise to additional cascades with respect to decays in the MSSM. The properties of this cascade have been analysed in ell5 , and in the following we will briefly discuss the branching ratios and the $`\stackrel{~}{B}`$ life times in the different parameter regimes:
a) $`\stackrel{~}{B}\stackrel{~}{S}\nu \overline{\nu }`$: This invisible process is mediated dominantly by sneutrino exchange. Since the sneutrino mass, as the mass of $`\stackrel{~}{B}`$, is essentially fixed by $`M_{1/2}`$ ell5 , the associated branching ratio varies in a predictable way with $`M_{\stackrel{~}{B}}`$: It can become up to 90% for $`M_{\stackrel{~}{B}}30`$ GeV, but decreases with $`M_{\stackrel{~}{B}}`$ and is maximally 10% for $`M_{\stackrel{~}{B}}>\mathrm{\hspace{0.33em}65}`$ GeV.
b) $`\stackrel{~}{B}\stackrel{~}{S}l^+l^{}`$: This process is mediated dominantly by the exchange of a charged slepton in the s-channel. If the lightest stau $`\stackrel{~}{\tau }_1`$ is considerably lighter than the sleptons of the first two generations, the percentage of taus among the charged leptons can well exceed $`\frac{1}{3}`$. If $`\stackrel{~}{\tau }_1`$ is lighter than $`\stackrel{~}{B}`$, it is produced on-shell, and the process becomes $`\stackrel{~}{B}\stackrel{~}{\tau }_1\tau \stackrel{~}{S}\tau ^+\tau ^{}`$. Hence we can have up to 100% taus among the charged leptons and the branching ratio of this channel can become up to 100%.
c) $`\stackrel{~}{B}\stackrel{~}{S}S`$: This two-body decay is kinematically allowed if both $`\stackrel{~}{S}`$ and $`S`$ are sufficiently light. (A light $`S`$ is not excluded by Higgs searches at LEP1, if its coupling to the $`Z`$ is too small ell4 .) However, the coupling $`\stackrel{~}{B}\stackrel{~}{S}S`$ is proportional to $`\lambda ^2`$, whereas the couplings appearing in the decays a) and b) are only of $`O(\lambda )`$. Thus this decay can only be important for $`\lambda `$ not too small. In ell5 , we found that its branching ratio can become up to 100% in a window $`10^3<\lambda <\mathrm{\hspace{0.33em}10}^2`$. Of course, $`S`$ will decay immediately into $`b\overline{b}`$ or $`\tau ^+\tau ^{}`$, depending on its mass. (If the branching ratio $`Br(\stackrel{~}{B}\stackrel{~}{S}S)`$ is substantial, $`S`$ is never lighter than $`5`$ GeV.) If the singlet is heavy enough, its $`b\overline{b}`$ decay gives rise to 2 jets with $`B`$ mesons, which are easily detected with $`b`$-tagging. In any case, the invariant mass of the $`b\overline{b}`$ or the $`\tau ^+\tau ^{}`$ system would be peaked at $`M_S`$, making this signature easy to search for.
d) $`\stackrel{~}{B}\stackrel{~}{S}\gamma `$: This branching ratio can be important if the mass difference $`\mathrm{\Delta }M=M_{\stackrel{~}{B}}M_{\stackrel{~}{S}}`$ is small ($`<\mathrm{\hspace{0.33em}5}`$ GeV).
Further possible final states like $`\stackrel{~}{B}\stackrel{~}{S}q\overline{q}`$ via $`Z`$ exchange have always branching ratios below 10%. (The two-body decay $`\stackrel{~}{B}\stackrel{~}{S}Z`$ is never important, even if $`\mathrm{\Delta }M`$ is larger than $`M_Z`$: In this region of the parameter space $`\stackrel{~}{\tau }_1`$ is always the NLSP, and thus the channel $`\stackrel{~}{B}\stackrel{~}{\tau }_1\tau `$ is always preferred.)
The $`\stackrel{~}{B}`$ life time depends strongly on the Yukawa coupling $`\lambda `$, since the mixing of the singlino $`\stackrel{~}{S}`$ with gauginos and higgsinos is proportional to $`\lambda `$. Hence, for small $`\lambda `$ (or a small mass difference $`\mathrm{\Delta }M`$) the $`\stackrel{~}{B}`$ can be so long lived that it decays only after a macroscopic length of flight $`l_{\stackrel{~}{B}}`$. An approximate formula for $`l_{\stackrel{~}{B}}`$ (in meters) is given by
$`l_{\stackrel{~}{B}}[m]210^{10}{\displaystyle \frac{1}{\lambda ^2M_{\stackrel{~}{B}}[GeV]}},`$ (2)
and $`l_{\stackrel{~}{B}}`$ becomes $`>1`$ mm for $`\lambda <\mathrm{\hspace{0.33em}6}10^5`$.
To summarize, the following unconventional signatures are possible within the CNMSSM, compared to the MSSM:
a) additional cascades attached to the original vertex (but still missing energy and momentum): one or two additional $`l^+l^{}`$, $`\tau ^+\tau ^{}`$ or $`b\overline{b}`$ pairs or photons, with the corresponding branching ratios depending on the parameters of the model.
b) one or two additional $`l^+l^{}`$ or $`\tau ^+\tau ^{}`$ pairs or photons with macroscopically displaced vertices, with distances varying from millimeters to several meters. These displaced vertices do not point towards the interaction point, since an additional invisible particle is produced.
More details on the allowed branching ratios and life times can be found in ell5 , applications to sparticle production processes et LEP 2 are published in ell6 , and differential (spin averaged) cross sections of the $`\stackrel{~}{B}\stackrel{~}{S}`$ decay are available upon request.
## III Report on the gluino-LSP scenario
H Baer, K. Cheung, J.F. Gunion
### III.1 Introduction
This contribution will present a brief overview of the results of Ref. bcg . Most GUT scale boundary conditions (e.g. the mSUGRA universal boundary conditions in which the gaugino masses are assumed to have a common value at $`M_U`$) lead to the gluino being much heavier than the lightest neutralino. However, there are several models in which the gluino is heavy but is yet the lightest supersymmetric particle, denoted $`\stackrel{~}{g}`$-LSP.
(a) A $`\stackrel{~}{g}`$-LSP can arise in models in which the gaugino masses are given by one-loop corrections plus a contribution from Green-Schwarz mixing (parameterized by $`\delta _{GS}`$) nonuniv ; guniondrees2 . An example is the O-II string model in the limit where all SUSY breaking arises from the size-modulus field and none from the dilaton. At $`M_U`$ one has:
$$M_3:M_2:M_1\stackrel{OII}{}(3+\delta _{GS}):(1\delta _{GS}):(\frac{33}{5}\delta _{GS}),$$
(3)
and after evolution down to 1 TeV or below a heavy gluino is the LSP when $`\delta _{GS}3`$ (a preferred range for the model).
(b) In the GMSB context, the possibility of a heavy $`\stackrel{~}{g}`$-LSP has been stressed in Ref. raby . In the model constructed, the $`\stackrel{~}{g}`$ is the LSP as a result of mixing between the Higgs fields and the messenger fields.
In fact, there are three significant indications that a light gluino is to be preferred over the universal gaugino mass result of $`m_{\stackrel{~}{g}}7m_{\stackrel{~}{\chi }_1^0}`$.
$``$ The first such hint relates to the magnitude of $`\alpha _s`$ predicted by requiring precise gauge coupling unification at $`M_U10^{16}\mathrm{GeV}`$. For universal masses at $`M_U`$, gauge coupling unification typically requires $`\alpha _s(m_Z)0.1250.13`$ when sparticle masses are $`<1\mathrm{TeV}`$, a value that is uncomfortably high relative to the best fit value of $`\alpha _s(m_Z)<0.12`$. Although unified gauge couplings at $`M_U`$ can be made consistent with $`\alpha _s(m_Z)<0.12`$ if the sparticle masses important in the gauge coupling running are all $`>10\mathrm{TeV}`$ (for which fine tuning is regarded as a problem), a much more interesting possibility is that discussed in Ref. lrshif . There it is shown that if the gluino mass is substantially below, or at least comparable to the $`\stackrel{~}{\chi }_1^\pm `$ mass (these determine the two most critical thresholds in gauge coupling running), then $`\alpha _s(m_Z)<0.12`$ is much easier to achieve.
$``$ The second hint is from fine tuning. As discussed most recently in Ref. kaneking , the most severe problem in fine tuning arises from the fact that the magnitude of the standard measure of fine tuning contains a term proportional to $`M_3^2(M_U)`$ with a very large numerical coefficient, much larger than the (possibly canceling) terms proportional to $`M_1^2(M_U)`$ and $`M_2^2(M_U)`$. (Very roughly, the relative size of these coefficients are determined by the relative size of the corresponding gauge couplings squared, although $`\mathrm{tan}\beta `$ is also an important ingredient.) The fine tuning problem is greatly relaxed if $`M_3(M_U)`$ is substantially smaller than $`M_{1,2}(M_U)`$. For example, using the $`\mathrm{tan}\beta =2.5`$ numerical coefficients given in kaneking , one measure of fine tuning is (all parameters at $`M_U`$ with $`\widehat{M}_i`$ and $`\widehat{m}_i`$ denoting $`M_i(M_U)/m_Z`$ and $`m_i(M_U)/m_Z`$)
$$\mathrm{\Delta }_\mu 23.8\widehat{M}_3^21.3\widehat{M}_2^2+0.01\widehat{M}_1^2+1.66\widehat{M}_2\widehat{M}_3+0.206\widehat{M}_1\widehat{M}_3+0.045\widehat{M}_1\widehat{M}_2+\mathrm{},$$
(4)
Since the $`\widehat{M}_2^2`$ coefficient is not very large, taking $`\widehat{M}_3=0`$ (leading to a gluino LSP) avoids significant fine-tuning problems. One can even arrange for the $`\widehat{M}_3`$ and $`\widehat{M}_2`$ terms to cancel ($`\widehat{M}_3\widehat{M}_2/5`$). Thus, the $`\widehat{M}_3`$ values predicted for O-II model boundary conditions in the favored $`\delta _{GS}3,4,5`$ range would lead to a considerable relaxation of the the fine tuning problem and a gluino that is much lighter than normally anticipated.
$``$ The third hint relates to vacuum stability. This has been recently reviewed in Ref. cim , where references to earlier work can be found. The most serious problem associated with many of the soft-supersymmetry-breaking boundary conditions motivated by string theory is the presence of directions in field space for which the effective potential is unbounded from below (UFB). Dangerous charge and color breaking (CCB) minima are also possible. The strongest constraint typically arises from the UFB-3 direction, which involves the fields $`H_2,\nu _{L_i},e_{L_j},e_{R_j}`$ with $`ij`$. After minimizing the effective potential, the latter three fields can be expressed in terms of $`H_2`$. The value of the potential in the UFB-3 direction is then given by:
$$V_{UFB3}=(m_{H_2}^2+m_{L_i}^2)|H_2|^2+\frac{|\mu |}{\lambda _{e_j}}(m_{L_j}^2+m_{e_j}^2+m_{L_i}^2)|H_2|\frac{2m_{L_i}^4}{g^{\mathrm{\hspace{0.17em}2}}+g^2},$$
(5)
where $`\lambda _{e_j}`$ is the leptonic Yukawa coupling of the $`j`$-th generation. We must have $`V_{UFB3}(Q=\widehat{Q})>V_{\mathrm{real}\mathrm{min}}=\frac{1}{8}(g^{\mathrm{\hspace{0.17em}2}}+g^2)(v_2^2v_1^2)^2`$, where $`\widehat{Q}\mathrm{Max}[\lambda _{\mathrm{top}}|H_2|,m_{\stackrel{~}{t}}]`$. The problem arises when $`m_{H_2}^2`$ is negative (as happens under RGE electroweak symmetry breaking) and is most restrictive for $`j=3`$ for which $`\lambda _{e_j}`$ is largest. One finds that any significant amount of supersymmetry breaking from the dilaton leads to violation of the condition. The modulus-dominated (or one-loop) limit boundary conditions of Eq. 3 (or extremely close thereto) open up some allowed parameter space. This happens as follows. The RGE for $`m_{H_2}^2`$ is (using $`t=\mathrm{ln}(Q)`$) $`dm_{H_2}^2/dt=12(\lambda _{\mathrm{top}}/4\pi )^2m_{\stackrel{~}{t}}^2+\mathrm{}`$. The smaller $`m_{\stackrel{~}{t}}`$, the less negative $`m_{H_2}^2`$ is driven in evolving down from $`M_U`$, and, since
$$\frac{dm_{\stackrel{~}{q}}^2}{dt}=\left(\frac{1}{2\pi }\right)\left(\frac{16}{3}M_3^2\alpha _3+\mathrm{}\right),$$
(6)
smaller $`M_3`$ implies smaller $`m_{\stackrel{~}{t}}`$. Meanwhile, the positive terms in Eq. (5) have evolution
$$\frac{dm_e^2}{dt}=\left(\frac{1}{2\pi }\right)\left(\frac{12}{5}M_1^2\alpha _1+\mathrm{}\right),\frac{dm_L^2}{dt}=\left(\frac{1}{2\pi }\right)\left(3M_2^2\alpha _2\frac{3}{5}M_1^2\alpha _1+\mathrm{}\right),$$
(7)
implying larger $`m_e^2`$ and $`m_L^2`$ in $`V_{UFB3}`$ for larger $`M_{1,2}`$. Thus, non-universal boundary conditions with small $`M_3`$ relative to $`M_{1,2}`$ are crucial in the string model context and, more generally, are quite useful in satisfying the UFB-3 constraint.
It is often stated that a stable $`\stackrel{~}{g}`$-LSP is ruled out by virtue of relic density constraints, especially the non-observation of anomalous isotopes. However, such constraints are inevitably model dependent. In Ref. bcg it is shown that if the annihilation cross section for gluinos is nonperturbatively enhanced near threshold (many models of this type exist) then the relic density of gluinos could be very small. If, in addition, they did not cluster with nucleons (e.g. if they are concentrated at galaxy cores), then relic constraints would not rule out this scenario. Alternatively, the reheating required to avoid the Polonyi problem would also effectively eliminate the relic gluinos. It is also possible that the gluino could decay but with a lifetime so long that it is effectively stable in the detector. This is possible if there is a very weak violation of R-parity or in gauge-mediated-SUSY-breaking (GMSB) models where $`\stackrel{~}{g}g\stackrel{~}{G}`$ (where the $`\stackrel{~}{G}`$ is the gravitino) can be very suppressed by a large supersymmetric breaking scale. Thus, it is important to consider how to place constraints on a detector-stable gluino (for which we use the generic $`\stackrel{~}{g}`$-LSP notation) using accelerator experiments. We will focus on constraints that arise by looking directly for the $`\stackrel{~}{g}`$’s themselves. That is, we do not include processes where other supersymmetric particles are produced and then decay into $`\stackrel{~}{g}`$’s.
### III.2 Behavior of a $`\stackrel{~}{g}`$-LSP in a Detector
As soon as a $`\stackrel{~}{g}`$-LSP is produced in a detector, it picks up a gluon or quark-antiquark combination to form an ‘R-hadron’; $`R^0=\stackrel{~}{g}g`$ is likely to be the lightest state, but color-singlet $`\stackrel{~}{g}q^{}\overline{q}`$ states could have very similar mass, and if the difference in mass between such states and the $`R^0`$ were $`<m_\pi `$ they would be pseudo-stable in the detector. The behavior of a $`\stackrel{~}{g}`$-LSP in a typical detector depends very much upon whether the dominant R-hadron fragment is charged (probability $`P`$) or neutral (probability $`1P`$). Simple quark counting models suggest that $`P<1/2`$. The important ingredients in determining the energy deposited by the $`\stackrel{~}{g}`$-LSP in the detector are:
* the hadronic interaction length, $`\lambda _T`$, as determined by $`\sigma _T`$;
* the average energy deposited per hadronic collision, $`\mathrm{\Delta }E`$, as a function of the $`\stackrel{~}{g}`$’s velocity;
* the amount of ionization energy deposited between hadronic interactions and how the calibrated detector measures this energy;<sup>1</sup><sup>1</sup>1For example, in iron a given amount of $`E_{\mathrm{ionization}}`$ is translated into measured energy of $`E_{\mathrm{measured}}=rE_{\mathrm{ionization}}`$ with $`r1.6`$.
* the thickness (measured most conveniently by the number of hadronic interaction lengths) of various components of the detector.
One should picture the $`\stackrel{~}{g}`$-LSP as emerging from the hard production process as a neutral (charged) R-hadron with probability $`1P`$ ($`P`$). At each hadronic interaction the light quarks and gluons are presumed to be stripped away and the $`\stackrel{~}{g}`$-LSP again fragments into a neutral or charged R-hadron with probability $`1P`$ or $`P`$, respectively. Thus, for any $`P0,1`$, the charge of the R-hadron between hadronic collisions fluctuates. In Ref. bcg , several models for $`\lambda _T`$ (i.e. for the $`\stackrel{~}{g}`$-LSP total cross section) and for $`\mathrm{\Delta }E`$ are considered. For the most likely case of $`P<1/2`$, one finds that the energy deposited by the $`\stackrel{~}{g}`$-LSP is dominated by hadronic energy deposits rather than by ionization energy deposits. Further, for $`P<1/2`$ only a small fraction of the gluino’s energy is actually deposited. The $`\stackrel{~}{g}`$-LSP behaves like a bowling ball moving through a sea of ping-pong balls. The result is that the $`\stackrel{~}{g}`$-LSP generally exits the detector, thereby leading to missing energy aligned with a soft jet. For processes of interest, the $`\stackrel{~}{g}`$’s are always produced in pairs and are seldom back to back. As a result, the net missing energy is usually large and not aligned with any one of the jets observed in the detector. Thus, the crucial signal for $`P<1/2`$ is jets + missing energy. This signal is also generally quite useful even for $`P>1/2`$ since the momentum of a $`\stackrel{~}{g}`$-jet is never properly determined even if the energy deposited via ionization is large. (In fact, the ionization energy deposit is generally overestimated, and for $`P1`$ the ‘measured’ gluino momentum can even exceed its true momentum in the OPAL analysis procedure used later.)
### III.3 Constraints from LEP
Gluinos can be directly produced via two processes: $`e^+e^{}q\overline{q}\stackrel{~}{g}\stackrel{~}{g}`$ cer ; css ; mts , which can take place at tree-level, and $`e^+e^{}\stackrel{~}{g}\stackrel{~}{g}`$ nper ; krol ; css , which takes place via loop diagrams (involving squarks and quarks). The latter process is very model dependent and can be highly suppressed. Thus, we focus on the $`q\overline{q}\stackrel{~}{g}\stackrel{~}{g}`$ final state. An explicit calculation of the cross section for this final state reveals that LEP2 running will yield rather few $`q\overline{q}\stackrel{~}{g}\stackrel{~}{g}`$ events unless $`m_{\stackrel{~}{g}}`$ is quite small. However, the number of $`Z`$’s accumulated during LEP1 running is sufficiently large that a significant number of $`Zq\overline{q}\stackrel{~}{g}\stackrel{~}{g}`$ events would be expected for $`m_{\stackrel{~}{g}}<25\mathrm{GeV}`$.
The only relevant LEP1 experimental analysis is the OPAL opal search for pair production of neutralinos, $`Z\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$, with $`\stackrel{~}{\chi }_2^0q\overline{q}\stackrel{~}{\chi }_1^0`$, in the $`\mathrm{jets}+/p_T`$ channel that is potentially relevant for the $`q\overline{q}\stackrel{~}{g}\stackrel{~}{g}`$ final state. Typically, $`q\overline{q}\stackrel{~}{g}\stackrel{~}{g}`$ events give $`n(\mathrm{jets})=2`$, 3, or 4, depending upon the amount of energy deposition by the $`\stackrel{~}{g}`$-jets. After implementing the OPAL procedures in a detailed Monte Carlo simulation of $`q\overline{q}\stackrel{~}{g}\stackrel{~}{g}`$ production (including a parameterization of experimental resolutions and a Peterson form peterson for $`\stackrel{~}{g}`$ R-hadron fragmentation), we find that a $`\stackrel{~}{g}`$-LSP is excluded for $`3<m_{\stackrel{~}{g}}<24\mathrm{GeV}`$ if $`P<1/2`$. (For $`P>1/2`$, the excluded range only extends to $`23\mathrm{GeV}`$.) This is illustrated in Fig. 1 for our favored choices of $`\lambda _T`$ and $`\mathrm{\Delta }E`$ model. (The excluded mass range is quite insensitive to these choices.) We note that for $`P>1/2`$, a $`\stackrel{~}{g}`$-LSP can also be excluded by OPAL over much the same mass range by virtue of no excess of heavily ionizing tracks having been seen.
### III.4 Constraints from Tevatron Run I and Prospects for Run II
At a hadron collider, gluinos are produced via $`gg,q\overline{q}\stackrel{~}{g}\stackrel{~}{g}`$. Initial state radiation in association with the hard process yields additional jets in the final state. In Ref. bcg , we explored the limits that can be placed on a $`\stackrel{~}{g}`$-LSP using the jets$`+/p_T`$ analysis by CDF cdfcuts ; cdffinal of a portion of their Run I data. We performed a Monte Carlo simulation of $`\stackrel{~}{g}\stackrel{~}{g}`$ events for the CDF jets$`+/p_T`$ analysis cuts using ISAJET, supplemented by a routine that models the behavior of the $`\stackrel{~}{g}`$-LSP’s in the CDF detector for a given choice of $`P`$, $`\lambda _T`$ and $`\mathrm{\Delta }E`$ model. For $`P>0`$, one must discard events that contain a ‘muonic’ jet (i.e. a jet that has substantial ionization energy and minimal hadronic energy, as defined in the CDF analysis).
In Fig. 2 (upper window) the predicted jets$`+/p_T`$ cross section from $`\stackrel{~}{g}\stackrel{~}{g}`$ production (assuming $`P=1/2`$ and after all cuts and efficiencies) is compared to the 95% CL upper limit obtained by analyzing $`L=19\mathrm{pb}^1`$ of Run I data. The range $`20<m_{\stackrel{~}{g}}<130\mathrm{GeV}`$ is clearly excluded. For $`P0`$, the upper limit of the excluded mass range increases slowly to $`150\mathrm{GeV}`$. These results are quite independent of $`\lambda _T`$ and the $`\mathrm{\Delta }E`$ model. Note that the $`130150\mathrm{GeV}`$ lower limit on $`m_{\stackrel{~}{g}}`$ obtained is substantially below the lower limit that Run I data places on $`m_{\stackrel{~}{g}}`$ in a typical MSSM model. For easy comparison, the figure shows the cross section (after cuts) resulting from gluino pair production in the MSSM model considered in Ref. cdfcuts with $`m_{\stackrel{~}{q}}=1000\mathrm{GeV}`$, $`\mu =400\mathrm{GeV}`$ and $`\mathrm{tan}\beta =4`$; one sees that Run I data yields a 95% CL limit of roughly $`m_{\stackrel{~}{g}}>210\mathrm{GeV}`$. This is because the $`\stackrel{~}{g}`$-LSP scenario yields fewer jets (in particular, none from $`\stackrel{~}{g}`$ decay) as compared to a canonical MSSM scenario. For $`P3/4`$, the ionization energy deposited by a $`\stackrel{~}{g}`$ jet increases significantly, and for some $`\lambda _T`$ and $`\mathrm{\Delta }E`$ model choices the hadronic energy deposit is sufficiently small that one or both of the $`\stackrel{~}{g}`$ jets are often declared to be ‘muonic’ and the event discarded. For such choices, the current jets$`+/p_T`$ analysis does not constrain the $`\stackrel{~}{g}`$-LSP scenario. (Note that a modified analysis in which muonic jets are not discarded would, and is highly recommended.) However, the complementary CDF search for events with heavily ionizing tracks does exclude $`50\mathrm{GeV}m_{\stackrel{~}{g}}`$ (up to at least $`200\mathrm{GeV}`$ for any $`P3/4`$).
The expected extension of the excluded range of $`m_{\stackrel{~}{g}}`$ that will result for $`P=1/2`$ by analyzing Run II jets$`+/p_T`$ data in exactly the same way is also shown in Fig. 2. For a requirement of $`S/B>0.2`$ (as possibly needed for a reliable signal in the presence of systematic uncertainties), the limits obtained (for $`L>0.5\mathrm{fb}^1`$) will require $`m_{\stackrel{~}{g}}160180\mathrm{GeV}`$ for $`P1/2`$ (and also for $`P1/2`$ for many $`\lambda _T`$ and $`\mathrm{\Delta }E`$ model choices). If systematics could be controlled so that a signal with $`S/B<10\%`$ becomes reliable, the lower limits would be increased by about $`30\mathrm{GeV}`$. This, of course, is still substantially lower than the $`m_{\stackrel{~}{g}}`$ lower bound that can be achieved in the reference MSSM model for the same $`S/B`$ criterion (e.g. $`250\mathrm{GeV}`$ for $`S/B>0.2`$). It is worth emphasizing that these Run II limits do not disappear for large $`P`$ even for those $`\lambda _T`$ and $`\mathrm{\Delta }E`$ model choices that yield the largest probability for $`\stackrel{~}{g}`$ jets to be declared ‘muonic’. For example, for $`P=3/4`$, the worst choice would still result in excluding $`m_{\stackrel{~}{g}}<130\mathrm{GeV}`$ using Run II jets$`+/p_T`$ data and Run I analysis procedures. We anticipate that a more optimized analysis procedure (in particular not throwing away muonic jets) will do even better.
### III.5 Discussion and Conclusions
To summarize, accelerator data places quite significant constraints on a gluino LSP. Currently, for any reasonable value of the probability $`P`$ for $`\stackrel{~}{g}`$ charged R-hadron fragmentation ($`P1/2`$), $`3<m_{\stackrel{~}{g}}<130150\mathrm{GeV}`$ is excluded at 95% CL by a combination of OPAL LEP1 and CDF Run I jets$`+/p_T`$ analyses. For the theoretically much less likely $`P3/4`$ range, there is a window $`23<m_{\stackrel{~}{g}}<50\mathrm{GeV}`$ that (depending upon the hadronic path length of the $`\stackrel{~}{g}`$-LSP in the detector and the average energy deposited in each hadronic collision) might not be excluded by the jets$`+/p_T`$ analyses and would also not be excluded by the OPAL and CDF searches for heavily ionizing tracks. However, it is apparent that more optimized CDF procedures are capable of easily excluding this window. The increase in the lower bound on $`m_{\stackrel{~}{g}}`$ that will result from Run II Tevatron jets$`+/p_T`$ data will be limited by the level of systematic uncertainty in the absolute normalization of the background level.
For completeness, we also considered the scenario in which the gluino is not the LSP, but rather the NLSP (next-to-lightest supersymmetric particle), with the gravitino ($`\stackrel{~}{G}`$) being the (now invisible) LSP. Such a situation can arise in GMSB models, including that of Ref. raby . In this scenario, $`\stackrel{~}{g}g\stackrel{~}{G}`$. Early universe/rare isotope limits are then irrelevant. Further, the decay will be prompt from the detector point of view if $`m_{\stackrel{~}{G}}`$ is in the $``$ few eV region. (If the scale of supersymmetry breaking is so large that the $`\stackrel{~}{g}g\stackrel{~}{G}`$ decay lifetime is long enough that most $`\stackrel{~}{g}`$’s exit the detector before decaying, then we revert to the earlier $`\stackrel{~}{g}`$-LSP results.) For a $`\stackrel{~}{g}`$-NLSP, we find that the OPAL jets$`+/p_T`$ analysis excludes $`m_{\stackrel{~}{g}}26\mathrm{GeV}`$. The CDF Run I analysis excludes $`m_{\stackrel{~}{g}}240\mathrm{GeV}`$ (down to very low values), while Run II data can be expected to exclude at the very least $`m_{\stackrel{~}{g}}280\mathrm{GeV}`$ (assuming $`S/B>0.2`$ is required — better if smaller $`S/B`$ can be excluded).
Given that there are arguments in favor of a light gluino, it is unwise to simply assume that the gluino cannot be the lightest or next-to-lightest supersymmetric particle. Fortunately, present and future experiments can exclude or find such a gluino for a significant range of $`m_{\stackrel{~}{g}}`$.
## IV Detecting a highly degenerate lightest neutralino and lightest chargino at the Tevatron
J.F. Gunion, S. Mrenna
For some choices of soft SUSY–breaking parameters, the LSP is a stable neutralino $`\stackrel{~}{\chi }_1^0`$, the NLSP is a chargino $`\stackrel{~}{\chi }_1^\pm `$ almost degenerate in mass with the LSP ($`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}m_{\stackrel{~}{\chi }_1^\pm }m_{\stackrel{~}{\chi }_1^0}m_\pi `$few GeV), and all other sparticles are relatively heavy. In this case, detection of sparticles using the usual, mSUGRA–motivated signals will be difficult, since the visible decay products in $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0+\mathrm{}`$ will be very soft, and alternative signals must be considered. In this note, we summarize the viability of signatures at the Tevatron based on highly–ionizing charged tracks, disappearing charged tracks, large impact parameters, missing transverse energy and a jet or a photon, and determine the mass reach of such signatures assuming that only the $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\chi }_1^0`$ are light. If $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ is sufficiently big that $`c\tau (\stackrel{~}{\chi }_1^\pm )<`$few cm and there are no other light superparticles, there is a significant possibility that the limits on $`m_{\stackrel{~}{\chi }_1^\pm }`$ based on LEP2 data cannot be extended at the Tevatron. If $`c\tau (\stackrel{~}{\chi }_1^\pm )>`$few cm, relatively background–free signals exist that will give a clear signal of $`\stackrel{~}{\chi }_1^\pm `$ production (for some range of $`m_{\stackrel{~}{\chi }_1^\pm }`$).
### IV.1 Introduction
In mSUGRA models, gaugino masses are assumed to be universal at $`M_U`$, leading to $`M_3:M_2:M_16:2:1`$ at the TeV energy scale, implying relatively large mass splitting between the lightest chargino and the lightest neutralino (most often the LSP). However, various attractive models exist for which $`M_2<M_1`$ at the TeV scale, which results in $`m_{\stackrel{~}{\chi }_1^\pm }m_{\stackrel{~}{\chi }_1^0}`$ if $`|\mu |M_{1,2}`$ (as is normally the case for correct EWSB). In particular, $`M_2<M_1`$ when the gaugino masses are dominated by or entirely generated by loop corrections. The first model of this type to receive detailed attention was the O-II superstring model proposed in Ref. ibanez and studied in Refs. cdg1 ; cdg2 ; cdg3 . For further review, see Ref. bcgrunii2 . More recently, the one-loop boundary conditions have arisen in the context of the conformal anomaly murayama ; randall .
In the O-II model, $`M_{1,2,3}`$ are determined both by the one-loop beta functions and by the Green-Schwarz mixing parameter, $`\delta _{GS}`$ (required in general to cancel anomalies). At the scale $`M_U`$ ($`m_Z`$), the O–II model with $`\delta _{GS}=4`$ yields $`M_3:M_2:M_1=1:5:10.6`$ ($`6:10:10.6`$); $`\delta _{GS}=0`$ (equivalent to the simplest version of the conformal anomaly approach) gives $`M_3:M_2:M_1=3:1:33/5`$ ($`3:0.3:1`$). After radiative corrections, $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ is near in value to $`m_\pi `$ for $`\delta _{GS}=0`$, increasing to $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}12`$ GeV (depending on $`|\mu |`$) for $`\delta _{GS}=4`$. Since the typical values of $`|\mu |`$ required by RGE electroweak symmetry breaking are large, the higgsino $`\stackrel{~}{\chi }_2^\pm `$, $`\stackrel{~}{\chi }_3^0`$ and $`\stackrel{~}{\chi }_4^0`$ states are very heavy.
When $`M_2<M_1|\mu |`$, the $`Z\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0`$, $`Z\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$, $`Z\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$, and $`W^\pm \stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_2^0`$ couplings are all small, while the $`Z,\gamma \stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ and $`W^\pm \stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0`$ couplings are full strength. Only cross sections induced by the latter can have large rates.
The most critical ingredients in the phenomenology of such models are the lifetime and the decay modes of the $`\stackrel{~}{\chi }_1^\pm `$, which in turn depend almost entirely on $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ when the latter is small. The $`c\tau `$ and branching ratios of the $`\stackrel{~}{\chi }_1^\pm `$ as a function of $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ have been computed in Ref. cdg3 . A tabulation of $`c\tau `$ values for the range of $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ of interest in this report is given in Table 2. For $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<m_\pi `$, $`c\tau `$ is at least several meters; once $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>m_\pi `$, $`c\tau `$ drops quickly. For $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<m_\pi `$, $`\stackrel{~}{\chi }_1^\pm e\nu _e\stackrel{~}{\chi }_1^0`$ dominates, while for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>m_\pi `$ the $`\stackrel{~}{\chi }_1^\pm \pi ^\pm \stackrel{~}{\chi }_1^0`$ mode turns on and is dominant for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<800\mathrm{MeV}`$. For still larger $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$, multi-pion modes become important merging eventually into $`\stackrel{~}{\chi }_1^\pm \mathrm{}\nu _{\mathrm{}}\stackrel{~}{\chi }_1^0,q^{}\overline{q}\stackrel{~}{\chi }_1^0`$.
We now give a brief review of the results of Ref. hithip for the case where we assume that only the $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\chi }_1^0`$ are light. We discuss the types of signals that will be important for different ranges of $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$. (The modifications that arise if the $`\stackrel{~}{g}`$ is also light, with $`m_{\stackrel{~}{g}}>m_{\stackrel{~}{\chi }_1^\pm }m_{\stackrel{~}{\chi }_1^0}`$, are discussed in Ref. hithip .) For this discussion, we ask the reader to imagine a canonical detector (e.g. CDF or DØ at Run II) with the following components ordered according to increasing radial distance from the beam.
(I) An inner silicon vertex (SVX) detector extending radially from the beam axis. The CDF Run II vertex detector has layers at $`r1.6`$, 3, 4.5, 7, 8.5 and $`11\text{cm}`$ (the first and second layers are denoted L00 and L0, respectively) extending out to $`z=\pm 45`$ cm cdfextra . The DØ SVX has 4 layers (but 2 are double–sided), with the first at 2.5 cm and the last at 11 cm.
(II) A central tracker (CT) extending from $`15\text{cm}`$ to $`73\text{cm}`$ (DØ) or from roughly $`20\text{cm}`$ to $`130\text{cm}`$ (CDF).
(III) A thin pre–shower layer (PS).
(IV) An electromagnetic calorimeter (EC) and hadronic calorimeter (HC).
(V) The inner–most muon chambers (MC), starting just beyond the HC.
(VI) Both CDF and DØ will have a precise time–of–flight measurement (TOF) for any charged particle that makes it to the muon chambers.
It is important to note that the SVX, CT and PS can all give (independent) measurements of the $`dE/dx`$ from ionization of a track passing through them. This will be important to distinguish a heavily–ionizing chargino (which would be $``$ twice minimal ionizing \[2MIP\] for $`\beta \gamma 0.85`$) from an isolated minimally ionizing particle \[1MIP\]. For example, at DØ the rejection against isolated 1MIP tracks will be $`\mathrm{few}\times 10^3`$, $`\mathrm{few}\times 10^3`$, and $`10^1`$ for tracks that pass through the SVX, CT and PS, respectively, with an efficiency of 90% for tracks with $`\beta \gamma <0.85`$ glandsberg . At CDF, the discrimination factors are roughly similar dstuart . Because of correlations, one cannot simply multiply these numbers together to get the combined discrimination power of the SVX, CT and PS for an isolated track that passes through all three; for such a track with $`\beta \gamma <0.85`$, the net discrimination factor would probably be of order $`\mathrm{few}\times 10^5`$. At LEP/LEP2, the detector structure is somewhat different and important features will be noted where relevant. We now list the possible signals.
(a) LHIT and TOF signals:
For $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<m_\pi `$, a heavy chargino produced in a collision travels a distance of order a meter or more and will often penetrate to the muon chambers. If it does, the chargino may be distinguished from a muon by heavy ionization in the SVX, CT and PS. There should be no hadronic energy deposits associated with this track, implying that the energy deposited in the hadronic calorimeter should be consistent with ionization energy losses for the measured $`\beta `$. This type of long, heavily–ionizing track signal will be denoted as an LHIT signal.
If the chargino penetrates to the muon chambers, its large mass will also be evident from the time delay of its TOF signal. This delay can substitute for the heavy ionization requirement. The passage of the chargino through the muon chamber provides an adequate trigger for the events. In addition, the chargino will be clearly visible as an isolated track in the CT, and this track could also be used to trigger the event. In later analysis (off–line even), substantial momentum can be required for the track without loss of efficiency. (The typical transverse momentum of a chargino when pair–produced in hadronic collisions is of order 1/2 the mass.)
After a reasonable cut on the $`p_T`$ of the chargino track, the LHIT and TOF signals will be background free.
(b) DIT signals:
For $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ above but near $`m_\pi `$, the chargino will often appear as an isolated track in the central tracker but it will decay before the muon chamber. (The appropriate mass range for which this has significant probability is $`m_\pi <\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<145\mathrm{MeV}`$, for which $`c\tau >17\text{cm}`$.) As such a chargino passes part way through the calorimeters beyond the CT, it will deposit little energy. In particular, any energy deposit in the hadronic calorimeter should be no larger than that consistent with ionization energy deposits for the $`\beta `$ of the track as measured using ionization information from the SVX+CT+PS. (In general, the chargino will only deposit ionization energy up to the point of its decay. Afterwards, the weakly–interacting neutralino will carry away most of the remaining energy, leaving only a very soft pion or lepton remnant.) Thus, we require that the track effectively disappear once it exits the CT. (The point at which the ionization energy deposits end would typically be observable in a calorimeter with sufficient radial segmentation, but we do not include this in our analysis.) Such a disappearing, isolated track will be called a DIT. The DIT will have substantial $`p_T`$, which can be used to trigger the event. A track with large $`p_T`$ from a background process will either be a hadron, an electron or a muon. The first two will leave large deposits in the calorimeters (EC and/or HC) and the latter will penetrate to the muon chamber. Thus, the signal described is very possibly a background–free signal. If not, a requirement of heavy ionization in the SVX, CT and PS will certainly eliminate backgrounds, but with some sacrifice of signal events. Thus, we will also consider the possibility of requiring that the DIT track be heavily ionizing. In the most extreme case, we require that the average ionization measured in the SVX, CT and PS correspond to $`\beta <0.6`$ ($`\beta \gamma <0.75`$), which signal is denoted by DIT6. For a DIT signal, this is a very strong cut once $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ is large enough that the average $`c\tau `$ is smaller than the radius of the CT. This is because rather few events will have both large enough $`\beta \gamma c\tau `$ to pass all the way through the CT and small enough $`\beta `$ to satisfy the heavy ionization requirement.
(c) STUB and KINK signals, including STUB$`+/E_T`$, or SMET signal:
For $`145\mathrm{MeV}<\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<160\mathrm{MeV}`$, $`17\text{cm}>c\tau >7\text{cm}`$. For such $`c\tau `$, the probability for the chargino to pass all the way through the central tracker will be small. The chargino will be most likely to pass all the way through the SVX and decay somewhere in the CT. Such a short SVX track we term a STUB. It will not be associated with any calorimeter energy deposits. At a hadron collider, the primary difficulty associated with a STUB signal is that it will not provide its own Level–1 trigger. We have found that it is most efficient ($`ϵ10\%`$) to trigger the event by requiring substantial missing transverse energy ($`/E_T`$). Once an interesting event is triggered, off–line analysis will provide a measurement of the ionization deposited by the STUB in the SVX.<sup>2</sup><sup>2</sup>2Note that for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>160\mathrm{MeV}`$ (for which $`c\tau <7`$ cm), requiring heavy ionization, i.e. small $`\beta `$, begins to significantly conflict with the requirement that the chargino pass all the way through the SVX. For smaller $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$, the $`c\tau `$ of the chargino is larger and this conflict is not very severe. Although we believe that the STUB signals will be background–free without a $`\beta `$ cuts, we have also considered discovery reach after imposing a $`\beta <0.6`$ cut. Altogether, we will define 4 STUB–based signals: (a) SNT – a STUB track only, with no other trigger; (b) SNT6 – a STUB track with $`\beta <0.6`$ and no other trigger. (c) SMET – a STUB track in an event with $`/E_T>35\mathrm{GeV}`$; and, (d) SMET6 – a STUB track with $`\beta <0.6`$ in an event with $`/E_T>35\mathrm{GeV}`$. Only (c) and (d) are possible using the triggering design planned by CDF and DØ for Run II.
In addition to the STUB, most of the soft charged pions from charginos that decay after passing through the vertex detector will be seen in the tracker. Typically, the soft pion track that intersects the STUB track will do so at a large angle, a signature we call a KINK. We have not explored this in detail, but believe that a KINK requirement in association with the STUB signals defined above would lead to little loss of signal and yet make the signals background–free with high certainty.
(d) HIP signals:
For $`160\mathrm{MeV}<\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<190\mathrm{MeV}`$, $`7\text{cm}>c\tau >3\text{cm}`$. Some of the produced charginos will decay late compared to $`c\tau `$ and yield a STUB signature of the type discussed just above. More typically, however, the $`\stackrel{~}{\chi }_1^\pm `$ will pass through two to three layers of the SVX. The $`\stackrel{~}{\chi }_1^\pm `$ track will then end and turn into a single charged pion with substantially different momentum. Both the sudden disappearance of and the lack of any calorimeter energy deposits associated with the $`\stackrel{~}{\chi }_1^\pm `$ track will help to distinguish it from other light–particle tracks that would normally register in all layers of the SVX and in the calorimeters.
For $`160\mathrm{MeV}<\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<190\mathrm{MeV}`$, $`p_\pi ^{}77130\mathrm{MeV}`$. The corresponding transverse impact parameter resolution of the SVX, $`\sigma _b`$, is approximately $`300170\mu `$m (taking $`p_\pi ^Tp_\pi ^{}`$ and applying the $`1\sigma `$ values from Fig. 2.2 of cdfextra when L00 is included), and is much smaller than the typical impact parameter (which is a sizeable fraction of $`c\tau >3\text{cm}`$). In this range, the KINK formed by the $`\stackrel{~}{\chi }_1^\pm `$ track and the soft $`\pi `$ should be visible. In addition, the layers that the $`\stackrel{~}{\chi }_1^\pm `$ passes through will provide an ionization estimate for $`\beta `$ that could be used to help eliminate backgrounds. However, we have not pursued either of these possibilities since in the end the STUB signals are still viable in this $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ range and are probably superior.
For $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>230\mathrm{MeV}`$, $`c\tau <1.6\text{cm}`$ and the typical $`\stackrel{~}{\chi }_1^\pm `$ will not even pass through the innermost SVX layer unless $`\beta `$ is very large . However, $`p_\pi ^{}>180\mathrm{MeV}`$ and the impact parameter resolution for the single emitted pion moves into the $`<150\mu `$m range. For example, if $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}=240,300,500,1000\mathrm{MeV}`$, $`c\tau 1.2,0.37,0.09,0.007\text{cm}`$ while $`p_\pi ^T195,265,480,990\mathrm{MeV}`$ yields $`1\sigma `$ impact parameter resolutions of $`\sigma _b120,90,50,25\mu `$m. We thus considered the signal based on events defined by the $`\gamma +/E_T`$ trigger and the presence of one or more large–$`b`$ charged pions.
Unlike the previous signals, the HIP signal has a large background even after requiring the $`\pi `$’s impact parameter $`b`$ to satisfy $`b>5\sigma _b`$, where $`\sigma _b`$ is the resolution. After imposing isolation criteria and low-$`p_T`$ for the $`\pi `$, the main background arises from production of long–lived baryons (e.g. $`\mathrm{\Sigma }^\pm ,\mathrm{}`$) that decay to a $`\pi `$ and a nucleon.
Once $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>1\mathrm{GeV}`$, a HIP signature will not be useful and we must consider the chargino decay to be prompt. This is because the largest possible impact parameter is only a few times the $`1\sigma `$ value for the resolution and we will be dominated by fakes. This leads us to one of two completely different types of signal.
(e) $`\gamma \mathbf{+}\mathbf{/}E_T`$ and jets+$`\mathbf{/}E_T`$ signals:
For some interval of $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ (e.g. $`200\mathrm{MeV}<\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<300\mathrm{MeV}`$ at the DELPHI LEP/LEP2 detector — see later — or, perhaps, $`1\mathrm{GeV}<\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<1020\mathrm{GeV}`$ at the Tevatron) the decay products (hadron(s) or $`\mathrm{}\nu `$) produced along with the $`\stackrel{~}{\chi }_1^0`$ will be too soft to be distinctively visible in the main part of the detector and at the same time high–impact–parameter tracks associated with chargino decay will not be apparent. One will then have to detect chargino production as an excess of events with an isolated photon or missing energy above a large $`\gamma +/E_T`$ or jet(s)+$`/E_T`$ background. We find that, despite its lower rate, after appropriate cuts the $`\gamma +/E_T`$ channel is superior to the jet(s)+$`/E_T`$ signals. For some values of the chargino mass and $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$, an excess in these channels could confirm the SVX signals discussed earlier.
(f) standard mSUGRA signals:
For large enough $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$, the extra lepton or hadron tracks from $`\stackrel{~}{\chi }_1^\pm `$ decay will be sufficiently energetic to be detected and will allow identification of chargino production events when associated with a photon or missing energy trigger. A detailed simulation is required to determine exactly how large $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ needs to be for this signal to be visible above backgrounds. At LEP/LEP2, backgrounds are sufficiently small that the extra tracks are visible for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>300\mathrm{MeV}`$ in association with a photon trigger while standard mSUGRA searches based on missing energy and jets/leptons require $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>3\mathrm{GeV}`$. At a hadron collider we estimate that $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>1020`$ GeV will be necessary to produce leptons or jets sufficiently energetic to produce a distinctive event assuming a missing energy trigger.
### IV.2 Collider Phenomenology of degenerate models
Although our main focus will be on Tevatron Run II, it is useful to summarize which of the above signals have been employed at LEP2 and the resulting constraints on the degenerate scenarios we are considering.
#### IV.2.1 Lepton Colliders
As discussed above and in Refs. cdg1 ; cdg2 ; cdg3 , collider phenomenology depends crucially on $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$. Most importantly, SUSY detection depends on which aspects (if any) of the $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ final state are visible. If the $`\stackrel{~}{\chi }_1^\pm `$ decay products are soft, $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ production may be indistinguishable from the large $`e^+e^{}e^+e^{}\gamma \gamma e^+e^{}+\mathrm{soft}`$ background. Tagging $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ production using a photon from initial or final state radiation (ISR) is necessary cdg1 . Even with an ISR tag, the $`\stackrel{~}{\chi }_1^+`$ and $`\stackrel{~}{\chi }_1^{}`$ might be invisible because of the softness of their decay products and the lack of a vertex detector signal. In this case, $`\gamma \stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ production is observable only as a $`\gamma /M`$ signature, where $`/M=\sqrt{(p_e^{}+p_{e^+}p_\gamma )^2}`$. Even after requiring $`/M>m_Z`$, the $`\gamma \nu \overline{\nu }`$ background is large. For $`M_2<M_1|\mu |`$, Ref. cdg1 found that at LEP2 with $`L=125\mathrm{pb}^1`$ per experiment, no improvement over the $`m_{\stackrel{~}{\chi }_1^\pm }<45\mathrm{GeV}`$ limit coming from LEP1 $`Z`$–pole data was possible. The experimental situation is greatly improved if LHIT and/or KINK signals can be employed, or if the soft pions from the $`\stackrel{~}{\chi }_1^\pm `$ decays in $`\gamma \stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ events can be detected. All of this is most clearly illustrated by summarizing the analysis from DELPHI at LEP2 delphideg .
* When $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<200\mathrm{MeV}`$, the charginos are sufficiently long–lived to produce either an LHIT or a KINK signal for $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ production. No additional trigger is required for either signal. As a result, DELPHI is able to exclude $`m_{\stackrel{~}{\chi }_1^\pm }`$ out to nearly the kinematic limit (currently 90 GeV).
* When $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>`$ 3 GeV, the decay products of the $`\stackrel{~}{\chi }_1^\pm `$ become easily visible, and the standard mSUGRA search results apply; the $`\stackrel{~}{\chi }_1^\pm `$ is excluded out to the kinematic limit (90 GeV for the data sets analyzed), except for the case of a relatively light sneutrino, for which the $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ cross section is smaller and the limit does not extend past 75 GeV.
* For $`200\mathrm{MeV}<\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<3\mathrm{GeV}`$, the chargino tracks are not long enough to use the KINK signature, and the chargino decay products are too soft to provide a clear signature on their own. As proposed in Ref. cdg1 , DELPHI employs an ISR photon tag. In order to essentially eliminate the $`\gamma \nu \overline{\nu }`$ background, the event is required to contain soft charged tracks consistent with the isolated pions expected from the chargino decays. DELPHI observed no events after all cuts. For $`M_2<M_1|\mu |`$, and a heavy (light) sneutrino, this excludes $`m_{\stackrel{~}{\chi }_1^\pm }<62\mathrm{GeV}`$ ($`49\mathrm{GeV}`$) for $`0.3<\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<3\mathrm{GeV}`$ ($`0.5<\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<3\mathrm{GeV}`$). The gap from $`0.20.3\mathrm{GeV}`$ ($`0.20.5\mathrm{GeV}`$) arises because of the low efficiency for detecting very soft pions.<sup>3</sup><sup>3</sup>3With the ISR tag, the $`\gamma \gamma `$ background is completely negligible.
Thus, there is a gap from just above $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}200\mathrm{MeV}`$ to at least $`300\mathrm{MeV}`$ for which the chargino is effectively invisible. DELPHI finds that the $`\gamma /M`$ signature, discussed earlier, is indeed insufficient to improve over the $`m_{\stackrel{~}{\chi }_1^\pm }>45\mathrm{GeV}`$ limit from $`Z`$ decays. We are uncertain whether DELPHI explored the use of high–impact–parameter tracks in their vertex detector (in association with the ISR trigger) to improve their sensitivity (by sharply reducing the $`\gamma \nu \overline{\nu }`$ background) in these gap regions.
#### IV.2.2 Hadron Colliders
At hadron colliders, typical signatures of mSUGRA are tri–lepton events from neutralino–chargino production, like–sign di–leptons from gluino pair production, and multi-jets$`+/E_T`$ from squark and gluino production. The tri–lepton signal from $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_2^0`$ production and the like–sign di–lepton signal from $`\stackrel{~}{g}\stackrel{~}{g}`$ production are both suppressed when $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ is small by the softness of the leptons coming from the $`\stackrel{~}{\chi }_1^\pm `$ decay(s). In $`M_2<M_1|\mu |`$ scenarios, the tri–lepton signal is further diminished by the suppression of the $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_2^0`$ cross section. Provided that $`m_{\stackrel{~}{g}}`$ is light enough, the most obvious signal for SUSY in degenerate models is jet(s) plus missing energy, as studied in Refs. cdg2 ; hithip . However, it is entirely possible that the gluino is much heavier than the light $`\stackrel{~}{\chi }_1^\pm ,\stackrel{~}{\chi }_1^0`$ states and that the $`\stackrel{~}{g}\stackrel{~}{g}`$ production rate (at the Tevatron at least) will be quite suppressed. In this case, the ability to detect events in which the only directly produced SUSY particles are light neutralino and chargino states could prove critical. In what follows, we assume that the sfermion, gluino and heavier chargino and neutralino states are sufficiently heavy that their production rates at the Tevatron are not useful, and investigate methods to probe $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ and $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0`$ production at the Tevatron. The possible signals were summarized earlier. Details regarding cuts and triggering appear in Ref. hithip .
We performed particle level studies using either the processes contained in the PYTHIA 6.125 event generator or by adding external processes (several of the $`\gamma +X`$ processes considered here) into PYTHIA. A calorimeter is defined out to $`\eta =4.4`$ with a Gaussian $`E_T`$ resolution of $`\sigma _{E_T}=80\%/\sqrt{E_T}`$. Jets with $`E_T>5`$ GeV and $`R=0.5`$ are reconstructed to define $`/E_T`$. Non–Gaussian contributions will be estimated as described later. Charged track momenta and impact parameters $`b`$ are unsmeared, but the effects of detector resolution on $`b`$ are included.
We find that there is a natural boundary near a mass splitting of $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}300\mathrm{MeV}`$, below which one or more of the background–free signals are viable but above which one must contend with large backgrounds.
Region (A) For $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ values $`<200300\mathrm{MeV}`$, one considers the background–free signals summarized above, which will have the most substantial mass reach in $`m_{\stackrel{~}{\chi }_1^\pm }`$. The $`L=2\mathrm{fb}^1`$ and $`L=30\mathrm{fb}^1`$ 95% CL (3 events, no background) limits on $`m_{\stackrel{~}{\chi }_1^\pm }`$ deriving from these signals are summarized in Fig. 3. We give a brief verbal summary.
$`𝚫𝒎_{\stackrel{\mathbf{~}}{𝝌}_\mathrm{𝟏}}\mathbf{<}𝒎_𝝅`$: For such $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$, the average $`c\tau `$ of the chargino is of order a meter or more. The LHIT and TOF signals are prominent, but the DIT and STUB signals appear if $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ is not extremely small. The relative weight between these signals is determined by the exponential form of the $`c\tau `$ distribution in the chargino rest frame and the event–by–event variation of the boosts imparted to the chargino(s) during production.
* The LHIT signature can probe masses in the range $`260325`$ ($`380425)`$ GeV for $`L=2\mathrm{fb}^1`$ ($`30\mathrm{fb}^1`$), the lower reach applying for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}m_\pi `$ and the highest reach applying for any $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<125\mathrm{MeV}`$. The reach of the TOF signature is nearly identical to that of the LHIT signature.
* The DIT signature has a reach of 320 (425) GeV for $`120\mathrm{MeV}\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}m_\pi `$,<sup>4</sup><sup>4</sup>4We did not study lower $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ values since they are highly improbable after including radiative correction contributions to $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$. and, in particular, is more efficient than the LHIT and TOF signals for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}m_\pi `$. The DIT signature reach drops by about 20 GeV with a $`\beta <0.6`$ cut (DIT6) designed to require that the chargino track be heavily–ionizing.
* The STUB signature with no additional trigger (SNT) can reach to $`340`$ (450) GeV for $`120\mathrm{MeV}\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}m_\pi `$, which mass reach drops by $`1020`$ GeV if $`\beta <0.6`$ is required. However, neither DØ nor CDF can use STUB information at Level–1 in their current design.
* With the addition of a standard $`/E_T`$ trigger, the resulting STUB signature (SMET) will be viable with the present detectors, reaching to $`240260`$ ($`350375`$) GeV for $`120\mathrm{MeV}\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}m_\pi `$, which numbers drop by about 10 GeV if $`\beta <0.6`$ is required (SMET6).
$`𝒎_𝝅\mathbf{<}\mathbf{}𝚫𝒎_{\stackrel{\mathbf{~}}{𝝌}_\mathrm{𝟏}}\mathbf{<}\mathbf{}\mathrm{𝟐𝟎𝟎}\mathbf{}\mathrm{𝟑𝟎𝟎}\mathrm{𝐌𝐞𝐕}`$:
* The LHIT and TOF signatures disappear, since almost all produced charginos decay before reaching the MC or TOF.
* The DIT signature remains as long as the $`\beta <0.6`$ (heavily–ionizing) requirement is not necessary to eliminate backgrounds. If we require $`\beta <0.6`$, there is a mismatch with the requirement that the chargino pass through the CT – once $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ is above $`145\mathrm{MeV}`$, the entire signal is generated by large boosts in the production process, which is in conflict with requiring small $`\beta `$.
* The SNT signature probes $`m_{\stackrel{~}{\chi }_1^\pm }<300\mathrm{GeV}`$ ($`<400\mathrm{GeV}`$) for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}m_\pi `$ and $`L=2\mathrm{fb}^1`$ ($`L=30\mathrm{fb}^1`$). For $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ as large as $`300\mathrm{MeV}`$, it alone among the background–free channels remains viable, probing $`m_{\stackrel{~}{\chi }_1^\pm }<70\mathrm{GeV}`$ ($`<95\mathrm{GeV}`$). Certainly, it would extend the $`90\mathrm{GeV}`$ limit obtained by DELPHI at LEP2 that applies for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<200\mathrm{MeV}`$ and the $`45\mathrm{GeV}`$ limit from LEP data that is the only available limit for $`200\mathrm{MeV}\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}300\mathrm{MeV}`$. But, as stated above, the SNT signature will not be possible without a Level–1 SVX trigger.
* The STUB+$`/E_T`$, SMET and SMET6 signatures are fully implementable at Run II and have a reach that is only about 20 GeV lower than their SNT and SNT6 counterparts.
Region (B) For $`300\mathrm{MeV}<\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}<600\mathrm{MeV}`$, the high–impact–parameter (HIP) signal (a $`\gamma +/E_T`$ tag for events yields the smallest backgrounds) is very useful despite the large background from production of $`\mathrm{\Sigma }^\pm ,\mathrm{}`$ hadrons. The luminosity required to achieve 95% CL exclusion or $`5\sigma `$ discovery was evaluated in Ref. hithip , requiring also that $`S/B>0.2`$. We find that one can achieve a 95% CL lower bound of $`95\mathrm{GeV}`$ ($`75\mathrm{GeV}`$) on $`m_{\stackrel{~}{\chi }_1^\pm }`$ for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}=300\mathrm{MeV}`$ ($`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}=600\mathrm{MeV}`$) for $`L=30\mathrm{fb}^1`$. This would represent some improvement over the $`60\mathrm{GeV}`$ lower bound obtained in the current DELPHI analysis of their LEP2 data for this same range of $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ if the sneutrino is heavy. (If the $`\stackrel{~}{\nu }`$ is light, then there is no useful LEP2 limit if $`300\mathrm{MeV}\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}500\mathrm{MeV}`$, but LEP data requires $`m_{\stackrel{~}{\chi }_1^\pm }>45\mathrm{GeV}`$.) With only $`L=2\mathrm{fb}^1`$ of data, the HIP analysis would only exclude $`m_{\stackrel{~}{\chi }_1^\pm }<68\mathrm{GeV}`$ ($`<53\mathrm{GeV}`$) for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}=300\mathrm{MeV}`$ ($`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}=600\mathrm{MeV}`$).
Region (C) For $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>600\mathrm{MeV}`$, up to some fairly large value (we estimate at least 10 to 20 GeV), the chargino decay products are effectively invisible at a hadron collider and the most useful signal is $`\gamma +/E_T`$. However, this signal at best probes $`m_{\stackrel{~}{\chi }_1^\pm }<60\mathrm{GeV}`$ (for any $`L>2\mathrm{fb}^1`$), whereas the DELPHI analysis of their LEP2 data already excludes $`m_{\stackrel{~}{\chi }_1^\pm }60\mathrm{GeV}`$ for $`500\mathrm{MeV}\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}3\mathrm{GeV}`$ (if the sneutrino is heavy — only $`48\mathrm{GeV}`$ if the sneutrino is light) and $`m_{\stackrel{~}{\chi }_1^\pm }90\mathrm{GeV}`$ for $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>3\mathrm{GeV}`$.
An overall summary of the signals and their mass reach at the Tevatron for detecting $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ and $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0`$ production in the $`M_2<M_1|\mu |`$ scenario appears in Table 3. Clearly, the very real possibility that $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}300\mathrm{MeV}`$ would present us with a considerable challenge.
For purposes of comparison, we note that in an mSUGRA scenario the tri–lepton signature from $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_2^0`$ production allows one to probe chargino masses up to about $`160\mathrm{GeV}`$ for $`L=30\mathrm{fb}^1`$ when the scalar soft–SUSY–breaking mass is large trilepref .
Finally, we wish to note that the precise values of $`m_{\stackrel{~}{\chi }_1^\pm }`$ and $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ will be of significant theoretical interest. $`m_{\stackrel{~}{\chi }_1^\pm }`$ will be determined on an event–by–event basis if the chargino’s momentum and velocity can both be measured. This will be possible for the LHIT, TOF, DIT, and STUB signals by combining tracking information with ionization information. (Note that, in all these cases, accepting only events roughly consistent with a given value of $`m_{\stackrel{~}{\chi }_1^\pm }`$ will provide further discrimination against backgrounds.) However, for the HIP and $`\gamma +/E_T`$ signals $`m_{\stackrel{~}{\chi }_1^\pm }`$ can only be estimated from the absolute event rate. As regards $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$, it will be strongly constrained by knowing which signals are present and their relative rates. In addition, if the the soft charged pion can be detected, its momentum distribution, in particular the end–point thereof, would provide an almost direct determination of $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$.
For large $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ ($`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}>2030\mathrm{GeV}\mathrm{?}`$), one should explore the potential of the tri–lepton signal coming from $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_2^0`$ production. However, this is a suppressed cross section when both the lightest neutralino and lightest chargino are wino–like. Standard mSUGRA studies do not apply without modification; the cross section must be rescaled and the lepton acceptance recalculated as a function of $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$. A detailed study is required to determine the exact mass reach as a function of $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$.
Of course, additional SUSY signals will emerge if some of the squarks, sleptons and/or sneutrinos are light enough (but still heavier than the $`\stackrel{~}{\chi }_1^\pm `$) that their production rates are substantial. In particular, leptonic signals from the decays \[e.g. $`\stackrel{~}{\mathrm{}}_L^\pm \mathrm{}^\pm \stackrel{~}{\chi }_1^0`$ or $`\stackrel{~}{\nu }_{\mathrm{}}\mathrm{}^\pm \stackrel{~}{\chi }_1^{}`$\] would be present.
Given the possibly limited reach of the Tevatron when the lightest neutralino and chargino are nearly degenerate, it will be very important to extend these studies to the LHC. A particularly important issue is the extent to which the large $`c\tau `$ tails of the $`\stackrel{~}{\chi }_1^\pm `$ decay distributions can yield a significant rate in the background–free channels studied here. Hopefully, as a result of the very high event rates and boosted kinematics expected at the LHC, the background–free channels will remain viable for significantly larger $`\mathrm{\Delta }m_{\stackrel{~}{\chi }_1}`$ and $`m_{\stackrel{~}{\chi }_1^\pm }`$ values than those to which one has sensitivity at the Tevatron. In this regard, a particularly important issue for maximizing the mass reach of these channels will be the extent to which tracks in the silicon vertex detector and/or in the central tracker can be used for triggering in a high–luminosity enviroment.
While finalizing the details of this study, other papers randall2 ; wells appeared on the same topic. Some of the signatures discussed here are also considered in those papers. Our studies are performed at the particle level and contain the most important experimental details.
## V Superheavy Supersymmetry
S. Ambrosanio, J. D. Wells
In the vast space of all viable physics theories, supersymmetry (SUSY) is not a point. Any theory can be “supersymmetrized” almost trivially, and the infinite array of choices for spontaneous SUSY breaking just increases the scope of possibilities in the real world. One thing that appears necessary, if SUSY has anything to do with nature, is superpartners for the standard model particles that we already know about: leptons, neutrinos, quarks, and gauge bosons. These superpartners must feel SUSY breaking and a priori can have arbitrary masses as a result.
Phenomenologically, the masses cannot be arbitrary. There are several measurements that have been performed that effectively limit what the SUSY masses can be. First, there are direct limits on $`Z\mathrm{SUSY}`$, for example, that essentially require all superpartners to be above $`m_Z/2`$. Beyond this, collider physics limits become model dependent, and it is not easy to state results simply in terms of the mass of each particle. Second, comparing softly broken SUSY model calculations with flavor changing neutral current (FCNC) measurements implies that superpartner masses cannot be light and arbitrary. And finally, requiring that the $`Z`$ boson mass not result from a fine-tuned cancellation of big numbers requires some of the particles masses be near $`m_Z`$ (less than about $`1\text{ TeV}`$, say).
Numerous explanations for how the above criteria can be satisfied have been considered. Universality of masses, alignment of flavor matrices, flavor symmetries, superheavy supersymmetry, etc., have all been incorporated to define a more or less phenomenologically viable explanation of a softly broken SUSY description of nature.
In this contribution, we would like to summarize some of the basic collider physics implications of superheavy supersymmetry (SHS) at the Tevatron. Our understanding is that analyses of all the specific processes that are mentioned here in principle are being pursued within other subgroups. Therefore, our goal in this submission is to succinctly explain what SHS is and how some of the observables being studied within other contexts could be crucial to SHS. We also hope that by enumerating some of the variations of this approach that this contribution could help us anticipate and interpret results after discovery of SUSY, and help distinguish between theories. The idea we are discussing goes under several names including “decoupling supersymmetry”, “more minimal supersymmetry”, “effective supersymmetry”, “superheavy supersymmetry”, etc. The core principle cohen is that very heavy superpartners do not contribute to low-energy FCNC or CP violating processes and therefore cannot cause problems. Furthermore, no fancy symmetries need be postulated to keep experimental predictions for them under control.
On the surface, it appears that decoupling superpartners is completely irrelevant for the Tevatron. After all, Tevatron phenomenology is limited to what the Tevatron can produce. Superheavy superpartners, which we define to be above at least $`20\text{ TeV}`$, are of course not within reach of a $`2\text{ TeV}`$ collider. However, not all sparticles need be superheavy to satisfy constraints. In fact, the third generation squarks and sleptons need not be superheavy to stay within the boundaries of experimental results on FCNC and CP violating phenomena. As an all important bonus, the third family squarks and sleptons are the only ones that contribute significantly at one loop to the Higgs potential mass parameters. By keeping the third generation sfermion light, we simultaneously can maintain a “natural” and viable lagrangian even after quantum corrections are taken into account.
In short, the first-pass description of SHS is to say that, in absence of any alignment, special symmetry or other mechanism yielding flavor-horizontal degeneracy, all particles which are significantly coupled to the Higgs states should be light, and the rest heavy. The gluino does not by itself contribute to FCNC, nor does it couple directly to the Higgs bosons and so it could be heavy or light. However, the gauginos usually have a common origin, either in grand unified theories (GUTs), theories with gauge-mediated supersymmetry breaking (GMSB), or superstring theories, and so it is perhaps more likely that the gluino is relatively light with its other gaugino friends, the bino and the wino. Furthermore, the $`H_d`$ could be superheavy as well, but that is not as relevant for Tevatron phenomenology. Therefore, we can summarize the “Basic Superheavy Supersymmetry” (BSHS) spectrum:
($`\stackrel{>}{}20\text{ TeV}`$): $`\stackrel{~}{Q}_{1,2}`$, $`\stackrel{~}{u}_{1,2}^c`$, $`\stackrel{~}{d}_{1,2}^c`$, $`\stackrel{~}{L}_{1,2}`$, $`\stackrel{~}{e}_{1,2}^c`$;
($`\stackrel{<}{}1\text{ TeV}`$): $`\stackrel{~}{Q}_3`$, $`\stackrel{~}{t}^c`$, $`\stackrel{~}{B}`$, $`\stackrel{~}{W}`$, $`H_u`$, $`\mu `$ (higgsinos);
(either light or heavy): $`\stackrel{~}{b}^c`$, $`\stackrel{~}{L}_3`$, $`\stackrel{~}{\tau }^c`$, $`\stackrel{~}{g}`$, $`H_d`$.
Specific models of SUSY breaking will put the “unconstrained” fields in either the “superheavy” or “light” categories.
Any question about relative masses within each category above can not be answered within this framework. In fact, that is one of the theoretically pleasing aspect of this approach: no technical details about the spectrum need be assumed to have a viable theory. Another nice feature is that the mass pattern for the scalar partners across generations is somewhat opposite to that of the SM fermions. This might well inspire a profound connection between the physics of flavor and SUSY breaking. A possible theoretical explanation of such a large mass hierarchy in the scalar sector is that it could be a result of new gauge interactions carried by the first two generations only, and which could be, e.g., involved in a dynamical breaking of SUSY. For Tevatron enthusiasts, it is a frustrating model, since we do not even know what phenomenology should be studied because things will change drastically depending on the relative ordering of states in the “light” category.
However, there are several features about the BSHS spectrum which are interesting not because of the phenomena that it predicts at the Tevatron, but rather for what it does not predict. For example, $`\stackrel{~}{q}_{1,2}\stackrel{~}{g}`$ and $`\stackrel{~}{q}_{1,2}\stackrel{~}{q}_{1,2}^{}`$ production is not expected at the Tevatron. This is a potentially large source of events in other scenarios, such as minimal supergravity (mSUGRA), but is not present here. A more predictive feature is the expectation of many bottom quarks and $`\tau `$ leptons in the final state of SUSY production. For example, $`p\overline{p}\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^\pm `$ will not be allowed to cascade decay through $`\stackrel{~}{e}_L`$ for example, but may have hundred percent branching fractions to $`\tau `$ final states. Therefore, while the “golden tri-lepton” signals are generally suppressed in these models, efforts to look for specific $`3\tau `$ final states are relatively more important to study in the context of SHS compared to other models. Furthermore, light $`\stackrel{~}{t}`$ and $`\stackrel{~}{b}`$ production either directly or from gluino (chargino, stop) decays is of added interest in the BSHS spectrum, and may lead to high multiplicity $`b`$-jet final states. In short, drawing production and decay diagrams for all possible permutations of the BSHS spectrum always yields high multiplicity $`\tau `$ or $`b`$-jet final states. From the BSHS perspective, preparation and analysis for $`\tau `$ and $`b`$-jet identification is of primary importance. For instance, while detection of selectrons and smuons would exclude BSHS, detection of many staus and no $`\stackrel{~}{e}`$ or $`\stackrel{~}{\mu }`$ would be a good hint for it (although one could think of other SUSY scenarios where the $`m_{\stackrel{~}{e}}m_{\stackrel{~}{\tau }}`$ splitting is rather large, due e.g. to large values of $`\mathrm{tan}\beta `$). An interesting place to look for violations of $`e\tau `$ universality is $`\chi _1^\pm `$ or $`\chi _2^0`$ branching fractions, after gaugino-pair ($`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ or $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_2^0`$) production.
There are two main problems with the BSHS spectrum. The heavy particles can generate a disastrously large hypercharge Fayet-Iliopoulos term proportional to $`g_1^2`$Tr($`Ym^2`$). In universal scalar mass scenarios these terms are proportional to Tr($`Y`$) which is zero because of the gravity–gravity–U(1)<sub>Y</sub> anomaly cancellation. In minimal GMSB scenarios $`m^2Y^2+\mathrm{}`$, and so Tr($`Ym^2`$) = Tr($`Y^3`$) + $`\mathrm{}`$ vanishes because of the U(1)$`{}_{}{}^{3}{}_{Y}{}^{}`$ and SU(N)–SU(N)–U(1)<sub>Y</sub> anomaly cancellation. No such principle exists in the BSHS ansatz given above, and so the Tr($`Ym^2`$) is generically a problem. Barring the possibility of miraculous cancellations, we can cure the “Tr($`Ym^2`$) problem” by postulating that the superheavy masses follow a GMSB hierarchy, or that the superheavy states come in complete multiplets of SU(5), and the masses of all states within an SU(5) representation are degenerate or nearly degenerate. We will consider both possibilities in the following. These requirements may lower the stock of “superheavy supersymmetry” ideas for some, or it may change how one perceives model building based on decoupling superpartners, but it has no direct effect on Tevatron phenomenology.
The superheavy states are inaccessible anyway, so how they arrange their masses in detail is of little consequence to us here. On the other hand, the generic pattern and theoretical principles beyond this arrangement may affect the light sector of the model as well, both directly and indirectly through higher-order mass corrections. Indeed, another more serious problem, which has direct consequence to Tevatron phenomenology is related to new two-loop logarithmic contributions to the light scalar masses in SHS arkani . For example, the relevant renormalization group equation has a term
$$\frac{d\stackrel{~}{m}_{\mathrm{light},\mathrm{f}}^2}{d\mathrm{ln}Q}\underset{i}{}\alpha _i^2C_i^f\stackrel{~}{m}_{\mathrm{heavy}}^2+\mathrm{}$$
(8)
where $`C_i^f`$ are Casimirs for $`f`$, $`i`$ labels the indices of the SM gauge groups, and $`m_{\mathrm{heavy}}^2`$ is the characteristic superheavy mass scale. This renormalization group equation begins its running at the scale where SUSY breaking is communicated to the superpartners. In supergravity, this is the Planck scale, and so the shift in light superpartner masses is proportional to the right side of eq. 8 multiplied by a large logarithm, of order $`\mathrm{ln}M_{\mathrm{Planck}}/m_Z`$. This term is so large that in order to keep, e.g., the top squark mass squared from going negative, it must have a mass greater than several TeV at the high scale arkani . (Similar problems occur for the other “light scalars” which could potentially put us in a charge or color breaking vacuum.) Even though the top squark mass can be tuned to be light at the $`Z`$ scale, the renormalization group effects of the heavy top-squark at the high scale feed into the Higgs sector and results in a fine-tuned Higgs potential. Since fine-tuning is a somewhat subjective criteria, this problem may not be fundamental.
A healing influence on the above two-loop malady is to make the SUSY breaking transmission scale much lower than the Planck scale. This reduces the logarithm and allows for a more natural Higgs potential without large cancellations. The most successful low-energy SUSY breaking idea is GMSB giudice . There, the relevant scale is not tied to gravity ($`M_{\mathrm{Planck}}`$), but rather to the scale of dynamical SUSY breaking. Transmission of this breaking to superpartner masses can take place at scales as low as $`m_{\mathrm{heavy}}`$ in this scheme.
With some thought about the BSHS spectrum and the troubles that could arise theoretically from it, we seem to be converging on something that looks more or less like GMSB. In fact, we can think of the input parameters for our converging model to be the input parameters of minimal GMSB giudice , which are
$$\mathrm{\Lambda },M,N_{\mathrm{mess}},\mathrm{sign}(\mu ),\mathrm{tan}\beta ,\mathrm{and}\sqrt{F_0}$$
(9)
where $`\mathrm{\Lambda }`$ sets the overall mass scale of the superpartners, $`M`$ is the messenger scale, $`N_{\mathrm{mess}}`$ characterizes the number of equivalent $`5+\overline{5}`$ messenger representations, and $`\sqrt{F_0}`$ determines the interactions of the goldstino with matter. Then we add to these parameters,
$$a_{1,2}=\frac{\stackrel{~}{m}_{f_{1,2}}^2(M)}{\stackrel{~}{m}_{f_3}^2(M)}$$
(10)
where we define $`\stackrel{~}{m}_{f_3}^2(M)`$ to be the minimal GMSB values of the sfermion masses at the messenger scale excluding D-terms ($`f=\stackrel{~}{Q}`$, $`\stackrel{~}{d}^c`$, $`\stackrel{~}{L}`$, $`\stackrel{~}{e}^c`$). The two $`a_{1,2}`$ parameters with the parameters of eq. 9 completely specify a gauge-mediated inspired superheavy SUSY (GMSS) model. (Another similar parameter might be introduced for the Higgs $`H_d`$ if this is heavy, but this is less relevant to Tevatron phenomenology). We suggest that analyses can use these input parameters to make experimental searches and studies of SHS. Adding some family dependent discrete symmetries on the superpartners and messengers would allow such a model to arise in a similar way as ordinary gauge-mediated models. Recall also that in gauge mediation the Tr($`Ym^2`$) problem can be solved by the triple gauge anomaly rather than by the gravity-gauge anomaly requirement as would be the case if we had heavy sparticles come in degenerate remnants of $`\overline{5}`$ and/or $`10`$ representation, as a result of the presence of an approximate global SU(5) symmetry.
The psychological disadvantage of this GMSS model is that it is overkill on the FCNC problem. Gauge mediation cures this problem by itself, and there might not be strong motivation to further consider mechanisms that suppress it. However, gauge mediation does not automatically solve the CP problem, and so the heavy first two generations may help ameliorate it to some degree. As an aside, the above discussion can be reinterpreted as a powerful motivation for GMSB. We started with no theory principles but rather only experimental constraints and with some basic reasoning were drawn naturally to gauge mediation. However, we know of no compelling theoretical reason why $`a_{1,2}1`$. We only know that if the heavy spectrum follows a minimal gauge-mediated hierarchy, then the “Tr($`Ym^2`$)” problem can be solved. (However, it is possible to construct a more complex gauge-mediated model that does not satisfy Tr$`(Ym^2)=0`$.) Gauge-mediation, of course, is not necessarily the only way to transmit low-energy SUSY breaking. From a phenomenological point of view, one should be open to a more general low-energy SUSY breaking framework.
It must be said that in some cases, even when SUSY breaking is transmitted at low scales as in GMSS, one still could have a hard time avoiding color- and charge-breaking vacua. Indeed, the contribution from the superheavy states in eq. 8 can still be large when loops from all the scalars of the first two generations add up.
As anticipated, another possibility to cure the “Tr($`Ym^2`$) problem” and the “two-loop problem” is with the hybrid multi-scale SUSY models (HMSSM) hybrid , using the “approximate global SU(5)” pattern:
The first two generations of the $`10`$ representation of SU(5) ($`\stackrel{~}{Q}_{1,2}`$, $`\stackrel{~}{u}_{1,2}^c`$, $`\stackrel{~}{e}_{1,2}^c`$) are superheavy ($`\stackrel{~}{m}_{10_{1,2}}`$), while the rest of the sparticles are light and approximately degenerate.
In HMSSM-IIa all three generations of the $`\overline{5}`$ representation of SU(5) ($`\stackrel{~}{d}_{1,2,3}^c`$, $`\stackrel{~}{L}_{1,2,3}`$) are superheavy ($`\stackrel{~}{m}_{\overline{5}_{1,2,3}}`$), while the rest of the sparticles are light. In HMSSM-IIb just the first two generations of the $`\overline{5}`$ are superheavy ($`\stackrel{~}{m}_{\overline{5}_{1,2}}`$).
In these models, one attempts a solution of the FCNC problem by using a combination of some decoupling (superheavy scalars) and some degeneracy. A theoretical motivation for this could be that due to an approximate SU(5) global symmetry of the SUSY breaking dynamics, only some of the quark/leptons superfields with the same SU(3)$``$SU(2)$``$U(1) quantum numbers are involved in the SUSY breaking sector, carry an additional quantum number under a new “strong” horizontal gauge group and are superheavy. The other superfields instead couple only weakly (but in a flavor-blind way) to SUSY breaking and are light and about degenerate.
Actually, these “hybrid” models present many advantages compared to other SHS realizations. The reduced content of the superheavy sector considerably weakens the “two-loop” problem, since the negative contribution to the light scalar masses squared is less important. This is especially true for the HMSSM-II, and in particular the IIb version. Actually, it is in this case possible to raise the $`m_{\mathrm{heavy}}`$ scale up to $`40\text{ TeV}`$, in a natural way. Most problems with FCNC phenomena come from $`LR`$ operators, and since these operators remain suppressed, the hybrid models are phenomenologically viable and attractive versions of superheavy supersymmetry.
The resulting spectrum is different than GMSS and BSHS in that some of first two generation states are now allowed to be light. For example, in the HMSSM-I model, the $`\stackrel{~}{L}`$ sleptons can be light, and on-shell decays of winos into $`L+\stackrel{~}{L}`$ can allow the trilepton signal $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_2^03l`$ to have near $`100\%`$ branching fraction. This is not possible in the BSHS spectrum. Also, it may be useful to study the total rate of jets plus missing energy and the kinematics of the events to discern that only $`\stackrel{~}{d}^c`$, $`\stackrel{~}{s}^c`$ and 3rd-generation squarks are light, and the remaining squarks are heavy. More detailed phenomenological studies might start from observing that in the “hybrid” case too, one still needs low-energy SUSY breaking to deal with the “two-loop problem”. Again, a GMSB-inspired spectrum for the light sector corrected by the (here reduced) presence of the heavy scalars seems relevant as a starting point. In this case, a parameterization along the lines described above for the GMSS would involve new additional parameters such as $`\stackrel{~}{m}_{10_{1,2}}`$ for the HMSSM-I or $`\stackrel{~}{m}_{\overline{5}_{1,2(,3)}}`$ for the HMSSM-IIa(,b), plus possibly an analogous parameter for $`H_d`$.
Whether the spectrum is more minimal GMSB-like or is better described by the “hybrid models”, there is one feature in common. Due to the “two-loop problem”, SHS appears more natural with low-energy supersymmetry breaking, independent of how the SUSY breaking and transmission are accomplished (minimal gauge-mediation ideas or otherwise). This implies that the lightest superpartner is the gravitino rather than the neutralino, as e.g. in mSUGRA.
Depending on the details of SUSY breaking and the transmission of that breaking to superpartners (e.g., whether $`\sqrt{F_0}`$ in eq. 9 is much larger or smaller than about $`100\text{ TeV}`$), the next-to-lightest superpartner (NLSP) will either decay promptly in the detector, or decay with a long lifetime outside the detector. This may very well dominate the phenomenological implications of the model. Another important feature is the identity of the NLSP. It is well known that, e.g. in a GMSB-like spectrum, the best candidates are the $`\stackrel{~}{\chi }_1^0`$ and the lightest stau $`\stackrel{~}{\tau }_1\stackrel{~}{\tau }_R`$. In SHS, a scenario with a neutralino NLSP, with associated decays such as $`\stackrel{~}{\chi }_1^0\gamma \stackrel{~}{G}`$ possibly inside the detector, is still an important possibility. In this case, multiple high-$`p_T`$ photons are the tags to spectacular events. On the other hand, in the GMSS model and in the HMSSM models, the $`\stackrel{~}{\tau }_R`$ is always part of the light scalar sector. In addition, here the negative contributions from the heavy scalars to its mass will tend to lower it compared to the neutralino mass.
Further, in many realizations of SHS the big mass hierarchy between the Higgses $`H_u`$ and $`H_d`$ can trigger very large $`𝒪(m_{\mathrm{heavy}}/M_Z)`$ values of $`\mathrm{tan}\beta `$ (which might provide a reasonable explanation of the large $`m_t/m_b`$ ratio without fine-tuning). As a side result, after $`LR`$ mixing, the $`\stackrel{~}{\tau }_1`$ mass might turn out to be even lighter relative to the other scalars and the neutralino than in GMSB models. Hence, we believe that the possibility of a $`\stackrel{~}{\tau }_1`$ NLSP in SHS deserves very serious consideration. If this is the case, the NLSP is charged and might live beyond the detector if $`\sqrt{F_0}`$ is relatively large. Then stable charged particle tracks in the calorimeter will be tags to even more spectacular events dimo . Many of the results in the gauge-mediation literature will directly apply for discovery. After discovery, the particles that come along with the spectacular stable charged tracks (SCTs) or the high-$`p_T`$ photons can then be studied to find out with great confidence the light particle content of the theory, that could distinguish between the superheavy and the “traditional” models.
As an example of distinguishing phenomenology, we can define $`R_{\mathrm{}^{}\mathrm{}}`$ to be
$$R_{\mathrm{}^{}\mathrm{}}\underset{\stackrel{\mathrm{}=e,\mu }{\mathrm{}^{}=e,\mu }}{}\frac{\sigma [2\mathrm{SCTs}+\mathrm{}^+\mathrm{}^{}]}{\sigma [2\mathrm{SCTs}+X]}.$$
(11)
From total SUSY production in HMSSM-IIb one expects $`R_{\mathrm{}^{}\mathrm{}}<1/10`$ since $`\stackrel{~}{e}`$ and $`\stackrel{~}{\mu }`$ cannot participate in the decays. Most events will then have $`X=\tau _{\mathrm{hard}}^+\tau _{\mathrm{hard}}^{}`$ accompanying the 2 SCTs. However, in minimal GMSB the $`\stackrel{~}{e}`$ and $`\stackrel{~}{\mu }`$ are present in the low-energy spectrum, and so $`\stackrel{~}{\chi }^+e^+\nu _e\tau _{\mathrm{soft}}^\pm \stackrel{~}{\tau }^{}`$ may proceed with large branching fraction. Although a precise number depends mainly on the number of messenger representations and $`\mathrm{tan}\beta `$, $`R_{\mathrm{}^{}\mathrm{}}`$ could be greater than $`1/2`$ in GMSB. More generally, the unusually large $`LR`$ mass hierarchies that are typical of “hybrid” models may allow identification of observables suitable for discerning superheavy supersymmetry from other more conventional forms of supersymmetry at the Tevatron.
## VI Study of missing-$`E_T`$ plus jet signals at the Tevatron Run II
A. Brignole, A. Castro, F. Feruglio, G.F. Giudice, M. Mangano,
R. Rattazzi, J.D. Wells, F. Zwirner
### VI.1 Introduction
Final states with large amounts of missing transverse energy ($`\text{E}\text{/}_T`$) are of particular interest for searches of physics beyond the Standard Model in hadronic collisions. The canonical signatures include, in addition to the $`\text{E}\text{/}_T`$, the presence of 3 or more jets. In this case, the signal may come, for example, from the associated production of gluinos or scalar quarks, which undergo cascade decays to stable, invisible neutralinos and to several hadronic jets. It has been recently pointed out that signatures with just one jet can be of interest for several classes of processes beyond the Standard Model. This is the case, for example, of the associated production of a jet with a pair of very light gravitinos brig , or of a jet with one of the light Kaluza-Klein modes of the graviton from large extra dimensions savas ; grw ; KK .
The signature of one jet plus $`\text{E}\text{/}_T`$ has been neglected experimentally for a long time. The parton-level studies presented in refs. brig ; grw ; KK indicate however that this signal can be usefully exploited for searches of new physics, in spite of the potentially large backgrounds. The purpose of this contribution is to provide a realistic estimate of the discovery potential for this channel, once a realistic simulation of the detector performance is included. We shall estimate the expected backgrounds remaining after proper signal selection criteria are applied, using current CDF measurements to extract the absolute normalization of the backgrounds. We then extrapolate these estimates to the 2 TeV, 2 fb<sup>-1</sup> scenario expected at the Tevatron in the coming years, and extract the number of events which can be excluded at the 95%CL from an observation consistent with the anticipated backgrounds. This result can then be used to study the potential for discovery of arbitrary processes leading to jet+$`\text{E}\text{/}_T`$ final states. We present explicit results for the case of a very light-gravitino, studied in ref. brig , and of Kaluza-Klein graviton excitations, studied in ref. grw .
### VI.2 The strategy
The topology we are interested in is the one with a high-$`E_T`$ jet plus a large missing transverse energy ($`\text{E}\text{/}_T`$). In this study we derive a set of simple requirements which select effectively such a topology. This selection is based on simple calorimetric and tracking informations. For this reason the full simulation of the CDF detector cdf operating at the Tevatron of Fermilab during Run I is appropriate enough. For the time being we use no specific lepton identification, so we do not correct the $`\text{E}\text{/}_T`$ measurement for the presence of high-$`P_T`$ muons but only for jet energy mismeasurement.
We rely on simple requests on the minimum energy of the leading jet (clustering cone 0.4, $`E_T80`$ GeV, $`|\eta |0.7`$), on the missing transverse energy ($`\text{E}\text{/}_T100`$ GeV) and on the azimuthal angle between the missing transverse energy direction and the closest among all jets reconstructed with $`E_T15`$ GeV and $`|\eta |2.4`$, $`\mathrm{\Delta }\varphi 90^{}`$. Finally we reject events having one or more isolated tracks with $`P_T15`$ GeV/c. Tracks are defined as isolated if the sum of the $`P_T`$ of all other tracks within $`\mathrm{\Delta }R=0.4`$ from the candidate is smaller than 10 GeV/c. The first requisite favors the selection of events where a large $`\text{E}\text{/}_T`$ is associated to a high-$`E_T`$ jet, without being affected much by the presence of additional radiation. The cut $`\text{E}\text{/}_T100`$ GeV is a pre-requisite which reduces the number of background processes which need to be considered. The third cut suppresses strongly the instrumental backgrounds, due to fluctuations in the jet energy measurement, which concentrate at low $`\mathrm{\Delta }\varphi `$. The cut on the isolated tracks helps in rejecting $`W/Z+jet`$ events where leptons from the bosons appear as well-isolated tracks. While these cuts are appropriate to estimate the potential for this measurement in Run II, it goes without saying that it will be possible to improve on them and to identify optimal selection criteria once the data are available.
### VI.3 Standard Model backgrounds
The main Standard Model sources of large $`\text{E}\text{/}_T`$ are $`W+jet`$, $`Z+jet`$, $`t\overline{t}`$, $`WW`$, $`WZ`$ and $`ZZ`$ production. As a first step we evaluate their production cross sections at $`\sqrt{s}=2`$ TeV by extrapolating the values measured/calculated at 1.8 TeV with the help of the PYTHIA PYTHIA generator: $`\sigma (\sqrt{s}=2\mathrm{TeV})=\sigma (\sqrt{s}=1.8\mathrm{TeV})\times \frac{\sigma _{PYTHIA}(2\mathrm{TeV})}{\sigma _{PYTHIA}(1.8\mathrm{TeV})}`$. We generate with PYTHIA Monte Carlo samples of all these processes with a full simulation of the CDF detector used for Run I. In the case of the leading source of backgrounds, namely the $`W/Z`$+jet processes, we compare our results with the study of $`W/Z`$+jet production documented in ref. cdf\_jet , and rescale the PYTHIA result by factors of the order of 1.2. It should be noted that this $`K`$ factor is extracted from the study of jets with transverse energies in the 20 GeV, while we shall be interested here in jets and $`\text{E}\text{/}_T`$ in the range of 100 GeV and above. We shall therefore assume for this study that the relative factor between the data and PYTHIA does not depend on the $`E_T`$ thresholds. Ultimately, with the Run II data available, it will be possible to determine the absolute normalization of the $`W/Z`$+jet backgrounds by using the sample of $`(Z\mathrm{}^+\mathrm{}^{})`$+jet, and by accurately measuring the efficiency of the isolated-track cuts on the $`W`$+jet events. The accuracy of these determinations will be limited by the available statistics. Our preliminary studies indicate that it should be possible to determine the absolute rate of the $`W/Z`$+jet backgrounds to within 10%.
For the Tevatron Run II, with $`\sqrt{s}=2`$ TeV and an expected integrated luminosity of 2 fb<sup>-1</sup>, we then expect $`5200\pm 520`$ $`W+jet`$ and $`5600\pm 560`$ $`Z+jet`$ events passing our selection. We notice that $`W`$+jet and $`Z`$+jet events have a similar rate. $`Z`$+jet events come mainly from the $`Z\nu \overline{\nu }`$ decay ($`98\%`$), while all $`W`$+jet events are from leptonic decays: $`We\nu `$ ($`27\%`$), $`W\mu \nu `$ ($`31\%`$) and $`W\tau \nu `$ ($`42\%`$). In addition we notice that the irreducible background ($`Z\nu \overline{\nu }+jet`$) represents $`50\%`$ of the total background, while the remaining backgrounds (mainly $`W\mathrm{}\nu +jet`$) might be reduced by a more refined selection which tries to identify the high-$`P_T`$ leptons from the $`W`$.
The other background processes have much smaller cross sections. For $`t\overline{t}`$ we use the CDF value of $`7.6_{1.5}^{+1.8}`$ pb cdf\_top . For the diboson production at 1.8 TeV we use the following values ohn : $`\sigma _{WW}=9.5`$ pb, $`\sigma _{WZ}=2.5`$ pb and $`\sigma _{ZZ}=1.1`$ pb, to which we associate a systematic uncertainty of the order of $`30\%`$ from different PDF choices. While this uncertainty is probably over conservative, its impact on the final result is marginal, since the overall contribution of these processes is negligible, as will be shown.
Such a calculation provides the following values for $`\sqrt{s}=2`$ TeV: $`\sigma (t\overline{t})=10.1_{2.0}^{+2.4}`$ pb, $`\sigma (WW)=11.1\pm 3.3`$ pb, $`\sigma (WZ)=3.1\pm 0.9`$ pb and $`\sigma (ZZ)=1.3\pm 0.4`$ pb. For the Tevatron Run II, we then expect $`60\pm 25`$ $`t\overline{t}`$ and $`110\pm 45`$ diboson events passing our selection, accounting for an additional 25% uncertainty on the selection efficiency. This represents of the order of 2% of the total background.
Figs. 4 and 5 show how our selection suppresses the production rates of $`W+jet`$ ($`W\mathrm{}\nu `$) and $`Z+jet`$ ($`Z\nu \overline{\nu }`$) events. Fig. 6 shows, for all the sources separately, the production rates expected as a function of the $`\text{E}\text{/}_T`$ threshold. From these rates we extract the number of non-SM events which can be excluded at Run II under the assumption that the observed rates will agree with the SM expectations. We count the number of background events expected above a given $`\text{E}\text{/}_T`$ cut and evaluate the excess which can be excluded at the 95 % confidence level. When the number of background events is much smaller than 1 this corresponds to a limit of 3 signal events. The result is presented in Fig. 7, which shows how the result depends upon the residual uncertainty on the absolute background rate.
### VI.4 The case of $`\stackrel{~}{G}\stackrel{~}{G}+jet`$ production
All realistic models with global supersymmetry must contain the goldstino, a massless and neutral spin-$`\frac{1}{2}`$ particle associated with the spontaneous breaking of supersymmetry. When gravitation is introduced, and supersymmetry is realized locally, the spin-$`2`$ graviton is accompanied by the gauge particle of supersymmetry, the spin-$`\frac{3}{2}`$ gravitino. After spontaneous supersymmetry breaking, the gravitino acquires a mass by absorbing the would-be goldstino. The signals of pair-production of squarks, gluinos, charginos and neutralinos have been extensively discussed at this workshop, in the case of a heavy as well as of a light gravitino. In the latter case, the gravitino is the Lightest Supersymmetric Particle, (LSP), and goes undetected giving rise to $`\text{E}\text{/}_T`$.
If the gravitino, $`\stackrel{~}{G}`$, is very light (much lighter than $`10^4`$ eV/c<sup>2</sup>), and all the other supersymmetric particles are above the production threshold, supersymmetry may still be seen at the Tevatron by looking at final states including gravitinos and ordinary particles only. In particular, we can have the following processes: $`q\overline{q}\stackrel{~}{G}\stackrel{~}{G}g`$, $`qg\stackrel{~}{G}\stackrel{~}{G}q`$, $`\overline{q}g\stackrel{~}{G}\stackrel{~}{G}\overline{q}`$ and $`gg\stackrel{~}{G}\stackrel{~}{G}g`$, which all lead to a topology with a single jet + $`\text{E}\text{/}_T`$. In this scenario the main parameters upon which these processes depend are the supersymmetry-breaking scale $`\sqrt{F}`$ and the CMS energy $`\sqrt{s}`$. We recall that the gravitino mass $`m_{\stackrel{~}{G}}`$ is related to $`F`$ via $`m_{\stackrel{~}{G}}M_{Pl}=F/\sqrt{3}`$, where $`M_{Pl}=2.4\times 10^{18}`$ GeV is the reduced Planck mass.
For the production of very light gravitinos ($`m_{\stackrel{~}{G}}10^4`$ eV/c<sup>2</sup>) we consider the processes $`p\overline{p}\stackrel{~}{G}\stackrel{~}{G}g,\stackrel{~}{G}\stackrel{~}{G}q`$ leading to a single jet + large $`\text{E}\text{/}_T`$ topology. These processes are simulated at $`\sqrt{s}=2`$ TeV in a private version of HERWIG V5.6 HERWIG , which reproduces the calculations of brig , in the limit in which the supersymmetric particles of the Minimal Supersymmetric Standard Model and other exotic particles are heavy. The production cross section depends on the supersymmetry-breaking parameter $`F`$: $`\sigma \frac{1}{F^4}`$. For the generation we use $`\sqrt{F}=290`$ GeV as justified in the next section. We choose a renormalization/factorization scale $`\mu =E_T`$ and the MRSD- MRSD set of structure functions. In order to reduce the computing time we require the hard scatter momentum $`P_{Thard}70`$ GeV/c: this choice is justified by the fact that events generated with $`P_{Thard}<70`$ GeV/c have $`\text{E}\text{/}_T`$ below 100 GeV and do not pass our cuts. With such choice of parameters, the production cross section amounts to $`2.62\pm 0.02`$ pb for $`\text{E}\text{/}_T100`$ GeV. To populate the $`\text{E}\text{/}_T`$ region above 300 GeV we generate also events with $`P_{Thard}>200`$ and 300 GeV/c. All samples are fed to the CDF detector simulation.
Fig. 8 shows the $`\stackrel{~}{G}\stackrel{~}{G}+jet`$ production cross section as a function of the $`\text{E}\text{/}_T`$ threshold, at parton level, after the event reconstruction and after our selection. Fig. 9 shows the effect of applying reconstruction and selection cuts to the signal, relative to the naive results which would be obtained at the parton level. The loss in efficiency is induced partly by the shower evolution of the initial and final state, which slightly reduces the amount of $`\text{E}\text{/}_T`$ relative to the parton-level configuration, and partly by other purely detector-driven effects, such as the cuts on isolated tracks. It is reasonable to assume that the efficiencies presented in this figure apply to the case of the signals studied in ref. KK , since the nature of both the hadronic initial and final states are very similar to those encountered in the light-gravitino case. Fig. 9 can therefore be used for realistic estimates of the signal deterioration in the case of ref. grw ; KK , and, when used in conjunction with the results of Fig. 7, can provide a realistic estimate of the exclusion potential of Run II for signals of millimeter-scale extra-dimensions.
### VI.5 Comparison $`\stackrel{~}{G}\stackrel{~}{G}+jet`$ production $`vs`$ Standard Model backgrounds
In this Section we use the estimated background rate (see Fig. 10) to derive the sensitivity to the scale $`\sqrt{F}`$ in the case of a very light gravitino. The best lower limit on $`\sqrt{F}`$ which can be obtained with our simple selection is $`370`$ GeV if we require $`\text{E}\text{/}_T100`$ GeV (see Fig. 11). The corresponding lower limit on the gravitino mass is $`3.3\times 10^5`$ eV/c<sup>2</sup>. We do not use lower thresholds because the background increases strongly below 100 GeV due to the presence of instrumental backgrounds. This calculation is repeated to account for the $`10\%`$ global uncertainty on the background estimate. In this case the best lower limit on $`\sqrt{F}`$ is $`290`$ GeV, reached for $`\text{E}\text{/}_T200`$ GeV, and the signal and background rates are similar in size (see Fig. 10). This limit corresponds to $`m_{\stackrel{~}{G}}2.0\times 10^5`$ eV/c<sup>2</sup>. The limit has become much worse because of the large systematic uncertainty on the background which dominates over the statistical fluctuations up to $`\text{E}\text{/}_T200`$ GeV. Finally, we remark that the existence of a very low supersymmetry- breaking scale may also show up indirectly at the Tevatron Run II, in the form of anomalous four-fermion interactions involving standard quarks and leptons brig2 .
### VI.6 Production of jet plus Gravitons
In ordinary gravity, the cross-section for $`p\overline{p}\mathrm{jet}+`$ graviton at the Tevatron is about $`10^{26}\mathrm{fb}`$. Needless to say, ordinary quantum gravity effects have no hope of being seen at the Tevatron. However, the scale of quantum gravity may be as low as the weak scale, rather than the Planck scale, if gravity propagates in a higher dimensional space with $`D`$ total space-time dimensions savas . The $`\delta =D4`$ extra spatial dimensions must be compactified with a large radius $`R`$ given by
$$M_{\mathrm{Pl}}^2=R^\delta M_D^{2+\delta },$$
(12)
where $`M_D`$ is the characteristic quantum gravity scale. The probability of producing a single light Kaluza-Klein excitation of the graviton at a collision with typical energy $`E`$ is of the order of $`E^2/M_{\mathrm{Pl}}^2`$. This small probability is compensated by the large number of available graviton excitations. Indeed, the number of Kaluza-Klein excitations of the graviton lighter than $`E`$ is equal to
$$N_G(ER)^\delta =E^\delta \frac{M_{\mathrm{Pl}}^2}{M_D^{2+\delta }}.$$
(13)
With $`\delta =2`$ extra spatial dimensions and $`M_D1\mathrm{TeV}`$, the multiplicity of graviton states in typical Tevatron processes is about $`10^{2529}`$. Therefore, the final cross-section may be roughly comparable to SM cross-sections, and signs of quantum-gravity effects may appear savas .
Single jet plus gravitons production at the Tevatron has been calculated and compared to SM backgrounds, using a parton-level Monte Carlo grw ; KK . With the efficiencies presented in sections II and III, we can now improve the parton-level results, relating them to a more realistic event simulation that includes initial state radiation, jet clustering, and jet energy reconstruction effects.
Our parton-level calculation is very similar to that described in ref. grw , except that here we follow the kinematic acceptance criteria put forth in section II above. We retain the $`\text{E}\text{/}_T>150\mathrm{GeV}`$ requirement in ref. grw in order to facilitate a comparison between that publication and the present analysis. After the parton-level simulation, we multiply by a factor of $`0.55`$ for the relative efficiency between full reconstruction and the parton-level results. We expect the relative efficiency results for the gravitino signal to be very similar to that of the gravitons signal considered here, and so this number is read directly from Fig. 9b with $`\text{E}\text{/}_T^{\mathrm{min}}=150\mathrm{GeV}`$.
In Fig. 12 we plot the total number of signal events expected in $`2\mathrm{fb}^1`$ of integrated luminosity, after all efficiencies are taken into account, for both $`\delta =2`$ and $`\delta =4`$. The b curve indicates that we integrate the perturbative jet plus gravitons amplitudes for all values of $`\widehat{s}`$ (partonic center of mass energy), and the a curve indicates that we integrate the perturbative jet plus gravitons amplitudes only for $`\widehat{s}<M_D^2`$ and set the amplitude to zero for $`\widehat{s}>M_D^2`$. This is a way to parameterize our ignorance about the non-perturbative quantum-gravity regime $`\widehat{s}>M_D^2`$. We expect the full, unitarized cross-section in a complete quantum-gravity theory to be somewhere between the a curve and the b curve. The limits on the quantum-gravity scale $`M_D`$ which can be obtained from $`2\mathrm{fb}^1`$ of data are presented in Table 4. The limits are derived under the more conservative assumption of curve a, i.e. assuming $`\sigma (\widehat{s}>M_D^2)=0`$. In Table 4 we also show the original estimate of the limits based on the parton results of ref. grw . As expected, the limits based on the full jet reconstruction simulation with 10% background uncertainty are somewhat lower than the limits based on parton-level Monte Carlo simulation. Nevertheless, probes of gravity above the TeV mass scale can be accomplished with $`2\mathrm{fb}^1`$ integrated luminosity. Furthermore, as the background becomes better understood, the range of $`M_D`$ that can be discovered quickly rises, as illustrated in Table 4. Finally we remark that the existence of a low quantum-gravity scale can also be tested at the Tevatron Run II by studying the effects of virtual graviton exchange in SM processes grw ; hew , as reviewed in these proceedings hewp .
### VI.7 Conclusions
We have studied the production of events with a high-$`E_T`$ jet plus large $`\text{E}\text{/}_T`$ at 2 TeV with the help of Monte Carlo simulations for the production of very-light gravitinos and of Standard Model backgrounds. We have defined a set of simple selection criteria which are quite efficient on the signal ($`50\%`$ for $`\text{E}\text{/}_T100`$ GeV) and reduce strongly the backgrounds (by a factor $`6000`$). Comparing the estimated background to the expected signal as a function of $`\sqrt{F}`$, and assuming a conservative 10% uncertainty on the absolute background rate, we derive a 95% C.L. lower limit on the supersymmetry-breaking scale $`\sqrt{F}290`$ GeV for the Tevatron Run II with 2 fb<sup>-1</sup>. This limit corresponds to a lower limit on the gravitino mass of $`2.0\times 10^5`$ eV/c<sup>2</sup> and is expected to improve once we reduce the systematic uncertainty on the background estimate. We have also derived limits on the quantum-gravity scale $`M_D`$ in theories with $`\delta `$ large extra dimensions. The Tevatron Run II with 2 fb<sup>-1</sup> can reach the 95% C.L. lower limits on $`M_D`$ presented in Table 4.
## VII Physics Implications of a Perturbative Superstring Construction
M. Cvetič, L. Everett, P. Langacker, and J. Wang
### VII.1 Introduction
Predictions from superstring theory provide natural possible extensions of the MSSM. However, there are several problems to be resolved in attempting to connect string theory to the observable, low energy world. First, many models can be derived from string theory, and there is no dynamical principle as yet to select among them. Furthermore, no fully realistic model, i.e., a model which contains just the particle content and couplings of the MSSM, has been constructed. In addition, there is no compelling scenario for how to break supersymmetry in string theory, and so soft supersymmetry breaking parameters must be introduced into the model by hand, just as in the MSSM.
We adopt a more modest strategy and consider a class of quasi-realistic models constructed within weakly coupled heterotic string theory. In addition to the necessary ingredients of the MSSM, such models generically contain an extended gauge structure that includes a number of $`U(1)`$ gauge groups and “hidden” sector non-Abelian groups, and additional matter fields (including a number of SM exotics and SM singlets). They predict gauge couplings unification (sometimes with non-standard Kač-Moody level) at the string scale $`M_{String}5\times 10^{17}`$ GeV. The most desirable feature of models in this class is that the superpotential is explicitly calculable; in particular, the non-zero Yukawa couplings are $`𝒪(1)`$, and can naturally accommodate the radiative electroweak symmetry breaking scenario. In addition, string selection rules can forbid gauge-allowed terms, in contrast to the case in general field-theoretic models.
In addressing the phenomenology of these models, there are two complementary approaches. The first is the ”bottom-up” approach, in which models with particle content and couplings motivated from quasi-realistic string models are studied to provide insight into the new physics that can emerge from string theory (such as additional $`Z^{}`$ gauge bosons)cl ; SY ; cdeel ; cceel1 ; lw . In this work we adopt the second (“top-down”) approach, and analyze a prototype string model (Model 5 of chl ) in detail. The analysis of this class of string models (done in collaboration with G. Cleaver and J. R. Espinosa) proceeds in several stages, which will be briefly summarized below and is documented in cceel2 ; cceel3 ; cceelw1 ; cceelw2 . We then focus on the main results: the determination of the mass spectrum and trilinear couplings at the string scale, the renormalization group analysis, the low energy gauge symmetry breaking patterns and the mass spectrum of the model at the electroweak scale.
Our analysis shows that the prototype model is not fully realistic. In particular, many of the SM exotics remain massless in the low energy theory. However, we find there are other general features of the model which have interesting phenomenological implications, including an additional low-energy $`U(1)^{}`$ gauge group, $`R`$ parity violating couplings, “mixed” effective $`\mu `$ terms, and extended chargino, neutralino, and Higgs sectors (with patterns of mass spectra that differ substantially from the case of the MSSM).
In section II, we discuss the generation of the effective mass terms and the trilinear couplings associated with the flat direction. In section III, we present the effective couplings and the implications of the effective theory along a particular flat direction as an illustrative example. We conclude in section IV.
### VII.2 Flat Directions and Effective Couplings
The model we have chosen as a prototype model to analyze is Model 5 of chl . Prior to vacuum restabilization, the model has the gauge group
$$\{SU(3)_C\times SU(2)_L\}_{\mathrm{obs}}\times \{SU(4)_2\times SU(2)_2\}_{\mathrm{hid}}\times U(1)_A\times U(1)^6,$$
(14)
and a particle content that includes the following chiral superfields in addition to the MSSM fields:
$`6(1,2,1,1)+(3,1,1,1)+(\overline{3},1,1,1)+`$
$`4(1,2,1,2)+2(1,1,4,1)+10(1,1,\overline{4},1)+`$
$`8(1,1,1,2)+5(1,1,4,2)+(1,1,\overline{4},2)+`$
$`8(1,1,6,1)+3(1,1,1,3)+42(1,1,1,1),`$ (15)
where the representation under $`(SU(3)_C,SU(2)_L,SU(4)_2,SU(2)_2)`$ is indicated. The SM hypercharge is determined as a linear combination of the six non-anomalous $`U(1)`$’s.
As the first step of the analysis, we address the presence of the anomalous $`U(1)_A`$ generic to this class of models. The standard anomaly cancellation mechanism generates a nonzero Fayet-Iliopoulos (FI) term of $`𝒪(M_{String})`$ to the $`D`$ term of $`U(1)_A`$. The FI term would appear to break supersymmetry at the string scale, but certain scalar fields are triggered to acquire large VEV’s along $`D`$ and $`F`$ flat directions, such that the new “restabilized” vacuum is supersymmetric. The complete set of $`D`$ and $`F`$ flat directions involving the non-Abelian singlet fields for Model 5 was classified in cceel2 .
In a given flat direction, the rank of the gauge group is reduced, and effective mass terms and trilinear couplings may be generated from higher order terms in the superpotential:
$`W_M`$ $`=`$ $`{\displaystyle \frac{\alpha _{K+2}}{M_{Pl}^{K1}}}\mathrm{\Psi }_i\mathrm{\Psi }_j\mathrm{\Phi }^K`$ (16)
$`W_3`$ $`=`$ $`{\displaystyle \frac{\alpha _{K+3}}{M_{Pl}^K}}\mathrm{\Psi }_i\mathrm{\Psi }_j\mathrm{\Psi }_k\mathrm{\Phi }^K,`$ (17)
in which the fields which are in the flat direction are denoted by $`\mathrm{\Phi }`$, and those which are not by $`\{\mathrm{\Psi }_i\}`$. Hence, some fields acquire superheavy masses and decouple. The effective Yukawa couplings of the remaining light fields are typically suppressed <sup>5</sup><sup>5</sup>5However, in the prototype model considered the effective trilinear couplings arising from fourth order terms are comparable in strength to the original Yukawas. compared with Yukawa couplings of the original superpotential (the $`\alpha _K`$ coefficients are in principle calculable; for details, see cew ).
This procedure has been carried out for the prototype model in cceel2 ; cceelw1 . We carry out the analysis of the implications of the model for the flat directions that break the maximal number of $`U(1)`$’s, leaving $`U(1)_Y`$ and $`U(1)^{^{}}`$ unbroken. The list of matter superfields and their $`U(1)_Y`$ and $`U(1)^{}`$ charges are presented in Table 1.
### VII.3 Example: Low Energy Implications of a Representative Flat Direction
We choose to present the analysis of the model along a particular flat direction <sup>6</sup><sup>6</sup>6Other flat directions involve other interesting features, such as fermionic textures, baryon number violation, and the possibility of intermediate scale $`U(1)^{}`$ breaking.. The flat direction we consider is the $`P_1^{}P_2^{}P_3^{}`$ direction (in the notation of cceel2 ; cceelw1 ; cceelw2 ), which involves the set of fields $`\{\phi _2,\phi _5,\phi _{10},\phi _{13},\phi _{27},\phi _{29},\phi _{30}\}`$.
Along this flat direction, the effective mass terms which involve the observable sector fields and the non-Abelian singlets <sup>7</sup><sup>7</sup>7We refer the reader to cceelw1 ; cceelw2 for further details of the model, such as the couplings involving the hidden sector fields. are given by
$`W_M`$ $`=`$ $`gh_f\overline{h}_b\phi _{27}+gh_g\overline{h}_d\phi _{29}+{\displaystyle \frac{\alpha _4^{(1)}}{M_{Pl}}}h_b\overline{h}_b\phi _5\phi _{10}+{\displaystyle \frac{\alpha _4^{(2)}}{M_{Pl}}}h_b\overline{h}_b\phi _2\phi _{13}`$ (18)
$`+`$ $`{\displaystyle \frac{g}{\sqrt{2}}}(e_d^ce_b+e_g^ce_a)\phi _{30}+{\displaystyle \frac{g}{\sqrt{2}}}(\phi _1\phi _{15}+\phi _4\phi _9)\phi _{10}+{\displaystyle \frac{g}{\sqrt{2}}}(\phi _7\phi _{16}+\phi _9\phi _{12})\phi _2`$
$`+`$ $`{\displaystyle \frac{g}{\sqrt{2}}}(\phi _6\phi _{26}+\phi _8\phi _{23}+\phi _{14}\phi _{17})\phi _{29}+{\displaystyle \frac{\alpha _4^{(3)}}{M_{Pl}}}\phi _{21}\phi _{25}\phi _{27}\phi _{29}.`$
The effective trilinear couplings involving all fields which couple directly to the observable sector fields are given by:
$`W_3`$ $`=`$ $`gQ_cu_c^c\overline{h}_c+gQ_cd_b^ch_c+{\displaystyle \frac{\alpha _4^{(4)}}{M_{Pl}}}Q_cd_d^ch_a\phi _{29}+{\displaystyle \frac{g}{\sqrt{2}}}e_a^ch_ah_c+{\displaystyle \frac{g}{\sqrt{2}}}e_f^ch_dh_c`$ (19)
$`+`$ $`{\displaystyle \frac{\alpha _5^{(1)}}{M_{Pl}^2}}e_h^ch_eh_a\phi _5\phi _{27}+{\displaystyle \frac{\alpha _5^{(2)}}{M_{Pl}^2}}e_e^ch_eh_a\phi _{13}\phi _{27}+gh_b^{}\overline{h}_c\phi _{20}.`$
In the observable sector, the fields which remain light include both the usual MSSM states and exotic states such as a fourth ($`SU(2)_L`$ singlet) down-type quark, extra fields with the same quantum numbers as the lepton singlets, and extra Higgs doublets. There are other massless states with exotic quantum numbers (including fractional electric charge) that also remain in the low energy theory. The $`U(1)^{}`$ charges of the light fields are family nonuniversal (and hence is problematic with respect to FCNC).
There are some generic features of the superpotential which are independent of the details of the soft supersymmetry breaking parameters. In addition to a large top-quark Yuakwa coupling ($`𝒪(1)`$) which is necessary for radiative electroweak symmetry breaking, the couplings indicate $`tb`$ and (unphysical) $`\tau \mu `$ Yukawa unification, with the identification of the fields $`\overline{h}_c`$, $`h_c`$ as the standard electroweak Higgs doublets. There is no elementary or effective canonical $`\mu `$-term involving $`\overline{h}_c`$ and $`h_c`$, but rather non-canonical effective $`\mu `$ terms involving additional Higgs doublets. Finally, there is also a possibility of lepton- number violating couplings; thus this model violates $`R`$\- parity, and has no stable LSP.
With the knowledge of the massless spectrum at the string scale, the gauge coupling beta-functions can be determined, and the gauge couplings can then be run from the string scale (where they are predicted to unify) to the electroweak scale. We determine the gauge coupling constant $`g=0.80`$ at the string scale by assuming $`\alpha _s=0.12`$ (the experimental value) at the electroweak scale, and evolving $`g_3`$ to the string scale. We then use this value as an input to determine the electroweak scale values of the other gauge couplings by their (1-loop) renormalization group equations (RGE’s). The low energy values of the gauge couplings are not correct due to the exotic matter and non-standard $`k_Y=11/3`$ for this model; however, it is surprising that $`\mathrm{sin}^2\theta _W0.16`$ and $`g_2=0.48`$ are not too different from the experimental values $`0.23`$ and $`0.65`$, respectively.
The string-scale values of the Yukawa couplings of (19) are calculable (with the knowledge of the VEV’s of the singlet fields in the flat direction). Utilizing the RGE’s, we can also determine the low energy values of the Yukawa coupling constants. The running of the gauge couplings and the Yukawa couplings are shown in Fig. 13.
To address the gauge symmetry breaking scenarios for this model, we introduce soft supersymmetry breaking mass parameters and run the RGE’s from the string scale to the electroweak scale. While the qualitative features of the analysis are independent of the details of the soft breaking, we choose to illustrate the analysis with a specific example with a realistic $`ZZ^{}`$ hierarchy. General considerations cl ; cdeel ; cceel1 ; lw and an inspection of the $`U(1)^{}`$ charges of the light fields indicate that, in this example, the $`U(1)^{}`$ breaking is at the electroweak (TeV) scale. Due to the lack of a canonical effective $`\mu `$ term between $`\overline{h}_c`$, $`h_c`$ and a singlet, an extended Higgs sector is required, with an additional Higgs doublet and singlet ($`\overline{h}_c`$, $`h_c`$, $`h_b^{}`$, and $`s\phi _{20}`$). The symmetry breaking is characterized by a large ($`𝒪(\mathrm{TeV})`$) value of the SM singlet VEV, with the electroweak symmetry breaking at a lower scale due to accidental cancellations.
We now present the mass spectrum for a concrete numerical example of this scenario, which requires mild tuning of the soft supersymmetry breaking mass parameters at the string scale. The initial and final values of the parameters for this example are listed in Table 2.
* Fermion Masses: The masses for the $`t`$, $`b`$, $`\tau `$, and $`\mu `$ are due to Yukawa couplings of the original superpotential, as shown in (19). With the identification of $`Q_c`$ as the quark doublet of the third family and $`h_d`$, $`h_a`$ as the lepton doublets of the third and second families, respectively, $`m_t=156`$ GeV, $`m_b=83`$ GeV, $`m_\tau =32`$ GeV, and $`m_\mu =27`$ GeV. The ratio $`m_b/m_\tau `$ is larger than in the usual $`b\tau `$ unification because of the ratio $`1:1/\sqrt{2}`$ of the Yukawa couplings at the string scale, and is probably inconsistent with experiment mbmtau (of course, the high values for $`m_b`$, $`m_\tau `$, and $`m_\mu `$ are unphysical). Finally, $`u`$, $`d`$, $`c`$, $`s`$, and $`e^{}`$ remain massless.
* Squarks/Sleptons: To ensure a large $`M_Z^{}`$ in this model, the squark and slepton masses have values in the several TeV range, with $`m_{\stackrel{~}{t}L}=2540`$ GeV, $`m_{\stackrel{~}{t}R}=2900`$ GeV; $`m_{\stackrel{~}{b}L}=2600`$ GeV, $`m_{\stackrel{~}{b}R}=2780`$ GeV; $`m_{\stackrel{~}{\tau }L}=2760`$ GeV, $`m_{\stackrel{~}{\tau }R}=3650`$ GeV; $`m_{\stackrel{~}{\mu }L}=2790`$ GeV, and $`m_{\stackrel{~}{\mu }R}=3670`$ GeV.
* Charginos/Neutralinos: The positively charged gauginos and higgsinos are $`\stackrel{~}{W}^+`$, $`\stackrel{~}{\overline{h}}_c`$, $`\stackrel{~}{\overline{h}}_a`$, and the negatively charged gauginos and higgsinos are $`\stackrel{~}{W}^{}`$, $`\stackrel{~}{h}_c`$, $`\stackrel{~}{h}_b^{}`$. There is one massless chargino, and the other two have masses $`m_{\stackrel{~}{\chi }_1^\pm }=591`$ GeV, and $`m_{\stackrel{~}{\chi }_2^\pm }=826`$ GeV.
The neutralino sector consists of $`\stackrel{~}{B}^{}`$, $`\stackrel{~}{B}`$, $`\stackrel{~}{W_3}`$, $`\stackrel{~}{\overline{h}}_c^0`$, $`\stackrel{~}{h}_c^0`$, $`\stackrel{~}{h}_b^{}_{}{}^{}0`$, $`\stackrel{~}{\phi }_{20}^{^{}}`$, and $`\stackrel{~}{\overline{h}}_a^0`$. The mass eigenvalues are $`m_{\stackrel{~}{\chi }_1}^0=963`$ GeV, $`m_{\stackrel{~}{\chi }_2}^0=825`$ GeV, $`m_{\stackrel{~}{\chi }_3}^0=801`$ GeV, $`m_{\stackrel{~}{\chi }_4}^0=592`$ GeV, $`m_{\stackrel{~}{\chi }_5}^0=562`$ GeV, $`m_{\stackrel{~}{\chi }_6}^0=440`$ GeV, $`m_{\stackrel{~}{\chi }_7}^0=2`$ GeV, and $`m_{\stackrel{~}{\chi }_8}^0=0`$.
In both cases, the absence of couplings of the Higgs field $`\overline{h}_a`$ in the superpotential leads to a massless chargino and neutralino state. The absence of an effective $`\mu `$ term involving $`h_c`$ leads to an additional global $`U(1)`$ symmetry in the scalar potential, and an ultralight neutralino pair in the mass spectrum.
* Exotics: There are a number of exotic states, including the $`SU(2)_L`$ singlet down-type quark, four $`SU(2)_L`$ singlets with unit charge (the $`e`$ and extra $`e^c`$ states), and a number of SM singlet ($`\phi `$) states, as well as exotics associated with the hidden sector. The scalar components of these exotics are expected to acquire TeV-scale masses by soft supersymmetry breaking. However, there is no mechanism within our assumptions to give the fermions significant masses.
* Higgs Sector: The non-minimal Higgs sector of three complex doublets and one complex singlet leads to additional Higgs bosons compared to the MSSM. Four of the fourteen degrees of freedom are eaten to become the longitudinal components of the $`W^\pm `$, $`Z`$, and $`Z^{}`$; and the global $`U(1)`$ symmetry is broken, leading to a massless Goldstone boson (which, however, acquires a small mass at the loop level) in the spectrum.
The spectrum of the physical Higgs bosons after symmetry breaking consists of two pairs of charged Higgs bosons $`H_{1,2}^\pm `$, four neutral CP even Higgs scalars $`(h_i^0,i=1,2,3,4)`$, and one CP odd Higgs $`A^0`$, with masses $`m_{H_1^\pm }=10`$ GeV, $`m_{H_2^\pm }=1650`$ GeV, $`m_{h_1^0}=33`$ GeV, $`m_{h_2^0}=47`$ GeV, $`m_{h_3^0}=736`$ GeV, $`m_{h_4^0}=1650`$ GeV, and $`m_{A^0}=1650`$ GeV.
In this model, the bound on the lightest Higgs scalar is different that the traditional bound in the MSSM. It is associated with the breaking scale of the additional global $`U(1)`$ symmetry; since this scale is comparable to the electroweak scale, not only one but two Higgs scalars will be light in the decoupling limit. In particular, the lightest Higgs mass satisfies the (tree-level) bound comelli
$`m_{h_1^0}^2{\displaystyle \frac{G^2}{4}}v_1^2+g_1^{}^2Q_1^2v_1^2=(35\mathrm{GeV})^2.`$ (20)
### VII.4 Conclusions
The purpose of this work has been to explore the general features of this class of quasi-realistic superstring models through a systematic, “top-down” analysis of a prototype model. The results of the investigation of the low energy implications of the mass spectrum and couplings predicted in a subset of the restabilized vacua of this model demonstrate that in general, the TeV scale physics is more complicated than that of the MSSM.
In particular, we have found that non-canonical couplings, such as mixed effective $`\mu `$ terms and $`R`$\- parity violating operators are typically present in the superpotential. In some other cases, there are possibilities for potentially interesting fermion textures. The model is also characterized by the presence of extra matter in the low energy theory such as SM exotics, extra $`Z^{}`$ gauge bosons with TeV scale masses, and additional charginos, neutralinos, and Higgs bosons with patterns of masses that differ substantially from the MSSM.
The particular model we studied is not realistic, in part due to the presence of (ultra-light or massless) extra matter. However, due to the large number of possible models that can be derived from string theory, this result does not invalidate the potential viability of string models, or the motivation for investigating their phenomenological implications <sup>8</sup><sup>8</sup>8Progress has been made in exploring models in which the exotic matter decouples above the electroweak scale; e.g. see faraggi98 .. We stress that the features of this model are likely to be generic to this class of quasi-realistic models based on weakly coupled heterotic string theory, and thus warrant further consideration.
| $`(SU(3)_C,SU(2)_L,`$ | | $`6Q_Y`$ | $`100Q_Y^{}`$ |
| --- | --- | --- | --- |
| $`SU(4)_2,SU(2)_2)`$ | | | |
| (3,2,1,1): | $`Q_a`$ | 1 | 68 |
| | $`Q_b`$ | 1 | 68 |
| | $`Q_c`$ | 1 | $``$71 |
| ($`\overline{3}`$,1,1,1): | $`u_a^c`$ | $``$4 | 6 |
| | $`u_b^c`$ | $``$4 | 6 |
| | $`u_c^c`$ | $``$4 | $``$133 |
| | $`d_a^c`$ | 2 | $``$3 |
| | $`d_b^c`$ | 2 | 136 |
| | $`d_c^c`$ | 2 | $``$3 |
| | $`d_d^c`$ | 2 | $``$3 |
| (1,2,1,1): | $`\overline{h}_a`$ | 3 | $``$74 |
| | $`\overline{h}_b`$ | 3 | 65 |
| | $`\overline{h}_c`$ | 3 | 204 |
| | $`\overline{h}_d`$ | 3 | 65 |
| | $`h_a`$ | $``$3 | 74 |
| | $`h_b`$ | $``$3 | $``$65 |
| | $`h_c`$ | $``$3 | $``$65 |
| | $`h_d`$ | $``$3 | $``$65 |
| | $`h_e`$ | $``$3 | $``$204 |
| | $`h_f`$ | $``$3 | $``$65 |
| | $`h_g`$ | $``$3 | $``$65 |
| (3,1,1,1): | $`𝒟_a`$ | $``$2 | $``$136 |
Table Ia: List of non-Abelian non-singlet observable sector fields in the model with their charges under hypercharge and $`U(1)^{}`$.
| | $`6Q_Y`$ | $`100Q_Y^{}`$ | | $`6Q_Y`$ | $`100Q_Y^{}`$ |
| --- | --- | --- | --- | --- | --- |
| $`e_{a,c}^c`$ | 6 | $``$9 | $`e_b^c`$ | 6 | $``$9 |
| $`e_{d,g}^c`$ | 6 | 130 | $`e_e^c`$ | 6 | 130 |
| $`e_f^c`$ | 6 | 130 | $`e_h^c`$ | 6 | 130 |
| $`e_i^c`$ | 6 | $``$9 | $`e_{a,b}`$ | 6 | $``$130 |
| $`e_c`$ | 6 | $``$130 | $`e_{d,e}`$ | $``$6 | 9 |
| $`e_f`$ | $``$6 | $``$269 | | | |
| $`\phi _1`$ | 0 | 0 | $`\phi _{2,3}`$ | 0 | 0 |
| $`\phi _{4,5}`$ | 0 | 0 | $`\phi _{6,7}`$ | 0 | 0 |
| $`\phi _{8,9}`$ | 0 | 0 | $`\phi _{10,11}`$ | 0 | 0 |
| $`\phi _{12,13}`$ | 0 | 0 | $`\phi _{14,15}`$ | 0 | 0 |
| $`\phi _{16}`$ | 0 | 0 | $`\phi _{17}`$ | 0 | 0 |
| $`\phi _{18,19}`$ | 0 | $``$139 | $`\phi _{20,21}`$ | 0 | $``$139 |
| $`\phi _{22}`$ | 0 | $``$139 | $`\phi _{23}`$ | 0 | 0 |
| $`\phi _{24}`$ | 0 | 0 | $`\phi _{25}`$ | 0 | 139 |
| $`\phi _{26}`$ | 0 | 0 | $`\phi _{27}`$ | 0 | 0 |
| $`\phi _{28,29}`$ | 0 | 0 | $`\phi _{30}`$ | 0 | 0 |
Table Ib: List of non-Abelian singlet fields in the model with their charges under hypercharge and $`U(1)^{}`$.
| | $`M_Z`$ | $`M_{String}`$ | | $`M_Z`$ | $`M_{String}`$ |
| --- | --- | --- | --- | --- | --- |
| $`g_1`$ | 0.41 | 0.80 | $`M_1`$ | 444 | 1695 |
| $`g_2`$ | 0.48 | 0.80 | $`M_2`$ | 619 | 1695 |
| $`g_3`$ | 1.23 | 0.80 | $`M_3`$ | 4040 | 1695 |
| $`g_1^{}`$ | 0.43 | 0.80 | $`M_1^{}`$ | 392 | 1695 |
| $`\mathrm{\Gamma }_{Q1}`$ | 0.96 | 0.80 | $`A_{Q1}`$ | 3664 | 8682 |
| $`\mathrm{\Gamma }_{Q2}`$ | 0.93 | 0.80 | $`A_{Q2}`$ | 4070 | 9000 |
| $`\gamma _{Q3}`$ | 0.27 | 0.08 | $`A_{Q3}`$ | 5018 | 1837 |
| $`\mathrm{\Gamma }_{l1}`$ | 0.30 | 0.56 | $`A_{l1}`$ | $``$946 | 4703 |
| $`\mathrm{\Gamma }_{l2}`$ | 0.36 | 0.56 | $`A_{l2}`$ | $``$707 | 4532 |
| $`\mathrm{\Gamma }_{l3}`$ | 0.06 | 0.05 | $`A_{l3}`$ | 4613 | 4425 |
| $`\mathrm{\Gamma }_{l4}`$ | 0.11 | 0.13 | $`A_{l4}`$ | 4590 | 4481 |
| $`\mathrm{\Gamma }_s`$ | 0.22 | 0.80 | $`A`$ | 1695 | 12544 |
| $`m_{Q_c}^2`$ | $`(2706)^2`$ | $`(2450)^2`$ | $`m_{d_d}^2`$ | $`(4693)^2`$ | $`(2125)^2`$ |
| $`m_{u_c}^2`$ | $`(2649)^2`$ | $`(2418)^2`$ | $`m_{d_c}^2`$ | $`(2734)^2`$ | $`(2486)^2`$ |
| $`m_{\overline{h}_c}^2`$ | $`(1008)^2`$ | $`(5622)^2`$ | $`m_{h_b^{}}^2`$ | $`(826)^2`$ | $`(2595)^2`$ |
| $`m_{\phi _{20^{}}}^2`$ | $`(518)^2`$ | $`(6890)^2`$ | $`m_{\phi _{22^{}}}^2`$ | $`(3031)^2`$ | $`(11540)^2`$ |
| $`m_{h_a}^2`$ | $`(3626)^2`$ | $`(3982)^2`$ | $`m_{h_c}^2`$ | $`(224)^2`$ | $`(5633)^2`$ |
| $`m_{h_d}^2`$ | $`(3666)^2`$ | $`(4100)^2`$ | $`m_{h_e}^2`$ | $`(4274)^2`$ | $`(4246)^2`$ |
| $`m_{e_a}^2`$ | $`(2770)^2`$ | $`(3564)^2`$ | $`m_{e_f}^2`$ | $`(2780)^2`$ | $`(3958)^2`$ |
| $`m_{e_e}^2`$ | $`(4195)^2`$ | $`(4254)^2`$ | $`m_{e_h}^2`$ | $`(4259)^2`$ | $`(4236)^2`$ |
Table II: $`P_1^{}P_2^{}P_3^{}`$ flat direction: values of the parameters at $`M_{String}`$ and $`M_Z`$, with $`M_Z^{}=735`$ GeV and $`\alpha _{ZZ^{}}=0.005`$. All mass parameters are given in GeV.
## VIII The signature at the Tevatron for the light doubly charged Higgsino of the supersymmetric left-right model
B. Dutta, R. N. Mohapatra and D. J. Muller
### VIII.1 Introduction
Supersymmetric left-right models (SUSYLR) where the $`SU(2)_R`$ gauge symmetry is broken by triplet Higgs fields $`\mathrm{\Delta }^c`$ with $`BL=2`$ have many attractive features: 1) they imply automatic conservation of baryon and lepton number moh ; 2) they provide a natural solution to the strong and weak CP problems of the MSSM rasin ; 3) they yield a natural embedding of the see-saw mechanism for small neutrino masses gell where the right-handed triplet field ($`\mathrm{\Delta }^c`$) that breaks the $`SU(2)_R`$ symmetry also gives heavy mass to the right-handed Majorana neutrino needed for implementing the see-saw mechanism.
Recently it has been shown that the doubly charged components of the triplet Higgs fields are massless unless there are some higher dimensional operators (HDO) kuchi ; aulakh ; chacko ; goran2 . This is independent of how supersymmetry breaking is transmitted to the visible sector (i.e., whether it is gravity mediated or it is gauge mediated) and also of whether the hidden sector supersymmetry breaking scale is above or below the $`W_R`$ scale. In the presence of HDO’s, they acquire masses of order $`v_R^2/M_{\mathrm{Pl}}`$. Since the measurement of the Z-width at LEP and SLC implies that such particles must have a mass of at least 45 GeV, this puts a lower limit on the $`W_R`$ scale of about $`10^{10}`$ GeV or so. For $`W_R`$ near this lower limit, the masses of the doubly charged particles are in the 100 GeV range. The rest of the particle spectrum below the $`W_R`$ scale can be the same as that of the MSSM with a massive neutrino or it can have an extra pair of Higgs doublets in the 10 TeV range depending on the structure of the model.
The ordering of the sparticles masses and the existence of the doubly-charged Higgs fields in the SUSYLR model makes the SUSY signature distinctive from other SUSY models. Mass spectra for this type of model have been studied dm . It was found that one of the Higgs fields has a coupling with the third generation charged leptons which reduces the third generation charged slepton masses. Because of this, in SUSYLR models with GMSB the lighter stau is predominantly the NLSP whenever the deltino is too massive to play that role. As a result, the decay chains of the SUSY particles typically lead to the lighter stau. The $`\stackrel{~}{\tau }_1`$ then decays into a $`\tau `$ lepton and a gravitino ($`\stackrel{~}{G}`$) which escapes the detector undetected (leading to missing energy). Since the gravitino mass is on the order of eV, the emitted $`\tau `$ will have high $`p_T`$ enhancing its detection possibility. Moreover, pair production of the light doubly charged Higgsinos always produces four $`\tau `$ leptons. When the $`\stackrel{~}{\tau }_1`$ is the NLSP, this occurs through $`\stackrel{~}{\mathrm{\Delta }}^{c\pm \pm }\stackrel{~}{\tau }_1^\pm \tau ^\pm `$ followed by $`\stackrel{~}{\tau }_1\tau \stackrel{~}{G}`$. When the $`\stackrel{~}{\mathrm{\Delta }}^{c\pm \pm }`$ is the NLSP, this occurs through the stau mediated decay $`\stackrel{~}{\mathrm{\Delta }}^{c\pm \pm }\tau ^\pm \tau ^\pm \stackrel{~}{G}`$. One can get a similar signal in supergravity motivated LR models with the gravitinos replaced by the lightest neutralino (its mass is greater than 33 GeV), which will constitute the missing energy.
Signals involving two or more high $`p_T`$ $`\tau `$ leptons are also important signals for conventional GMSB models as the lighter stau is frequently the NLSP for these models as well. In the SUSYLR model, however, we find that the production of the deltino can greatly enhance the signal. In addition, since the deltino decays into like sign $`\tau `$ leptons, we find that the distribution in angle between same sign $`\tau `$ leptons can be used to distinguish this model from other GMSB models.
### VIII.2 Sparticle Masses and Production
The particle content of this model above the LR scale includes $`\varphi (2,2,0)`$, $`\mathrm{\Delta }(3,1,2)`$, $`\overline{\mathrm{\Delta }}(3,1,2)`$, $`\mathrm{\Delta }^c(1,3,2)`$, $`\overline{\mathrm{\Delta }}^c(1,3,2)`$ and a singlet where the numbers in the parentheses refer to their transformation properties under $`SU(2)_L\times SU(2)_R\times U(1)_{BL}`$. After integrating out the fields at the left-right scale, we are left with the following additional part to the MSSM:
$$W=M_\mathrm{\Delta }\mathrm{\Delta }^c\overline{\mathrm{\Delta }}^{c++}+f_il^cl^c\mathrm{\Delta }^c$$
(21)
where we have assumed that f is diagonal. The PSI experiment psi has put an upper bound on the product of the first two generation couplings of $`f_1f_2<1.2\times 10^3`$. The magnitude of $`f_3`$ is unrestrained. The term $`M_\mathrm{\Delta }`$ originates from the nonrenormalizable terms.
The model considered here involves gauge mediated supersymmetry breaking. In GMSB type models the SUSY breaking is communicated to the observable sector by the SM gauge interactions. We choose GMSB since the lighter stau will then have a mass that is almost always below that of the lightest neutralino due to the presence of the additional coupling $`f_3`$. The lighter stau then decays to a $`\tau `$ lepton and a gravitino. Since the gravitino is very light, the $`\tau `$ lepton will typically be very energetic.
There are a number of potential SUSY production mechanisms here. Given the current lower bounds on the various sparticle masses and the hierarchy of sparticle masses in GMSB models, the important SUSY production mechanisms will typically include EW gaugino production. At the Tevatron, chargino pair ($`\chi _1^+\chi _1^{}`$) production takes place through s-channel $`Z`$ and $`\gamma `$ exchange and $`\chi _2^0\chi _1^\pm `$ production is through s-channel $`W`$ exchange. Squark exchange via the t-channel also contributes to both processes, but the contributions are expected to be negligible since the squark masses are large in GMSB models. The production of $`\chi _1^0\chi _1^\pm `$ is suppressed due to the smallness of the coupling involved.
In addition to these usual SUSY production mechanisms of the MSSM, we also have deltino pair ($`\stackrel{~}{\mathrm{\Delta }}^{c++}\stackrel{~}{\mathrm{\Delta }}^c`$) production. This proceeds through s-channel $`Z`$ and $`\gamma `$ exchange. Given that the $`\stackrel{~}{\mathrm{\Delta }}^{c\pm \pm }`$ can be relatively light, it can be a very important SUSY production mode. In fact, it frequently is the dominant mode. The masses of some of the particles of interest are given in Fig. 14 and Fig. 15. In Fig. 14 we take $`M_{\stackrel{~}{\mathrm{\Delta }}}(M)=90`$ GeV (97 GeV at the weak scale), but the masses of the gauginos and sleptons (with the exception of the stau) do not vary much with the messenger scale deltino mass (M is the scale at which the soft breaking masses are introduced in the observable sector) . Fig. 15 gives the masses of the delta boson and the deltino. The deltino mass is not very sensitive to the value of $`\mathrm{\Lambda }`$, while the delta boson mass is highly dependent on $`\mathrm{\Lambda }`$ due to the contributions from the messenger scale loops.
The cross sections for the more traditional SUSY production modes are given in Fig. 16. We also have deltino pair production; the cross sections for which are tabulated in Table 5. Since the deltino mass does not vary much over the values of $`\mathrm{\Lambda }`$ ($`\mathrm{\Lambda }`$ is related to the SUSY breaking scale) considered, the cross section for deltino pair production does not vary much either. This cross section is high enough for all the deltino masses considered that deltino pair production is always an important SUSY production mode.
### VIII.3 Tau Jet Analysis
We now give an account of the possible $`\tau `$-jet signatures for SUSY production at the Tevatron in the context of the left-right GMSB model.
This analysis is performed in the context of the Main Injector (MI) and TeV33 upgrades of the Tevatron collider. The center of mass energy is taken to be $`\sqrt{s}=2`$ TeV and the integrated luminosity is taken to be 2 fb<sup>-1</sup> for the MI upgrade and 30 fb<sup>-1</sup> for the TeV33 upgrade.
In performing this analysis, the cuts employed are that final state charged leptons must have $`p_T>10`$ GeV. Jets must have $`E_T>10`$ GeV and $`|\eta |<2`$. In addition, hadronic final states within a cone size of $`\mathrm{\Delta }R\sqrt{(\mathrm{\Delta }\varphi )^2+(\mathrm{\Delta }\eta )^2}=0.4`$ are merged to a single jet. Leptons within this cone radius of a jet are discounted. For a $`\tau `$-jet to be counted as such, it must have $`|\eta |<1`$. The most energetic $`\tau `$-jet is required to have $`E_T>20`$ GeV. In addition, a missing transverse energy cut of $`\text{/}E_T>30`$ GeV is imposed.
In Table 6 we give the inclusive $`\tau `$-jet production cross sections for a messenger scale deltino mass of 90 GeV, respectively. We include only up to four $`\tau `$-jets as the cross sections for more than four $`\tau `$-jets are small. We see that before cuts the production of two and three $`\tau `$-jets are dominant, but the four $`\tau `$-jet cross section is also significant at slightly over 100 fb. After the cuts are applied, however, the situation changes substantially. The one $`\tau `$-jet mode is now dominant, but the cross section for two $`\tau `$-jets is not far below and the three $`\tau `$-jets cross section is not insignificant. For $`\mathrm{\Lambda }=35`$ TeV the cross section for inclusive production of three $`\tau `$-jets is 32.3 fb. For an integrated luminosity of 2 fb<sup>-1</sup> (the approximate initial value at Run II), this corresponds to about 65 events. For 30 fb<sup>-1</sup>, the number of observable events is $`970`$.
In comparison to the GMSB model with the MSSM symmetry, the two $`\tau `$-jets and the three $`\tau `$-jets cross sections are considerably higher in this model. In the GMSB model with MSSM symmetry, the two $`\tau `$-jets cross section can be seen at RUN II, but not the three $`\tau `$ jets mdn .
### VIII.4 Angular Distributions
The excess of $`\tau `$-jets expected in this model does not constitute an unequivocal signal for this model. $`\tau `$-jets are part of the signatures for other models including the minimal GMSB model when the lighter stau is the NLSP. The question then arises as to whether there is any way to distinguish this model from the minimal GMSB model. A possible distinguishing characteristic is the distribution in angle between the two highest $`E_T`$ $`\tau `$-jets when they come from same sign $`\tau `$-jets.
Consider deltino pair production. The deltino tends to decay to like sign $`\tau `$ leptons. This occurs directly when the deltino is the NLSP and so decays via the three-body decay $`\stackrel{~}{\mathrm{\Delta }}^{\pm \pm }\tau ^\pm \tau ^\pm \stackrel{~}{G}`$. When the two-body decay of the deltino $`\stackrel{~}{\mathrm{\Delta }}^{\pm \pm }\stackrel{~}{\tau }_1^\pm \tau ^\pm `$ occurs, then the second like sign $`\tau `$ lepton comes from the subsequent decay of the stau. In the rest frame of the deltino, the $`\tau `$ leptons are widely distributed. In the lab frame, however, the deltinos are quite energetic and have a large velocity, especially if their masses are small. As a consequence of this, the decay products of the deltino tend to be collimated in the direction in which the deltino was moving. Thus when the two most energetic $`\tau `$-jets have the same sign in deltino pair production, the angle between them tends to be smaller than when the two most energetic $`\tau `$-jets have opposite sign charges.
Fig. 17 gives the distribution in angle between the two most energetic $`\tau `$-jets for deltino pair production. This example is for a weak scale deltino mass of about 97 GeV. We can see that the distribution in angle for like sign $`\tau `$-jets, which is given in Fig. 17(a), peaks at about $`40^{}`$. Fig. 17(b) gives the distribution in angle between the two most energetic $`\tau `$-jets when they come from opposite sign $`\tau `$ leptons. In stark contrast to the previous case, here the peak occurs at $`110^{}`$.
The situation changes as the deltino mass gets larger. This is in part due to the fact that the deltino pair production cross section gets smaller and so production of charginos and neutralinos can have a larger impact on the distributions. In addition, a larger deltino mass means the deltinos will typically be moving slower.
Fig. 19 shows the angular distributions for combined $`\chi _2^0\chi _1^\pm `$ and $`\chi _1^+\chi _1^{}`$ production for the weak scale $`\chi _2^0`$ mass is $`100`$ GeV. The distribution for same sign $`\tau `$-jets is given in Fig 5(a). We see that the peak occurs at about $`110^{}`$. In $`\chi _2^0\chi _1^\pm `$ production, one of the same sign $`\tau `$-jets generally comes from the chargino and the other from the neutralino. We now consider the angular distribution for opposite sign $`\tau `$-jets which are given in Fig. 5(b). In $`\chi _2^0\chi _1^\pm `$ production, opposite sign $`\tau `$-jets frequently come from the neutralino, while in $`\chi _1^+\chi _1^{}`$ production one of the $`\tau `$-jets comes from one of the charginos and the other $`\tau `$-jet comes from the other chargino. Since there is a strong possibility that the opposite sign $`\tau `$-jets come from the same particle ($`\chi _2^0`$), the distribution should peak at a lower angle than for same sign $`\tau `$-jets. We see from the figure that the peak occurs at about $`85^{}`$. This feature does not change with the increase in the gaugino mass.
### VIII.5 Conclusion
In conclusion, we have found that the doubly charged Higgs bosons of LR models can be potentially observable at Run II of the Tevatron through the production of $`\tau `$-jets. In a GMSB type theory, SUSYLR models typically produce large numbers of two and three $`\tau `$-jet final states. This large $`\tau `$-jet signal is also due in large part to pair production of the doubly charged Higgsino. It is also due to the relatively low mass of the lighter stau (which is frequently the NLSP) in these models, which is due to the additional coupling $`f`$. We have also shown that the distribution in angle between the two highest $`E_T`$ $`\tau `$-jets is different from other models which do not have this doubly charged Higgsino.
## IX Indirect Signals for Extra Dimensions
J.L. Hewett
One manifestation of theories of low scale quantum gravity is the existence of a Kaluza Klein (KK) tower of massive gravitons which can interact with the SM fields on the wall. In this section we examine the indirect effects of these massive gravitons being exchanged in Drell-Yan production. We consider a novel feature of this theory, namely, the contribution of gluon-gluon initiated processes to lepton pair production.
The effective theory below the effective Planck scale in the bulk, $`M_{eff}`$, consists of the SM fields on the wall and gravity which propagates in the full $`4+n`$ bulk. The interactions of these fields are given by
$$d^{4+n}xT^{\widehat{\mu }\widehat{\nu }}\frac{h_{\widehat{\mu }\widehat{\nu }}(x^\mu ,x^a)}{M_{eff}^{n/2+1}},$$
(22)
where $`T^{\widehat{\mu }\widehat{\nu }}`$ is the symmetric, conserved stress-energy tensor in the bulk and $`h_{\widehat{\mu }\widehat{\nu }}`$ is the graviton field-strength tensor, which can be decomposed into spin-2, 1, and 0 fields. Here the indices $`\widehat{\mu }`$ extend over the full $`4+n`$ dimensions, $`\mu `$ over the $`3+1`$ dimensions on the wall, and $`a`$ over the $`n`$ bulk dimensions. The interactions with the SM matter fields are obtained by decomposing the above into the 4-dimensional states. The bulk fields $`h_{\widehat{\mu }\widehat{\nu }}`$ appear as Kaluza-Klein towers in the 4-dimensional space arising from a Fourier analysis over the cyclic boundary conditions of the compactified dimensions. Performing this decomposition, we immediately see that $`T_{\mu a}=0`$ and hence the spin-1 KK states don’t interact with the wall fields. The scalar, or dilaton, states couple proportionally to the trace of the stress-energy tensor. For interactions with fermions, this trace is linear in the fermion mass, while for gauge bosons it is quadratic in the boson mass. Hence, the dilaton does not contribute to the processes under consideration here.
We thus only have to consider the interactions of the KK spin-2 gravitons with the SM fields. All the gravitons in the KK tower, including the massless state, couple in an identical manner. Hence we may use the couplings to matter as obtained in the case of linearized general relativitygr . In this linearized theory, the matrix element for $`q\overline{q}\mathrm{}^+\mathrm{}^{}`$ generalized for the case of $`n`$ massive graviton exchanges can be written as
$$=\frac{1}{M_{Pl}^2}\underset{n}{}\frac{T_{\mu \nu }^eP^{\mu \nu \lambda \sigma }T_{\lambda \sigma }^f}{sm_{gr}^2[n]},$$
(23)
where the sum extends over the KK modes. $`P_{\mu \nu \lambda \sigma }`$ represents the polarization sum of the product of two graviton fields and is given in gr . The terms in the polarization sum that are quadratic and quartic in the transferred momentum do not contribute to the above matrix element since $`T_{\mu \nu }`$ is conserved. Likewise, the terms which go as $`\eta _{\mu \nu }\eta _{\lambda \sigma }`$ lead to terms proportional to $`T_\mu ^{e\mu }T_\lambda ^{f\lambda }`$ which vanish in the limit of zero electron mass. The remaining terms are $`P_{\mu \nu \lambda \sigma }=\frac{1}{2}[\eta _{\mu \lambda }\eta _{\nu \sigma }+\eta _{\mu \sigma }\eta _{\nu \lambda }\eta _{\mu \nu }\eta _{\lambda \sigma }]`$ and are exactly those present in the massless graviton case; they are thus universally applicable to all of the states in the KK tower. Since the spacing of the KK states is given by $`1/r`$, the sum over the states in (23) above can be approximated by an integral which is log divergent for $`n=2`$ and power divergent for $`n>2`$. A cut-off must then be applied to regulate these ultraviolet divergences, and is generally taken to be the scale of the new physics. For $`n>2`$ it can be shownsumm that the dominant contribution to this integral is of order $`M_{Pl}^2/M_s^4`$, where we have taken the cut-off to be the string scale, while for $`n=2`$ this result is multiplied by a factor of order $`\mathrm{log}(M_s^2/E^2)`$, where $`E`$ is the center-of-mass energy of the process under consideration. The exact computation of this integral can only be performed with some knowledge of the full underlying theory. Combining these results yields the matrix element
$``$ $`=`$ $`{\displaystyle \frac{\lambda }{M_s^4}}\{\overline{q}(p_1)\gamma _\mu q(p_2)\overline{\mathrm{}}(p_3)\gamma ^\mu \mathrm{}(p_4)(p_2p_1)(p_4p_3)`$
$`\overline{q}(p_1)\gamma _\mu q(p_2)\overline{\mathrm{}}(p_3)\gamma _\nu \mathrm{}(p_4)(p_2p_1)^\nu (p_4p_3)^\mu \}.`$
Here, the momentum flow is defined with $`p_{1,2}`$ into the vertex and $`p_{3,4}`$ outgoing. Note that graviton exchange is $`C`$ and $`P`$ conserving, and is independent of the flavor of the final state. The coefficient $`\lambda `$ is of $`𝒪(1)`$ and cannot be explicitly calculated without knowledge of the full quantum gravity theory. It is dependent on the number of extra dimensions, how they are compactified, and is in principle a power series in $`s/M_s^2`$. However, we neglect this possible energy dependence in $`\lambda `$ and note that the limits obtained here, which go as $`|\lambda |^{1/4}`$, are only very weakly dependent on its precise value and hence on the specific model realization. In principle the sign of $`\lambda `$ is undetermined and we examine the constraints that can be placed on $`M_s`$ with either choice of signs.
The angular distribution for $`q\overline{q}\mathrm{}^+\mathrm{}^{}`$ with massless leptons is then calculated to be
$`{\displaystyle \frac{d\sigma }{dz}}`$ $`=`$ $`N_c{\displaystyle \frac{\pi \alpha ^2}{2s}}\{P_{ij}[A_{ij}^qA_{ij}^{\mathrm{}}(1+z^2)+2B_{ij}^qB_{ij}^{\mathrm{}}z]`$
$`{\displaystyle \frac{\lambda s^2}{2\pi \alpha M_s^4}}P_i\left[2z^3v_i^qv_i^{\mathrm{}}(13z^2)a_i^qa_i^{\mathrm{}}\right]`$
$`+{\displaystyle \frac{\lambda ^2s^4}{16\pi ^2\alpha ^2M_s^8}}[13z^2+4z^4]\},`$
where the indices $`i,j`$ are summed over $`\gamma `$ and $`Z`$ exchange, $`z=\mathrm{cos}\theta `$, $`P_{ij}`$ and $`P_i`$ are the usual propagator factors (defined in e.g., leptos ), $`A_{ij}^f=(v_i^fv_j^f+a_i^fa_j^f)`$ , and $`B_{ij}^f=(v_i^fa_j^f+v_j^fa_i^f)`$. Note that the total cross section from $`q\overline{q}`$ annihilation is unaltered by graviton exchanges, independently of fermion flavor up to terms of order $`s^4/M_s^8`$, and hence the angular distributions will be most sensitive to these new exchanges.
In addition to the $`q\overline{q}`$ channel, gravitons can also mediate gluon-gluon contributions to lepton pair production via $`s`$-channel exchange, as noted above. Such gluon initiated processes are a remarkable consequence of this theory and have the potential to modify the Drell-Yan spectrum in a unique manner. Following an analogous procedure as outlined above for the four-fermion case, the matrix element for $`gg\mathrm{}^+\mathrm{}^{}`$ via graviton exchanges is found to be
$``$ $`=`$ $`{\displaystyle \frac{\lambda }{4M_s^4}}\overline{\mathrm{}}(p^{})[(p^{}p)_\mu \gamma _\nu +(p^{}p)_\nu \gamma _\mu ]\mathrm{}(p)\{k_\alpha ^{}(k_\mu \eta _{\beta \nu }+k_\nu \eta _{\beta \mu })+k_\beta (k_\mu ^{}\eta _{\alpha \nu }+k_\nu ^{}\eta _{\alpha \mu })`$
$`\eta _{\alpha \beta }(k_\mu ^{}k_\nu +k_\mu k_\nu ^{})+\eta _{\mu \nu }(k^{}k\eta _{\alpha \beta }k_\beta k_\alpha ^{})kk^{}(\eta _{\mu \alpha }\eta _{\nu \beta }+\eta _{\mu \beta }\eta _{\nu \alpha }\}ϵ_g^\beta (k^{})ϵ_g^\alpha (k),`$
where the momentum flow is defined with both $`k,k^{}`$ flowing into the vertex and $`p,p^{}`$ being outgoing and $`ϵ`$ represents the gluon polarization vector. Because the graviton couplings and the summation over the KK tower of states for $`22`$ processes are universal, $`\lambda `$ is the same $`𝒪(1)`$ coefficient as in Eq. (IX). This matrix element yields the $`gg\mathrm{}^+\mathrm{}^{}`$ differential cross section for massless leptons
$$\frac{d\sigma }{dz}=\frac{\lambda ^2\widehat{s}^3}{64\pi M_s^8}(1z^2)(1+z^2),$$
(27)
which has a remarkably simple form. While the overall numerical coefficient appears to be very small, it must be compared to $`\alpha ^2`$ which appears in the usual contributions to Drell-Yan production. In addition, the large parton luminosity for gluons at higher energy colliders may also compensate for the small numerical factor. Since this cross section is also even in $`\mathrm{cos}\theta `$, the gluon-gluon contributions will only affect the total cross section and not the forward-backward asymmetry. Also note that the ambiguity in the sign of $`\lambda `$ does not affect the gluon-gluon contributions as they do not interfere with the $`q\overline{q}`$ initiated process.
The bin integrated lepton pair invariant mass distribution and forward-backward asymmetry $`A_{FB}`$ is presented in Fig. 21. In each case the solid histogram represents the SM expectations, and the ‘data’ points include the graviton exchanges with the error bars representing the statistics in each bin. The rapidity cuts, parton density parameterizations, and assumed integrated luminosity are as labeled, and we have summed over electron and muon final states. Here we present the sample case of $`M_s=800`$ GeV and the sign ambiguity in $`\lambda `$ is visible in the forward-backward asymmetry. Since the graviton exchanges only affect the invariant mass distribution at order $`\lambda ^2/M_s^8`$, we expect only minor modifications to this spectrum and we see that this holds true. For higher energy colliders which have a larger gluon parton luminosity, such as the LHC, large string scales do have a sizable effect on the $`M_{\mathrm{}\mathrm{}}`$ spectrumjlh . However, the deviations in $`A_{FB}`$, are more pronounced at the Tevatron, where even the two cases $`\lambda =\pm 1`$ are statistically distinguishable from each other for this sample case. The resulting 95% C.L. search reach are given in Fig. 23. We also find that present Tevatron data from Run I with 110 $`\mathrm{pb}^1`$ of integrated luminosity excludes a string scale up to 980 (920) GeV at 95% C.L. for $`\lambda =1(+1)`$.
## X Higgs Bounds in Three and Four Generation Scenarios
D. Dooling, K. Kang and S.K. Kang
### X.1 Introduction
The search for the Higgs boson being one of the major tasks along with that for supersymmetric sparticle and fourth generation fermions at future accelerators such as LEP200 and LHC makes it a theoretical priority to examine the bounds on the Higgs boson mass in the SM and its supersymmetric extension and to look for any distinctive features. The actual measurement of the Higgs boson mass could serve to exclude or at least to distinguish between the SM(3,4) and the MSSM(3,4) models for electroweak symmetry breaking. Recently, bounds on the lightest Higgs boson mass were calculated in 1 ; 2 ; 3 ; 4 ; 5 ; 6 ; 7 ; 8 ; 9 . It was found that for a measured $`M_H`$ lying in a certain mass range, both the SM vacuum stability lower bound and the MSSM upper bound are violated, thus shaking our confidence in these theories just as the final member of the mass spectrum is observed. One method of curing this apparent illness is to take a leap of faith by adding another fermion generation, to fortify these theories with another representation of the gauge group. This additional matter content, for certain ranges of its mass values, has the desired effect of raising the MSSM3 upper bound above that of the SM lower bound and avoids the necessity of being forced to introduce completely new physics.
In this work, we use the latest LEP Electroweak Working Group data as well as the most recent experimental lower limits on the masses of fourth generation fermions to see how the Higgs boson mass bounds are affected. With only the standard three generations, a violation of both the SM and the MSSM bounds occurs for a broad range of $`M_H`$ values, signalling the need for an additional generation. Our presentation is organized as follows. We first present a summary of our one-loop effective potential (EP) improved by the two-loop renormalization group equations (RGE) analysis. An expression for the Higgs boson mass is derived up to the next-to-leading logarithm order. Bounds on $`M_H`$ are obtained by imposing different boundary conditions on the Higgs self-coupling $`\lambda `$. Finally, we present our results and also obtain information on the possible mass range of the fourth generation leptons, $`M_L`$.
### X.2 Effective Potential Approach
As shown in 10 , in order to calculate the Higgs boson mass up to the next-to-leading logarithm approximation, we must consider the one-loop EP improved by two-loop RGE for the $`\beta `$ and $`\gamma `$ functions of the running coupling constants, masses and the $`\varphi `$ field for the Higgs boson 11 .
The two-loop RGE improved one loop EP of the SM4 is given by
$$V_1=V_{(0)}+V_{(1)}$$
(28)
where
$`V_{(0)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}m^2(t)\varphi _c^2(t)+{\displaystyle \frac{1}{24}}\lambda (t)\varphi _c^4(t)`$ (29)
$`V_{(1)}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{5}{}}}\left({\displaystyle \frac{\kappa _i}{64\pi ^2}}\right)h_i^4(t)\varphi _c^4(t)\left[\mathrm{ln}{\displaystyle \frac{h_i^2(t)\zeta ^2(t)}{2}}{\displaystyle \frac{3}{2}}\right]`$ (30)
Here $`\varphi _c`$ is the classical field corresponding to the physical Higgs boson $`\varphi `$, $`i=(t,T,B,N,E)`$ , $`\kappa _i=3`$ for $`i=(t,T,B)`$ and $`\kappa _i=1`$ for $`i=(N,E)`$ and $`h_i`$ is the Yukawa coupling of the i<sup>th</sup> fermion to the Higgs field. In addition,
$`\zeta (t)`$ $`=`$ $`\left({\displaystyle _0^t}\gamma _\varphi (t^{})𝑑t^{}\right)`$ (31)
$`\varphi _c(t)`$ $`=`$ $`\varphi _c\zeta (t)`$ (32)
and
$$\mu (t)=\mu \mathrm{exp}(t)$$
(33)
where $`\mu `$ is a fixed scale.
Starting with the above expression for the SM(3,4) EP, one may follow the analysis in 1 ; 4 and obtain the expression for the Higgs boson mass:
$`m_\varphi ^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}\lambda v^2+{\displaystyle \frac{\mathrm{}v^2}{16\pi ^2}}\{{\displaystyle \frac{\lambda ^2}{3}}+2\lambda (h_t^2+h_T^2+h_B^2)+{\displaystyle \frac{2\lambda }{3}}(h_E^2+h_N^2)`$ (34)
$`{\displaystyle \frac{\lambda }{2}}(3g_2^2+g_1^2)+{\displaystyle \frac{9}{8}}g_2^4+{\displaystyle \frac{3}{4}}g_1^2g_2^2+{\displaystyle \frac{3}{8}}g_1^4`$
$`6h_t^4\mathrm{ln}{\displaystyle \frac{h_t^2\zeta ^2}{2}}6h_T^4\mathrm{ln}{\displaystyle \frac{h_T^2\zeta ^2}{2}}6h_B^4\mathrm{ln}{\displaystyle \frac{h_B^2\zeta ^2}{2}}`$
$`2h_E^4\mathrm{ln}{\displaystyle \frac{h_E^2\zeta ^2}{2}}2h_N^4\mathrm{ln}{\displaystyle \frac{h_N^2\zeta ^2}{2}}\}+O(\mathrm{}^2)`$
where the three generation case is obtained by simply letting the fourth generation Yukawa couplings go to zero.
Following the method similar of Kodaira et al. 4 , we arrive at the appropriate energy scale at which to evaluate $`\frac{^2V}{\varphi _c(t)^2}`$ by requiring $`\frac{V}{\varphi _c(t_v)}=`$ 0 at the scale $`t_v`$ where $`\varphi _c(t_v)=v=(\sqrt{2}G_F)^{\frac{1}{2}}`$ = 246 GeV. $`t_v`$ is found to satisfy:
$$t_v=\mathrm{ln}\frac{v}{\mu }+_0^{t_v}\gamma _\varphi (t^{})𝑑t^{}.$$
We then evaluate the first and second derivatives of the EP at the scale $`t_v`$ where $`\varphi _c(t_v)=v`$. In the above relation satisfied by $`t_v`$, $`\mu `$ is a fixed, constant mass scale. In the MSSM theories, we will take $`\mu `$ to be $`M_{susy}=1`$ or 10 TeV, while the SM lower bounds will be derived after choosing $`\mu =\mathrm{\Lambda }=10^{19}`$ GeV.
We define the running Higgs mass as:
$$m_h^2(t)=\frac{m_h^2(t_v)\zeta ^2(t_v)}{\zeta ^2(t)}.$$
The physical, pole masses are related to the running masses via the following equations:
$$M_i=\left[1+\beta _i\frac{4\alpha _3(M_i)}{3\pi }\right]m_i(M_i)$$
$$M_H^2=m_h^2(t)+Re\mathrm{\Pi }(M_H^2)Re\mathrm{\Pi }(0)$$
where $`\beta _i`$ = 0 for i = (N,E) and = 1 for i = (t,T,B). $`\mathrm{\Pi }(q^2)`$ is the renormalized electroweak self-energy of the Higgs boson.
The method of solving the RGE and the appropriate boundary conditions for the couplings is explained in 1 . In this update, we use $`M_Z=91.1867`$ GeV and $`\alpha _3(M_Z)=.119`$.
### X.3 Bounds on $`M_H`$
Let us first discuss the procedure for determining a lower bound on the Higgs boson mass in the SM 5 ; 12 Working with the two-loop RGE requires the imposition of one-loop boundary conditions on the running parameters. As pointed out by Casas et al. 5 ; 7 , the necessary condition for vacuum stability is derived from requiring that the effective coupling $`\stackrel{~}{\lambda }(\mu )>`$ 0 rather than $`\lambda >0`$ for $`\mu (t)<\mathrm{\Lambda }`$, where $`\mathrm{\Lambda }`$ is the cut-off beyond which the SM is no longer valid. The effective coupling $`\stackrel{~}{\lambda }`$ in the SM4 is defined as:
$$\stackrel{~}{\lambda }=\frac{\lambda }{3}\frac{1}{16\pi ^2}\left\{\underset{i=1}{\overset{5}{}}2\kappa _ih_i^4\left[\mathrm{ln}\frac{h_i^2}{2}1\right]\right\}$$
where the three generation case is simply the same as the above expression without the fourth generation Yukawa coupling contributions. Choosing $`\mathrm{\Lambda }=10^{19}`$ GeV and $`M_{top}=172`$ GeV, we arrive at a vacuum stability lower bound on $`M_h`$ of $``$ 134 GeV for the SM with three generations. Allowing $`M_{top}`$ to be as large as 179 GeV increases the lower bound on $`M_H`$ to $``$ 150 GeV.
To compute the MSSM upper bound on $`M_H`$, we assume that all of the sparticles have masses $`O(M_{susy})`$ or greater and that of the two Higgs isodoublets of the MSSM, one linear combination is massive, also with a mass of $`O(M_{susy})`$ or greater, while the other linear combination, orthogonal to the first, has a mass of the order of weak-scale symmetry breaking. With these two assumptions, it is clear that below the supersymmetry breaking scale $`M_{susy}`$, the effective theory is the SM. This fact enables us to use the SM effective potential for the Higgs boson when we treat the lightest Higgs boson in the MSSM.
In the MSSM(3,4), the boundary condition for $`\lambda `$ at $`M_{susy}`$ is
$$\frac{\lambda }{3}(M_{susy})=\frac{1}{4}\left[g_1^2(M_{susy})+g_2^2(M_{susy})\right]\mathrm{cos}^2(2\beta )+\frac{\kappa _ih_i^4(M_{susy})}{16\pi ^2}\left(2\frac{X_i}{M_{susy}^2}\frac{X_i^4}{6M_{susy}^4}\right)$$
where $`\kappa _i`$ = 3 for $`i=(t,T,B)`$ and $`\kappa _i`$ = 1 for $`i=(N,E)`$ and $`X_i`$ is the supersymmetric mixing parameter for the ith fermion. Zero threshold corrections correspond to $`X_i`$ = 0. Maximum threshold corrections occur for $`X_i=6M_{susy}^2`$.
In Fig. (1) we present our numerical two-loop results for the lightest Higgs boson mass bounds in the SM and the MSSM3 as a function of the supersymmetric parameter $`\mathrm{cos}^2(2\beta )`$. The bottom two curves correspond to the MSSM3 upper bound for the two cases $`M_{top}=172`$ GeV and the slightly greater upper bound that results when $`M_{top}`$ = 179 GeV and with no threshold corrections. When the case of maximum threshold corrections is considered, these two curves are translated upwards by $``$ 55 GeV - 60 GeV, illustrating the strong dependence of the upper bound on the precise value of the threshold corrections. Yet even with such a dramatic increase in the upper bounds with increasing threshold corrections, we observe that the SM lower bound exceeds the MSSM upper bound for $`M_{top}=172`$ GeV and $`0<\mathrm{cos}^2(2\beta )<.2`$ for all values of the threshold correction contribution. Similarly, for $`M_{top}=179`$ GeV, the troublesome situation is only exacerbated, as the SM lower bound exceeds the MSSM upper bound for $`0<\mathrm{cos}^2(2\beta )<.38`$ independent of the threshold corrections.
In Fig.(2) we present the problem more clearly. Taking into account the present experimental lower limit on $`M_H`$ of $``$ 90 GeV at 95$`\%`$ CL, we find the value of the threshold correction that gives a smallest upper bound consistent with the experimental lower limit. Clearly, for this phenomenologically determined lower limit of the threshold contributions, there is a large area in $`M_H\times \mathrm{cos}^2(2\beta )`$ space that is inconsistent with both the SM and the MSSM. For $`M_{top}=172`$ GeV, the region 92 GeV $`<M_H<`$ 134 GeV invalidates both theories independent of $`\mathrm{cos}^2(2\beta )`$, while for $`M_{top}=179`$ GeV, the range of mutual invalidation is 92 GeV $`<M_H<`$ 150 GeV.
### X.4 Fourth Generation
To resolve the above conundrum, one would like to either raise the MSSM upper bounds, lower the SM lower bounds, or both. Adding more massive fermions to the theory only increases both bounds, so it is readily apparent that the way out of the area of inconsistency is to consider the MSSM4 and see if the additional matter of the MSSM4 results in MSSM4 upper bounds that exceed the SM3 lower bounds.
We now discuss restrictions on the possible fourth generation fermion masses 2 ; 13 ; 14 ; 15 . The close agreement between the direct measurements of the top quark at the Tevatron and its indirect determination from the global fits of precision electroweak data including radiative corrections within the framework of the SM imply that there is no significant violation of the isospin symmetry for the extra generation. Thus the masses of the fourth generation isopartners must be very close to degenerate 14 ; i.e.
$$\frac{M_T^2M_B^2}{M_Z^2}1,\frac{M_E^2M_N^2}{M_Z^2}1$$
Recently, the limit on the masses of the extra neutral and charged lepton masses, $`M_N`$ and $`M_E`$, has been improved by LEP1.5 to $`M_N>59`$ GeV and $`M_E>62`$ GeV. Also, CDF has yielded a lower bound on $`M_B`$ of $``$ 140 GeV.
In our previous work, we considered a completely degenerate fourth generation of fermions with mass $`m_4`$. We derived an upper bound on $`m_4`$ in the MSSM4 by demanding perturbative validity of all the couplings out to the GUT scale 16 . This constraint led to an upper bound on $`m_4`$ of $``$ 110 GeV. The above experimental lower limit on $`M_B`$ naturally forces us to now a consider a fourth generation where degeneracy only holds among the isodoublets separately. We therefore consider a fourth generation with masses $`M_L`$ and $`M_Q`$.
In Fig.(3), we present the SM lower bound, the MSSM4 upper bound with the fourth generation masses at their experimental lower limits and with fourth generation masses large enough to remove the problem area for all values of $`\mathrm{cos}^2(2\beta )`$. The MSSM bounds were calculated with no threshold corrections, and $`M_{top}`$ is fixed at 172 GeV. Fig.(4) shows the same information for $`M_{top}`$ = 179 GeV. The MSSM4 upper bounds are much more sensitive to $`M_Q`$ than they are to $`M_L`$. This qualitative behavior is readily understood from inspection of the equation for $`m_\varphi ^2`$. For this reason, it is necessary to increase $`M_Q`$ appropriately in order to generate a MSSM4 upper bound that is greater than the SM lower bound for all values of $`\mathrm{cos}^2(2\beta )`$. In fact, keeping $`M_Q`$ at 146 GeV and allowing $`M_L`$ to be 110 GeV does not resolve the problem. But increasing both $`M_Q`$ and $`M_L`$ as indicated in the figures does remove the problem. Because all of the bounds increase as $`M_L`$ and $`M_Q`$ increase, and because the upper bounds on $`m_4`$ from the previous work are saturated when the masses of the fourth generation reach some critical values from below, we can conclude that $`M_L`$ must still be $`<110`$ GeV. This conclusion follows because it is $`h_N`$ that violates perturbative validity, so in the non-degenerate case, it is $`M_L`$ that must still respect this upper bound if gauge coupling unification is still to be achieved in the MSSM4.
### X.5 Concluding remarks
In conclusion, we have studied the upper bounds on the lightest Higgs boson mass $`M_H`$ in the MSSM with four generations by solving the two-loop RGE’s and using the one-loop EP. We have considered a fourth generation of quarks and leptons with degenerate masses $`M_Q`$ and $`M_L`$. For certain values of $`M_Q`$ and $`M_L`$, the area of mutual inconsistency between the SM and MSSM3 Higgs mass bounds is found to be consistent with the MSSM4 upper bounds.
## XI Fourth Family Status and Prospects
D. McKay
### XI.1 The Bottom Line
In light of the general expectation from SUSY considerations and from precision electroweak fits that the lightest Higgs is light, I believe that the inclusion of the essential Higgs FCNC channel in the simulation of the various decay channels of $`b^{}`$ is the the fourth family project that deserves a strong effort among the BTMSSM searches. Analysis of the $`b^{}b+Z`$ mode is well advanced at both CDF and D0. The inclusion of the Higgs mode, with the subsequent $`Hb+\overline{b}`$ decay, will remove the last loophole in the $`b^{}`$ search up to the top mass. A good reason for stressing this issue is that it ties in directly to the crucial Higgs search mission at Run II, which makes the $`b^{}`$ search more subtle and more interesting.
A small but important project is the re-fitting of electroweak parameters to include relatively light new physics (50-150 GeV) effects. A preliminary look is shown in the “mass bounds” graph in Section 3.
More detail on these points is given in Sec. 3 below.
### XI.2 Words of Motivation
1) There is no theoretical reason against having a fourth family. Since there is no solid theoretical understanding of the family pattern, there is motivation to simply look to see if a new one shows up. Unless stated otherwise, fourth family means a family with a lepton doublet with mass greater than $`m_Z`$/2 but with an otherwise sequential, chiral set of leptons $`(\nu ^{},\tau ^{})`$ and quarks $`(t^{},b^{})`$. (It is assumed in all of the discussion that $`V_{b^{}t^{}}V_{b^{}t}V_{t^{}b}V_{b^{}c}`$
2) Though it is not in the vein of the fourth family as just defined, the proposal by Frampton and collaborators to solve the strong CP problem calls for heavy vector-like fermion states, and the search strategies are largely the same as for chiral lepton and quark states. ( see Frampton and Kephart, PRL 66, 166; Frampton and Ng, PRD 43, 3034).
3) In a four family extension of the standard model, Hung has recently argued that unification of the gauge couplings can be obtained at a high enough scale to avoid proton decay. The scale for the fourth family can be higher than in the MSSM, but it is still constrained to rather low values - lepton masses below $`m_Z`$ and quark masses at around $`m_t`$ \- by the requirement that all couplings stay perturbative up to the unification scale. If this motivation applies to SUSY, it applies to this four family SM. The hierarchy problem is left unsolved, of course. ( P.Q. Hung, PRL 80, 3000 (1998); see also P. Frampton and P.Q. Hung, Phys. Rev. D 58, 057704 (1998) for phenomenology of quasi-stable quarks in the $`m_t`$ range that could have escaped detection in the CDF and D0 analysis to date. For a study of four family consequences in the MSSM with SUGRA boundary conditions, see J. Gunion, D. McKay and H. Pois, PRD 53, 1616,(1996) ).
4) For those with a taste for dynamical breaking of electroweak symmetry, new heavy quark flavors in a fourth family may have some appeal, since the heavy quark condensate mechanism for EW breaking may have a chance if there are extra, heavy quarks in the range 150-230 GeV to help top do the job. (see for example Hung and Isidori, P.L. B402, 122).
5) Tying in directly with the Higgs boson search and BTMSSM theme of this workshop, D. Dooling, K. Kang and S.-K. Kang (hep-ph/9710258, to be published, and contribution to this workshop) have argued that the lower mass bound on the SM Higgs boson and the upper limit on the lightest Higgs in the MSSM are relaxed in four family extensions of these models. This consideration is especially timely as the limit on the SM Higgs mass has recently been pushed up to nearly 100 GeV at LEP II.
### XI.3 Experimental Facts that tightly constrain a sequential Fourth Family
1) The S parameter in precision electroweak analysis is now $`S=0.16\pm 0.14(0.10)`$ as listed by the PDG ’98 discussion of ”Electroweak model and constraints on new physics”(European Physical Journal C 3,1 (1998)). A new, heavy (several times $`M_Z`$) degenerate, chiral family adds $`2/3\pi `$ to S and this excludes an extra, heavy family of fermions at the 99.2% C.L. when S is allowed to be arbitrary and fixed by the values of the EW parameters fit at the Z (primarily). Requiring S to be positive with a given Higgs mass relaxes this bound somewhat. If some of the masses of the family members are in the neighborhood of $`m_Z`$, then the theoretical value of S is less than $`2/3\pi `$, and a more general analysis of the data is called for (I. Maksymyk, C.P. Burgess and D. London, Phys. Rev. D 50, 529 (1994)). In the “mass bounds” figure, the S parameter bound excludes all of the region to the right of the curve in this plot at the 95% C.L. $`m_Q/m_Z`$ is plotted on the y axis vs $`m_L/m_Z`$ on the x axis. The Q and L designate mass degenerate quark and lepton doublets, respectively. ($`S0.15`$) is the constraint. Clearly the (chiral) lepton doublet must be very light to satisfy this requirement. The constraint is insensitive to the quark masses because they nearly, but not quite, saturate the bound by themselves for large quark mass values, so they have little effect when $`m_Q/m_Z>1.5`$.
A completely consistent re-fit to the data with relatively light, fourth family doublets needs to be done and is in the works.
2) The direct searches for $`b^{}b+V(\gamma ,Z,g)`$ have severely narrowed the options below $`M_b^{}=140GeV`$. The D0 Collaboration ( PRL 78, 3817 (1997)) reported the $`b^{}b+\gamma `$ bound of $`m_b^{}<95GeV`$ at 95% C.L. 95 GeV is the limit for this mode, since the larger b+Z mode sets in here. Results of a study of the $`b^{}b+Z`$ mode were reported by Joao Guimaraes da Costa of CDF at the April, ’98 AAPT/APS meeting (see also the contribution by Dave Stuart to this workshop). A limit $`m_b^{}>140GeV`$ at 95% C.L. is achieved, assuming that $`\mathrm{\Gamma }(b^{}b+Z)=100\%`$. This bound is shown by the horizontal line in the “mass bounds” figure. A slip knot in this noose tightening around b’ could be the $`b^{}b+H`$ mode, which is likely to be as large as or larger than the b+Z mode. This would seriously degrade the $`b^{}`$ mass bounds quoted unless the Higgs final state were included in the tagged sample with favorable efficiencies. Failing this, roughly speaking, when the b+Higgs decay mode is accessible to $`b^{}`$, the $`b+\gamma `$ quoted number of expected events would be halved and the expected b+Z events would be quartered. But in fact, as argued by Guimaraes da Costa (private communication), if the Higgs + b final state is present, it will be included in the sample as long as the b + Z branching fraction is not zero so that the lepton tagging of the Z is operative. Moreover the efficiency for tagging the b’s from Higgs decay is good enough that the factor 4 decrease in events from the competing channel is compensated. The bottom line is that if $`m_Zm_H`$, the bound $`m_b^{}<140GeV`$ is still likely to obtain. The b+Z and b+H decay rates are roughly proportional to $`m_t^2m_t^{}^2`$, so their ratio is insensitive to $`m_t^{}`$, but when $`m_t^{}m_t`$ with $`m_b^{}<m_t^{}`$, then $`b^{}c+W`$ or $`b^{}c+W^{}`$ becomes potentially important. These issues are thoroughly considered in, for example, W-S. Hou and R.G. Stuart, Phys. Rev. D 43, 3669 (1991). What is new is the near consensus from SUSY and precision electroweak fits that the lightest Higgs mass is between $`m_Z`$ and $`m_t`$. I believe that the simulation of competing effects of the various decay channels of $`b^{}`$ and the onset of the $`t^{}`$ threshold and decays, with inclusion of the essential Higgs FCNC channel, is the project that deserves a strong effort in Run II. The reason for stressing this issue is that it ties in directly to the crucial Higgs search mission at Run II. A search for the FCNC decay modes of $`b^{}`$ or $`t^{}`$ must include the Higgs effects.
CDF has analyzed its data for the presence of long-lived charged 1/3 $`(b^{})`$ and 2/3 $`(t^{})`$ quarks and has reported mass limits of 195 GeV and 220 GeV at 95% C.L. This pretty analysis closes the door on the possibility that top search signatures missed stable $`(\gamma \tau >10^8s)`$ $`b^{}`$ or $`t^{}`$ ( Dave Stuart, contribution to this workshop).
3) Searches for N and L leptons by the LEPII detector groups limit masses of stable charged leptons to be above 90 GeV at 95% C.L.(eg. OPAL Collaboration, Phys.Lett.B, 433, 195 (1998). Analysis of LEPII data also shows that unstable lepton doublets can escape detection only if N and L are degenerate at the level of a GeV or so, with the neutral one lighter and stable enough to exit the detector. (see S. Thomas and J. Wells, hep-ph/9804359 for a technique to detect such states. The motivation is that vector-like doublets with radiative splitting on the order of several hundreds of MeV are typical in a variety of motivated extensions to SM, but the analysis applies equally well to “unmotivated” doublets).
### XI.4 Considerations on Unification and Perturbation Theory, SUSY Breaking and Other Theory Prejudices
As mentioned above, requiring that perturbation theory remain valid as Yukawa and gauge coupling constants are run to the grand unification scale tightly constrains new family masses. The Yukawa couplings blow up (Landau poles) sooner or later as they are run to higher and higher mass scales. The bigger they start (big masses at the weak scale), the faster they take off. Because of the many superpartners, this effect is exaggerated in the MSSM compared to the SM. Even the latter suffers rapid blow up when fourth family quark masses exceed 150 GeV or so (a recent study is reported in Yu. F. Pirogov and O.V. Zenin, hep-ph/9808396). In the context of a fourth family extension of the MSSM, the constraint is severe indeed. Setting the top mass in the range 170- 175 GeV and requiring that all couplings run to the GUT scale at a few times $`10^{16}`$ GeV, one finds that $`m_b^{}<100GeV`$ and $`m_t^{}<120GeV`$. Unless the lightest MSSM Higgs boson is at the Z mass or less, the D0 and CDF searches eliminate a $`b^{}`$ with a mass this low. In the framework of gravity mediated SUSY breaking and the SUGRA boundary conditions, giving up perturbative unification means giving up the predictions of the SUSY mass spectrum at low energy, since the mass parameters are all specified at the GUT scale. This is a punishing penalty to pay! In gauge mediated SUSY breaking, however, the mass parameters are set at the messenger scale - roughly 100 TeV. Within this picture of SUSY breaking, one may require perturbative running of the coupling constants $`onlytothemessengerscale`$, still retain predictive power over the SUSY mass spectrum, and buy a wider range of theoretically acceptable fourth family mass values. The result of a preliminary one-loop RG study is shown in the “Landau Pole Bound” curve on the “mass bounds” figure for the case $`tan\beta =1.3`$, $`M_{susy}=M_Z`$ and $`\alpha _s=0.119`$. Values above the curve are excluded. The constraints are still surprisingly tight! Only values of $`m_q^{}`$ up to about 1.9 $`m_Z`$, or just below $`m_t`$, are allowed with $`m_L`$ in the range 0.5 - 1.1 $`m_Z`$, as shown in the figure. Whether a viable SUSY spectrum and retention of spontaneous breaking of EW symmetry breaking can be achieved in a four family GMSB model upon running mass relations back down from the messenger scale remains to be seen. The “mass bounds” figure allows a small rectangular region at the lowest, degenerate values of the lepton doublet, L, after application of the mixture of theoretical and experimental constraints. The lower bound of 140 GeV from the $`b^{}b+Z`$ search should be set with light ($`m_h100150GeV`$) Higgs effects included, which brings us back to our “bottom line”.
In connection with the above discussion, it should be remarked that a clever re-identification of the $`t^{}`$ as the top signal and a proposed “hidden” top with $`m_tm_W`$ and a top decay to stop plus LSP in the fourth family extension of the MSSM is probably safely eliminated by the failure to find $`t\overline{t}`$ production at LEP II and the necessary light stop to go with it. (M. Carena, H. Haber and C. Wagner, hep-ph/9512446 and Nucl. Phys. B433, 195; for a stop limit of 70-80 GeV, see The ALEPH Collaboration, Phys. Lett. B434, 189(1998)).
This summary includes contributions through conversations, email, published papers and unpublished talks by the following people (but errors and distortions in this summary are strictly my doing):Jack Gunion, Wei-Shu Hou, Robin Stuart, Dave Gerdes, Joao Guimar$`\overline{a}`$es da Costa, Bhaskar Dutta, Dave Stuart, Kara Hoffman, Mel Shocket, Tom Diehl, Taka Yasuda, Kaori Maeshima, Howie Haber, P.Q. Hung, Paul Frampton, Kyungsik Kang and Sin Kyu Kang.
## XII Light Gluino Predictions for Jet Cross Sections <br>in Tevatron Run II
L. Clavelli
The CDF collaboration at Fermilab has published Abe1 a study of the jet inclusive transverse energy cross section in $`p\overline{p}`$ cross sections at $`1.8TeV`$ which suggest the possibility of anomalous behavior in both the low and high transverse energy regions. D0 has not published results in the low transverse energy region but has presented data at high transverse energy which appear to be consistent with either the CDF result or the standard model. The apparent anomaly at high $`E_T`$ seen by CDF could, therefore, be a statistical fluctuation. It has also been suggested Lai that these results are compatible with the standard model if the gluon distribution at high x is appreciably higher than expected on the basis of previous fits. On the other hand, the anomalous behavior observed by CDF in both the low and high $`E_T`$ regions is also consistent with that expected if the gluino of supersymmetry is light (below $`10GeV`$ in mass is sufficient) CT . Although all direct searches for a light gluino have turned up negative, many indirect indications of such a light color octet parton have been noted. A partial list is contained in the references of CT .
The measured inclusive cross section at center of mass energy $`\sqrt{s}`$ to produce a jet of transverse energy $`E_T`$ averaged over a certain rapidity interval is theoretically expected to have the form
$$d\sigma /dE_T=\alpha _s(\mu )^2s^{3/2}F(X_T,\frac{\mathrm{\Lambda }}{\sqrt{s}},\frac{m}{\sqrt{s}})+𝒪(\alpha _s^3)$$
(35)
Here $`\mu `$ is the scale parameter, $`X_T=2E_T/\sqrt{s},\mathrm{\Lambda }`$ is the QCD dimensional transmutation parameter, and m represents any of the masses of the strongly interacting particles in the theory. Taken to all orders the cross section is independent of $`\mu `$ but at finite order the theoretical result depends on $`\mu `$ which must therefore be treated as a parameter of the theory. The CDF best fits correspond to $`\mu =E_T/2`$. At high energy the scaling function F depends only on $`X_T`$ . The CDF data for this cross section compared to the next-to-leading order (NLO) QCD predictions are below unity at low $`E_T`$ and rise dramatically above unity at high $`E_T`$. In the Supersymmetry (SUSY) treatment of CT2 this behavior was attributed to three phenomena.
* With a light gluino the strong coupling constant runs more slowly being higher than the standard model at high $`\mu `$ and lower at low $`\mu `$.
* The production of gluino pairs increases F by a roughly uniform factor of 1.06 for all $`E_T`$
* A squark, if present, will cause a bump in the cross section at about $`m_{\stackrel{~}{Q}}/2`$.
The fit of CT2 used the CDF suggested value of $`\mu `$ and a value of $`\mathrm{\Lambda }`$ corresponding to $`\alpha _s(M_Z)=0.113`$ and a squark mass of about $`106GeV`$. The theoretical ratio of the SUSY prediction relative to the standard model prediction is relatively insensitive to higher order corrections since both will have roughly equal higher order enhancements. In this work the CTEQ3 parton distribution functions (pdf’s) were used. In a later study Bhatti , CDF considered the scaled ratio of the inclusive jet $`E_T`$ cross sections at $`630GeV`$ and $`1.8TeV`$.
$$r(X_T)=\frac{s^{3/2}d\sigma /dE_T(\sqrt{s}=630GeV)}{s^{3/2}d\sigma /dE_T(\sqrt{s}=1800GeV)}$$
(36)
Since at both energies, $`\sqrt{s}`$ is much greater than the QCD scale parameter $`\mathrm{\Lambda }`$ and all the quark masses of the standard model (except the top quark which contributes negligibly at these energies), the standard model prediction modulo residual corrections from higher order and from scaling violation in the pdf’s is just
$$r(X_T)=\frac{\alpha _s^2(\lambda X_T0.630GeV/2)}{\alpha _s^2(\lambda X_T1.8TeV/2)}$$
(37)
We have assumed here that the appropriate choice of $`\mu `$ is $`\lambda E_T`$ with $`\lambda =1/2`$ being the result of the CDF best fit to the $`1.8TeV`$ data. The full standard model prediction with corrections incorporated seriously overestimates the CDF data. In addition there is a possible structure in r that, if real, might suggest the existence of a strongly interacting particle in the $`100GeV`$ region with a production cross section many times larger than that of top. As always, there is the possibility that the anomaly is due to systematic errors although it would be surprising if such errors induced structure in $`E_T`$. In fact, the D0 experiment does not confirm the existence D0 of structure in r suggesting, therefore, an explanation in terms of systematic errors. Although the systematic errors could easily affect the normalization of the r parameter, it would be surprising if they affect the point to point errors. These systematic errors derive primarily from the lower energy (630 GeV) data and hence the existence or non-existence of structure should be definitively resolved by comparing the ratio of the $`2TeV`$ data which will be available beginning in the year 2000 with the $`1.8TeV`$ data. Although the energy step is small, the greatly increased luminosity in run II coupled with the small systematic errors in the $`1.8TeV`$ data should guarantee sufficient sensitivity to settle the question.
The features observed by CDF in the scaling ratio are those expected in the light gluino scenario CT2 . The slower fall-off of $`\alpha _s`$ predicts that the r parameter should be generally lower than the standard model expectations in agreement with the data. In addition a squark in the $`100GeV`$ mass range would provide a bump in each cross section at roughly fixed $`E_T=m_{\stackrel{~}{Q}}/2`$. This would lead to a dip-bump structure separated by a factor of 1.8/0.63 in $`X_T`$ in qualitative agreement with the CDF data. If the bump had occurred at lower $`X_T`$ than the dip there would have been no possibility of a fit in any model where the structure was attributed to a new particle. Reference CT2 provided two fits to the CDF data. The first used the CTEQ3 pdf’s and the scale choice $`\mu =E_T/2`$ with a squark mass of $`130GeV`$. In the CTEQ3 pdf’s there are, of course, no initial state gluinos so the cross section bump derives from the reaction
$$qgq\stackrel{~}{g}\stackrel{~}{g}$$
(38)
with an intermediate squark in the $`q\stackrel{~}{g}`$ channel. The dynamics are such that the initial state gluon splits into two dominantly collinear gluinos one of which interacts with the initial state quark to produce the intermediate squark. Other non-resonant light gluino contributions to the cross section come from the parton level processes
$$q\overline{q}\stackrel{~}{g}\stackrel{~}{g}gg\stackrel{~}{g}\stackrel{~}{g}$$
(39)
If the gluino is light it should have a pronounced presence in the proton dynamically generated from the gluon splitting discussed above. Two groups BB ; RV have analyzed deep inelastic scattering allowing for a light gluino and presented fits to the gluino pdf as well as modifications of the other pdf’s due to the gluino presence. In ref. CT2 we compared the scaling violation using the Rückl-Vogt pdf’s with that using the CTEQ3 set. With intrinsic gluinos there are extra contributions to the jet inclusive cross sections from the processes
$$g\stackrel{~}{g}g\stackrel{~}{g}q\stackrel{~}{g}q\stackrel{~}{g}\stackrel{~}{g}\stackrel{~}{g}\stackrel{~}{g}\stackrel{~}{g}\stackrel{~}{g}\stackrel{~}{g}gg\stackrel{~}{g}\stackrel{~}{g}q\overline{q}$$
(40)
The second process replaces the higher order reaction of Eq. (38) and provides a direct channel pole at the mass of the squark leading to a peak in the transverse energy cross sections. We treat the squark as a resonance in the quark-gluino channel. Each of these reactions of course is subject to higher order corrections but these tend to cancel in the scaling ratio and in the ratio of the SUSY transverse energy cross section to that of the standard model. In this second fit the scale $`\mu `$ was chosen to be the parton-parton CM energy.
The purpose of the current work is to return to the inclusive jet transverse energy cross section and seek a combined fit to this plus the scaling curve allowing for intrinsic gluinos in the proton. Fitting both the scaling curve $`r(X_T)`$ and the $`1.8TeV`$ cross section is equivalent to fitting the transverse energy cross section at both $`1.8TeV`$ and $`0.63TeV`$. Using the parameters of this combined best fit we then present the predictions for the $`E_T`$ cross section at $`2TeV`$ CM energy of run II and the scaling curve for $`1.8TeV/2TeV`$. The primary parameters of the combined fit are the scale $`\mu `$, the QCD $`\mathrm{\Lambda }`$ parameter or equivalently $`\alpha _s(M_Z)`$, and the squark mass $`m_{\stackrel{~}{Q}}`$. We find the optimal values
$$\mu =0.6E_T\alpha _s(M_Z)=0.116m_{\stackrel{~}{Q}}=133GeV$$
(41)
In the fit we estimate NLO corrections by the K factor $`1+10\alpha _s(\mu )/\pi `$ and we simulate resolution smearing by increasing the width of the squark by a factor of 2 from its SUSY QCD prediction $`2\alpha _sm_{\stackrel{~}{Q}}/3`$. In addition it is known that the systematic errors in the $`630GeV`$ data form a fairly broad band Geer . We therefore allow the scaling data to float by a uniform factor near unity. The results are presented in figures 1-4.
Figure 1 shows the fit to the $`1.8TeV`$ jet inclusive $`E_T`$ cross section averaged over the CDF rapidity range $`0.1<|\eta |<0.7`$. In order to compare with the data of ref. Bhatti , the light gluino prediction is plotted relative to the QCD prediction given to us in a private communication by the author of that reference. At high $`E_T`$ the fit goes through the lower range of the CDF errors which suggests it is also consistent with the D0 data. The fit qualitatively reproduces the dip at low $`E_T`$ and shows a peak at low $`E_T`$ due to the $`133GeV`$ squark. Figure 2 shows the scaling function $`r(X_T)`$ as given in the light gluino scenario with a $`133GeV`$ squark and as given by the standard model. The data has been moved up by a uniform factor of 1.2 which is consistent with the effect of systematic errors in the $`630GeV`$ data. The height and width of the dip-bump structure is in qualitative agreement with the expectations of the light gluino plus $`133GeV`$ squark model. One might expect that a full simulation including hadronization and detector acceptance could somewhat shift this mass.
If a squark exists at $`133GeV`$ it should be apparent in the $`e^+e^{}`$ annihilation cross section through the quark-squark-gluino final state CCFHPY . The L3 data L3 shows what is possibly an upward statistical fluctuation in the hadronic cross section in $`e^+e^{}`$ annihilation in the $`130GeV`$ region. Since the gluino decays are expected to leave very little missing energy, the quark-squark-gluino final state might also explain an apparent surplus in the visible energy cross section at high $`E_{vis}`$ L32 . In addition, a SUSY symmetry breaking scale of $`133GeV`$ would, in the light gluino scenario, predict stop quarks in the region just above the top and could explain some anomalies in the top quark events CG and lead to an enhancement in the deep inelastic cross section at high $`Q^2`$ and high hadronic mass CC . If there is indeed a light gluino and a squark in the $`100135GeV`$ region, a dip-bump structure should also be found at LEP II in the scaling ratio of the inclusive dijet cross section in $`e^+e^{}`$ annihilation.
$$r(M^2/s)=\frac{s^2d\sigma /dM^2(\sqrt{s}=E_1)}{s^2d\sigma /dM^2(\sqrt{s}=E_2)}$$
(42)
where both $`E_1`$ and $`E_2`$ are above the squark mass. Since the squark decays in the present model into quark plus gluino, the excess should be in the four-jet sample but should not appear in the pair production of two high mass states.
In figure 3, we show the predictions for the jet inclusive $`E_T`$ cross section in $`p\overline{p}`$ collisions at the energy $`2TeV`$ relative to the standard model expectations. The curve shows a pronounced peak at $`m_{\stackrel{~}{Q}}/2`$ and is generally $`5`$ to $`10\%`$ below unity due to the slower running of $`\alpha _s`$ in the light gluino case and to the scaling violations in the parton distribution functions. In Figure 4 the scaling ratio
$$r(X_T)=\frac{s^{3/2}d\sigma /dE_T(\sqrt{s}=1.8TeV)}{s^{3/2}d\sigma /dE_T(\sqrt{s}=2TeV)}$$
(43)
is plotted for the case of light gluino plus $`133GeV`$ squark and for the case of light gluino but no squark present (non-resonant solid line). The dash-dotted curve gives the prediction of the standard model.
Although we have not attempted to estimate hadronization corrections nor resolution smearing (apart from doubling the squark width), we expect that run II will be sensitive to the predicted peaks if they exist and will therefore either discover or rule out a squark in the $`100GeV`$ mass region in conjunction with a light gluino. With additional information on dijet mass and angular distributions TC ; T ; DHR , the Run II measurements are sensitive to a light gluino with a squark up to $`1TeV`$. Since most of the value of SUSY would be lost with squarks so high in mass, Run II should definitively settle the question as to whether the light gluino indications including those referenced in CT are the first signs of SUSY or merely an amazing string of coincidences attributable to systematic errors.
## XIII Detection of Long-Lived Particles in Run II with DØ
D. Cutts and G. Landsberg
### XIII.1 Physics motivation
Long-lived neutral or charged massive particles appear in many extensions of the MSSM. There are two main scenarios which result in a long lifetime of some of the SUSY partners. In the models that predict a degenerate mass spectrum of SUSY particles, light SUSY partners might be stable or long-lived, since the phase space would not allow for strong decays modes. Some models (see sections
It is particularly interesting if the NLSP is long-lived, since the observation of such an NLSP might not be doable via standard decay signatures expected for short-lived SUSY particles. As an example, many MSSM extensions predict degenerate mass spectrum of the neutralinos, in which case the following radiative decay could be a dominant decay mode of the second-lightest neutralino:
$$\stackrel{~}{\chi _2^0}\stackrel{~}{\chi }_1^0\gamma .$$
(44)
Being electromagnetic, this decay mode is suppressed, and for fine splitting between the masses of the two lightest neutralinos, the decay constant could be small enough to result in a long-lived $`\stackrel{~}{\chi }_2^0`$.
Among other models which allow for a long-lived NLSP, are the GMSB scenarios in which a neutral NLSP (usually a neutralino) radiatively decays into a gravitino LSP:
$$\stackrel{~}{\chi }_1^0\gamma \stackrel{~}{G},$$
(45)
or a charged NLSP (usually the right-handed stau, $`\stackrel{~}{\tau }_R`$) decays into a $`\tau `$ and a gravitino NLSP:
$$\stackrel{~}{\tau }_R\tau \stackrel{~}{G}$$
(46)
In the GMSB scenarios the gravitino mass is given by (see, e.g., Feng ):
$$M_{\stackrel{~}{G}}=\sqrt{\frac{8\pi }{3}}\frac{F}{M_P},$$
where $`F`$ is the vacuum expectation value of the dynamical supersymmetry breaking, and $`M_P`$ is Planck mass. Since most of the GMSB models predict $`F10^[14]`$ GeV<sup>2</sup>, and $`M_P10^{19}`$ GeV, the gravitino is expected to be very light. The fact that the gravitino interacts with matter only weakly, could make the decays (45) and (46) very slow. For example, the $`c\tau `$ for the decay (46) is given by Feng :
$$c\tau 10\text{km}\times \beta \gamma \times \left[\frac{\sqrt{F}}{10^7\text{GeV}}\right]^4\times \left[\frac{100\text{GeV}}{m_{\stackrel{~}{\tau }_R}}\right]^5,$$
and could be very large.
Other models (see, e.g. Section II) consider cascade decays of the SUSY particles that can also result in displaced vertices. Massive charged long-lived SUSY particles are expected in SUSY models with chargino LSP Feng1 where the lightest neutralino and chargino are nearly degenerate in mass. Similar degeneracy automatically occurs in superstring-inspired models of SUSY breaking Randall . Weakly interacting bound states of gluino LSP also act as long-lived massive particles (see, e.g. Section III).
Generally speaking, in different SUSY models, the $`c\tau `$ of the charged or neutral long-lived particles can be from subatomic distances to many kilometers.
### XIII.2 Detection of the Delayed Decays
Detection and identification of long-lived particles depend on the decay path and the charge of this particle. A typical collider detector, e.g. the upgraded DØ apparatus D0Upgrade , has an inner Silicon Microstrip Tracker (SMT), capable of identifying secondary vertices from long-lived particle decays. The silicon detector is surrounded by a less precise Central Fiber Tracker (CFT), which is enclosed in the calorimeter. The calorimeter is surrounded by the muon system. Typical outer radii of the vertex detector, central tracking detector, calorimeter, and the muon system are $`10`$, 100, 200, and 1000 cm, respectively.
In the case of the charged long-lived particles, one can identify the secondary decay vertex in the silicon vertex detector for $`\gamma c\tau `$ between $`0.1`$ cm and $`10`$ cm. Since we expect SUSY particles to be heavy, in what follows we will assume $`\gamma 1`$, i.e. $`\gamma c\tau c\tau `$. Both the CDF and DØ experiments will be capable of identifying such secondary vertices and possibly trigger on them. For 10 cm $`c\tau 100`$ cm, the charged long-lived particle will predominantly decay inside the outer tracking volume. The resulting characteristic signature is a kink in the outer tracker and large $`dE/dx`$ in the silicon vertex tracker and the inner layers of the outer tracker, typical for a slow-moving charged particle. While triggering on the kinks in the outer tracker won’t be possible in Run II, there is a good chance that such kinks can be found offline, if the event is accepted by one of the standard triggers. Since a massive slow moving particle still has a significant momentum, one would likely trigger on such events using a single high-$`p_T`$ track trigger, or a designated $`dE/dx`$ trigger. For 100 cm $`c\tau 200`$ cm, the charged long-lived particle will decay inside the calorimeter, giving a jet-like energy deposition inside it. Identification of these particles will therefore rely on the fact that such a jet has only one track pointing to it (similar to one-prong $`\tau `$-decays); moreover, this track will have high $`dE/dx`$. In the case of the DØ detector, an additional $`dE/dx`$ measurement in the preshower detector can be used to aid the identification of such particles both offline and at the trigger level. A single high-$`p_T`$ track trigger, a designated $`dE/dx`$ trigger, or a $`\tau `$-trigger could be used to trigger on such events. Finally, for $`c\tau 200`$ cm, the long-lived particle will look stable from the point of view of the detector, and its identification will rely on a high $`dE/dx`$ track in the silicon vertex detector, the outer tracker, the calorimeter, and the muon system (if the $`c\tau `$ exceeds it outer radius). An additional time-of-flight (TOF) information from a designated TOF system (CDF) or muon system scintillators (DØ) could be used to trigger on such events and identify them offline. To summarize, both the CDF and DØ detectors will have very good capabilities for searches for charged long-lived particles with lifetimes from a few tens of picoseconds to infinity. Identification of the slow-moving particles using $`dE/dx`$ techniques in the case of the DØ detector is discussed in more detail in Section XIII.4.
The situation is much more complicated for a neutral weakly interacting particle. First of all, if such a particle decays outside of the calorimeter, it can not be detected, and will look like a missing $`E_T`$ in the event, i.e. like the LSP. It’s unlikely that one could identify the presence of a long-lived NLSP in such events, but it still might be possible to identify these events as the non-SM ones. If events with such a signature are found in Run II, one could conceivably install a wall of scintillating detectors far away from the main detector volume and try to look for a photon from the radiative decay in this additional scintillator Gunion . For the purpose of this report, however, we will focus only on the case of $`c\tau <2`$ m, which roughly coincides with the outer radius of the DØ calorimeter. The detection technique for such decays relies heavily on the fine longitudinal and transverse segmentation of the preshower detector and the calorimeter, which is an essential and unique feature of the DØ detector.
The signature for the radiative decay of a long-lived particle (e.g., (45) or (44)), is a production of a photon with a non-zero impact parameter. If there was a way to identify the point of photon origin, one could single out such a delayed radiative decay corresponding to a very distinct and low background topology. As is shown in Section XIII.3, the DØ calorimeter information, combined with the preshower information, can be used to achieve a very precise determination of the photon impact parameter. This technique would allow to identify long-lived neutral particles with 10 cm $`c\tau 100`$ cm, in the case of the upgraded DØ detector. The detectable $`c\tau `$ range can be further extended by a factor of two by looking for the photons from the radiative decay in the DØ hadronic calorimeter. The signature for such photons would look like a “hot cell,” i.e. an isolated energy deposit in one or two hadronic calorimeter cells. The reason that the EM energy deposition in the calorimeter is so isolated is the fact that each hadronic calorimeter cell contains many radiation lengths and completely absorbs an EM shower. The main background for this signature is production of high-$`p_T`$ $`K_S^0`$-mesons which decay into a pair of $`\pi ^0`$’s corresponding to the 4$`\gamma `$ final state. Since the $`K_S^0`$ is significantly boosted, these four photons are highly collimated and will be identified as a single EM shower in the calorimeter. For a typical $`K_S^0`$ momentum of 25 GeV, which corresponds to $`20`$ GeV $`E_T`$ of the resulting $`4\gamma `$ system, the $`\gamma c\tau `$ is about 130 cm, i.e. most of the $`K_S^04\gamma `$ decays will occur in the hadronic calorimeter. This background, however, can be well predicted and also has a well defined $`r`$-dependence, where $`r`$ is the radial distance of the “hot cell” from the detector center. We therefore expect this background to be under control in Run II. Thus, the upgraded DØ detector will have unique capabilities for searching for neutral long-lived particles with 10 cm $`c\tau 2`$ m.
Both the CDF and DØ could also explore the 0.1 cm $`c\tau 10`$ cm range by looking for a conversion of the photon from a radiative decay in the silicon vertex detector. If the photon converts, one can determine its direction and the impact parameter fairly precisely by looking at the tracks from the $`e^+e^{}`$-pair. One, however, would pay a significant price for being able to explore this region of $`c\tau `$, since the probability of photon conversion in the silicon vertex detector is only a few per cent. CDF could in principle use the conversion technique to explore higher values of the $`c\tau `$ by looking for conversions in the outer tracking chamber. However, due to the low conversion probability the sensitivity of the CDF detector in this region is by far lower than that of the DØ detector.
Apart from being crucial for the study of physics with delayed radiative decays, the DØ detector ability for photon pointing is very attractive from other points of view. In the high luminosity collider environment the average number of interactions per crossing exceeds one. (It can be as high as five, for the 396 ns bunch spacing expected at the beginning of Run II.) Therefore, each event will generally have several primary vertices with only one being from the high-$`p_T`$ interaction of physics interest, and the others being due to minimum bias $`p\overline{p}`$ collisions. The presence of multiple vertices creates a problem in choosing the right one for determination of the transverse energies of the objects. Photon pointing can solve this problem. It is especially important at the trigger level when the information about all the objects in the events is not generally available, and therefore high-$`p_T`$ objects which produce tracks in the tracking chambers can not always be used to pinpoint the hard scattering vertex. In some cases, for example the $`\gamma +E\text{/}_T`$ final state, there are no objects with tracks at all, so there is no way to determine what vertex the photon originated from without utilizing the calorimeter-based pointing. Not knowing which vertex is the one from high-$`p_T`$ collision results in the object $`E_T`$ mismeasurement, which is especially problematic for missing transverse energy calculations. Indeed, $`E\text{/}_T`$ calculations rely heavily on the vertex position and picking the wrong one may result in significant missing transverse energy calculated in the event which in fact does not have any physics sources of real $`E\text{/}_T`$. Not does only this affect physics analyses for topologies with $`E\text{/}_T`$, especially the $`\gamma +E\text{/}_T`$ one, but it also results in a worsening of the $`E\text{/}_T`$ resolution, hence a slower turn-on of the $`E\text{/}_T`$-triggers and ergo higher trigger rates. It is, therefore, very important to have a way of telling the high-$`p_T`$ primary vertex, and photon pointing is the only way of doing this in the $`\gamma +E\text{/}_T`$ case.
Subsequent sections contain technical detail on the high-$`dE/dx`$ and delayed photon identification.
### XIII.3 Photon Pointing at DØ
The importance of photon pointing was appreciated in some Run I analyses, particularly the studies of $`Z(\nu \nu )\gamma `$ Zg and $`\gamma \gamma `$ monopoles . We have utilized the fine longitudinal segmentation of the DØ EM calorimeter (which has four longitudinal layers) by calculating the c.o.g. of the shower in all four layers independently and then fitting the four spatial points to a straight line in order to determine the impact parameter and the $`z`$-position of the photon point-of-origin EMVTX (see Fig. 33). The algorithm used for the c.o.g. finding is based on a logarithmic weighting of the energy deposition in the EM calorimeter cells which belong to the EM shower. The spatial resolution of the c.o.g. finding algorithm in four calorimeter layers averaged over central (CC) and forward (EC) rapidity range, as well as the geometrical parameters of the calorimeter layers are given in Table 7. This study resulted in an algorithm, EMVTX EMVTX , that has been used in several DØ analyses involving photons Zg ; monopoles .
In order to study the improvement of photon pointing made possible by the utilization of the fine spatial resolution of the preshower detector, we have written a toy Monte Carlo (MC) simulation package which takes into account detector geometry, position error in calorimeter and preshower, as well as the primary vertex distribution. First, we compare the resolution obtained from the toy MC with the actual distributions obtained from $`We\nu `$ events collected in Run I. For electrons from $`W`$-events it is possible to determine the point of origin by using track information, which is quite precise. The difference between the $`z`$ position of the vertex obtained from tracking and from electron pointing as well as the signed impact parameter for the electron obtained by pointing are shown in Fig. 34, for central electrons. (The positive sign corresponds to the impact parameter to the right of the center of the detector observed from the EM cluster location.) The distributions are fitted with Gaussian functions with the widths of $`\sigma _z=14.0`$ cm and $`\sigma _r=9.5`$ cm, as determined from toy MC. The data agrees well with the MC predictions and also appears Gaussian. An analogous comparison for forward electrons is shown in Fig. 35. The corresponding resolutions are $`\sigma _z=17.0`$ cm and $`\sigma _r=4.5`$ cm. The resolution on impact parameter improves in the forward region because the physical size of the calorimeter cells becomes smaller with an increase in $`|\eta |`$. In the $`z`$-direction, this effect is compensated by the small angle of the cluster pointing, which dilutes the $`z`$-resolution.
The central and forward preshower detectors of the DØ Upgrade provide a precision measurement of the photon cluster position. The resolution for a typical EM shower perpendicular to the preshower strip is $`1`$ mm. Taking into account the crossing angle between the $`u`$\- and $`v`$-planes in the preshower detectors, the following resolutions in $`r\varphi `$ and $`z`$ ($`r`$) can be obtained: $`\sigma _{r\varphi }=1.5`$ mm, $`\sigma _z=2.5`$ mm (CPS) and $`\sigma _{r\varphi }=2.5`$ mm, $`\sigma _r=2.5`$ mm (CPS). As one can see, the preshower cluster position measurement is superior to that obtained from the EM calorimeter; its position is also spatially separated from that in the calorimeter, as seen in Fig. 36.
With the additional position measurement coming from the preshower detectors, the following pointing resolutions can be obtained for central and forward photons in Run II: $`\sigma _z=2.2`$ cm, $`\sigma _r=1.4`$ cm (CC, see Fig. 37) and $`\sigma _z=2.8`$ cm, $`\sigma _r=1.2`$ cm (EC, see Fig. 38). A very significant improvement in photon pointing (by a factor of six) is achieved by utilizing the additional position measurement provided by preshower detectors.
The implementation of the photon-pointing algorithm can be done as early as at the trigger Level 3. An approximate algorithm that uses only the c.o.g. of the EM shower in the preshower and in the third layer of the EM calorimeter could be used to decrease the amount of calculations and, hence, the decision-making time. Monte Carlo simulations show that the impact parameter and $`z`$-resolutions for a simplified algorithm are only 10% worse than those obtained by complete five-point fit, which is quite satisfactory for trigger purposes. Having the precise vertex information from photon-pointing for photon triggers at Level 3, we will also recalculate the $`E\text{/}_T`$based on this vertex, and that would significantly improve the turn-on of the $`E\text{/}_T`$part of less inclusive triggers which would require an EM cluster and $`E\text{/}_T`$.
### XIII.4 Detection of Slow-Moving Massive Charged Particles in DØ
As described in Section XIII.2, long-lived charged particles have a variety of characteristics which enable their identification, particularly in the analysis stage when the event’s complete data set is available. Depending on the lifetime, a combination of time-of-flight, ionization ($`dE/dx`$) in several different detectors, and the muon-like penetration of a high momentum, isolated track all can be employed, with the additional presence of a kink where there is a decay within the detector volume. Even though discovery at the analysis stage seems possible, triggering is really crucial, if heavy stable particles produced at the Tevatron are to be detected. In this section we discuss the possibilities for triggering, and in particular, the use of $`dE/dx`$ in the hardware trigger as a tool for detecting these objects.
The DØ trigger system is hierarchical, with 3 levels, each level passing a small subset of the events it examines to the next level for further analysis. Thus Level 1 has a high input rate and examines a limited amount of information in making its decision, while subsequent levels have progressively lower input rates and spend longer analyzing more data. At Level 3, all the data digitized for the event is available and the selection is made running software algorithms written in high level code and derived from the offline analysis. We expect that the massive stable particles will be sufficiently rare, and their characteristics clear, such that separating candidates at Level 3 will be straight forward. However, we need to understand how to identify these events in the hardware triggers so that they survive to Level 3.
For Level 1 several tools are available to detect a massive stable particle. The Central Fiber Tracker provides track candidates binned in momentum to which it can apply an isolation criteria Yuri . The CFT has 80 trigger segments independently processed by the trigger. For a heavy stable particle one can select a high $`p_T`$ track with no other tracks in the home and adjacent segments, imposing an isolation of $`\pm 6.75^{}`$ in azimuthal angle Yuri . Association of a muon with such a track would provide an efficient Level 1 trigger for massive stable particles, produced centrally and in isolation.
An additional tool KJ at Level 1 is the measurement of the time of flight (TOF) from scintillators associated with the muon detector and used primarily to reduce background in the large area muon detector from cosmic ray and beam associated accidentals. These counters cover much of the area just outside the calorimeter, inside the first (“A”) layer of muon detectors, and completely outside the detector, on top of the muon “C” layer planes. Electronics associated with the scintillation counters provides both a trigger gate and a TOF gate, thus giving time windows relative to the interaction time. Scintillator hits received within the TOF gate will have the time of flight measurement digitized and saved with the data, assuming the event otherwise passes Level 1. Typically the TOF gate will be set sufficiently wide to accept slow moving particles, namely at least 100 nsec. To contribute to the Level 1 trigger, however, the scintillator hit must occur within the trigger gate, which is necessarily much shorter (of order 25 nsec) because of the high rate of accidentals, particularly in the “A” layer counters. It may be possible to run with a considerable wider trigger gate for the “C” layer counters, given the lower expected accidental rate in these counters. In this case, TOF from the muon system would provide a useful tool for Level 1 triggering on slow particles. Given a Level 1 trigger, generated either through TOF or through an isolated stiff CFT track, the TOF data will be a useful tool for the Level 2 and Level 3 triggers.
Beyond Level 1, the new Silicon Microstrip Tracker provides interesting possibilities for triggering on slow moving particles. Since energy loss of a slow moving particle drops as $`1/\beta ^2`$ as a function of its velocity, the excellent $`dE/dx`$ energy resolution of the silicon chip provides a good handle for offline identification. For example, the energy deposited in one silicon layer is shown in Fig. 39, as measured in a test beam Maria . More importantly for the detection of slow moving particles, there are possibilities to exploit these measurements at the trigger level, with recent approval of the DØ Silicon Tracker Trigger (STT) STT as a component of the Level 2 trigger system. The hardware design for the STT may allow for the inclusion of several additional backplane lines to carry $`dE/dx`$ information along with the other data associated with each cluster of hits Uli . We have studied the basic capabilities of such trigger hardware to explore its potential for slow particle identification.
Our simulation of the STT slow particle trigger is based on a simple Monte Carlo generator which returns an ADC value appropriate for the spectrum shown in Fig. 39. We assume that two backplane lines per SMT hit are available, and we encode each ADC value into these two bits of data, or four bins. After trying various values for the three bin edges, we have chosen to define the bins by those ADC values below which lie 67%, 95%, and 99% of the data. These cutoffs correspond to 38, 75, and 105 ADC counts (see Fig. 39). The STT Level 2 trigger processor finds tracks using four layers of silicon; so, in our model, the hardware would provide four samples of this two bit $`dE/dx`$ data for each track. At this point the trigger would use some algorithm to combine the four samplings most advantageously. The major concern for the correct identification of a slow moving particle is the likelihood of false signals from a minimum ionizing particle (MIP), due to an occasional response in the very long tail of the Landau energy loss distribution. Based on a few studies, our preliminary suggestion is simply to sum the three lowest values, rejecting the largest ADC count of the four. The data then would provide a parameter, related to a $`dE/dx`$ of the particle in the silicon tracker, which we call “slowness,” and which has 10 possible values (0..9). From our simulation we derive the distribution in slowness for a $`\beta =1`$ particle (equivalent to the test beam pion), as shown in Fig. 40.
The ADC response to the passage of a slow moving particle will be similar to the ADC distribution of Fig. 39 scaled by the factor $`1/\beta ^2`$. However, because of the very different kinematics (the slow moving particle is massive, typically 150 GeV) the energy loss distribution will have a less pronounced Landau tail. We make a very conservative assumption and use only the Gaussian component in generating ADC values for massive slow moving particles. Several resulting distributions in our $`dE/dx`$ trigger parameter, “slowness,” for particles with $`\beta =0.4`$ and $`\beta =0.7`$, are included in Fig. 40. There is a clear separation in this parameter compared to the $`\beta =1`$ distribution.
To estimate the effectiveness of the STT $`dE/dx`$ trigger we consider a selection which tags as a slow particle those whose “slowness” is greater than or equal to some value. We vary this selection to study the efficiency for slow particles (at highest possible $`\beta `$) while maintaining a strong rejection against $`\beta =1`$ particles. Figure 41 shows the acceptance as a function of a particle’s velocity, for events with slowness $`>4`$. The appearance of a few $`\beta =1`$ particles but not other high $`\beta `$ events reflects the conservative use, for massive particles, of a Gaussian ADC response, rather than the Landau distribution, which is used to generate the ADC values from a MIP.
Despite the intent to trigger only on centrally produced objects, tracks will tend to be inclined to the silicon, increasing and broadening the $`dE/dx`$ response. In fact, the STT hardware will search for track candidates from adjacent barrels; so that tracks may be inclined to the normal by as much as $`45^{}`$ although the STT will have no information about this angle. We have modeled the effect of inclined tracks by scaling the appropriate ADC response by $`1/\mathrm{cos}\theta `$, where $`\mathrm{cos}\theta `$ is generated uniformly between .707 and 1.0, or by $`\eta `$, where $`\eta `$ is generated uniformly between 0 and 0.88. The differences between smearing based on $`\mathrm{cos}\theta `$ and that based on $`\eta `$ are small. We present here results using $`\mathrm{cos}\theta `$ smearing for the centrally produced massive stable particles and $`\eta `$ smearing for the $`\beta =1`$ background. As seen in Fig. 41, the effect of the variation in ADC response due to track inclination is only a small reduction in the rejection for $`\beta =1`$ particles; moreover, this effect helpfully raises the cutoff in $`\beta `$ for massive particles. Overall, this smearing does not appear to affect the STT “slowness” trigger significantly. Including track inclination, the study suggests that the STT Level 2 $`dE/dx`$ trigger would provide a fully efficient tag for slow particles up to $`\beta =0.7`$ with an acceptance of $`\beta =1`$ particles less than $`2\times 10^3`$.
Because of its excellent rejection for $`\beta =1`$ particles, the Level 2 $`dE/dx`$ trigger will be a good means to select slow moving particles, independent of other criteria such as TOF or track isolation. However, if the particle is sufficiently long-lived to traverse the entire detector, it may be possible to relax the $`dE/dx`$ requirement. We have studied the effect of modifying the hardware ADC sampling, to explore widening the acceptance of heavy particles in $`\beta `$ at the expense of $`\beta =1`$ rejection. There is good physics motivation in doing so, as some GMSB models predict the production of heavy stable particles with $`\beta `$ in the range 0.8-0.9 Jianming . Using ADC bins with edges corresponding to 50%, 68%, and 90% of the distribution, we find good acceptance for massive particles at high $`\beta `$, as shown in Fig. 42, with rejection factors between 20 and 100, for $`\beta =1`$ particles.
The above study illustrates the potential of a STT $`dE/dx`$ trigger. It does assume that the ADC distribution for a MIP is as seen in the test beam, and that differences in the silicon chip response over the detector won’t significantly affect this distribution. The study suggests that track inclination may not be a serious problem. Further, the “slowness” flag derived with the STT from $`dE/dx`$ can in Level 2 be combined with other information (as of a straight, non-oblique and isolated muon) to provide a global Level 2 trigger. In summary, it seems promising that the STT could provide very useful $`dE/dx`$ information early in a slow particle’s lifetime, which can be combined with other data to create an efficient trigger for these interesting objects.
## XIV Searches for Beyond the MSSM Phenomena at CDF
M. Chertok
Two remarks motivate the discussion to follow. First, if the LSP is charged and has an appreciable lifetime (or is stable), it can be detected in a search for charged massive particles (CHAMPs). Second, while the MSSM makes no predictions regarding the possibility of extra quark families, there is no theoretical reason against them mckay .
CDF has performed searches for these phenomena using data taken with the Run I detector det . Upgrades of both the detector and Tevatron are currently underway tdr . These will provide substantial enhancements for these and other searches in Run II, scheduled to begin in 2000. During this run, CDF II will collect roughly 2 $`\mathrm{fb}^1`$ of data at $`\sqrt{s}=2\mathrm{TeV}`$, corresponding to twenty times the present statistics. The 10% increase in energy corresponds to a 40% increase in the $`t\overline{t}`$ yield, and similarly will aid new phenomena searches.
Although the LEP data at the $`Z^0`$ pole exclude extra fermion generations with light neutrinos LEP , models including fourth family quarks have received recent theoretical attention mckay ; four\_theory . CDF performs three complementary searches for such quarks, described below.
### XIV.1 Search for long-lived parent of the $`Z^0`$
CDF performs a general search for long-lived particles decaying to $`Z^0`$ bosons lxy . One example is a $`b^{}`$ quark decaying via a FCNC to a $`b`$ quark and a $`Z^0`$. This decay can dominate (depending on the $`b^{}`$ mass) but may lead to a long lifetime. Another example is low-energy symmetry breaking models with SUSY ambrosanio that predict the decay $`\stackrel{~}{\chi }_1^0Z^0\stackrel{~}{G}`$ with a long lifetime due to the small coupling constant of the gravitino.
The technique for this search is to select $`e^+e^{}`$ pairs from $`Z^0`$ decay and search the transverse decay length distribution for evidence of $`Z^0`$ bosons originating at a displaced vertex. Electrons ($`e^+`$ or $`e^{}`$) are required to satisfy $`E_T>20\mathrm{GeV}`$, $`p_T>15\mathrm{GeV}/c`$, and $`|\eta |<1`$, with the pair reconstructing to the $`Z^0`$ mass: $`|M_{ee}M_{Z^0}|<15\mathrm{GeV}/c^2`$. Quality cuts are imposed to reduce the effects of misreconstructed tracks. The $`e^+`$ and $`e^{}`$ are required to originate from a common vertex, and nearly collinear tracks are removed by requiring $`|\mathrm{\Delta }\varphi \pi |>0.02`$. A high-purity sample of $`J/\psi \mu ^+\mu ^{}`$ is used to check that these cuts do not introduce lifetime biases. In 90 $`\mathrm{pb}^1`$, 703 events pass the full selection.
The reconstructed transverse decay length $`L_{xy}`$ distribution from these events is clustered around the origin, as expected, and can be modeled with a central gaussian and tails due to tracking errors. The negative $`L_{xy}`$ region is used to estimate the background from tracking errors and select an $`L_{xy}`$ cut which is then used to search for a signal in the positive $`L_{xy}`$ region. A cut of $`L_{xy}>1mm`$ is chosen, which corresponds to an expectation of $``$ one event based on the central gaussian. There is one candidate event with $`L_{xy}>+1mm`$. As the data are consistent with background levels, limits are derived as shown in Figure 43. The left figure shows the model independent cross section times branching ratio limit as a function of the lifetime, as well as a limit (insert) for the case of $`b^{}bZ^0`$. In this case two jets with $`E_T>10\mathrm{GeV}`$ and $`|\eta |<2`$ are additionally required and the $`L_{xy}`$ cut is made at $`0.1mm`$. Also shown (right) is the resulting limit on the $`b^{}`$ mass as a function of the lifetime.
Plans for the continuation of this analysis in Run II include the addition of the muon channel and an attempt to perform the search using only the outer tracking chamber without the SVX. The former would increase the cross section sensitivity by approximately a factor of two while the latter would flatten out the limit to $`30`$ cm. Folding in the other factors for Run II improvements and integrated luminosity, the present cross section limit minimum of 0.5 pb for the model independent case should be reduced to below 5 fb. For the $`b^{}`$ model, masses up to the top quark mass should be excluded in Run II.
### XIV.2 Search for $`b^{}\overline{b}^{}Z^0Z^0b\overline{b}(\mathrm{}^+\mathrm{}^{})(q\overline{q})(b\overline{b})`$
If a fourth generation $`b^{}`$ quark is lighter than both $`t`$ and $`t^{}`$, then the allowed CC decay $`b^{}cW^{}`$ is doubly Cabbibo suppressed and the FCNC decay $`b^{}bZ^0`$ dominates if $`m_b^{}>m_{Z^0}+m_b`$. At the Tevatron, $`b^{}`$ would be pair produced with a cross section like that for top with $`b^{}\overline{b}^{}Z^0Z^0b\overline{b}(\mathrm{}^+\mathrm{}^{})(q\overline{q})(b\overline{b})`$ as one decay pattern. CDF performs a search for such a $`b^{}`$ by requiring two electrons consistent with $`Z^0e^+e^{}`$, with $`E_T(e_1)>20\mathrm{GeV}`$, $`E_T(e_2)>10\mathrm{GeV}`$, and $`75\mathrm{GeV}/c^2M(e^+e^{})105\mathrm{GeV}/c^2`$. In 87 $`\mathrm{pb}^1`$, 6548 events satisfy this electron selection. Three or more jets with $`|\eta |<2`$ are then required. For the search region $`M(b^{})<130\mathrm{GeV}/c^2`$, two jets must satisfy $`E_T>15\mathrm{GeV}`$ while the third can have $`E_T>7\mathrm{GeV}`$. For $`M(b^{})>130\mathrm{GeV}/c^2`$, all jets must satisfy $`E_T>15\mathrm{GeV}`$. Also, the sum of the jet $`E_T`$ for jets with $`E_T>15\mathrm{GeV}`$ must scale with the $`b^{}`$ mass in the following way: $`E_T(E_T(j)>15\mathrm{GeV})>M(b^{})60\mathrm{GeV}`$, motivated by a study of the backgrounds. Finally, one $`b`$ tag (displaced vertex) using the SVX detector information is required for the event to pass. This final requirement removes the 31 events remaining after the sum $`E_T`$ requirement.
The $`b^{}`$ signal is generated using HERWIG, with masses in the range $`100\mathrm{GeV}/c^2M(b^{})170\mathrm{GeV}/c^2`$ and passed through the CDF detector simulation program. The efficiency for $`Z^0e^+e^{}`$ is about 60% and the efficiency for the 3 jet cut ranges from 25% to 66% with $`M(b^{})`$. The main background process for $`b^{}bZ^0`$ comes from $`Z^0+nj`$ events, where $`n3`$, and are modeled using VECBOS and HERWIG. These events are passed through the CDF detector simulation program and filtered with the same event selection as for the signal.
Presently, this analysis is being finalized and the muon channel is being added. Including this additional acceptance along with the improvements expected in Run II described above, the ultimate sensitivity of the search will reach $`M(b^{})M(t)`$. It should be noted that the decay $`b^{}bH`$ must be considered in this regime. However, if $`Hb\overline{b}`$ is appreciable, much of the signal can be recovered as long as $`b^{}bZ^0`$ remains large enough to provide triggerable leptons from $`Z^0\mathrm{}^+\mathrm{}^{}`$.
### XIV.3 Search for long-lived charged massive particles
Several models outside the MSSM predict new charged massive particles (CHAMPs) with appreciable lifetimes. These include fourth generation leptons and quarks, weak $`R_p`$ violation, and gauge mediated SUSY breaking models. A strongly interacting CHAMP would have a large cross section (like that for top), which allows large masses to be probed with little background. CDF performs such a search using as a reference model a fourth generation quark with fragmentation to an integer charged meson within a jet.
The technique for this search is to use the large ionization energy loss, $`\mathrm{d}E/\mathrm{d}x`$, in the tracking chambers to tag massive (and therefore slow-moving) particles. Starting with 90 $`\mathrm{pb}^1`$ of data from Run I, various tracking cuts for the SVX and central tracker (CTC) are applied to select high quality tracks. Candidate tracks are required to have $`|p|>35\mathrm{GeV}/c`$ and $`|\eta |<1`$ and to pass ionization cuts in both the SVX and CTC which correspond to requiring $`\beta \gamma <0.85`$. Finally, a mass is calculated from the momentum and $`\mathrm{d}E/\mathrm{d}x`$ and a sliding cut, $`M_{\mathrm{d}E/\mathrm{d}x}>0.6M_{CHAMP}`$, is applied for each assumed $`M_{CHAMP}`$.
The background to this search is due to particles which fake a large ionization signal, mostly from particles whose tracks overlap. The probability of a track faking a signal in both the SVX and CTC is quite low: of the 20K events passing the selection cuts, $`12.1\pm 1.8`$ tracks from fakes are expected with $`\beta \gamma <0.85`$. Since the background is higher at low mass, a tighter cut of $`\beta \gamma <0.7`$ is used for the search region $`M_{CHAMP}<100\mathrm{GeV}/c^2`$. With this cut, $`2.5\pm 0.8`$ tracks from fakes are expected. These background levels agree well with the data as shown in Figure 44. The overall efficiency for this search ranges from 0.75% to 3% as a function of the mass of the CHAMP and is about twice as high for assumed charge $`+2/3`$ as for charge $`1/3`$ due to fragmentation effects. Systematic uncertainties for this analysis are dominated by the uncertainty on interactions of a massive quark in the calorimeter which contribute between 13 and 20% depending upon the assumed quark charge. Figure 44 also shows the resulting cross section limits for this analysis. Comparing these curves to the cross section prediction from PYTHIA, mass limits are obtained at $`190\mathrm{GeV}/c^2`$ for $`b^{}`$ and $`220\mathrm{GeV}/c^2`$ for $`t^{}`$.
In Run II, the CHAMP search will be enhanced by the larger integrated luminosity, higher cross section ($`40\%`$), and improved detector acceptance ($`80\%`$) due to the new tracking chambers. Moreover, there is a proposal to include a time of flight system to the CDF detector. This will greatly help searches for weakly interacting CHAMPs, and will improve the acceptance for analysis described here by $`50\%`$. Combining these factors, the Run II cross section limit for strongly interacting CHAMPs should reduce to 20 fb.
## XV New Gauge Bosons at the Tevatron Run II
T. Rizzo
In this section, we provide an overview of the capabilities of the Tevatron to both discover and explore the couplings of new gauge bosons during Run II. The problems associated with identifying a $`Z^{}`$ once discovered at the Tevatron and issues related to gauge kinetic mixing are also discussed.
### XV.1 Conventional $`Z^{}`$ and $`W^{}`$ Search Reaches
From the string point of view, new gauge bosons are perhaps one of the most natural extensions of the MSSMlang . The conventional approach in searching for new gauge bosons at hadron colliders is via the Drell-Yan channel where a resonant, on-shell particle is produced which subsequently decays to lepton pairs and is easily observable over any continuum backgroundover . The relevant parameter is thus the product of the production cross section times the leptonic branching fraction, $`\sigma B`$. Given a particular extended gauge model with a fixed set of couplings and a set of PDF’s there is very little theoretical uncertainty associated with the calculation of this quantity. This implies that search or exclusion reaches are relatively straightforward to establish under the assumption that the new particle can only decay to SM fermion pairs. (It has been shown that if new decay modes are responsible for decreasing the leptonic branching fraction of a $`Z^{}/W^{}`$ by a factor of 2 it results in a search reach degradation of only $`5060`$ GeVover at Run II.)
Using the canonical $`Z^{}`$ from $`E_6`$ as an example and remembering that in this case the fermionic couplings depend on a parameter $`\theta `$, we show in Fig. 45 both the current exclusion reach from Run I as well as the search reach anticipated from Run II. As the $`Z^{}`$ couplings vary the search reach also varies over a respectable range of $`150`$ GeV. We also see from this figure that this spread of values is quite typical given the fairly wide set of $`Z^{}`$ models. Fig. 45 also shows that for a rather wide range of models, including the $`E_6`$ and several Left-Right cases, the search reach scales almost linearly with the log of the integrated luminosity over the range of values relevant for future Tevatron running. In fact a fit to the curves in Fig. 45 reveals that the mass reach scales with luminosity in a quite model-independent manner as $`MM_0+91.3\mathrm{log}L2.68\mathrm{log}^2L`$ GeV with $``$ in $`fb^1`$ with $`M_0`$ being the only model-dependence.
The situation with $`W^{}`$ search reaches is somewhat different as the canonical example is the Left-Right Modelover . Here not only can the overall $`W^{}`$ coupling strength, $`g_R`$, differ from that of the $`W`$ in the SM, $`g_L`$, but the $`W^{}`$ production cross section may be modulated by a distinct, right-handed CKM matrix, $`V_R`$, and its leptonic decay may involve a massive neutrino, $`N`$, which can decay in the detector (and thus not appear as missing $`p_t`$, as has been studied by D0d0 ). The variation of the search reach with $`V_R`$, whose structure is a priori unknown, assuming $`g_L=g_R`$ and massless neutrinos in the final state is shown in Fig. 46. While the reach in the case $`V_R=V_L`$ is $`1040`$ GeV, a serious degradation is experiences as one scans over other possible forms of $`V_R`$ and is easily reduced by more than $`2030\%`$ and perhaps as much as $`50\%`$. Fig. 46 also shows the almost linear scaling of the $`W^{}`$ search reach with the log of the integrated luminosity in the most naive case; a fit gives to the curve yields $`MM_0^{}+76.3\mathrm{log}L1.86\mathrm{log}^2L`$ GeV for the reach.
A last point to remember regarding the $`W^{}`$ reach is that if $`N`$ is heavier than the $`W^{}`$ then the Drell-Yan search becomes useless and other modes such as $`W^{}WZ`$ or $`jj`$ must be employedother . By scaling existing searches for dijet mass bumps to the energy and luminosity of Run II we have estimated that a $`W^{}`$ can be discovered in this mode up to masses of 880(970) GeV for a luminosity of 2(30) $`fb^1`$.
### XV.2 $`Z^{}`$ Coupling Determinations and Gauge Kinetic Mixing
If new gauge bosons do indeed exist not far beyond the present limits then they will be easily be discovered during Run II. It then becomes mandatory to address the next question–what $`Z^{}`$ is it? To do this one needs to measure as many of the various $`Z^{}`$ properties as possible. While such discussions have taken place for the LHCover , little has been done to address these issues at the Tevatron. We would suppose that if a $`700800`$ GeV $`Z^{}`$ were discovered accumulating additional luminosity would be easily justified and that at least a few hundreds of events would eventually be collected in both the $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ channels. Even with this number of events only a few ‘on-peak’ observables will have statistical errors less than $`10\%`$ and are thus useful for coupling determinations: ($`i`$) the familiar forward-backward asymmetry of the final leptons, $`A_{fb}`$, ($`ii`$) the longitudinal polarization of one of the $`\tau `$’s in $`Z\tau ^+\tau ^{}`$, $`P_\tau `$, and ($`iii`$) the relative rate for $`Z^{}b\overline{b}`$ compared to $`\mathrm{}^+\mathrm{}^{}`$, $`R_{bl}`$x2bb . To see that indeed these three variables, if reasonably well measured, will be able to separate many of the more popular $`Z^{}`$ models, we compare their correlated values in Fig. 47. Note that the couplings in some cases, such as the Left-Right Model, depend upon the values of a single continuous parameter and so their predictions lie along specific curves and not at unique points. It is clear that a determination of only one of these variables will not be sufficient and generally all three are necessary for good model separation.
For extended gauge models based on GUTS with only an additional $`U(1)^{}`$ factor, such as those arising in the $`E_6`$ case, the results shown in Fig. 47 can be too optimistic due to the presence of gauge kinetic mixing(GKM)gkm . GKM is an induced mixing between the $`Z^{}`$ field strength and that associated with the SM hypercharge that arises due to vacuum polarization-like graphs. At the GUT scale, such graphs cancel since complete GUT matter representations exist. Once the GUT symmetry is broken the matter content of the theory below that GUT scale no longer lies in complete representations. In, e.g., the MSSM only the Higgs doublet superfields remain light whereas their associated color triplet partners remain at the GUT scale. The existence of incomplete representations then leads to GKM via the RGE’s and results in a modification of the naive expectations for the $`Z^{}`$ couplings at the TeV scale. In the case of conventional $`E_6`$, the $`Z^{}`$ couples to a charge $`Q_\theta ^{}`$, which is dictated by group theory and the value of the $`\theta `$ mixing parameter. GKM modifies this coupling as $`Q_\theta ^{}\lambda (Q_\theta ^{}+\delta \sqrt{\frac{3}{5}}\frac{Y}{2})`$ where $`Q_{em}=T_3+\frac{Y}{2}`$. Given a specific matter content of the low energy theory below the GUT scale the parameters $`\lambda `$ and $`\delta `$ become calculable via the RGE’s. As far as $`Z^{}`$ physics at the Tevatron is concerned it is quite fortunate that they cannot take on arbitrary values; for the case of the $`\eta `$-type model($`\theta 37.76^o`$) made popular by string theory, $`\delta =\frac{1}{3}`$ leads to leptophobia. In this case the $`Z^{}`$ does not couple to leptons and the standard Drell-Yan search technique fails. This is evident by the ‘hole’ in the search reach shown in Fig. 48. Given a reasonable set of assumptions the matter content of the extended $`E_6`$ model below the GUT scale is not arbitrary, there being only(!) 68 possible sets of superfields to considergkm . For each set the allowed ranges of $`\lambda `$ and $`\delta `$ are calculable, as shown in Figs. 4 and 5 of Ref.gkm , from which $`Z^{}`$ search reaches can be calculated. From this analysis it can be shown that at one loop $`0.286\delta 0.250`$ so that exact leptophobia does not happen. From a complete scan of the parameter space the worse case scenario is found to occur for an $`\eta `$-type model with $`\lambda =0.862`$ and $`\delta =0.286`$ which yields a search reach of 482(736) GeV assuming a luminosity of 2(30) $`fb^1`$.
GKM also affects the $`E_6`$ predictions for $`P_\tau `$, $`R_{bl}`$ and $`A_{fb}`$ since these now depend on the additional parameters $`\delta ,\lambda `$, making model separation more difficult. Fig. 49 shows the result of a Monte Carlo scan of the allowed space of $`\theta \delta \lambda `$ values; we see that the locus of $`E_6`$ points no longer lie along a pair of lines but are now in fairly broad (statically) shaded regions. These plots explicitly show the extreme importance of having multiple observables available for model separation. While the models labelled by A, B, L, U, and $`\mathrm{W}_3`$ remain easily distinguishable from one another and $`E_6`$ if all three observables are well measured, there is now a reasonably large overlap between the Left-Right case and $`E_6`$. Similarly, a $`Z^{}`$ with SM couplings(S) is now no longer separable from $`E_6`$ even when all three observables are employed. Having more observables available would clearly prove useful.
## XVI Conclusions
We hope that the reader has been impressed with the variety of different possible phenomenological manifestations for supersymmetry. We have seen that some versions of supersymmetry could present significant challenges, while other versions will be either discovered or eliminated very soon during Run II. It seems very likely that at least a few of the supersymmetric particles will be observed during Run II.
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# Fine structure of excitons in a quantum well in the presence of a nonhomogeneous magnetic field
## I Introduction
The trapping and guiding of atoms has been the subject of a number of experimental and theoretical works during the last few years. Among the various methods that have been applied to the trapping of atoms, the magnetic one has been clearly recognized as the most powerful. The possibility of exciton trapping in semiconductors has been actively pursued for many years. Recently, nonhomogeneous stress was used in order to create an energy minimum for the trapping of excitons in GaAs quantum wells. Several works have argued that the utilization of homogeneous magnetic fields during the trapping of excitons in semiconductors quantum wells will enhance Bose-Einstein effects in a gas of excitons. In particular, Ref. reported a strong increase of the exciton confinement properties above a critical threshold of uniform magnetic field. An alternative approach to localize excitons is by using nonhomogeneous magnetic fields. In a preceding paper, we discussed the exciton lateral confinement in a GaAs two-dimensional electron gas (2DEG) in the presence of nonhomogeneous magnetic fields. We showed that excitons can be trapped in magnetic field inhomogeneities, with a confinement degree that depends strongly on their angular momentum. Experimentally, nonhomogeneous magnetic fields were used to control the confinement of excitons in quantum wells, where spatially resolved photoluminescence (PL) and PL-excitation spectroscopy was utilized to measure the field gradient effect in the exciton lateral motion.
With the continuing improvements in the experimental realization of spatially nonhomogeneous magnetic fields, several magnetic structures which can be used for the trapping of particles have been proposed, such as: magnetic quantum dots produced by a scanning tunneling microscope lithographic technique, magnetic superlattices created through the patterning of ferromagnetic materials integrated with semiconductors, type-II superconducting materials deposited on conventional heterostructures, and nonplanar 2DEG systems grown by molecular-beam epitaxy. In particular, a magnetic dipole type of profile, which experimentally can be created by ferromagnetic materials deposited on top of semiconductor heterostructures, has attracted a great deal of interest. This system is essentially different from the others since the local nonhomogeneous magnetic field has zero average magnetic field strength. Furthermore, the magnetized disks can create much stronger magnetic field inhomogeneities than superconducting one.
In our previous paper, we neglected the spin of the electron and the heavy-hole and considered only their orbital motion. We showed that this could already induce a weak confinement of the excitons for special configurations of the nonhomogeneous magnetic field profile. In the trapping of atoms, their spin/total angular momentum is the key quantity which leads to their confinement in a nonhomogeneous magnetic field. Therefore, we expect that in the case of excitons, the inclusion of the spin of the particles will lead to an increased confinement.
In this work, we report on a detailed analysis of the trapping of excitons in a GaAs/Al<sub>x</sub>Ga<sub>1-x</sub>As quantum well in a nonhomogeneous magnetic field taking into account spin effects. A magnetized disk deposited on top of a quantum well with a homogeneous magnetic field applied parallel to its growth direction $`z`$, i.e., perpendicular to the magnetic disk, generates the nonhomogeneous magnetic field profile used in the present calculations. This system gives rise to a magnetic dipole type of profile in the $`xy`$ plane of the exciton motion. The effects of the well confinement and the importance of the orbital and spin Zeeman splitting, and the diamagnetic contributions to the exciton trapping energy and the wave function are analysed.
The paper is organized as follows. In Sec. II the theoretical procedure used to obtain the exciton Hamiltonian in a quantum well in nonhomogeneous magnetic fields is given. The circular magnetic trap used in this work is also presented in Sec. II. The calculation method to obtain the exciton trapping energy, the exciton effective mass and its effective confinement potential are discussed in Sec. III. Our numerical results for GaAs/Al<sub>x</sub>Ga<sub>1-x</sub>As quantum wells are presented and discussed in Sec. IV. These theoretical results strongly suggest that it is possible to realize experimentally exciton trapping using nonhomogeneous magnetic fields. Our results are summarized in Sec. V.
## II Theoretical Model
### A Exciton Hamiltonian
The exciton Hamiltonian describing the electron and heavy-hole motion in a quantum well in nonhomogeneous magnetic fields can be written as follows:
$$H=H^{2D}(𝐫_e,𝐫_h)+W(r,z_{e,h})+H^{}(z_{e,h})+H^{m_z}(𝐫_e,𝐫_h),$$
(1)
where
$`H^{2D}(𝐫_e,𝐫_h)`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m_{e,}^{}}}\left\{i_e^{2D}+{\displaystyle \frac{e}{\mathrm{}c}}𝐀(𝐫_e)\right\}^2+{\displaystyle \frac{\mathrm{}^2}{2m_{h,}^{}}}\left\{i_h^{2D}{\displaystyle \frac{e}{\mathrm{}c}}𝐀(𝐫_h)\right\}^2{\displaystyle \frac{\gamma e^2}{\epsilon r}},`$ (3)
$`W(r,z_{e,h})`$ $`=`$ $`{\displaystyle \frac{\gamma e^2}{\epsilon r}}{\displaystyle \frac{e^2}{\epsilon \sqrt{r^2+(z_ez_h)^2}}},`$ (4)
$`H^{}(z_{e,h})`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m_{e,}^{}}}{\displaystyle \frac{^2}{z_e^2}}+V_e(z_e){\displaystyle \frac{\mathrm{}^2}{2m_{h,}^{}}}{\displaystyle \frac{^2}{z_h^2}}+V_h(z_h),`$ (5)
$`H^{m_z}(𝐫_e,𝐫_h)`$ $`=`$ $`\mu _B{\displaystyle \underset{i=x,y,z}{}}g_{e,i}S_{e,i}B_i(r_e)2\mu _B{\displaystyle \underset{i=x,y,z}{}}\left(k_iJ_{h,i}+q_iJ_{h,i}^3\right)B_i(𝐫_h){\displaystyle \frac{2}{3}}{\displaystyle \underset{i=x,y,z}{}}c_iS_{e,i}J_{h,i}.`$ (6)
The in-plane Hamiltonian \[$`H^{2D}(𝐫_e,𝐫_h)`$\] describes the electron and hole motion in the $`xy`$ plane in a nonhomogeneous magnetic field, which is represented by the vector potential $`𝐀(𝐫)`$. We will make use of a variational approach to assume that the difference between the 2D and 3D Coulomb interactions \[$`W(r,z_{e,h})`$\] can be made very small by the choice of an optimum value for $`\gamma `$, which is a variational parameter that is calculated from the average of $`W(r,z_{e,h})`$ over the exciton wave function. The perpendicular contribution \[$`H^{}(z_{e,h})`$\] describes the exciton confinement in the quantum well, i.e., in the $`z`$ direction, which is not affected by the magnetic field. The exciton spin interaction with the magnetic field is described by the spin Hamiltonian \[$`H^{m_z}(𝐫)`$\]. The terms in the latter correspond to the electron and heavy hole Zeeman spin interaction, and to the spin-spin coupling energy, respectively.
In Eq. (1), $`m_{e,}^{}`$, $`m_{h,}^{}`$ ($`m_{e,}^{}`$, $`m_{h,}^{}`$) are the perpendicular (in-plane) electron and heavy hole effective masses,$`S_{e,i}`$ and $`J_{h,i}`$ ($`g_{e,i}`$ and $`k_i,q_i`$) are related to the electron and heavy hole spin (Luttinger Zeeman splitting constants), respectively, $`m_z=S_{e,i}+J_{h,i}=\pm 1,\pm 2`$, is the total spin angular momentum,$`\mu _B=e\mathrm{}/2m_{e,}^{}c`$ is the Bohr magneton, and $`c_i`$ is the spin-spin coupling constant related to the zero field spin interaction. We will assume an isotropic dispersion for electrons and holes. Notice that $`𝐫_e`$, $`𝐫_h`$ are 2D coordinates and give the $`xy`$ plane position of the electron and hole, respectively.
To simplify the in-plane \[$`H^{2D}(𝐫_e,𝐫_h)`$\] and spin \[$`H^{m_z}(𝐫_e,𝐫_h)`$\] Hamiltonian, we introduce the exciton relative and center-of mass coordinates, $`𝐫=𝐫_e𝐫_h`$, and $`𝐑=\left(m_e^{}𝐫_e+m_h^{}𝐫_h\right)/M`$, respectively, with $`M=(m_e^{}+m_h^{})`$ the exciton mass. Following Ref. , we apply the adiabatic approach in which we assume that the exciton relative motion is fast as compared to the center-of-mass motion. Following the approach of Freire et al., the $`H^{2D}(𝐫_e,𝐫_h)`$ Hamiltonian can be separated in a center-of-mass and a relative Hamiltonian, i.e., $`H^{2D}(𝐫_e,𝐫_h)=H^{CM}(𝐑)+H_\gamma ^r(𝐫,𝐑,_R)`$, with $`H^{CM}(𝐑)=(\mathrm{}^2/2M)_R^2`$, and:
$`H_\gamma ^r(𝐫,𝐑,_R)={\displaystyle \frac{\mathrm{}^2}{2\mu }}_r^2{\displaystyle \frac{\gamma e^2}{\epsilon r}}+W_1+W_2,`$ (7)
where $`W_1=(e/2\mu c)\xi 𝐁(R)𝐋+(ie\mathrm{}/2Mc)\left[𝐁(R)\times _R_R\times 𝐁(R)\right]𝐫`$ and $`W_2=(e^2/8\mu c^2)B(𝐑)^2r^2`$ are terms related with the first and the second power in the magnetic field strength, respectively. In the above Hamiltonian, $`\mu =m_e^{}m_h^{}/M`$ is the exciton reduced mass, $`\xi =\left(m_h^{}m_e^{}\right)/M`$, and $`𝐋=𝐫\times \left(i\mathrm{}_r\right)`$ is the exciton angular momentum associated to the relative motion.
In the adiabatic approximation, the magnetic field in the spin Hamiltonian \[see $`H^{m_z}(𝐫_e,𝐫_h)`$ in Eq. (6)\] can be expanded to zero order in the relative coordinates, i.e., $`𝐁(𝐫_{e,h})=𝐁(𝐑)`$. With this assumption, the exciton spin Hamiltonian can be written as:
$$H^{m_z}(𝐑)=\mu _B\underset{i=1}{\overset{3}{}}\left[g_{e,i}S_{e,i}\frac{1}{3}g_{h,i}J_{h,i}\right]B_i(R)\frac{2}{3}\underset{i=1}{\overset{3}{}}c_iS_{e,i}J_{h,i},$$
(8)
where we introduced the following definitions for the exciton heavy hole $`g`$ factor $`g_{h,x}=3q_x`$, $`g_{h,y}=3q_y`$, $`g_{h,z}=6k_z+13.5q_z`$.
Following the adiabatic approach, the total exciton wave function can be written as $`\mathrm{\Psi }^{m_z}(𝐑,𝐫,z_{e,h})=\mathrm{\Phi }(𝐫)\psi (𝐑)F(z_{e,h})^{m_z}(𝐑)`$, where $`\mathrm{\Phi }(𝐫)`$ \[$`\psi (𝐑)`$\] is the exciton wave function associated to the relative (center-of-mass) motion, and $`^{m_z}(𝐑)`$, $`F(z_{e,h})`$ is the exciton wave function corresponding to the spin and to the confinement in the quantum well, respectively. Please notice that in the case of the exciton motion in a homogeneous magnetic field, the spin interaction with the magnetic field can be solved separately from the center-of-mass and relative motion coordinates. Therefore, we assume here that there is no coupling between the wave function related to the exciton spin contribution and those of the exciton center-of-mass and relative motion. Finally, we obtain the following Schrödinger equations:
$`\left\{H^{CM}(𝐑)+E^r(\gamma ,𝐑,_R)+E^{m_z}(𝐑)+E^{}E\right\}\psi (𝐑)=0,`$ (10)
$`\left\{H_\gamma ^r(𝐫,𝐑,_R)E^r(\gamma ,𝐑,_R)\right\}\mathrm{\Phi }(𝐫)=0,`$ (11)
$`\left\{H^{m_z}(𝐑)E^{m_z}(𝐑)\right\}^{m_z}(𝐑)=0,`$ (12)
$`\left\{H^{}(z_{e,h})E^{}\right\}F(z_{e,h})=0.`$ (13)
The variational parameter ($`\gamma `$) is chosen in such a way to minimize the expectation value of $`W(r,z_{e,h})`$. The optimum condition is found when
$$\mathrm{}E^{^{}}=\mathrm{\Phi }(𝐫)F(z_{e,h})\left|W(r,z_{e,h})\right|\mathrm{\Phi }(𝐫)F(z_{e,h})=0.$$
(14)
The above equation determines $`\gamma _{\mathrm{min}}`$ which is inserted into the exciton relative motion energy $`E^r(\gamma ,𝐑,_R)`$.
### B The Nonhomogeneous Magnetic Field
Here, we are interested in the possibility of exciton trapping by using a confinement potential which is created by a circular nonhomogeneous magnetic field. Experimentally, these magnetic field profiles can, e.g., be created by the deposition of nanostructured ferromagnetic disks (or superconducting disks) on top of a semiconductor heterostructure with a homogeneous magnetic field applied perpendicular to the $`xy`$ plane. This produces a nonhomogeneous perpendicular magnetic field in the plane. In this work, we consider a magnetized disk on top of a GaAs/Al<sub>x</sub>Ga<sub>1-x</sub>As quantum well, which creates a magnetic dipole type of profile. A sketch of the experimental setup is shown in Fig. 1(a). This nonhomogeneous magnetic field profile is given by the following equation:
$`B_z(R)`$ $`=`$ $`B_a+2B_0^D{\displaystyle \frac{a(a+R)}{R\sqrt{(a+R)^2+d^2}}}\left\{E\left(p^2\right)+\left(1{\displaystyle \frac{p^2}{2}}\right)K\left(p^2\right)\right\}+B_0^D{\displaystyle \frac{a\left(a^2R^2+d^2\right)}{R^2\sqrt{aR}}}p^3`$ (15)
$`\times `$ $`\left\{{\displaystyle \frac{}{p^2}}E(p^2){\displaystyle \frac{1}{2}}K(p^2)+\left(1{\displaystyle \frac{p^2}{2}}\right){\displaystyle \frac{}{p^2}}K(p^2)\right\},`$ (16)
where $`p=2\sqrt{aR}/\sqrt{\left(a+R\right)^2+d^2}`$, $`B_0^D=h/a`$ is the strength of the magnetization of the disk, $`B_a`$ is the uniform applied field, $`h`$ is the disk thickness, $`a`$ is the disk radius, $`d`$ the distance of the magnetic disk to the quantum well, $``$ the magnetization (in units of $`\mu _0/4\pi `$, where $`\mu _0`$ is the permeability of free space ), R the radial coordinate in the $`xy`$ plane, and $`K(x)`$ $`\left[E(x)\right]`$ is the elliptic integral of first (second) type. The magnetic field profile $`B_z(R)`$ is shown in Fig. 1(b), for $`a=2`$ $`\mu `$m, $`d=0.2`$ $`\mu `$m, $`B_0^D=0.05`$ T, and for different values of the external applied magnetic field $`B_a=0`$, $`0.25`$ T, and $`0.35`$ T. Notice that for this magnetic field profile, the average magnetic field strength is zero, which give us the additional possibility to apply a background field $`B_a`$ to shift this magnetic field profile up and down. We will show that this results into two different effective confinement profiles for the exciton.
## III Effective Mass and Confinement Potential
In order to estimate the exciton confinement, we defined the trapping energy $`E_T`$ as the energy difference between an excitonic state in the homogeneous applied field $`B_a`$ and the corresponding state in the nonhomogeneous magnetic field:
$$E_T=E^r(\gamma ,B_a)+E^{m_z}(B_a)+E^{}E.$$
(17)
Thus, the important terms (effective potentials) which determine $`E_T`$ in Eq. (10) are only those appearing in Eqs. (11-13) which are $`B(R)`$-dependent. Therefore, the Hamiltonian corresponding to the exciton confinement in the quantum well does not have to be solved in order to obtain the trapping energy. Indeed the corresponding result will cancel out in the $`E_T`$ calculation. All the well confinement related contribution to the exciton trapping is then determined through the $`\gamma `$ parameter \[see Eq. (11) and Eq. (14)\]. It is worthwhile to point out that if we are interested in calculating the exciton binding energy, all these terms \[Eqs. (11-13)\] have to be included.
Following Ref. , the exciton relative motion \[see Eq. (11)\] can be solved by perturbation techniques, where all the $`B(R)`$-dependent terms are treated as perturbations. The corresponding eigenvalues are:
$`E^r(\gamma ,R,_R)`$ $`={\displaystyle \frac{e\mathrm{}}{2\mu c}}\xi m_rB_z(R)+{\displaystyle \frac{e^2}{8\mu c^2}}\beta _{m_r}^{n_r}\gamma ^2B_z(R)^2+{\displaystyle \frac{e^2\mu }{2M^2c^2}}\alpha _{m_r}^{n_r}\gamma ^4_R\left[B_z(R)^2_R\right],`$ (18)
with the following wave function:
$$\mathrm{\Phi }_0^{n_r,m_r}(𝐫)=A_{n_r,m_r}\left(\frac{2\gamma r}{a_B^{}(n_r1/2)}\right)^{|m_r|}\mathrm{exp}\left(im_r\phi \frac{\gamma r}{a_B^{}(n_r1/2)}\right)L_{n_r+|m_r|1}^{2|m_r|}\left(\frac{2\gamma r}{a_B^{}(n_r1/2)}\right),$$
(19)
where $`n_r`$, $`m_r`$ are quantum numbers of the exciton relative motion, $`a_B^{}=\epsilon \mathrm{}^2/\mu e^2`$ is the effective Bohr radius, $`A_{n_r,m_r}`$ is the normalization constant, $`L_n^{(\alpha )}(x)`$ is the generalized Laguerre polynomial, and $`\alpha _{m_r}^{n_r}`$ and $`\beta _{m_r}^{n_r}`$ (in units of $`a_{B}^{}{}_{}{}^{4}`$ and $`a_{B}^{}{}_{}{}^{2}`$, respectively) are constants related to the relative quantum numbers $`n_r`$, $`m_r`$. In writing the above energy for the relative motion \[Eq. (18)\] we have neglected the field independent term, $`\gamma ^2R_y^{}/(n_r1/2)^2`$ (where $`R_y^{}=\mu e^4/2\epsilon ^2\mathrm{}^2`$ is the effective Rydberg), which does not give any contribution to the exciton trapping energy. In solving the exciton relative motion we assumed that the magnetic field intensity is such that we are in the weak-field regime, i.e., $`\mathrm{}\omega _c^{}<2R_y^{}`$ where $`\omega _c^{}=eB/\mu c`$ is the cyclotron-resonance frequency. For GaAs this implies $`B<5`$ T.
For the sake of simplicity, instead of solving Eq. (14) to obtain $`\gamma `$, we used the results of Stébé and A. Moradi and of Andreani and Pasquarello to determine the zero magnetic field binding energy of excitons in GaAs/Al<sub>0.3</sub>Ga<sub>0.7</sub>As quantum wells. We found that the gamma parameter only slightly influences the trapping energy and this as a coefficient in two smaller terms (diamagnetic and mass correction) in the exciton relative motion energy \[see Eq. (18)\]. For a complete description of the dependence of the $`\gamma `$ parameter with the quantum well width and also with an applied homogeneous magnetic field (which is negligible in the weak field regime) we refer to Ref. .
The eigenenergy of the spin Hamiltonian for a magnetic field parallel to the z direction can be straightforwardly calculated from Eq. (12):
$$E^{m_z}(R)=\pm \frac{1}{2}\sqrt{\mu _B^2\left[\left(1\right)^{m_z+1}g_{e,z}+g_{h,z}\right]^2B_z(R)^2+\left[c_x\left(1\right)^{m_z+1}c_y\right]^2},$$
(20)
In the above energy, we neglected the $`B(R)`$-independent term, i.e., the e-h exchange energy in the z-direction $`\left(1\right)^{m_z+1}c_z/2`$, which does not contribute to the calculus of $`E_T`$. The spin wave function is a linear combination \[$`|^{m_z}(𝐑)=(|+\pm |)/\sqrt{2}`$\] of the exciton spin states $`m_z`$:
$$|^{m_z}(R)=\frac{||m_z|+(Q\pm \sqrt{1+Q^2})||m_z|}{\sqrt{2\left(1+Q^2\pm Q\sqrt{1+Q^2}\right)}},$$
(21)
where $`Q=\mu _B\left[\left(1\right)^{m_z+1}g_{e,z}+g_{h,z}\right]B_z(R)/\left[c_x\left(1\right)^{m_z+1}c_y\right]`$.
We can now use all the results of the exciton relative motion, spin interaction and well confinement \[Eqs. (18), (20), and (14), respectively\] to solve the exciton center-of-mass equation \[see Eq. (10)\]. Now, the eigenenergies $`E^r(\gamma ,R,_R)`$ and $`E^{m_z}(R)`$ act like an effective potential and an effective mass in the center-of-mass equation of motion:
$`\{{\displaystyle \frac{\mathrm{}^2}{2M}}_R\{1{\displaystyle \frac{e^2\mu }{\mathrm{}^2Mc^2}}\alpha _{m_r}^{n_r}\gamma ^4B_z(R)^2\}_R+{\displaystyle \frac{e^2}{8\mu c^2}}\beta _{m_r}^{n_r}\gamma ^2B_z(R)^2+{\displaystyle \frac{e\mathrm{}}{2\mu c}}\xi B_z(R)m_r`$ (22)
$`\pm {\displaystyle \frac{1}{2}}\sqrt{\mu _B^2\left[\left(1\right)^{m_z+1}g_{e,z}+g_{h,z}\right]^2B_z(R)^2+\left[c_x\left(1\right)^{m_z+1}c_y\right]^2}E\}\psi (R)=0.`$ (23)
Our magnetic field is parallel to the z-direction and has $`\phi `$-symmetry, i.e., $`B_z(𝐑)=B_z(R)`$ \[see Eq. (16)\]. Then, we can use the cylindrical symmetry of our problem to write the exciton center-of-mass wave function as $`\psi (𝐑)=e^{im_R\phi }\psi (R)`$, where $`m_R`$ is the quantum number for the angular momentum of the exciton center-of-mass motion. Inserting the above wave function in Eq. (23) and writing the $``$ operator in cylindrical coordinates, Eq. (23) can be written as follows:
$$\left\{\frac{\mathrm{}^2}{2}\frac{1}{R}\frac{d}{dR}\left[\frac{R}{M^{eff}(R)}\frac{d}{dR}\right]+V^{eff}(R)E\right\}\psi (R)=0,$$
(24)
where
$$M^{eff}(R)=\frac{M}{1\frac{e^2\mu }{\mathrm{}^2Mc^2}\alpha _{m_r}^{n_r}\gamma ^4B_z(R)^2},$$
(25)
is the exciton effective mass, and
$`V^{eff}(R)`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2}}{\displaystyle \frac{1}{M^{eff}(R)}}{\displaystyle \frac{m_R^2}{R^2}}+{\displaystyle \frac{e^2}{8\mu c^2}}\beta _{m_r}^{n_r}\gamma ^2B_z(R)^2+{\displaystyle \frac{e\mathrm{}}{2\mu c}}\xi m_rB_z(R)`$ (25)
$`\pm `$ $`{\displaystyle \frac{1}{2}}\sqrt{\mu _B^2\left[\left(1\right)^{m_z+1}g_{e,z}+g_{h,z}\right]^2B_z(R)^2+\left[c_x\left(1\right)^{m_z+1}c_y\right]^2},`$ ()
is the effective confinement potential of the exciton center-of-mass motion. In the above equation, the different terms correspond to the centrifugal, diamagnetic, orbital momentum, and spin contribution, respectively. The centrifugal term is related with the exciton effective mass and the angular momentum of the center-of-mass motion, which we found to give the smallest contribution to the effective potential. The magnetic field squared dependence is present in the effective mass and in the diamagnetic term of the effective potential. Notice that all the direct contributions of the quantum well confinement is included through the $`\gamma `$ parameter, which only gives a small increase in the magnetic field squared intensity. As occurs in the trapping of atoms, the two more important terms are the angular momentum and spin contributions. The first is only important for the excited states. The sign in the spin contribution is related to the spin quantum state $`m_z=\pm 1,\pm 2`$, which characterizes the spin Zeeman splitting.
## IV Numerical Results and Discussion
We have calculated the trapping energy and wave function of excitons in a GaAs/Al<sub>0.3</sub>Ga<sub>0.7</sub>As quantum well in a nonhomogeneous dipole type of magnetic field. The numerical solution of Eq. (24) was obtained by using a discretization technique. We use electron and heavy-hole $`g`$ factors which depend on the quantum well width $`L`$, but we neglect its magnetic field dependence since only small magnetic fields are considered. We take advantage of the system symmetry to assume that $`c_x=c_y`$. The electron and heavy hole mass, and the dieletric constant used in this work are the same as in Ref. ($`m_e^{}/m_0=0.067`$, $`m_h^{}/m_0=0.34`$, and $`\epsilon =12.5`$). The confinement energy level of the exciton in the quantum well does not directly enter into the trapping energy, but it enters indirectly through the $`\gamma `$ parameter which depends on the confinement state of the exciton and on its relative motion state. As an example we took a ferromagnetic disk with radius $`a=2`$ $`\mu `$m, which is placed a distance $`d=0.2`$ $`\mu `$m above the quantum well.
The effective potential and effective mass for the exciton ground state is shown in Fig. 2 as a function of the radial coordinate $`R`$, for a quantum well width of $`L=90`$ Å, for the situation in which an uniform applied field of (a) $`B_a=0.35`$ T and (b) $`B_a=0.25`$ T is presented, where the strength of the disk magnetization is $`B_0^D=0.05`$ T. The Zeeman effect can be seen in the effective potential by the shift in the corresponding spin states $`m_z=\pm 1`$ (dashed and dotted curves), as compared to the spinless state (solid curves). Notice that there is an interesting change of both sign and curvature in the effective potential associated to the spin states $`m_z=\pm 1`$, which can be observed by comparing the dotted curves ($`m_z=+1`$) with the dashed curves ($`m_z=1`$) depicted in Fig. 2. This occurs due to the fact that the diamagnetic contribution to the effective potential \[see Eq. (()III)\] is usually smaller than the Zeeman contribution. Then, the latter dictates the behavior of the effective confinement potential, changing the confinement of the trapped exciton from a centered structure to a ring-like structure.
The quantum well confinement effects on the exciton motion are very important not because of the confinement itself, but due to the large dependence of the exciton $`g`$ factor on the quantum well width \[see inset of Fig. 3(a)\]. The exciton trapping energy dependence on the strength of the magnetization of the disk ($`B_0^D`$), for a homogeneous applied field of $`B_a=0.25`$ T, for quantum wells widths of $`L=50`$ Å and $`L=100`$ Å is shown in Fig. 3(a) and Fig. 3(b), respectively. In both figures we showed the results for the case when the exciton center-of-mass quantum number is $`n_R=1`$ with $`m_R=0`$ (curves) and $`m_R=1`$ (symbols), and when the relative quantum numbers are ($`n_r=1`$, $`m_r=0`$), for the ground state level of the exciton well confinement, and for spin quantum number $`m_z=0,\pm 1`$ (i.e., $`\sigma ^\pm `$ polarized states), $`\pm 2`$ (i.e., the dark excitons). The results for the dark exciton are only shown for reference. They are not optically active and they will not be further considered. Notice that the exciton $`g`$ factor increases the trapping energy as much as by a factor of 20 and that the behavior of the Zeeman splitting is strongly influenced by the nonhomogeneous magnetic field \[compare the dashed and dotted curves in Fig. 3(a) and 3(b)\]. Also notice that for large $`B_0^D`$, the trapping energy starts to decrease because the nonhomogeneous field created by $`B_0^D`$ can now be comparable to the homogeneous applied field $`B_a`$, which decreases the confinement region in the effective potential \[see Fig. 2(b) and Fig. 1(b)\]. The energy of the $`m_R=1`$ states \[the centrifugal term in the effective potential of Eq. (()III)\] are only slightly different from the energy of the correspondent non-excited state.
The dependence of the exciton ground state trapping energy on the quantum well width and on the magnetization strength $`B_0^D`$ is shown in Figs. 4, 5(a), and 5(b), for the spin quantum numbers $`m_z=0`$, $`1`$, and $`+1`$, respectively, and for an applied field $`B_a=0.35`$ T. For thin wells (widths smaller than approximately 60 Å) the exciton trapping energy is strongly sensitive to the nonhomogeneous magnetic field and to the size of the well. As the well width increases, the trapping energy becomes basically independent of the well confinement, and the magnetic field is only a small perturbation \[see Figs. 5(a) and 5(b)\]. There is a minimum in the exciton trapping energy when the quantum well width is near 100 Å which is due to the fact that the $`g`$ factor is approximately equal to zero \[see inset of Fig. 3(a)\]. The spin Zeeman splitting as due to the nonhomogeneous magnetic field does not follow the same behavior of shifting up and down the energy as occurs in the homogeneous field case \[compare Figs. 4, 5(a) and 5(b)\]. The $`\sigma ^{}`$ polarized state \[Fig. 5(a)\] is always shifted up as compared to the spinless situation (Fig. 4), but the $`\sigma ^+`$ polarized state can be shifted down, but also up \[see Fig. 5(b) or dotted curves in Fig. 3(a) and 3(b)\], depending on the relation between the quantum well width (exciton $`g`$ factor) and the magnetization strength $`B_0^D`$. This is due to the change in the confinement region of the effective potential related with the $`m_z=+1`$ spin state, as previously described in our discussion of Fig. 2. It is quite remarkable that the exciton spin interaction with the nonhomogeneous magnetic field can be responsible for increases in the exciton trapping energy as large as a factor of 100, as compared to the spinless situation of Fig. 4. It is important to highlight that the considered magnetic fields and quantum well widths are in the range which are currently available experimentally.
The trapping energy of the exciton excited states as a function of the strength of the disk magnetization $`B_0^D`$ are shown in Fig. 6 for ($`n_R=1`$, $`m_R=0`$) with relative quantum number $`n_r=2`$, for the first excited state of the exciton well confinement, for a quantum well width $`L=50`$ Å, for $`B_a=0.25`$ T, with the following relative angular quantum numbers: (a) $`m_r=1`$ , (b) $`m_r=0`$, and (c) $`m_r=+1`$, and for $`B_a=0.25`$ T, for (d) $`m_r=1`$ , (e) $`m_r=0`$, and (f) $`m_r=+1`$. The interaction of the exciton angular and spin momentum is responsible for several interesting effects, mainly in the case of negative $`B_a`$ where the Zeeman splitting exhibits different behavior for each angular quantum number of the relative motion. The trapping energy for a positive applied field always increases linearly with increasing nonhomogeneous magnetic field intensity $`B_0^D`$, which is not the case for the results with negative $`B_a`$. The linear increase is due to the fact that the angular momentum term is the dominant term in the exciton Hamiltonian, and it has a linear dependence on the magnetic field. The decrease of the trapping energies in Figs. 6(a) and 6(b) can be explained by the competition between $`B_a`$ and $`B_0^D`$, as discussed previously in connection with Fig. 3.
The contour plot of the conditional probability to find the electron somewhere in the $`xy`$ plane for the case of the exciton ground state $`|\psi (𝐑)\mathrm{\Phi }(𝐫)|^2`$, in the presence of an applied homogeneous magnetic field $`B_a=0.25`$ T, is shown in Fig. 7 for the spin quantum number $`m_z=1`$ and in Fig. 8 for $`m_z=+1`$. In both pictures, we considered a quantum well width $`L=90`$ Å and the strength of the disk magnetization was $`B_0^D=0.05`$ T. The hole (indicated by the symbol in the figures with a cross in the middle) is fixed in the position $`𝐫_h=(0,0)`$ in Figs. 7(a) and 8(a), and in $`𝐫_h=(0.5,0)`$ $`\mu `$m in Figs. 7(b) and 8(b). It can clearly be seen that the spin orientation changes the exciton confinement from a centered structure to a ring like structure (compare Figs. 7 and 8). Also notice that the electron tries to follow the hole but this is partially worked against by the attraction to the minimum of the effective potential (see Fig. 2) in the case of a centered wave function (see Fig. 7). In the ring-like structure, the electron moves to the hole nearest position (see Fig. 8).
## V Conclusions
The possibility of carrier trapping in a well-defined confinement region in a semiconductor heterostructure open the possibility to numerous studies for electronic phase transitions, confined exciton gases and excitonic molecules. The trapping of excitons in semiconductor systems of reduced dimensionality have also been proposed as a method of observing Bose-Einstein condensation of excitons. The experimental efforts to observe exciton condensation in semiconductors were concentrated on the analysis of the PL line shape and the transport of excitons. Previous PL experiments (see, e.g., Ref. and references therein) of the exciton trapping by using nonhomogeneous magnetic fields have shown that excitons are driven to regions of minimum field, but they could not detect the direct effect of the field gradient on the exciton motion. They concluded that the experimental conditions had to be improved by taking better samples which have larger exciton mobility, and a larger nonhomogeneous magnetic field strength was necessary.
We have investigated the exciton trapping in a quantum well in the presence of a nonhomogeneous dipole type of magnetic field. The effects of the well confinement and of the exciton interaction with the nonhomogeneous magnetic field taking into account spin states were analyzed and discussed. As compared to our previous results in which the spin of the exciton was not taken into account, we found a substantial increase of the trapping energy of the excitons by using nonhomogeneous magnetic fields. We also showed that the obtained energies, as well as the considered magnetic field strength and well widths are currently experimental obtainable.
Our results show that the trapping energy is strongly dependent on the quantum well width and shape of the nonhomogeneous magnetic field profile, but that a increase in the magnetic field intensity is not always related with a stronger confinement (as occurs in the uniform magnetic field situation), as previously discussed in Sec. IV. Our main result is that the spin of the exciton is responsible for an increases in the trapping energy, which can be as large as a factor of 100. This can be used to increase the experimental conditions by choosing a suitable set of parameters (quantum well width and magnetic field profile) in order to maximize the exciton trapping. Furthermore, the spin interaction with the magnetic dipole type of profile can change the exciton spatial localization of a centered structure to a ring-like structure, increasing the energy level corresponding to the $`m_z=+1`$ spin state, which suggests that the two lines of the $`\sigma ^\pm `$ polarized states will be very close in energy in the PL spectra, which may make it difficult for its identification. For the case of narrow wells, exciton localization due to quantum well width fluctuations will also be present. This effect was not considered in the present paper but should be easily distinguished experimentally with the present exciton trapping by nonhomogeneous magnetic field, by changing the strength of the magnetization of the magnetized disk.
Our results open up a new path for experiments on the confinement of excitons, which should reveal new kinds of exciton trapping and increase the knowledge about the exciton interaction with nonhomogeneous magnetic fields.
## VI Acknowledgments
This research was supported by the Flemish Science Foundation (FWO-VI), the IUAP (Belgium), the ”Onderzoeksraad van de Universiteit Antwerpen”, and by the Inter-university Micro-Electronics Center (IMEC, Leuven). J. A. K. Freire was supported by the Brazilian Ministry of Culture and Education (MEC-CAPES) and F. M. Peeters was supported by the FWO-Vl. V. N. Freire and G. A. Farias would like to acknowledge the partial financial support received from CNPq, the Funding Agency of the Ceará State in Brazil (FUNCAP), and the Brazilian Ministry of Planning through FINEP.
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# A Fourier-Mukai approach to spectral data for instantons
## 1 Introduction
The mathematical study of gauge theory was born some 25 years ago, and one of the first main results is just as old, namely the ADHM construction of instantons over the Euclidean 4-space. Since then, several different types of Euclidean instantons have been studied from various points of view: monopoles, calorons, Higgs pairs, doubly-periodic instantons, to name a few. The common feature in the study of such objects is the so-called Nahm transform. It relates instantons on $`^4`$ which are invariant under a subgroup of translations $`\mathrm{\Lambda }^4`$ to instantons on the dual $`(^4)^{}`$ which are invariant under the dual subgroup of translations $`\mathrm{\Lambda }^{}`$ (see \[15, Section 7\] for a detailed exposition). The Nahm transform has many interesting properties, but perhaps its greatest virtue is that it often converts the difficult problem of solving nonlinear PDE’s into a simpler problem involving only ODE’s or just vector spaces and linear maps between them.
More recently, the so-called Fourier-Mukai transforms have generalised the Nahm transform, bringing it to the realms of algebraic geometry and derived categories. In this paper, we shall study a particular version of such transforms. It is defined for torsion-free sheaves $`E`$ over relatively minimal elliptic surfaces $`X\stackrel{\pi }{}B`$ with a section, transforming them into torsion sheaves on $`J_X\stackrel{\widehat{\pi }}{}B`$, the relative Jacobian of $`X`$, which are supported over the spectral curves studied by Friedman-Morgan-Witten for physical reasons. These constructions are briefly reviewed in Sections 2 and 3, and we observe that there is a correspondence between instantons on a 4-torus, spectral curves with line bundles on it, and instantons on the dual 4-torus.
The bulk of the paper is contained in Section 5, where we compare the $`\mu `$-stability of $`E`$ with the concept of stability introduced by Simpson for its transform. In particular, we show that $`E`$ is $`\mu `$-stable and locally-free if and only if its transform is stable in the sense of Simpson.
As an application of these ideas, we shed new light into the moduli space of irreducible $`SU(r)`$-instantons over elliptic K3 and abelian surfaces, showing in Sections 6 and 7 that it has the structure of a complex Lagrangian fibration. The case of rank 2 instantons on an elliptic surface is treated in Friedman’s book . We extend the results to higher rank and provide some more details about the structure of the fibration. The proofs are greatly simplified using the Fourier-Mukai technology and this enables us to go further in our description of the moduli spaces.
#### Notation.
The elliptic surface $`X`$ will be polarised by $`\mathrm{}`$ which can be any polarisation (for the product of elliptic curves it is convenient to choose the sum of the elliptic curves). If $`F_x`$ denotes the fibre of $`X\stackrel{\pi }{}B`$ passing through $`xX`$, then $`\widehat{\pi }^1(\pi (x))=\widehat{F}_{\pi (x)}`$ is the dual of $`F_x`$. Let $`J_X`$ denote the dual elliptic surface parametrising flat line bundles on the fibres of $`\pi `$. It also fibres over $`B`$ and has a canonical section given by the trivial bundles. We shall denote its fibres by $`\widehat{F}_b`$, for $`bB`$.
Given a sheaf $`E`$ on $`X`$, we write its Chern character as a triple:
$$ch(E)=(\mathrm{rank}(E),c_1(E),ch_2(E)),$$
where $`ch_2(E)=\frac{1}{2}c_1(E)^2c_2(E)`$.
By $`𝒫X\times _BJ_X`$ we mean the relative Poincaré sheaf parametrising flat line bundles on the smooth fibres. Given $`pJ_X`$, let $`\iota _p`$ be the inclusion of $`F_{\widehat{\pi }(p)}`$ inside $`X`$. We define:
$$E(p)=\iota _p^{}(E𝒫_p)$$
where $`𝒫_p`$ denotes the torsion sheaf on $`X`$ corresponding to $`pJ_X`$.
Finally, $`_X(r,k)`$ (or just $``$) will denote the moduli space of $`\mu `$-stable locally-free sheaves over $`X`$, with Chern classes $`(r,0,k)`$, where $`r,k>1`$. Note that this coincides with the moduli space of $`SU(r)`$ instantons of charge $`k`$.
## 2 Spectral Data for instantons
Let $``$ be a vector bundle over the elliptic surface $`X`$, and let $`A`$ be an anti-self-dual $`SU(r)`$ connection on $``$. Then $`A`$ induces a holomorphic structure $`\overline{}_A`$ on $``$; we denote by $`E`$ the associated holomorphic vector bundle.
Given a holomorphic vector bundle $`EX`$ as above, we say that $`E`$ is generically fibrewise semistable if its restriction to a generic elliptic fibre is semistable; $`E`$ is said to be fibrewise semistable if its restriction to every elliptic fibre is semistable. We shall also say that $`E`$ is regular if $`h^0(F_{\widehat{\pi }(p)},E(p))1`$ for all $`pJ_X`$. Observe that regularity is a generic condition.
Now assume that $`E`$ is a regular holomorphic vector bundle over $`X`$; we define the instanton spectral curve with respect to the projection $`X\stackrel{\pi }{}B`$ in the following manner:
$$S=\{pJ_X|h^1(F_{\widehat{\pi }(p)},E(p))0\}$$
(1)
It is not difficult to see that $`S\stackrel{\pi }{}B`$ is a branched $`r`$-fold covering map.
To define the second part of the instanton spectral data, recall that $`\chi (E(p))=0`$ for every $`pJ_X`$. We define a line bundle $``$ on $`S`$ by attaching the vector space $`H^1(F_{\widehat{\pi }(p)},E(p))`$ to the point $`pS`$. Alternatively, consider the diagram:
(2)
where $`\tau _1`$ is the inclusion map and $`\tau _2,\sigma `$ are the obvious projections. We define $``$ as follows:
$$=R^1\sigma _{}(\tau _2^{}E\tau _1^{}𝒫)$$
(3)
Geometrically, it is not difficult to see that $``$ can be regarded as a bundle of cokernels associated to a family of coupled Dirac operators parametrised by $`S`$. For the simple case $`X=T\times ^1`$, see . This follows from the natural identification between $`H^1(F_{\widehat{\pi }(p)},E(p))`$ and the cokernel of the Dolbeault operator $`\overline{}_A|_{F_{\widehat{\pi }(p)}}`$.
As we will see in Section 3 below, the above construction is invertible, and the original holomorphic bundle $`E`$ can be reconstructed from its spectral pair $`(S,)`$. However, the spectral pair contains only the holomorphic information, so some extra information is needed in order to reconstruct the original instanton connection. One possibility is to construct a connection on $`S`$, as it was done in ; see also for a related problem. In this paper we will follow a different path, which will be described in Section 5 below.
## 3 Spectral Data and the Relative <br>Fourier-Mukai Transform
Given a sheaf $`E`$ on the relatively minimal elliptic surface $`XB`$ with a section, we define its relative Fourier-Mukai transform to be the complex of sheaves $`\mathrm{\Phi }(E)`$ on $`J_X`$ given by:
$$\mathrm{\Phi }(E)=R\widehat{\pi }_{}\left(\pi ^{}E𝒫\right)$$
(4)
where $`X\stackrel{\pi }{}X\times _BJ_X\stackrel{\widehat{\pi }}{}J_X`$ are the projection maps and $`𝒫`$ is the relative Poincaré sheaf as before.
We say that $`E`$ is $`\mathrm{\Phi }`$-WIT<sub>n</sub> if $`\mathrm{\Phi }^i(E)=0`$ unless $`i=n`$, where $`\mathrm{\Phi }^i(E)`$ denotes the $`i`$th homology of the complex representing $`\mathrm{\Phi }(E)`$ which is well defined up to isomorphism.
It can be shown that (4) gives an equivalence between $`D(X)`$ and $`D(J_X)`$, the derived categories of bounded complexes of coherent sheaves. Its inverse functor is given by:
$$\widehat{\mathrm{\Phi }}(L)=R\pi _{}\left(\widehat{\pi }^{}L𝒫^{}\right)$$
(5)
so that $`\widehat{\mathrm{\Phi }}^0(\mathrm{\Phi }^1(E))=E`$ and $`\mathrm{\Phi }^1(\widehat{\mathrm{\Phi }}^0(L))=L`$
This is one of a series of Fourier-Mukai transforms for elliptic surfaces parameterized by SL$`{}_{2}{}^{}()`$. For a description of these and their invertibility, see . The rank and fibre degree of the transform $`\mathrm{\Phi }(E)`$ is related to the rank and fibre degree of $`E`$ by multiplication by the element of SL$`{}_{2}{}^{}()`$. The case we shall be interested in this paper is given by the matrix
$$\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\mathrm{SL}_2()$$
Let us now recall the details of the definition of stability (due to Mumford and Takemoto) of torsion-free sheaves:
###### Definition.
A torsion-free sheaf $`E`$ over a polarised surface $`(X,\mathrm{})`$ is said to be $`\mu `$-stable (respectively, $`\mu `$-semistable) if for all proper subsheaves $`F`$ of $`E`$, $`\mu (F)<\mu (E)`$ (respectively, $`\mu (F)\mu (E)`$), where $`\mu ()=c_1()\mathrm{}/r()`$ is called the slope.
Recall that since we are working over surfaces, we can assume that the destabilizing subsheaf is locally-free and that the quotient $`E/F`$ is torsion-free. Moreover, we also may choose $`F`$ to be $`\mu `$-stable.
If $`E`$ is $`\mu `$-stable and locally-free, that is if $`E`$ is a holomorphic vector bundle arising from an irreducible instanton connection via the Hitchin-Kobayashi correspondence, it is easy to see that $`E`$ is $`\mathrm{\Phi }`$-WIT<sub>1</sub>. Hence $`\mathrm{\Phi }(E)=\mathrm{\Phi }^1(E)`$ is not a complex of sheaves, but simply a sheaf on $`J_X`$. This also applies if $`E`$ is not necessarily locally-free but is still torsion-free, as we shall see this in Section 5.
The following Proposition brings together the instanton spectral data described in the previous Section with the Fourier-Mukai transform.
###### Proposition 1.
Let $``$ be a vector bundle with an irreducible $`SU(r)`$ instanton connection over an elliptic surface $`X`$. Let $`E`$ be the associated $`\mu `$-stable, regular holomorphic vector bundle with $`c_1(E)=0`$. Then $`\mathrm{\Phi }^1(E)`$ is a torsion sheaf supported over the instanton spectral curve $`SJ_X`$. Moreover, the restriction of $`\mathrm{\Phi }^1(E)`$ to its support is naturally isomorphic to $``$.
###### Proof.
Let $`j:SJ_X`$ be the inclusion map. This fits into a commuting diagram of maps introduced in (2):
Then
$`j_{}`$ $``$ $`R(j_{}\sigma _{})(\tau _2^{}E\tau _1^{}𝒫)`$
$``$ $`R\widehat{\pi }_{}R\tau _1(\tau _2^{}E\tau _1^{}𝒫)`$
$``$ $`R\widehat{\pi }_{}(𝒫R\tau _1\tau _2^{}E).`$
But $`\tau _2=\pi \tau _1`$ (see diagram 2), so $`R\tau _1\tau _2^{}E\pi ^{}E𝒪_{X\times _BS}`$. Hence, we have a natural map $`\mathrm{\Phi }^1(E)j_{}`$. Since the fibres of $`\mathrm{\Phi }^1(E)`$ and $`j_{}`$ are naturally isomorphic, this must be an isomorphism. Hence, $`j^{}\mathrm{\Phi }^1(E)=j^{}j_{}=`$ as required. ∎
## 4 Instantons on flat 4-tori
Let $`𝕋`$ be a flat 4-dimensional torus, with a fixed complex structure. Denote by $`\widehat{𝕋}`$ the corresponding dual torus, which inherits a flat metric and a complex structure from $`𝕋`$ polarized by the associated Kähler class.
We start by considering an irreducible $`SU(r)`$ bundle $`𝕋`$ with $`c_1()=0`$ and $`c_2()=k`$, plus an instanton connection $`A`$. Equivalently, we can adopt an algebraic geometric point of view and look at the associated $`\mu `$-stable holomorphic vector bundle $`E𝕋`$ of rank $`r`$ and the same Chern classes. Stability now implies the irreducibility hypothesis, which in turns implies that $`h^0(𝕋,E)=h^2(𝕋,E)=0`$, so that $`h^1(𝕋,E)=k`$.
The Nahm transformed bundle $`\widehat{E}\widehat{𝕋}`$ is then constructed as follows. Let $`P𝕋\times \widehat{𝕋}`$ be the Poincaré line bundle, and consider the natural projection maps:
$$𝕋\stackrel{p}{}𝕋\times \widehat{𝕋}\stackrel{\widehat{p}}{}\widehat{𝕋}$$
(6)
Then $`\widehat{E}=R^1\widehat{p}_{}(p^{}EP^{})`$. The fibre of $`\widehat{E}`$ over $`\xi \widehat{𝕋}`$ can be canonically identified with $`\widehat{E}_\xi =H^1(𝕋,EL_\xi )`$, where $`L_\xi 𝕋`$ denotes the flat holomorphic line bundle associated with $`\xi \widehat{𝕋}`$.
Conversely, it can be shown that $`E=R^1p_{}(\widehat{p}^{}\widehat{E}P)`$, with fibres canonically identified with $`E_z=H^1(\widehat{𝕋},\widehat{E}L_z)`$, where $`L_z\widehat{𝕋}`$ denotes the flat holomorphic bundle associated with $`z𝕋`$.
###### Theorem 2.
$`\widehat{E}`$ is a $`\mu `$-stable holomorphic vector bundle of rank $`k`$, flat determinant and $`c_2(\widehat{E})=r`$. Therefore, since it is invertible, the Nahm transform is a bijective correspondence $`_𝕋(r,k)_{\widehat{𝕋}}(k,r)`$.
For the proof, we refer to . The algebro-geometric version is given in . In particular, in the $`k=1`$ case, the Theorem tells us that there are no stable holomorphic bundles with this Chern character, thus precluding the existence of instantons with unit charge. For larger $`k`$, this result is not really helpful as it only converts the problem of constructing instantons/stable bundles over $`𝕋`$ into the same problem over $`\widehat{𝕋}`$. A better understanding of the moduli space of $`SU(r)`$ instantons on $`𝕋`$ with $`k>1`$ requires a new tool, namely the Relative Fourier-Mukai transform introduced in the previous Section.
Now assume that $`𝕋=V\times W`$, i.e. $`𝕋`$ is given by the product of two elliptic curves $`V`$ and $`W`$. We shall denote by $`\widehat{V}`$ and $`\widehat{W}`$ the respective Jacobian curves, so that the dual torus is given by $`\widehat{𝕋}=\widehat{V}\times \widehat{W}`$. This case is particularly interesting because we have two fibration structures leading apparently to two distinct spectral curves, and we would like to understand the relation between them.
Regarding $`𝕋`$ as an elliptic surface over $`V`$, consider the diagram:
(7)
where $`\pi _V`$ is the projection onto the first and second factors, while $`\widehat{\pi }_V`$ is the projection onto the first and third factors. Let also $`𝐏_WW\times \widehat{W}`$ be the Poincaré line bundle. Then
$$=R^1\widehat{\pi }_V(\pi _V^{}E𝐏_W)$$
(8)
is a torsion sheaf on $`V\times \widehat{W}`$, supported over the instanton spectral curve $`S`$.
On the other hand, regard $`\widehat{𝕋}`$ as an elliptic surface over $`\widehat{W}`$ and consider the Nahm transformed bundle $`\widehat{E}\widehat{𝕋}`$. Based on the diagram:
(9)
where $`\pi _W`$ is the projection onto the second and third factors, while $`\widehat{\pi }_W`$ is the projection onto the first and third factors. Let also $`𝐏_VV\times \widehat{V}`$ be the Poincaré line bundle. Then
$$𝒩=R^1\widehat{\pi }_W(\pi _W^{}E𝐏_V)$$
(10)
is a torsion sheaf on $`V\times \widehat{W}`$, supported over the dual instanton spectral curve $`R`$.
Since the bundles $`E`$ and $`\widehat{E}`$ are related via Nahm transform, it seems natural to ask how are the two sets of spectral data $`(S,)`$ and $`(R,𝒩)`$ are related to one another.
Regarding $`𝕋`$ as an elliptic surface over $`V`$, we have just seen that the spectral data $`(S,)`$ is encoded into the torsion sheaf $`\mathrm{\Phi }^1(E)`$ on $`V\times \widehat{W}`$. On the other hand, regarding $`\widehat{𝕋}`$ as an elliptic surface over $`\widehat{W}`$ let us denote the Fourier-Mukai transform of $`\widehat{E}`$ by $`\mathrm{\Psi }(\widehat{E})`$. It is a torsion sheaf on $`V\times \widehat{W}`$ encoding the dual spectral data $`(R,𝒩)`$. Let also $`\widehat{\mathrm{\Psi }}`$ denote the inverse of $`\mathrm{\Psi }`$.
We are finally in a position to prove our first main result:
###### Theorem 3.
The spectral pairs $`(S,)`$ and $`(R,𝒩)`$ are equivalent, in the sense that $``$ and $`𝒩`$ can be canonically identified as sheaves on $`J_X`$.
###### Proof.
If we define the functor $`:D(𝕋)D(\widehat{𝕋})`$ as follows:
$$(E)=R\widehat{p}_{}(p^{}EP^{})$$
where $`p`$ and $`\widehat{p}`$ are the projection maps in (6), then $`\widehat{E}=^1(E)`$. The functor $``$ is Mukai’s original Fourier transform introduced in . Using this notation, it is enough to show that $`\widehat{\mathrm{\Psi }}\mathrm{\Phi }=`$.
To see this it suffices to show that $`\widehat{\mathrm{\Psi }}(\mathrm{\Phi }(𝒪_x))=\widehat{\mathrm{\Psi }}(𝒫_x)=P_x`$ over $`\widehat{𝕋}`$. But $`𝒫_x𝒫_y^{}`$ is supported at the intersection of translates of $`V`$ and $`\widehat{W}`$. Hence, $`𝒫_x`$ is $`\mathrm{\Psi }`$-WIT<sub>0</sub> and $`\widehat{\mathrm{\Psi }}(𝒫_x)`$ is a flat line bundle. On the other hand, the properties of $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ imply that translating $`x`$ to $`y`$ twists $`\widehat{\mathrm{\Psi }}(𝒫_x)`$ by $`P_{yx}`$ and so by normalizing the Poincaré bundles appropriately, we have $`\widehat{\mathrm{\Psi }}(𝒫_x)=P_x`$. ∎
Similarly, one could regard $`𝕋`$ as an elliptic surface over $`W`$ and $`\widehat{𝕋}`$ as an elliptic surface over $`\widehat{V}`$. The analogue of Theorem 3 would again hold, so that the corresponding spectral data would also coincide as torsion sheaves in $`\widehat{V}\times W`$.
More generally, one has the following commuting diagram of derived categories:
(11)
Of course, the functor $`\mathrm{{\rm Y}}`$ is given by the Fourier-Mukai transform regarding $`𝕋`$ as an elliptic surface over $`W`$, while the functor $`\mathrm{\Xi }`$ is the Fourier-Mukai transform regarding $`\widehat{𝕋}`$ as an elliptic surface over $`\widehat{V}`$.
###### Remark.
One concludes from Theorem 3 that there is an equivalence between the following three objects: instantons on a flat 4-torus (i.e. $`\mu `$-stable bundles over $`𝕋`$), spectral data, and instantons on the dual 4-torus (i.e. $`\mu `$-stable bundles over $`\widehat{𝕋}`$). Such “circle of ideas” has been previously establish for monopoles , doubly-periodic instantons and periodic monopoles . Indeed, one expects that a similar scheme will hold for all translation invariant instantons on $`^4`$.
## 5 Preservation of Stability
We now aim at establishing the link between the stability of $`E`$ and the stability of $`L=\mathrm{\Phi }^1(E)`$. For torsion sheaves, we have a concept of stability due to Simpson :
###### Definition.
A torsion sheaf $`L`$ on a polarised variety $`(X,\mathrm{})`$ is said to be $`p`$-stable (respectively, $`p`$-semistable) if it is of pure dimension (i.e. the support of all subsheaves have the same dimension) and for all proper subsheaves $`M`$ of $`L`$ we have $`p(M,n)<p(L,n)`$ (respectively, $`p(M,n)p(L,n)`$) for all sufficiently large $`n`$, where $`p(L,n)`$ denotes the reduced Hilbert polynomial of $`L𝒪(n\mathrm{})`$.
In particular, a torsion sheaf which is supported on an irreducible curve and whose restriction has rank one is automatically $`p`$-stable with respect to any polarisation.
Now let $`XB`$ be a (relatively minimal) elliptic surface with a section $`\sigma `$; let $`E`$ be a torsion-free sheaf on $`X`$ with $`ch(E)=(r,0,k)`$.
###### Proposition 4.
Suppose that the restriction of $`E`$ to some smooth fibre $`F_s`$ is semistable. Then $`E`$ is $`\mathrm{\Phi }`$-WIT<sub>1</sub> and $`L=\mathrm{\Phi }^1(E)`$ has Chern character $`(0,k\widehat{f}+r\widehat{\sigma },0)`$, where $`\widehat{\sigma }`$ is the class given by the zero section of $`J_X`$. In particular, $`p(L,n)=n`$.
###### Proof.
First note that if $`E`$ is torsion-free then $`\mathrm{\Phi }^0(E)=\widehat{\pi }_{}\left(\pi ^{}E𝒫\right)`$ is also torsion-free. Hence its support consists is either empty or consists of the whole relative jacobian surface $`J_X`$. Moreover, if $`E`$ is not $`\mathrm{\Phi }`$-WIT<sub>1</sub> then the support of $`\mathrm{\Phi }^0(E)`$ is contained in the support of $`\mathrm{\Phi }^1(E)`$ since $`\chi (E|_{F_x})=0`$ for all $`xX`$.
However, $`\mathrm{\Phi }^1(E)`$ cannot be supported on the whole surface since it does not contain the whole fibre $`\widehat{F}_s`$ where $`E|_{F_s}`$ is semistable. Thus $`\mathrm{\Phi }^0(E)=0`$ and $`E`$ is $`\mathrm{\Phi }`$-WIT<sub>1</sub>.
Now $`\mathrm{\Phi }^1(𝒪_X)`$ is the trivial line bundle supported on the zero section of $`J_X`$ and $`\mathrm{\Phi }^0(𝒪_x)`$ is a flat line bundle supported on (a divisor in) $`\widehat{f}`$. Hence, since $`ch(E)=rch(𝒪_X)+kch(𝒪_x)`$, we conclude that:
$$ch(\mathrm{\Phi }^1(E))=rch(\mathrm{\Phi }^1(𝒪_X))+kch(\mathrm{\Phi }^0(𝒪_x))=(0,k\widehat{f}+r\widehat{\sigma },0)$$
as required. ∎
Note that the rank and fibre degrees of $`L`$ are given by $`r(L)=c_1(E)f=0`$ and $`c_1(L)\widehat{f}=r`$ from the definition of the relative transform.
Let $`S`$ denote the support of $`\mathrm{\Phi }^1(E)`$; the statement above does not imply that $`S`$ is a divisor in $`J_X`$, since it might contain some 0-dimensional components (i.e. isolated points).
###### Lemma 5.
If $`E`$ is generically fibrewise semistable, then $`S=\mathrm{supp}\mathrm{\Phi }^1(E)`$ has no 0-dimensional components. Moreover, $`L=\mathrm{\Phi }^1(E)`$ is of pure dimension one.
###### Proof.
Suppose that $`qS`$ is an isolated point. It cannot belong to a fibre $`\widehat{F}_x`$ such that $`E|_{F_x}`$ is semistable, since the restriction to every nearby fibre is also semistable and the restriction of $`E`$ to the fibres varies holomorphically.
So $`E|_{F_q}`$ is unstable. But $`S`$ contains all such fibres, since $`h^1(F_{\widehat{\pi }(p)},E(p))>0`$ for all $`p\widehat{\pi }^1(\pi (q))`$.
To obtain the second statement, we must show that $`L`$ has no proper subsheaves supported on points. For a contradiction, suppose it does, and let $`F`$ be a subsheaf of $`L`$ with 0-dimensional support. Then $`\widehat{\mathrm{\Phi }}^0(F)`$ is a torsion subsheaf of $`\widehat{\mathrm{\Phi }}^0(L)=E`$, thus contradicting the hypothesis that $`E`$ is torsion-free. ∎
Thus, we will call $`S`$ the spectral curve associated to the generically fibrewise semistable torsion-free sheaf $`E`$, since it generalizes our previous definition for regular locally-free sheaves. As before, let us denote by $``$ the restriction of $`\mathrm{\Phi }^1(E)`$ to its support; in general, it is a coherent sheaf on $`S`$. Note that $`\mathrm{deg}()=0`$.
###### Remark.
If the restriction of $`E`$ to some smooth fibre has no sections, then semicontinuity implies that $`E`$ is generically fibrewise semistable, and hence $`\mathrm{\Phi }`$-WIT<sub>1</sub>. It is somewhat surprising that the spectral data can actually be defined in an interesting, meaningful way under such mild conditions.
Summing up, we also conclude:
###### Corollary 6.
If $`E`$ is a torsion-free and $`\mathrm{\Phi }`$-WIT<sub>1</sub>, then $`\mathrm{\Phi }(E)=\mathrm{\Phi }^1(E)`$ is a torsion sheaf of pure dimension one.
Conversely, we also have:
###### Lemma 7.
If $`L`$ is a $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> torsion sheaf of pure dimension one on $`J_X`$ with $`ch(L)=(0,k\widehat{f}+r\widehat{\sigma },0)`$, then $`\widehat{\mathrm{\Phi }}^0(L)`$ is torsion-free.
###### Proof.
Suppose that the support of $`L`$ decomposes as $`\mathrm{\Sigma }+F`$, where $`F`$ is the sum of fibres. When this happens, we have a subsheaf $`K`$ of $`L`$ supported on $`F`$ with degree 0. Assume that $`F`$ is the maximal such effective subdivisor of $`S`$. Let $`Q=L/K`$. Then $`Q`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> and $`K`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>1</sub>. The resulting sequence after transforming with $`\widehat{\mathrm{\Phi }}`$ is then
$$0EE^{}𝒪_Z0$$
(12)
where $`E=\widehat{\mathrm{\Phi }}^0(L)`$, for some zero dimensional subscheme $`Z`$. ∎
###### Remark.
Notice that applying $`\mathrm{\Phi }`$ to the sequence (12) shows that if $`E`$ is torsion-free but not locally-free, then the support of $`L`$ contains fibres.
The following lemma due to Bridgeland characterizes $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> sheaves on $`J_X`$:
###### Lemma 8.
A sheaf $`L`$ on $`J_X`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> if and only if
$$\mathrm{Hom}(L,𝒫_x)=0xX$$
### 5.1 Suitable polarizations
Now let $`\mathrm{}`$ be a polarisation of the elliptic surface $`X`$, and let $`\widehat{\mathrm{}}`$ be the induced polarisation on $`J_X`$. If $`\mathrm{}`$ is arbitrary, it is not difficult to see that a $`\mu `$-unstable torsion-free sheaf on $`X`$ can have a $`p`$-stable transform.
Indeed, let $`X`$ be an elliptic surface whose fibres are all smooth, and let $`E`$ be the locally-free sheaf given by the following extension:
$$0𝒪(\sigma +df)E𝒪(\sigma df)0$$
Clearly, $`c_1(E)=0`$ and $`c_2(E)=\sigma ^2+2d`$. For $`d`$ sufficiently large, $`c_2(E)>0`$ and $`\mathrm{}(\sigma +df)=(\mathrm{}\sigma )+d(\mathrm{}f)>0`$, so $`E`$ is $`\mu `$-unstable with respect to $`\mathrm{}`$.
On the other hand, notice that the restriction of $`E`$ to each fibre is an extension of $`Q\mathrm{Pic}^d(T)`$ by its dual. Thus, the bundles obtained by the above extension are generically regular. This means $`\mathrm{\Phi }^1(E)`$ is supported on a smooth, irreducible curve, so that $`\mathrm{\Phi }^1(E)`$ is necessarily $`p`$-stable with respect to $`\widehat{\mathrm{}}`$.
Therefore, we can only expect the Fourier-Mukai transform to preserve stability if we restrict the choice of polarisation on $`X`$ in some convenient way.
###### Definition.
Let $`c`$ be a positive integer, and consider the set
$$\mathrm{\Xi }(c)=\{\xi \mathrm{Div}(X)|4c\xi ^2<0\mathrm{and}\xi \mathrm{mod2}=0\}$$
Let $`W^\xi `$ be the intersection of the hyperplane $`\xi ^{}`$ with the ample cone of $`X`$. A polarisation $`\mathrm{}`$ is said to be $`c`$-suitable if $`\mathrm{}W^\xi `$ and sign$`(\mathrm{}\xi )=`$sign$`(f\xi )`$ for all $`\xi \mathrm{\Xi }(c)`$.
It is easy to see that suitable polarizations exist for every $`c`$. The following important result is due to Friedman, Morgan & Witten :
###### Theorem 9.
Let $`E`$ be a torsion-free sheaf on $`X`$ with $`ch(E)=(r,0,k)`$. If $`E`$ is $`\mu `$-semistable with respect to a $`k`$-suitable polarisation $`\mathrm{}`$, then $`E`$ is generically fibrewise semistable.
### 5.2 Preservation of stability
It follows from Theorem 9 that if $`E`$ is a torsion-free sheaf on $`X`$ with $`ch(E)=(r,0,k)`$, which is $`\mu `$-semistable with respect to a $`k`$-suitable polarisation $`\mathrm{}`$, then $`L=\mathrm{\Phi }^1(E)`$ is a torsion sheaf of pure dimension one. Let us now consider its stability (in the sense of Simpson).
###### Proposition 10.
Let $`E`$ be a torsion-free sheaf on $`X`$ which is $`\mu `$-semistable with respect to a $`k`$-suitable polarisation $`\mathrm{}`$. Then $`L=\mathrm{\Phi }^1(E)`$ is $`p`$-semistable with respect to $`\widehat{\mathrm{}}`$.
###### Proof.
Suppose that
$$0MLN0$$
is a destabilizing sequence for $`L`$. We can assume $`M`$ is semistable and $`N`$ has pure dimension 1. Now $`p(M,n)=n+\alpha `$, for some $`\alpha >0`$. Then $`Hom(M,𝒫_x)=0`$ for all $`xX`$, since $`p(𝒫_x,n)=n`$. Thus $`M`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> by Lemma 8. Arguing as in Proposition 4, we have $`c_1(\widehat{\mathrm{\Phi }}^0(M))=\alpha f`$ hence $`\mu (\widehat{\mathrm{\Phi }}^0(M))>0`$. So $`\widehat{\mathrm{\Phi }}^0(M)`$ will destabilize $`E=\widehat{\mathrm{\Phi }}^0(L)`$ unless $`r(\widehat{\mathrm{\Phi }}^0(M))=r(E)`$. For this to be the case we must have $`c_1(M)\widehat{f}=c_1(L)\widehat{f}`$. Then $`N`$ is supported on a sum of fibres. But since $`N`$ must be $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> we see that it has non-negative degree on these fibres and so $`\alpha =0`$, thus contradicting the assumption on $`M`$. ∎
If we examine the proof more carefully, we can also see that the assumption that $`E`$ is $`\mu `$-stable implies that $`L`$ is $`p`$-stable unless its support decomposes as $`D+D^{}`$, where $`D`$ is the sum of fibres. When this happens, we have a subsheaf $`K`$ of $`L`$ supported on $`D`$ with degree 0. Assume that $`D`$ is the maximal such effective subdivisor of $`S`$. Let $`Q=L/K`$. Then $`Q`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> and $`K`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>1</sub>. The resulting sequence after transforming with $`\widehat{\mathrm{\Phi }}`$ is then $`EE^{}𝒪_Z`$ for some zero dimensional subscheme $`Z`$. Conversely, applying $`\mathrm{\Phi }`$ to this sequence shows that the support of $`L`$ contains fibres. We have therefore established:
###### Proposition 11.
Suppose $`E`$ is $`\mu `$-semistable with respect to a $`k`$-suitable polarisation $`\mathrm{}`$ with $`c_1(E)=0`$. If $`E`$ is $`\mu `$-stable and locally-free then $`L`$ is $`p`$-stable. Moreover, $`E`$ is locally-free if and only if $`L`$ is destabilized by a sheaf supported on a fibre.
We can now consider the opposite question: if we assume $`L`$ is $`p`$-semistable what can we say about its transform?
###### Proposition 12.
Suppose $`L`$ is a $`p`$-stable sheaf on $`J_X`$ with Chern character $`(0,k\widehat{f}+r\widehat{\sigma },0)`$, where $`r,k>1`$. Then $`L`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> and $`\widehat{\mathrm{\Phi }}^0(L)`$ is locally-free, $`\mu `$-stable with respect to a $`k`$-suitable polarisation and such that $`ch(\widehat{\mathrm{\Phi }}^0(L))=(r,0,k)`$.
###### Proof.
The first statement follows from lemma 8, which requires us to show that $`Hom(L,𝒫_x)=0`$ for all $`x`$. However, $`p(𝒫_x,n)=n`$ and any map $`L𝒫_x`$ would contradict the stability of $`L`$.
Now suppose that $`\widehat{\mathrm{\Phi }}^0(L)`$ is not $`\mu `$-stable and let $`A`$ be the destabilizing subsheaf, so that $`\mu (A)\mu (E)=0`$ Moreover, we may assume that $`A`$ is $`\mu `$-stable, thus $`A`$ is $`\mathrm{\Phi }`$-WIT<sub>1</sub>.
We may also assume that the quotient $`B=E/F`$ is $`\mu `$-stable. Then both $`A`$ and $`B`$ are $`\mathrm{\Phi }`$-WIT<sub>1</sub>. Since their transforms must have zero rank, $`A`$ and $`B`$ both have zero fibre degree and so $`p(\mathrm{\Phi }^1(A),n)n`$ and this contradicts the stability of $`L`$.
The fact that $`\widehat{\mathrm{\Phi }}^0(L)`$ is locally-free follows from the last part of Proposition 11. ∎
A similar argument shows that if $`L`$ is $`p`$-semistable and $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> then $`\widehat{\mathrm{\Phi }}^0(L)`$ must be $`\mu `$-semistable. Note, however, that it the semistability of $`L`$ does not imply that $`L`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub>. Examining the proof above we see that this happens precisely when $`L`$ is destabilized by mapping to a sheaf supported on a fibre. In particular, such sheaves are S-equivalent to sheaves which are $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> but whose transforms are not locally-free.
We summarize the results of this section in the following theorem
###### Theorem 13.
Suppose $`E`$ is a torsion-free sheaf with rank $`r`$, $`c_1(E)=0`$ and $`c_2(E)=k`$ on a relatively minimal elliptic surface $`X`$ over $`B`$. Let $`\omega `$ be a $`k`$-suitable polarisation on $`X`$. Then
1. $`E`$ is a $`\mu `$-stable locally free sheaf if and only if its transform is a $`p`$-stable torsion sheaf supported on a divisor in $`|k\widehat{f}+r\widehat{\sigma }|`$, where $`\widehat{\sigma }`$ is a section of $`J_X`$ and $`\widehat{f}`$ is a fibre class.
2. $`E`$ is $`\mu `$-stable properly torsion-free if and only if its transform is a $`p`$-semistable torsion sheaf supported on a reducible divisor in $`|k\widehat{f}+r\widehat{\sigma }|`$ which is destabilized only by a sheaf supported on a fibre.
3. $`E`$ is properly $`\mu `$-semistable if and only if its transform is a $`p`$-semistable torsion sheaf supported on a reducible divisor in $`|k\widehat{f}+r\widehat{\sigma }|`$.
One can also give criteria for the Gieseker stability of $`E`$ in terms of destabilizing properties of $`L`$, but these seem to be less useful.
###### Remark.
Recently, Hernández Ruipérez and Muñoz Porras and Yoshioka have independently obtained stability results close to those presented in this Section.
###### Remark.
It follows from the first item in Theorem 13 that the Fourier-Mukai functor $`\mathrm{\Phi }`$ induces a bijective map from $`_X(r,k)`$ onto $`𝒮_{J_X}(0,k\widehat{f}+r\widehat{\sigma },0)`$, the Simpson moduli space of p-stable torsion sheaves on $`J_X`$ with the given Chern classes. As we will observe in Section 7 below, this map is also a hyperkähler isometry when $`X`$ is an elliptic K3 or abelian surface (i.e. when $`X`$ is hyperkähler).
We can see that the geometry of the spectral data is easily linked to the sheaf theoretic properties of the original sheaf. This should make it very easy in practice to use the spectral data to analyse properties of the specimen sheaf. We shall see an example of this in the subsequent Sections where we use the spectral data to explore the fibration structure on the moduli spaces.
## 6 The Fibration Structure
From now on we assume that $`X`$ is either an elliptic $`K3`$ surface with a section or a product of elliptic curves. We have seen that there is a natural map $`\mathrm{\Pi }`$ from the moduli space of $`\mu `$-stable torsion-free sheaves $`_X^{\mathrm{TF}}(r,k)`$ to the set $``$ of spectral curves. In the case of the K3 surface the base $``$ is just the linear system $`|k\widehat{f}+r\sigma |`$ while for the abelian surface it can be expressed as the total space of the projectivized bundle $`^{}(𝒪(kf+r\widehat{\sigma }))`$, where $`^{}`$ is the Mukai transform $`D(\widehat{V}\times W)D(V\times \widehat{W})`$ (this can be factored as $`\mathrm{\Phi }\widehat{\mathrm{{\rm Y}}}=\widehat{\mathrm{\Psi }}\mathrm{\Xi }`$ from (11)). The fibre of $`\mathrm{\Pi }`$ over $`S`$ is given by suitable subspaces of Jac$`{}_{g(S)1}{}^{}(S)`$. The hypothesis on $`X`$ implies that for $``$ to be non-empty we must have both $`r`$ and $`k`$ at least 2.
###### Theorem 14.
The map $`_X(r,k)\stackrel{\mathrm{\Pi }}{}`$ is a surjective map of varieties.
###### Proof.
We must first show that $`\mathrm{\Pi }`$ is a well defined map of varieties. To see this we use the argument given by Friedman and Morgan in \[4, Chapter VII, Thm 1.14\]. The key idea is to observe that $`\mathrm{\Pi }`$ coincides (locally in the étale or complex topology) with the projection map from the universal sheaf corresponding to the relative Picard scheme of degree $`g1`$ line bundles on families of genus $`g`$ curves. The universal sheaf only exists locally in these topologies but this is enough to show that $`\mathrm{\Pi }`$ is holomorphic.
Given a spectral curve $`S`$, observe that its structure sheaf $`𝒪_S`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> since in the short exact sequence $`0\mathrm{\Lambda }^1𝒪𝒪_S0`$, $`\mathrm{\Lambda }^1`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>1</sub> and $`𝒪`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub>.
We aim to construct a $`g(S)1`$ degree line bundle $`L`$ over the curve $`S`$, which is stable as a torsion sheaf on $`J_X`$. That is, we need to choose $`g(S)1`$ points on $`S`$ and guarantee that the restriction of $`L`$ to any proper component $`S_i`$ of $`S`$ satisfies deg$`(L|_{S_i})>0`$. We do this by choosing a $`g(S_i)`$ points on each component $`S_i`$. This can always be done because if $`g(S_i)>0`$ then $`S_i`$ intersects $`SS_i`$ at least twice since $`k`$ and $`r`$ are both at least 2 and so $`g(SS_i)<g(S)g(S_i)`$. Without loss of generality, we assume that the set we have just chosen $`ZS`$ consists of distinct points away from the singularities of $`S`$. Therefore,
$$Ext^1(𝒪_Z,𝒪_S)=\underset{zZ}{}Ext^1(𝒪_z,𝒪_S)=\underset{zZ}{}T_zS$$
We pick a class in $`Ext^1(𝒪_Z,𝒪_S)`$ which is non-zero on each factor. This defines a torsion sheaf $`L`$ on $`J_X`$ given by this extension class. Then $`L`$ is stable and locally free on its support by the choice of $`Z`$.
By Proposition 12, $`L`$ is $`\widehat{\mathrm{\Phi }}`$-WIT<sub>0</sub> and $`E=\widehat{\mathrm{\Phi }}^0(L)`$ is $`\mu `$-stable vector bundle such that $`\mathrm{\Pi }(E)=S`$. ∎
From this proof we can also see that any $`E`$ can be written as an extension
$$0\widehat{\mathrm{\Phi }}^0(𝒪_S)E\widehat{\mathrm{\Phi }}^0(𝒪_Z)0.$$
We shall make use of this in the next Section. Note that such representation is not unique.
Such a fibration structure will exist for the moduli spaces over any elliptic surface with a section but this surjectivity result may not hold. For the rank 2 case see \[5, Thm. 37\] or \[4, Thm. 1.14\]. It is also the case that the fibres will not exist in the middle dimension which is the situation we wish to consider in the next Section.
## 7 The Hyperkähler Structure
When $`X`$ is K3 or a torus, it is well known that the moduli spaces of sheaves are hyperkähler. To see this is, note that by a result of Mukai’s , for each complex structure $`I`$ on $`X`$ we obtain a complex symplectic structure $`\mathrm{\Omega }_I`$ on $`_X`$. Then by Yau’s proof of the Calabi conjecture to obtain the full hyperkähler structure on $`_X`$.
In fact, the complex symplectic structures arise in a very natural way:
$$T_{[E]}\times T_{[E]}Ext^1(E,E)\times Ext^1(E,E)\stackrel{}{}Ext^2(E,E),$$
where the $`E`$ is an $`𝒪_X`$-module with respect to $`I`$. But we can express this even more simply by observing that $`Ext^i(E,E)=Hom_{D(X)}(E,E[i])`$. Then the cup product $``$ is just composition of maps in the derived category. Since $`\mathrm{\Phi }`$ is a functor, it must preserve these and so we see that $`\mathrm{\Phi }`$ induces a complex symplectomorphism:
$$_X(r,k)𝒮_{J_X}(0,k\widehat{f}+r\widehat{\sigma },0)$$
Since this happens for each complex structure, the moduli map induced by $`\mathrm{\Phi }`$ is actually a hyperkähler isometry.
We aim to show now that the fibration structure we have defined in the last Section on $``$ has Lagrangian fibres with respect to this complex symplectic structure.
###### Proposition 15.
If $`t_1`$ and $`t_2`$ are two tangent vectors to the fibre of $`\mathrm{\Pi }`$ then $`\mathrm{\Omega }_I(t_1,t_2)=0`$.
###### Proof.
By continuity it suffices to prove this when $`t_i`$ are defined over a point $`E`$ given by an extension:
$$0\widehat{\mathrm{\Phi }}(𝒪_S)E\widehat{\mathrm{\Phi }}(𝒪_Z)0.$$
where $`Z`$ consists of discrete points and $`S`$ is smooth. Deformations of $`E`$ arising from the fibre of $`\mathrm{\Pi }`$ are determined by deformations of $`Z`$ along $`S`$. Then $`Ext^2(E,E)Ext^2(L,L)`$ is generated by a non-zero vector in $`Ext^2(𝒪_Z,𝒪_Z)`$ and the fibre $`\mathrm{\Pi }_S`$ over $`S`$ has tangent space given by $`_{zZ}\lambda _z,`$ where $`\lambda _z`$ generates the tangent space to $`S`$ at $`z`$. But $`\lambda _z\lambda _z=0`$ in $`Ext^2(𝒪_Z,𝒪_Z)`$ for each $`z`$ and so $`\mathrm{\Omega }_I(t_1,t_2)=0`$ as required. ∎
In the case of a product abelian surface we have natural hyperkähler actions of the torus and its dual via translations and twisting by flat line bundles. These act naturally on the fibration structure and the resulting quotient is also a Lagrangian fibration.
We can say a little more in the particular case of the product of two elliptic curves. For $`𝕋=V\times W`$, observe that $`_𝕋(r,k)`$ has two such fibration structures:
(13)
where $`\widehat{}=\widehat{}(𝒪(k\widehat{f}+r\sigma ))`$. Note that the fibres are Lagrangian with respect to the same complex symplectic structure and the bases are biholomorphic since the underlying vector bundles can be canonically identified by pulling back along the isomorphism $`V\times \widehat{W}\widehat{V}\times W`$.
In summary we have proved:
###### Theorem 16.
If X is an elliptic K3 surface with a section or a product of elliptic curves, then the moduli space of instantons admits a fibration structure over a compact base (which is a projective space or a projective bundle over an abelian surface, respectively) and the fibres are Lagrangian with respect to the natural complex symplectic structure.
## Acknowledgements
The first named author would like to thank Yale University and the University of Pennsylvania for their hospitality. The second named author would like to thank Yale University and the Engineering and Physical Sciences Research Council of the UK for their support. We would also like to thank Tom Bridgeland, Robert Friedman and the referee for several useful comments.
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# Statistics of energy levels and zero temperature dynamics for deterministic spin models with glassy behaviour
## 1 Introduction
A main issue in glassy systems is the analogy between glass-forming liquids and discontinuous spin-glasses, first pointed out in the pioneering works by Kirkpatrick, Thirumalai and Wolynes \[KTW\]. In both cases the thermodynamical properties can be indeed related to the dynamical evolution in an energy landscape. In liquid theory one can define the notion of inherent structures \[SW\] (local minima of the potential energy, each one surronded by its attraction basin or valley) and configurational entropy, i.e. the logarithm of the number of these minima divided by the number of particles in the system. Then the low-temperature dynamical evolution can be described as a superposition of an intra-basin “fast” motion and a “slow” crossing of energy barriers. If the temperature of the system is small enough, namely less than the Mode Coupling critical temperature $`T_{MC}`$, the system gets trapped in one of the basins. Since the number of energy minima diverges exponentially with the size of the system, a thermodynamic transition can be associated with an entropy crisis: the Kauzmann temperature $`T_K`$ of the glassy transition corresponds to the vanishing of the configurational entropy. We refer the reader to \[MP\] for an overview on equilibrium thermodynamics of glasses.
Consider now the class of discontinuous spin glasses, i.e. the mean-field models involving a random $`p`$-spin interaction. Also these models show a dynamical transition at a temperature $`T_D`$ (corresponding to $`T_{MC}`$) where dynamical ergodicity breaks down; a thermodynamic entropy-driven transition takes place at a lower temperature $`T_{1RSB}`$ (corresponding to $`T_K`$), at which replica symmetry breaks down with a “one step” pattern. Here the local minima of the free energy correspond to the solutions of the mean field TAP equations. Anyway at temperature $`T=0`$ metastable states with respect to any dynamics reduce to $`1`$-spin-flip stable states.
The main gap in the analogy between structural glasses and discontinuous spin-glasses is that in the latter models, unlike the former, the couplings between spins are quenched random variables. A significant, recent step in filling this gap has been made by the introduction of the deterministic, i.e. non-random, spin models which show a dynamical phase transition with a discontinuous order parameter and an equilibrium phase transition at a lower temperature associated with the vanishing of the high-temperature entropy \[MPR, BM, BDGU, NM, PP\]. It is the high degree of frustation among the couplings, not the disorder, to generate a huge number of metastable states and thus the glassy behaviour. The discovery of these models proved that disorder is not necessary to reproduce a complex free energy landscape.
Metastables states in infinite-range disordered spin-glasses have been extensively studied, both in the SK model where the number of $`1`$-spin-flip stable states scales like $`\mathrm{exp}(0.1992N)`$ \[TE, BrM, DGGO, MPV\], $`N`$ being the size of the system (number of spins in the one-dimensional case), and in $`p`$-spin interaction spin-glasses \[OF\].
Here we deal with the same question in deterministic models. By probabilistic arguments we will obtain, for the models introduced in \[MPR\], a lower bound on the number of $`1`$-spin-flip stable states, which increases exponentially with the size of the system. Hence this deterministic model exhibits the main feature of glassy behaviour.
The paper is organized as follows: in Section 2 we review the basic properties of the model we consider, the sine model of \[MPR\], a deterministic, one dimensional chain of $`N`$ spins with long-range oscillating interaction. In Section 3 we study the limiting distribution of the rescaled energy density, showing that it gets $`\delta `$-distributed in the thermodynamic limit. This property, which holds for the Curie-Weiss case, is an indication of the mean-field nature of the model. In Section 4 we deal with the state space of the system, computing explicitly (for $`N`$ prime analytically, for other values of $`N`$ numerically) the distribution of the energy levels by flipping one spin at a time; among other things we show that there can be a large number of states with almost zero overlap with the ground state but very close in energy to it. Finally, in Section 5 we derive the main result of the paper, that is a lower exponential bound for the number of metastable states at temperature $`T=0`$. Acknowledgments: M. D.E. contributed to this paper during his visit to the School of Mathematics of the Georgia Institute of Technology, whose support and excellent working conditions he gratefully acknowledges.
C. G. wants to thank the Department of Physics of Brown University, where part of this work was performed, for the kind hospitality during his visit under the exchange program between the Bologna and Brown Universities .
## 2 Orthogonal interaction matrices and the sine model
The basic setup is a probability space $`(\mathrm{\Sigma }_N,_N,\mathrm{P}_N)`$. The sample space $`\mathrm{\Sigma }_N`$ is the configuration space, i.e. $`\mathrm{\Sigma }_N=\{1,1\}^N`$ whose elements are the sequences $`\sigma =(\sigma _1,\mathrm{},\sigma _N)`$ with $`\sigma _i=\pm 1`$; $`_N`$ is the finite algebra with $`2^{2^N}`$ elements, and the a priori (or infinite-temperature) probability measure $`\mathrm{P}_N`$ is given by
$$\mathrm{P}_N(C)=\frac{1}{2^N}\underset{\sigma C}{}1.$$
(2.1)
The Hamiltonian is the function on $`\mathrm{\Sigma }_N`$ defined as
$$H(\sigma )=\frac{1}{2}\underset{xy}{}J_{xy}\sigma _x\sigma _y=\frac{1}{2}<J\sigma ,\sigma >$$
(2.2)
where $`J=(J_{xy})`$ is a symmetric real orthogonal $`N\times N`$ matrix given from the outset. Although many of the results presented here will hold for a generic symmetric orthogonal matrix (e.g. of the form $`J=OLO^T`$ with $`L`$ a diagonal matrix whose elements are $`\pm 1`$ and $`O`$ a generic orthogonal matrix chosen at random w.r.t. the Haar measure on the orthogonal group) in what follows we shall examine the following particular example known as the sine model:
$$J_{xy}=\frac{2}{\sqrt{2N+1}}\mathrm{sin}\left(\frac{2\pi xy}{2N+1}\right).$$
(2.3)
which satisfies (we assume $`N`$ odd) <sup>1</sup><sup>1</sup>1One might also consider interaction matrices with zero diagonal terms, recovering orthogonality in large $`N`$ limit. This amounts to put the average energy equal to zero (instead of $`1/2`$) and may be convenient for particular purposes (see Section 3).
$$JJ^T=\mathrm{Id}\text{and}\underset{x=1}{\overset{N}{}}J_{xx}=\underset{x=1}{\overset{N}{}}J_{xx}^2=1.$$
(2.4)
This model has been introduced by Marinari, Parisi and Ritort as a deterministic system with high frustration (competiting interactions with different signs and strengths) able to reproduce the complex thermodynamical behaviour typical of spin glasses \[MPR\]. It has been investigated analytically in the high-temperature regime, (through an high-temperature expansion), and numerically also in the low-temperature phase (using Montecarlo annealing). The analytical study revealed the existence of a static phase transition at a temperature $`T_S=0.065`$ where the high-temperature entropy vanishes, while evidence of the existence of a higher temperature $`T_D=0.134`$ where the system undergoes a dynamical transition of second order (i.e. with a jump in the specific heat) with a discontinuous order parameter has been put forward by numerical analysis. It has also been shown, using the replica formalism, that most of the thermodynamical properties of this model are the same as those of a generic symmetric orthogonal matrix (the static transition corresponding to RSB while the dynamical transition being given by the so-called “marginality condition”).
## 3 The limiting distribution of the rescaled energy levels.
The knowledge of the eigenvalues of $`J`$ imposes simple bounds on the energy of any spin configuration. Indeed a state vector $`\sigma `$ can be decomposed into his projections on the various orthogonal eigenspaces relatives to different eigenvalues. Here, due to orthogonality, the possible eigenvalues are $`+1,1`$ so that
$$\frac{N}{2}H(\sigma )\frac{N}{2}.$$
(3.5)
Let us consider the rescaled and shifted Hamiltonian (representing the energy per site, or energy density of the model, plus the ‘zero point’ energy $`1/2`$)
$$h(\sigma )=\frac{H(\sigma )}{N}+\frac{1}{2}$$
(3.6)
which takes values in $`[0,1]`$. We shall show that in the limit $`N\mathrm{}`$ the energy density $`h`$ gets $`\delta `$-distributed at $`x=1/2`$. We point out that this property can be immediately proved for the Curie-Weiss model, thus indicating a mean field behaviour of the present model in the thermodynamic limit. To this end consider the partition function $`Z_N`$ at inverse temperature $`\beta `$:
$$Z_N(\beta )=\underset{\sigma \mathrm{\Sigma }_N}{}\mathrm{exp}(\beta H(\sigma ))=2^N\mathrm{E}_N(e^{\beta H}),$$
(3.7)
where $`\mathrm{E}_N`$ denotes the expectation wrt $`\mathrm{P}_N`$, and note that the characteristic function of $`h`$ can be written as
$$\mathrm{E}_N(e^{\lambda h})=e^{\lambda /2}\frac{Z_N(\lambda /N)}{2^N}.$$
(3.8)
This expression will prove useful to compute the limiting expression of the characteristic function of the energy density $`h`$ without knowing the expression of all its moments. To see this, we first decouple the spins as follows: let $`B`$ be an orthogonal matrix such that $`B^TJB=D`$ with $`D=\mathrm{diag}(d_1,\mathrm{},d_N)`$. Since $`\mathrm{det}J0`$ we have $`d_i0`$, $`i=1,\mathrm{},N`$, and $`\mathrm{det}J^1=_id_i^1`$. Let $`uIR^N`$ be such that $`\sigma =Bu`$. We have $`<J\sigma ,\sigma >=<Bu,JBu>=<u,Du>`$, and thus
$`\mathrm{exp}({\displaystyle \frac{\lambda }{2N}}<J\sigma ,\sigma >)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{exp}({\displaystyle \frac{\lambda }{2N}}d_iu_i^2)`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{exp}\left({\displaystyle \frac{x_i^2}{2}}+\sqrt{{\displaystyle \frac{\lambda d_i}{N}}}u_ix_i\right)𝑑x_i`$
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^{N/2}}}{\displaystyle _{IR^N}}\mathrm{exp}({\displaystyle \frac{1}{2}}<x,x>+\sqrt{{\displaystyle \frac{\lambda }{N}}}<u,D^{1/2}x>)dx`$
$`=`$ $`{\displaystyle \frac{\mathrm{det}J^{\frac{1}{2}}}{(2\pi \lambda )^{N/2}}}{\displaystyle _{IR^N}}\mathrm{exp}({\displaystyle \frac{1}{2\lambda }}<y,J^1y>+<\sigma ,{\displaystyle \frac{y}{\sqrt{N}}}>)dy,`$
which, together with (2.2), (3.7) and (3.8) yields
$$\mathrm{E}_N(e^{\lambda h})=e^{\lambda /2}\frac{\mathrm{det}J^{\frac{1}{2}}}{(2\pi \lambda )^{N/2}}_{IR^N}\mathrm{exp}(\frac{1}{2\lambda }<y,J^1y>+\underset{i}{}\mathrm{log}\mathrm{cosh}\frac{y_i}{\sqrt{N}})𝑑y$$
(As usual, the square roots appearing in the above formulas are only apparently ill defined: they disappear in the expansion because it contains only the even terms). The above integral can be evaluated by means of standard high-temperature expansion techniques which turn out to be considerably simpler if one assumes that $`_iJ_{ii}=0`$ (see \[PP\]). As we have already noted, this assumption amounts to fix at zero the mean value of the energy. Also, the division by $`N`$ of the argument of the partition function leads to a convergence domain which is increasing as $`N`$ itself. In this way, the asymptotic expression (for $`N\mathrm{}`$) of $`\mathrm{E}_N(e^{\lambda h})`$ can be written in the form
$$\mathrm{E}_N(e^{\lambda h})=e^{\lambda /2}e^{NG(\lambda /N)}\left(1+𝒪(N^1)\right)$$
(3.9)
where the function $`G(x)`$ is an effective specific free energy. For the orthogonal interaction matrix (2.3) one finds \[PP\]:
$$G(x)=\frac{1}{4}\left[\sqrt{1+4x^2}\mathrm{log}\left(\frac{1+\sqrt{1+4x^2}}{2}\right)1\right].$$
(3.10)
It has the following expansion in the vicinity of $`x=0`$:
$$G(x)=\frac{x^2}{4}+𝒪(x^3),$$
(3.11)
which, by the way, coincides with what one obtains for the SK model. This yields
$$e^{NG(\lambda /N)}=1+\frac{\lambda ^2}{4N}+𝒪\left(\frac{\lambda ^3}{N^2}\right).$$
(3.12)
Summarizing, we have found that for any fixed $`\lambda `$,
$$\mathrm{E}_N(e^{\lambda h})=e^{\lambda /2}\left[1+𝒪\left(\frac{1}{N}\right)\right],N\mathrm{}.$$
(3.13)
Using a well known theorem of probability theory \[Si\] which says that a distribution function $`G_N`$ converges weakly to $`G`$ if and only if $`\phi _N(\lambda )\phi (\lambda )`$ for any $`\lambda `$ (where $`\phi _N(\lambda )`$ and $`\phi (\lambda )`$ are the characteristic fcts of $`G_N`$ and $`G`$ respectively) and noting that $`\phi (\lambda )=e^{\lambda /2}`$ is the characteristic function of the distribution fct $`G(x)=\chi _{[\frac{1}{2},\mathrm{})}`$, we then conclude that the distribution of $`h`$ tends to $`\chi _{[\frac{1}{2},\mathrm{})}`$.
## 4 Flipping spins from the ground state and statistics of levels
As already noted in \[MPR\], for special values of $`N`$ the ground state, i.e. the configuration $`\sigma ^0\mathrm{\Sigma }_N`$ which minimizes the energy, can be explicitly constructed. Indeed, for $`N`$ odd such that $`p=2N+1`$ is prime of the form $`4m+3`$, where $`m`$ is an integer, let $`\sigma ^0`$ be the state given by the sequence of Legendre symbols, i.e.
$$\sigma _x^0=\left(\frac{x}{p}\right)=\{\begin{array}{cc}+1,\hfill & \text{if }x=k^2(modp),\hfill \\ 1,\hfill & \text{if }xk^2(modp),\hfill \end{array}$$
(4.14)
with $`k=1,2,\mathrm{},p1`$. Then (see the Appendix):
$$H(\sigma ^0)=\frac{N}{2}$$
(4.15)
A typical ground state for $`p`$ prime of the form $`4m+3`$ reflects the well known random distribution of the Legendre symbols (see Fig. 1, where a pair of ground states are shown for two different $`N`$ values). No structure is present at any scale. Nevertheless, denoting by $`m^0`$ the specific magnetization of the ground state, i.e.
$$m^0=\frac{1}{N}\underset{x=1}{\overset{N}{}}\sigma _x^0=\frac{1}{N}\underset{x=1}{\overset{N}{}}\left(\frac{x}{p}\right),$$
(4.16)
one observes that it tends to be a positive function of $`N`$, fluctuating around the value $`1/\sqrt{N}`$. To let the reader better appreciate this fact we plot in Fig. 2 the total magnetization $`Nm^0`$ versus $`N`$.
For the remaining part of the paper we will restrict to $`p=2N+1`$ prime, with $`p=3(\text{mod}4)`$. We point out that this set has measure zero as a subset of the natural numbers. However, we have strong numerical evidence that some relevant properties that we are going to discuss hereafter, such as the behaviour of $`m^0(N)`$, the statistics of energy levels and the number of metastable states (see below) are somehow generic in $`N`$.
Let $`\mathrm{\Omega }_s\mathrm{\Sigma }_N`$ be the subspace consisting of the $`\left(\genfrac{}{}{0pt}{}{N}{s}\right)`$ configurations obtained by starting from the ground state $`\sigma ^0`$ described above and flipping exactly $`s`$ different spins. Each point of $`\mathrm{\Omega }_s`$ can thus be identified with a $`s`$-dimensional vector $`\tau \{1,\mathrm{},N\}^s`$ of the form $`\tau =(x_1,\mathrm{},x_s)`$, with $`x_ix_j`$ for $`ij`$, which specifies the positions of the flipped spins along the chain of length $`N`$. We then define the ‘flipping’ map $`L_\tau :\mathrm{\Sigma }_N\mathrm{\Sigma }_N`$ as:
$$\left(L_\tau \sigma \right)_x=\{\begin{array}{cc}\sigma _x,\hfill & x\tau ,\hfill \\ \sigma _x,\hfill & x\tau .\hfill \end{array}$$
(4.17)
In this way we can write
$$\mathrm{\Omega }_s=\{L_\tau \sigma ^0\}_\tau .$$
(4.18)
The correspondence $`\tau \sigma `$ given by $`\sigma =L_\tau \sigma ^0`$ is plainly one-to-one. Therefore in the sequel we shall freely identify a state $`\sigma =L_\tau \sigma ^0`$ with the vector $`\tau `$. Alternatively, we can proceed as follows. Define the overlap $`q(\sigma )`$ of a given configuration $`\sigma \mathrm{\Sigma }_N`$ with respect to the ground state $`\sigma ^0`$ as:
$$q(\sigma )=\frac{1}{N}\underset{x=1}{\overset{N}{}}\sigma _x\sigma _x^0,$$
(4.19)
so that $`q(\sigma ^0)=1`$. Then
$$\mathrm{\Omega }_s=\{\sigma \mathrm{\Sigma }_N:q(\sigma )=12\frac{s}{N}\}$$
(4.20)
The following straightforward calculation yields the energy values on the space $`\mathrm{\Omega }_s`$: using the definition of $`L_\tau `$, the symmetry of $`J`$ and the fact that the ground state $`\sigma ^0`$ is an eigenvector of $`J`$ to the eigenvalue $`1`$ we have:
$`H(L_\tau \sigma ^0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{x\tau }{}}{\displaystyle \underset{y\tau }{}}J_{xy}\sigma _x^0\sigma _y^0+{\displaystyle \frac{1}{2}}{\displaystyle \underset{x\tau }{}}{\displaystyle \underset{y\tau }{}}J_{xy}\sigma _x^0\sigma _y^0`$ (4.21)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{x\tau }{}}{\displaystyle \underset{y\tau }{}}J_{xy}\sigma _x^0\sigma _y^0{\displaystyle \frac{1}{2}}{\displaystyle \underset{x\tau }{}}{\displaystyle \underset{y\tau }{}}J_{xy}\sigma _x^0\sigma _y^0`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{x,y=1}{\overset{N}{}}}J_{xy}\sigma _x^0\sigma _y^0+2{\displaystyle \underset{x\tau }{}}{\displaystyle \underset{y\tau }{}}J_{xy}\sigma _x^0\sigma _y^0`$
$`=`$ $`{\displaystyle \frac{N}{2}}+2{\displaystyle \underset{x\tau }{}}{\displaystyle \underset{y=1}{\overset{N}{}}}J_{xy}\sigma _x^0\sigma _y^02{\displaystyle \underset{x\tau }{}}{\displaystyle \underset{y\tau }{}}J_{xy}\sigma _x^0\sigma _y^0`$
$`=`$ $`{\displaystyle \frac{N}{2}}+2{\displaystyle \underset{x\tau }{}}(\sigma _x^0)^22{\displaystyle \underset{x\tau }{}}{\displaystyle \underset{y\tau }{}}J_{xy}\sigma _x^0\sigma _y^0`$
$`=`$ $`{\displaystyle \frac{N}{2}}+2s2{\displaystyle \underset{x\tau }{}}{\displaystyle \underset{y\tau }{}}J_{xy}\sigma _x^0\sigma _y^0.`$
Notice that for the Ising mean-field interaction: $`J_{xy}=1/N`$, one finds
$$H(L_\tau \sigma ^0)=\frac{N}{2}+2s2\frac{s^2}{N}=\frac{N}{2}\left(1\frac{2s}{N}\right)^2.$$
(4.22)
It is now possible to study the distribution of the energy levels on the individual subspaces $`\mathrm{\Omega }_s`$, where $`s=0,1,\mathrm{},N`$. Let $`p_s`$ be the probability distribution restricted on $`\mathrm{\Omega }_s`$, i.e.
$$p_s(C)=\left(\genfrac{}{}{0pt}{}{N}{s}\right)^1\underset{\sigma C\mathrm{\Omega }_s}{}1,$$
(4.23)
and let $`\mathrm{E}_s`$ denote the expectation wrt $`p_s`$. The $`n`$-th moment of the energy $`H`$ on the subspace $`\mathrm{\Omega }_s`$ is given by
$$\mathrm{E}_s(H^n)_{\mathrm{\Omega }_s}H^n(\sigma )𝑑p_s(\sigma )=\left(\genfrac{}{}{0pt}{}{N}{s}\right)^1\underset{\tau \mathrm{\Omega }_s}{}H^n(L_\tau \sigma ^0),$$
(4.24)
so that the $`n`$-th moment $`\mathrm{E}_N(H^n)`$ of the energy on the whole configuration space $`\mathrm{\Sigma }_N`$ is
$$\mathrm{E}_N(H^n)_{\mathrm{\Sigma }_N}H^n(\sigma )𝑑\mathrm{P}_N(\sigma )=\frac{1}{2^N}\underset{s=0}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{N}{s}\right)\mathrm{E}_s(H^n).$$
(4.25)
A tedious but straightforward calculation (see Appendix) yields the following expressions for the first two $`s`$-moments:
$$\mathrm{E}_s(H)=\frac{N}{2}\left(1\frac{2s}{N}\right)^2,$$
(4.26)
$$\sigma _s^2(H)\mathrm{E}_s(H^2)(\mathrm{E}_s(H))^2=\frac{4s(sN)(2s^22sN+N)}{N^2(N2)}$$
(4.27)
and consequently
$$\mathrm{E}_N(H)=\frac{1}{2},\mathrm{E}_N(H^2)(\mathrm{E}_N(H))^2=\frac{N1}{2}.$$
(4.28)
These results indicate that, at variance with the ferromagnetic case where the energy is constant on each subspace $`\mathrm{\Omega }_s`$, here there is a significant overlap between the distributions (for different $`s`$ values) of the energy when restricted to $`\mathrm{\Omega }_s`$. In particular, from the espression of $`\sigma _s^2`$ we see that there can be a large number of states having small overlap with the ground state but nevertheless with energy very close to $`N/2`$. For example we have $`\sigma _{N/2}^2N/2`$, indicating that the energy restricted to the subspace $`\mathrm{\Omega }_{N/2}`$ may fluctuate over the whole energy range. This simple phenomenon is intimately related to the existence of metastable states and it will prove crucial in the understanding of the zero temperature dynamics, as discussed below. Fig. 3 shows the distributions of the energy restricted to various subspaces $`\mathrm{\Omega }_s`$.
Another quantity of interest is the specific magnetization $`m(\sigma )`$ of an arbitrary state $`\sigma \mathrm{\Sigma }_N`$, given by:
$$m(\sigma )=\frac{1}{N}\underset{x=1}{\overset{N}{}}\sigma _x,$$
(4.29)
In particular, given $`\tau =(x_1,\mathrm{},x_s)`$, we have
$$m(L_\tau \sigma ^0)m(\tau )=m^0\frac{2}{N}\underset{x\tau }{}\left(\frac{x}{p}\right).$$
(4.30)
Clearly,
$$\frac{1}{\left(\begin{array}{c}N\\ s\end{array}\right)}\underset{\tau \mathrm{\Omega }_s}{}\left(\underset{x\tau }{}\left(\frac{x}{p}\right)\right)=\frac{1}{\left(\begin{array}{c}N\\ s\end{array}\right)}\underset{x=1}{\overset{N}{}}\left(\begin{array}{c}n1\\ s1\end{array}\right)\left(\frac{x}{p}\right)=sm^0,$$
(4.31)
i.e.
$$E_s(m)=\frac{1}{\left(\begin{array}{c}N\\ s\end{array}\right)}\underset{\tau \mathrm{\Omega }_s}{}m(\tau )=m^0\left(12\frac{s}{N}\right).$$
(4.32)
Moreover, we show in the Appendix that
$$\sigma _s^2(m)E_s(m^2)(E_s(m))^2=\frac{4s(Ns)}{N^3(N1)}+\frac{4s(s1)}{N^3(N1)}\underset{x=1}{\overset{p}{}}\left(\frac{x}{p}\right)d_N(x),$$
(4.33)
where $`d_N(x)`$ is the integer valued function giving the number of elements $`u\{1,\mathrm{},N\}`$ such that $`u^1x\{1,\mathrm{},N\}`$ (where $`u^1`$ denotes the inverse $`\text{mod }p`$ of $`u`$). As will be discussed in the Appendix, $`d_N(x)`$ takes values around $`p/4`$ with rather small fluctuations. Since $`_{x=1}^p\left(\frac{x}{p}\right)=0`$, the last term in (4.33) can beconsidered as a small correction to the constant value $`4s(Ns)/N^3(N1)`$.
## 5 Zero temperature dynamics and metastable states.
We first introduce the following discrete $`1`$flip dynamics, given by:
$$\sigma (t+1)=\{\begin{array}{c}L_{\omega (t)}\sigma (t),\text{if }H(L_\omega \sigma )<H(\sigma )\text{,}\hfill \\ \sigma (t),\text{otherwhise,}\hfill \end{array}$$
where, for each $`t`$, $`\omega (t)`$ is chosen randomly in $`\{1,\mathrm{},N\}`$ with uniform distribution. Choosing an initial condition $`\sigma (0)`$ at random with respect to $`\mathrm{P}_N`$, one obtains a random orbit $`\{\sigma (0),\sigma (1),\mathrm{},\sigma (\mathrm{})\}`$ for any realization $`\{\omega (t)\}_{1t\mathrm{}}`$ of length $`\mathrm{}`$. As a consequence of the previous analysis, we have the following remarks.
* On one hand, it may happen that starting from $`\sigma (0)`$ one reaches after $`t`$ iterations a state $`\sigma (t)\mathrm{\Omega }_s`$, of the form $`\sigma (t)=L_\tau \sigma ^0`$ for some $`\tau =(x_1,\mathrm{},x_s)`$, such that $`H(L_\tau \sigma ^0)<H(L_\omega L_\tau \sigma ^0)`$ for any $`\omega \{x_1,\mathrm{},x_s\}`$.
* On the other hand one can reach $`\sigma (t)\mathrm{\Omega }_s`$ such that for some $`\omega \{1,\mathrm{},N\}`$, $`L_\omega \sigma \mathrm{\Omega }_{s+1}`$ and $`H(L_\omega \sigma (t))<H(\sigma (t))`$.
Due to the above observations, the overlap function $`q(\sigma (t))`$ is in general not monotonically non-decreasing along a given random orbit (this at variance with the Ising mean field model). In particular there might be metastable states \[NS\]. Now, given $`\omega \{1,\mathrm{},N\}`$, we shall say that a configuration $`\sigma \mathrm{\Sigma }_N`$ is $`\omega `$-stable if
$$H(L_\omega \sigma )>H(\sigma ).$$
(5.34)
Moreover, we say that $`\sigma `$ is $`1`$-flip stable (or metastable) if it is $`\omega `$-stable $`\omega \{1,\mathrm{},N\}`$.
We denote by $`n(N)`$ the total number of such metastable states as a function of $`N`$. From (4.21) one readily obtains that
$$H(L_\omega \sigma )=H(\sigma )+\mathrm{\hspace{0.17em}2}\underset{x\omega }{}J_{x\omega }\sigma _x\sigma _\omega ,$$
(5.35)
so that $`\sigma `$ is $`\omega `$-stable if and only if
$$\left(J\sigma \right)_\omega \sigma _\omega >J_{\omega \omega }.$$
(5.36)
Summing over $`\omega `$ and using (2.4) we see that if $`\sigma `$ is $`1`$-flip stable then
$$<J\sigma ,\sigma >>\mathrm{\hspace{0.33em}1}.$$
(5.37)
Recalling the expression (2.2) of the Hamiltonian we see that a necessary condition for $`\sigma \mathrm{\Sigma }_N`$ to be $`1`$-flip stable is that
$$H(\sigma )<\mathrm{E}_N(H)=\frac{1}{2}$$
(5.38)
The main goal of this paper is to give an estimate of the number of metastable states for any given $`N`$. To this end we first performed some numerical investigations. For $`N30`$ we performed an exact enumeration of all configurations, whereas for larger $`N`$ we run the zero temperature dynamics described above (“deep quench”) for a number or realizations $`\{\omega (t)\}`$ as large as $`10^8`$ for bigger sizes, keeping track of the metastable states. As shown in Fig. 4, the growth of these states is exponential for generic values of the $`N`$.
The best numerical fit yields
$$n(N)Ce^{0.28N}.$$
(5.39)
We remark that the same behaviour has been observed in \[PP\] for the Random Orthogonal Model; for spin glasses see \[PP2\]. We now proceed to give a partial justification of this result by means of probabilistic arguments.
Let $`\tau =(x_1,\mathrm{},x_s)`$ and $`\omega \{1,\mathrm{},N\}`$ be given. Using (5.35) and $`J\sigma ^0=\sigma ^0`$, it is easy to see that if $`\omega \tau `$ (i.e. $`L_\omega \sigma \mathrm{\Omega }_{s+1}`$)
$$H(L_\omega \sigma )=H(\sigma )+\mathrm{\hspace{0.17em}2}\left(1J_{\omega \omega }2\underset{x\tau }{}J_{x\omega }\sigma _x^0\sigma _\omega ^0\right),$$
whereas, if $`\omega \tau `$ (i.e. $`L_\omega \sigma \mathrm{\Omega }_{s1}`$), we have
$$H(L_\omega \sigma )=H(\sigma )\mathrm{\hspace{0.17em}2}\left(1+J_{\omega \omega }2\underset{x\tau }{}J_{x\omega }\sigma _x^0\sigma _\omega ^0\right).$$
If we define
$$h(\tau ,\omega )=\underset{x\tau }{}J_{x\omega }\sigma _x^0\sigma _\omega ^0,$$
(5.40)
we then see that a configuration $`\sigma =L_\tau \sigma ^0`$ is $`\omega `$-stable if and only if
$$\{\begin{array}{ccc}h(\tau ,\omega )\hfill & <\hfill & \frac{1}{2}(1J_{\omega \omega }),\text{if }\omega \tau ,\hfill \\ & & \\ h(\tau ,\omega )\hfill & >\hfill & \frac{1}{2}(1+J_{\omega \omega }),\text{if }\omega \tau .\hfill \end{array}$$
(5.41)
We now dwell upon the problem of characterising the behaviour of the function $`h(\tau ,\omega )`$ so as the condition (5.41) can be effectively used to estimate the number of metasable states. Let us rewrite $`h(\tau ,\omega )`$ in the form
$$h(\tau ,\omega )=\frac{2}{\sqrt{p}}\underset{x\tau }{}\xi _x(\omega ),$$
(5.42)
where
$$\xi _x(\omega ):=\left(\frac{\omega x}{p}\right)\mathrm{sin}\left(\frac{2\pi \omega x}{p}\right).$$
(5.43)
Now, having fixed $`\tau `$ and $`x\tau `$, we can view the function $`\xi _x(\omega )`$ defined in (5.43) as a random variable uniformly distributed on $`\{1,\mathrm{},N\}`$ and taking values in $`[1,1]`$. Its mean $`\mu _x`$ and variance $`\sigma _x^2`$ are easily computed:
$$\mu _x=\frac{1}{N}\underset{\omega =1}{\overset{N}{}}\left(\frac{\omega x}{p}\right)\mathrm{sin}\left(\frac{2\pi \omega x}{p}\right)=\frac{\sqrt{p}}{2N},$$
(5.44)
and, using (2.4),
$$\sigma _x^2=\frac{1}{N}\underset{\omega =1}{\overset{N}{}}\mathrm{sin}^2\left(\frac{2\pi \omega x}{p}\right)\mu _x^2=\frac{p}{4N}\left(1\frac{1}{N}\right).$$
(5.45)
Here we want to study the behaviour of the sum $`\eta (\tau ,\omega ):=_{x\tau }\xi _x(\omega )`$. We remark that this sum, and thus $`h(\tau ,\omega )=2\eta (\tau ,\omega )/\sqrt{p}`$, has to be regarded as a r.v. defined on the product of two probability spaces: for each fixed $`\tau `$, it is the sum of the i.i.d.r.v.’s $`\xi _x(\omega )`$ on the space $`\{1,\mathrm{},N\}`$ with uniform distribution (this comes from the very definition of the zero temperature dynamics); on the other hand, for each fixed $`\omega `$, it can be regarded as a r.v. on $`\mathrm{\Omega }_s`$ viewed as a probability space endowed with the distribution $`p_s`$. Its mean is given by (recall that the symbol $`\mathrm{E}_s`$ denotes the expectation wrt $`p_s`$):
$`\mathrm{E}_s\eta `$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{N}{s}}\right)^1{\displaystyle \underset{\tau =(x_1,\mathrm{},x_s)}{}}{\displaystyle \underset{x\tau }{}}\xi _x(\omega )=\left({\displaystyle \genfrac{}{}{0pt}{}{N}{s}}\right)^1{\displaystyle \underset{x=1}{\overset{N}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{N1}{s1}}\right)\xi _x(\omega )`$
$`=`$ $`{\displaystyle \frac{s}{N}}{\displaystyle \underset{x=1}{\overset{N}{}}}\xi _x(\omega )={\displaystyle \frac{s\sqrt{p}}{2N}}.`$
which does not depend on $`\omega `$ and equals $`s`$ times $`\mu _x`$. Along the same lines one shows that
$$\mathrm{Var}_s\eta =\frac{sp}{4N}\left(1\frac{s}{N}\right).$$
(5.47)
Notice that unlike the means, here we have $`\mathrm{Var}_s\eta s\sigma _x^2`$. This discrepancy comes from the fact that, for any fixed $`\omega `$, the sequence $`\xi _{x_1},\xi _{x_2},\mathrm{},\xi _{x_s}`$ is a sequence of distinct (and ordered) elements so that by no means we can view $`\eta `$ as a sum of independent and identically distributed objects. Nonetheless, also supported by strong numerical evidence (see Fig. 5), we claim that a version of the central limit theorem is applicable so that when $`N\mathrm{}`$, $`s\mathrm{}`$ with $`s/N\lambda `$, we have
$$p_s\left(\alpha <\frac{\eta \mathrm{E}_s\eta }{\sqrt{\mathrm{Var}_s\eta }}<\beta \right)\frac{1}{\sqrt{2\pi }}_\alpha ^\beta e^{y^2/2}𝑑y.$$
(5.48)
Assuming the validity of (5.48), performing the change of variables $`y=(x\lambda )/\sqrt{\gamma }`$, with $`\gamma =\lambda (1\lambda )`$, and setting $`a=\lambda +\alpha \sqrt{\gamma }`$, $`b=\lambda +\beta \sqrt{\gamma }`$, we thus obtain an asymptotic gaussian distribution for $`h(,\omega )`$:
$$p_s\left(a<h(\tau ,\omega )<b\right)\frac{1}{\sqrt{2\pi \gamma }}_a^be^{(x\lambda )^2/2\gamma }𝑑x.$$
(5.49)
Note that the r.h.s. does not depend on $`\omega `$. One can actually say more: for any $`\stackrel{~}{\omega },\omega \{1,\mathrm{},N\}`$ we have $`h(\stackrel{~}{\omega }^1\tau ,\stackrel{~}{\omega }\omega )=h(\tau ,\omega )`$. Therefore the set of values of $`h(,\omega )`$ on $`\mathrm{\Omega }_s`$ does not depend on the choice of $`\omega `$, i.e. $`\left\{h(\tau ,\omega )\right\}_{\tau \mathrm{\Omega }_s}=\left\{h(\tau ^{},\omega ^{})\right\}_{\tau ^{}\mathrm{\Omega }_s}`$, for all $`\omega `$, $`\omega ^{}`$.
Having fixed an order for the lattice points $`(\omega _1,\mathrm{},\omega _N)`$, $`\omega _j\omega _k`$, we now consider the following quantities:
$$\pi _{s,N}(\omega _k)=p_s\left(\{\sigma \mathrm{\Omega }_s:H(L_{\omega _k}\sigma )>H(\sigma )\}\right),$$
(5.50)
the $`p_s`$-probability that a randomly chosen state $`\sigma \mathrm{\Omega }_s`$ is $`\omega _k`$-stable,
$`\pi _{s,N}(\omega _{k+1}|\omega _1,\mathrm{},\omega _k)`$ $`=`$ $`p_s(\{\sigma \mathrm{\Omega }_s:H(L_{\omega _{k+1}}\sigma )>H(\sigma )`$ (5.51)
$`|H(L_{\omega _j}\sigma )>H(\sigma ),\omega _j=\omega _1,\mathrm{}\omega _k\}),`$
the conditional $`p_s`$-probability that a randomly chosen state $`\sigma \mathrm{\Omega }_s`$ is $`\omega _{k+1}`$-stable given that it is $`\omega _j`$-stable for $`j=1,\mathrm{},k`$, and
$$\pi _{s,N}=p_s\left(\{\sigma \mathrm{\Omega }_s:H(L_\omega \sigma )>H(\sigma ),\omega =\omega _1,\mathrm{},\omega _N\}\right),$$
(5.52)
the $`p_s`$-probability that a randomly chosen state $`\sigma \mathrm{\Omega }_s`$ is $`1`$-flip stable (i.e. stable for all possible flipping). Notice that by condition (5.41) the last quantity can be written as
$$\pi _{s,N}=\left(\genfrac{}{}{0pt}{}{N}{s}\right)^1\underset{\genfrac{}{}{0pt}{}{(x_1,\mathrm{},x_s)=\tau }{(y_1,\mathrm{},y_{Ns})=\tau ^c}}{}\underset{i=1}{\overset{s}{}}\theta \left(h(\tau ,x_i)>\frac{1+J_{x_ix_i}}{2}\right)\underset{j=1}{\overset{Ns}{}}\theta \left(h(\tau ,y_j)<\frac{1J_{y_iy_i}}{2}\right).$$
(5.53)
The three quantities introduced above are related by the following identity:
$$\pi _{s,N}=\pi _s(\omega _1)\pi _s(\omega _2|\omega _1)\pi _s(\omega _3|\omega _1,\omega _2)\mathrm{}\pi _s(\omega _N|\omega _1,\mathrm{},\omega _{N1})$$
(5.54)
and the total number of $`1`$-flip stable states in $`\mathrm{\Sigma }_N`$ is, by definition,
$$n(N)=\underset{s=0}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{N}{s}\right)\pi _{s,N}.$$
(5.55)
We shall study the quantity $`n(N)`$ in the thermodynamic limit: $`N\mathrm{}`$, $`s\mathrm{}`$, with $`s/N\lambda `$ and $`0<\lambda <1`$. In this regime we write $`\pi _{s,N}\pi _\lambda `$ and apply Stirling’s formula to obtain
$$\left(\genfrac{}{}{0pt}{}{N}{s}\right)\frac{e^{NF(\lambda )}}{\sqrt{2\pi N\lambda (1\lambda )}},\text{with}F(\lambda )=\lambda \mathrm{log}\lambda (1\lambda )\mathrm{log}(1\lambda ).$$
(5.56)
Note that $`F(\lambda )`$ is concave and symmetric around $`\lambda =1/2`$, with $`F(1/2)=\mathrm{log}2`$. In this way we get for $`N`$ large and $`s\lambda N`$ with $`\lambda `$ ranging in the unit interval,
$$n(N)_0^1\sqrt{\frac{N}{2\pi \lambda (1\lambda )}}\mathrm{exp}\left[N\left(F(\lambda )+\frac{\mathrm{log}\pi _\lambda }{N}\right)\right]𝑑\lambda .$$
(5.57)
It thus remains to estimate the probability $`\pi _\lambda `$. Let us consider first the unconditioned probability (5.50). According to (5.41) and the total probability formula we have:
$`\pi _{s,N}(\omega _k)`$ $`=`$ $`{\displaystyle \frac{s}{N}}p_s\left(h(\tau ,\omega _k)>{\displaystyle \frac{1+J_{\omega _k\omega _k}}{2}}|\tau \omega _k\right)`$ (5.58)
$`+`$ $`\left(1{\displaystyle \frac{s}{N}}\right)p_s\left(h(\tau ,\omega _k)<{\displaystyle \frac{1J_{\omega _k\omega _k}}{2}}|\tau \ni ̸\omega _k\right).`$
Here $`s/N`$ and $`1s/N`$ are the probabilities that $`\tau \omega _k`$ and $`\tau \ni ̸\omega _k`$, respectively. We can easily compute the conditional expectations
$$\mathrm{E}(h|\tau \omega _k)=\left(\genfrac{}{}{0pt}{}{N1}{s1}\right)^1\frac{2}{\sqrt{p}}\underset{\tau \omega _k}{}\underset{x\tau }{}\xi _x(\omega _k)=\frac{s1}{N1}\left(\frac{sN}{N1}\right)J_{\omega _k\omega _k},$$
(5.59)
and
$$\mathrm{E}(h|\tau \ni ̸\omega _k)=\left(\genfrac{}{}{0pt}{}{N1}{s}\right)^1\frac{2}{\sqrt{p}}\underset{\tau \ni ̸\omega _k}{}\underset{x\tau }{}\xi _x(\omega _k)=\frac{s}{N1}\left(1J_{\omega _k\omega _k}\right).$$
(5.60)
In a similar way one can compute the variances $`\gamma _+(\omega _k)`$ and $`\gamma _{}(\omega _k)`$ conditioned to the events $`\{\tau \omega _k\}`$ and $`\{\tau \ni ̸\omega _k\}`$. For $`N`$ large and $`s\lambda N`$, retaining only terms $`𝒪(1)`$, one gets
$$\mathrm{E}(h|\tau \omega _k)\mathrm{E}(h|\tau \ni ̸\omega _k)\lambda ,\gamma _{}(\omega _k)\gamma _+(\omega _k)\gamma =\lambda (1\lambda ).$$
(5.61)
Moreover in the thermodynamic limit specified above we write $`\pi _{s,N}(\omega _k)\pi _\lambda (\omega _k)`$ and argue from (5.49) the following approximate expression for $`\pi _\lambda (\omega _k)`$:
$`\pi _\lambda (\omega _k)`$ $``$ $`{\displaystyle \frac{\lambda }{\sqrt{2\pi \gamma }}}{\displaystyle _{1/2}^{\mathrm{}}}e^{(x\lambda )^2/2\gamma }𝑑x+{\displaystyle \frac{(1\lambda )}{\sqrt{2\pi \gamma }}}{\displaystyle _{\mathrm{}}^{1/2}}e^{(x\lambda )^2/2\gamma }𝑑x`$ (5.62)
$`=`$ $`{\displaystyle \frac{1}{2}}+\left({\displaystyle \frac{1}{2}}\lambda \right)\mathrm{erf}\left({\displaystyle \frac{\frac{1}{2}\lambda }{\sqrt{2\gamma }}}\right)`$
where we have denoted the error function by
$$\mathrm{erf}(z)=\frac{2}{\sqrt{\pi }}_0^ze^{x^2}𝑑x.$$
(5.63)
It is not difficult to check that the r.h.s. of (5.62) is convex and symmetrix around $`\lambda =1/2`$, where it reaches its minimum value equal to $`1/2`$.
Let us now come to $`\pi _{s,N}`$. In principle this quantity is to be computed by specifying the whole set of constraints embodied in (5.53) or, which is the same, by computing the conditional probabilities appearing in eq.(5.54). However, this appears to be a difficult task. A first approach which drastically simplifies this task is to forget about the constraints implied by (5.53) and assume that (in the thermodynamic limit) the various $`\omega `$-stability conditions become mutually independent, that is $`\pi _\lambda (\omega _{k+1}|\omega _1,\mathrm{},\omega _k)=\pi _\lambda (\omega _{k+1})`$, for all $`k=1,\mathrm{},N1`$, so that
$$\pi _\lambda =\underset{k=1}{\overset{N}{}}\pi _\lambda (\omega _k).$$
(5.64)
Recalling eq. (5.57) and (5.62), one is led to the following expression for $`n(N)`$:
$$n(N)_0^1\sqrt{\frac{N}{2\pi \lambda (1\lambda )}}\mathrm{exp}\left(NG_1(\lambda )\right)𝑑\lambda ,$$
(5.65)
where
$$G_1(\lambda )=F(\lambda )+\mathrm{log}\frac{1}{2}+\left(\frac{1}{2}\lambda \right)\mathrm{erf}\left(\frac{\frac{1}{2}\lambda }{\sqrt{2\gamma }}\right)$$
(5.66)
We show the shape of the the function $`G_1(\lambda )`$ in Fig. 6.
Evaluating the integral (5.65) with the saddle-point method one gets
$$n(N)Ce^{0.14N}.$$
(5.67)
Notice that the exponent is the half of what is observed numerically (cfr (5.39)). In the remaining part of the paper we shall argue that (5.67) is indeed an estimate from below of the actual number of metastable states. The above discussion has been able to reproduce the exponential growth of the number of metastable states with the size of the system. To understand the discrepacy between the estimated exponent and the one measured numerically one should note that the nature of the interaction makes the conditional probabilities play a major role in the asymptotics of the number of metastable states. To be more precise, our approximation which assumes mutually independent individual $`\omega `$-stability events, i.e. $`\pi _\lambda (\omega _{k+1}|\omega _1,\mathrm{},\omega _k)\pi _\lambda (\omega _{k+1})`$, is actually reasonable only for small value of $`k`$ (this can be checked, for example, calculating the correlation functions). As numerical results shows, for large values of $`k`$ the specific form of the interactions make these events strongly dependent. In Fig. 7 we show the function $`P(k)`$ providing the average of $`\pi _\lambda (\omega _{k+1}|\omega _1,\mathrm{},\omega _k)`$ over a large sample of different permutations $`(\omega _1,\mathrm{},\omega _N)`$ of the lattice points. The conditional probabilities $`P(k)`$ grow monotonically, almost linearly, from the initial (unconditioned) value up to a number close to $`1`$. In other words, requiring that a large number $`k`$ of spins produce an $`\omega `$-stable state increases substantially the probability of doing the same for the remaining spins.
Another way of understanding the constructive effect of the correlations is the following. Consider again the function $`h(\tau ,\omega )`$. Having fixed $`\omega _{k+1}`$, we have already noticed that for $`s`$ and $`N`$ large enough the values of $`h(\tau ,\omega _{k+1})`$ with $`\tau \mathrm{\Omega }_s`$ are approximately distributed according to a gaussian probability density with mean $`\lambda =s/N`$ and variance $`\gamma =\lambda (1\lambda )`$, regardless of the particular value of $`\omega _{k+1}`$. Thus, in particular, the same distribution are expected to arise if one considers the values of $`h(\tau ,\omega _{k+1})`$ constrained to the subsets of configurations such that $`\omega _{k+1}\tau `$, or $`\omega _{k+1}\tau `$. On the other hand, if one picks $`\omega _1,\omega _2,\mathrm{},\omega _k`$ with $`\omega _j\omega _{k+1}`$, $`j=1,\mathrm{},k`$, and computes numerically the two conditional distributions of the values of $`h(\tau ,\omega _{k+1})`$ given $`\omega _1`$-stability, … , $`\omega _k`$-stablity (again with the constraints $`\omega _{k+1}\tau `$ or $`\omega _{k+1}\tau `$), one finds that their means move to opposite directions, thus increasing the probability of $`\omega _{k+1}`$-stability (see (5.58)). This is shown in Fig. 8, where a system of size $`N=21`$ and $`s=10`$ is considered. The two central distributions correspond to the unconditioned cases, namely the values of $`h(\tau ,\omega )`$ for $`\tau \mathrm{\Omega }_{10}`$ with the only constraints $`\omega \tau `$ or $`\omega \tau `$, respectively. Considering instead the values taken by $`h`$ on the states $`\tau `$ which, besides the constraints specified above, are stable with respect the first $`10`$ spins, one finds two distributions whose mean values have moved towards opposite directions. An averaged over $`\omega `$ has been performed.
## 6 Conclusion
We have investigated the statistical properties of energy levels and metastables states for a class of deterministic models, the most representative being the sine model\[MPR\], which have attracted much attention in recent years for their glassy behaviour despite the non-random nature of the interaction. We have pushed further on the analogy with glassy systems, proving a number of properties typical of disordered spin models. In particular, using number theoretic methods, we have described the energy (equivalently, free energy at $`T=0`$) landscape as a function of configurations with a fixed overlap with the ground state. The analysis revealed the existence of states very different from the ground state but with energy arbitrarily close to it: this corresponds to the “chaoticity” property of spin-glasses systems, well established in long range models. More importantly, some of these states can be local energy minima (equivalently, $`1`$-flip stable at $`T=0`$). They are expected to have a significant weight on the partition function in the low-temperature region, giving rise to the non-equibrium behaviour observed in annealing Montecarlo experiments. We have been able to estimate the approximate number of these energy minima. The analytic computations, combined with the numerical findings, strongly support the conclusion that the bound (5.67) estimates from below the number of metastable states $`n(N)`$, proving their exponential increase with the size of the system.
A number of basic questions about metastability arises now in a natural way, such as computing the energy density distributions of metastables states, studying energy barriers among them and their attraction basins. Stability of configurations with respect to the flip of an arbitrary number of spins is an interesting question as well. These problems are currently under investigation using the approach developed in this paper and will be addressed elsewhere.
## Appendix
* (Proof of (4.15)) Choose $`N`$ odd such that $`p=2N+1`$ is prime of the form $`4m+3`$, where $`m`$ is an integer. Denote by $`\sigma ^0`$ the spin configuration given by the sequence of Legendre symbols, i.e.
$$\sigma _x^0=\left(\frac{x}{p}\right)=\{\begin{array}{cc}+1,\hfill & \text{if }x=k^2(modp),\hfill \\ 1,\hfill & \text{if }xk^2(modp),\hfill \end{array}$$
with $`k=1,2,\mathrm{},p1`$. Let us show that $`\sigma ^0=(\sigma _1^0\mathrm{}\sigma _N^0)`$ is the ground state for the sine model or, which is the same, that $`\sigma ^0`$ is an eigenvector of $`J`$ with eigenvalue $`1`$. For basic results of number theory used in the proof see, for example, ref. \[Ap\].
$`(J\sigma ^0)_y`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{p}}}{\displaystyle \underset{x=1}{\overset{N}{}}}\mathrm{sin}\left({\displaystyle \frac{2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)`$
$`=`$ $`{\displaystyle \frac{2}{\sqrt{p}}}{\displaystyle \underset{x=1}{\overset{N}{}}}{\displaystyle \frac{1}{2i}}\left(\mathrm{exp}\left({\displaystyle \frac{i2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)\mathrm{exp}\left({\displaystyle \frac{i2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)\right)`$
changing $`xx`$ in the second summation
$`=`$ $`{\displaystyle \frac{1}{i\sqrt{p}}}\left[{\displaystyle \underset{x=1}{\overset{N}{}}}\mathrm{exp}\left({\displaystyle \frac{i2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right){\displaystyle \underset{x=N}{\overset{1}{}}}\mathrm{exp}\left({\displaystyle \frac{i2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)\right]`$
$`\text{using multiplicativity of Legendre symbols:}\left({\displaystyle \frac{x}{p}}\right)=\left({\displaystyle \frac{1}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)`$
and the fact that $`\left(\frac{1}{p}\right)=1`$ if $`p=3(mod4)`$
$`=`$ $`{\displaystyle \frac{1}{i\sqrt{p}}}\left[{\displaystyle \underset{x=1}{\overset{N}{}}}\mathrm{exp}\left({\displaystyle \frac{i2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)+{\displaystyle \underset{x=N}{\overset{1}{}}}\mathrm{exp}\left({\displaystyle \frac{i2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)\right]`$
$`\text{using the periodicity of Legendre’ symbols}:\left({\displaystyle \frac{x+p}{p}}\right)=\left({\displaystyle \frac{x}{p}}\right)`$
$`=`$ $`{\displaystyle \frac{1}{i\sqrt{p}}}\left[{\displaystyle \underset{x=1}{\overset{N}{}}}\mathrm{exp}\left({\displaystyle \frac{i2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)+{\displaystyle \underset{x=N+1}{\overset{2N}{}}}\mathrm{exp}\left({\displaystyle \frac{i2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)\right]`$
being $`\left(\frac{p}{p}\right)=0`$ by definition
$`=`$ $`{\displaystyle \frac{1}{i\sqrt{p}}}{\displaystyle \underset{x=1}{\overset{p}{}}}\mathrm{exp}\left({\displaystyle \frac{i2\pi xy}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)`$
using the separability for Gauss sums
$`=`$ $`{\displaystyle \frac{1}{i\sqrt{p}}}\left({\displaystyle \frac{y}{p}}\right){\displaystyle \underset{x=1}{\overset{p}{}}}\mathrm{exp}\left({\displaystyle \frac{i2\pi x}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)`$
evaluating the Gauss sum
$`=`$ $`{\displaystyle \frac{1}{i\sqrt{p}}}\left({\displaystyle \frac{y}{p}}\right)i\sqrt{p}=\sigma _y^0,`$
which is the desired property.
* (Proof of (4.26) and (4.27)) We sketch the basic steps of the calculation. Set
$$\alpha =\frac{N}{2}+2sc_j=\left(\begin{array}{c}Nj\\ sj\end{array}\right).$$
We then have
$$\mathrm{E}_s(H^n)=\frac{1}{c_0}\underset{\tau \mathrm{\Omega }_s}{}\left(\alpha 2\underset{x\tau }{}\underset{y\tau }{}J_{xy}\right)^n=\frac{1}{c_0}\underset{\tau \mathrm{\Omega }_s}{}\underset{k=0}{\overset{n}{}}\left(\begin{array}{c}n\\ k\end{array}\right)\alpha ^{nk}\left(2\underset{x\tau }{}\underset{y\tau }{}J_{xy}\right)^k.$$
For $`n=1`$ we get
$`{\displaystyle \underset{\tau \mathrm{\Omega }_s}{}}\left({\displaystyle \underset{x\omega }{}}{\displaystyle \underset{y\omega }{}}J_{xy}\right)`$ $`=`$ $`{\displaystyle \underset{x=1}{\overset{N}{}}}{\displaystyle \underset{yx,y=1}{\overset{N}{}}}c_2J_{xy}\sigma _x\sigma _y+{\displaystyle \underset{x=1}{\overset{N}{}}}c_1J_{xx}`$
$`=`$ $`{\displaystyle \underset{x=1}{\overset{N}{}}}{\displaystyle \underset{y=1}{\overset{N}{}}}c_2J_{xy}\sigma _x\sigma _y+{\displaystyle \underset{x=1}{\overset{N}{}}}(c_1c_2)J_{xx}=c_2(N1)+c_1,`$
whereas for $`n=2`$ we have
$`{\displaystyle \underset{\omega \mathrm{\Omega }_s}{}}\left({\displaystyle \underset{x\omega }{}}{\displaystyle \underset{y\omega }{}}J_{xy}\right)^2={\displaystyle \underset{z=1}{\overset{N}{}}}{\displaystyle \underset{uz;u=1}{\overset{N}{}}}J_{zu}\sigma _z\sigma _u[{\displaystyle \underset{xz,u;x=1}{\overset{N}{}}}{\displaystyle \underset{yz,u,x;y=1}{\overset{N}{}}}c_4J_{xy}\sigma _x\sigma _y`$
$`+`$ $`{\displaystyle \underset{yz,u;y=1}{\overset{N}{}}}c_3J_{zy}\sigma _z\sigma _y+{\displaystyle \underset{yz,u;y=1}{\overset{N}{}}}c_3J_{uy}\sigma _u\sigma _y+{\displaystyle \underset{xz,u;x=1}{\overset{N}{}}}c_3J_{xz}\sigma _x\sigma _z+{\displaystyle \underset{xz,u;x=1}{\overset{N}{}}}c_3J_{xu}\sigma _x\sigma _u`$
$`+`$ $`{\displaystyle \underset{xz,u;x=1}{\overset{N}{}}}c_3J_{xx}+c_2J_{zu}\sigma _z\sigma _u+c_2J_{zz}+c_2J_{uu}]+{\displaystyle }_{z=1}^NJ_{zz}[{\displaystyle \underset{xz;x=1}{\overset{N}{}}}{\displaystyle \underset{yz,x;y=1}{\overset{N}{}}}c_3J_{xy}\sigma _x\sigma _y`$
$`+`$ $`{\displaystyle \underset{xz;x=1}{\overset{N}{}}}c_2J_{xz}\sigma _x\sigma _z+{\displaystyle \underset{yz;y=1}{\overset{N}{}}}c_2J_{yz}\sigma _y\sigma _z+{\displaystyle \underset{xz;x=1}{\overset{N}{}}}c_2J_{xx}+c_1J_{zz}]`$
$`=`$ $`c_4(N1)(N3)+2c_3(N1)+2c_2(N1)+c_1,`$
which easily give the desired identities.
* (Proof of (4.33)) Let us first extend everything to the set $`\{1,2,\mathrm{},p1\}`$ which, $`p`$ being prime, is a number field. Here we can exploit the multiplicative structure of the field and of the ‘extended ground state’ $`\sigma _x^0=\left(\frac{x}{p}\right)`$, $`x=1,\mathrm{},q`$, with $`q=p1`$. With slight abuse of notation we shall use the same symbols $`\mathrm{\Omega }_s`$, $`p_s`$ and $`\mathrm{E}_s`$ to denote the corresponding extended quantities. It is immediate to see that $`m(\sigma ^0)=0`$ and
$$E_s(m)=\frac{1}{\left(\begin{array}{c}q\\ s\end{array}\right)}\underset{\tau \mathrm{\Omega }_s}{}\left(\frac{2}{q}\underset{x\tau }{}\left(\frac{x}{p}\right)\right)=0.$$
In order to calculate the second moment we consider
$`{\displaystyle \frac{1}{\left(\begin{array}{c}q\\ s\end{array}\right)}}{\displaystyle \underset{\tau \mathrm{\Omega }_s}{}}\left({\displaystyle \underset{x\tau }{}}\left({\displaystyle \frac{x}{p}}\right)\right)^2`$ $`=`$ $`{\displaystyle \frac{1}{\left(\begin{array}{c}q\\ s\end{array}\right)}}{\displaystyle \underset{\tau \mathrm{\Omega }_s}{}}\left({\displaystyle \underset{x\tau ^2}{}}\left({\displaystyle \frac{x}{p}}\right)\right)`$
$`=`$ $`{\displaystyle \frac{1}{\left(\begin{array}{c}q\\ s\end{array}\right)}}{\displaystyle \underset{x=1}{\overset{q}{}}}c_p(x)\left({\displaystyle \frac{x}{p}}\right),`$
where, for any given $`\tau \mathrm{\Omega }_s`$, $`\tau ^2`$ is the collection, with multiplicity, of all possible products $`x_jx_i`$, $`x_i,x_j\tau `$ (all the operations are $`\text{mod }p`$). For example, if $`\tau =\{x_1,x_2,x_3\}`$ then $`\tau ^2=\{x_1^2,x_2^2,x_3^2,x_1x_2,x_1x_3,x_2x_1,x_3x_1,x_2x_3,x_3x_2\}`$. Also, for any given $`x\{1,\mathrm{},q\}`$,
$$c_p(x)=\underset{\tau \mathrm{\Omega }_s}{}\{\text{number of times }x\tau ^2\}=\underset{u=1}{\overset{q}{}}\mathrm{}\{\tau |u\tau \text{ and }u^1x\tau \}.$$
In particular, if $`\left(\frac{x}{p}\right)=1`$ then $`uu^1x`$, $`u=1,\mathrm{},q`$, therefore
$$c_p(x)=\underset{u=1}{\overset{q}{}}\left(\begin{array}{c}q2\\ s2\end{array}\right)=q\left(\begin{array}{c}q2\\ s2\end{array}\right).$$
If instead $`\left(\frac{x}{p}\right)=1`$, then there exists $`\overline{u}`$ such that $`\overline{u}^2=x`$, i.e.
$$c_p(x)=\underset{u\pm \overline{u}}{}\left(\begin{array}{c}q2\\ s2\end{array}\right)+2\left(\begin{array}{c}q1\\ s1\end{array}\right)=(q2)\left(\begin{array}{c}q2\\ s2\end{array}\right)+2\left(\begin{array}{c}q1\\ s1\end{array}\right).$$
Putting everything together we get the following expression for the variance $`\sigma _s^2(m)`$:
$`\sigma _s^2(m)`$ $`=`$ $`{\displaystyle \frac{4}{q^2\left(\begin{array}{c}q\\ s\end{array}\right)}}{\displaystyle \underset{\tau \mathrm{\Omega }_s}{}}\left({\displaystyle \underset{x\tau }{}}\left({\displaystyle \frac{x}{p}}\right)\right)^2`$
$`=`$ $`{\displaystyle \frac{4}{q^2\left(\begin{array}{c}q\\ s\end{array}\right)}}\left[{\displaystyle \frac{q^2}{2}}\left(\begin{array}{c}q2\\ s2\end{array}\right)+{\displaystyle \frac{q}{2}}(q2)\left(\begin{array}{c}q2\\ s2\end{array}\right)+q\left(\begin{array}{c}q1\\ s1\end{array}\right)\right]`$
$`=`$ $`{\displaystyle \frac{4}{q\left(\begin{array}{c}q\\ s\end{array}\right)}}\left[\left(\begin{array}{c}q1\\ s1\end{array}\right)\left(\begin{array}{c}q2\\ s2\end{array}\right)\right]={\displaystyle \frac{4s(qs)}{q^2(q1)}}.`$
We now turn back to our the original lattice $`\{1,\mathrm{},N\}`$. Again we can write
$$\frac{1}{\left(\begin{array}{c}N\\ s\end{array}\right)}\underset{\tau \mathrm{\Omega }_s}{}\left(\underset{x\tau }{}\left(\frac{x}{p}\right)\right)^2=\frac{1}{\left(\begin{array}{c}N\\ s\end{array}\right)}\underset{x=1}{\overset{p}{}}c_N(x)\left(\frac{x}{p}\right).$$
In this case, however, the multiplicity function $`c_N(x)`$ can not be handled as easily as before. In particular, given $`x\{1,\mathrm{}N\}`$, we denote by $`\mathrm{\Gamma }(x)`$ the set $`=\{u_1,\mathrm{},u_{d_N(x)}\}`$ given by the $`u`$’s in $`\{1,\mathrm{},N\}`$ such that $`u^1x\{1,\mathrm{},N\}`$. The cardinality $`d_N(x)`$ of the set $`\mathrm{\Gamma }(x)`$ is a non trivial function of $`x`$. It is shown in Fig. 9 for $`1xN`$. If now $`\left(\frac{x}{p}\right)=1`$, then clearly
$$c_N(x)=d_N(x)\left(\begin{array}{c}N2\\ s2\end{array}\right).$$
On the other hand, if $`\left(\frac{x}{p}\right)=1`$ (i.e. $`\overline{u}^2=x`$), then (note that either $`\overline{u}\{1,\mathrm{},N\}`$ or $`\overline{u}\{1,\mathrm{},N\}`$)
$$c_N(x)=(d_N(x)1)\left(\begin{array}{c}N2\\ s2\end{array}\right)+\left(\begin{array}{c}N1\\ s1\end{array}\right).$$
We can then use these informations and write
$`{\displaystyle \frac{1}{\left(\begin{array}{c}N\\ s\end{array}\right)}}{\displaystyle \underset{\tau \mathrm{\Omega }_s}{}}\left({\displaystyle \underset{x\tau }{}}\left({\displaystyle \frac{x}{p}}\right)\right)^2`$ $`=`$ $`{\displaystyle \frac{1}{\left(\begin{array}{c}N\\ s\end{array}\right)}}{\displaystyle \underset{x=1}{\overset{p}{}}}c_N(x)\left({\displaystyle \frac{x}{N}}\right)`$
$`=`$ $`{\displaystyle \frac{1}{\left(\begin{array}{c}N\\ s\end{array}\right)}}[{\displaystyle \underset{\left(\frac{x}{p}\right)=1}{}}d_N(x)\left(\begin{array}{c}N2\\ s2\end{array}\right)+`$
$`+{\displaystyle \underset{\left(\frac{x}{p}\right)=1}{}}((d_N(x)1)\left(\begin{array}{c}N2\\ s2\end{array}\right)+\left(\begin{array}{c}N1\\ s1\end{array}\right))]`$
$`=`$ $`{\displaystyle \frac{1}{\left(\begin{array}{c}N\\ s\end{array}\right)}}[\left(\begin{array}{c}N2\\ s2\end{array}\right){\displaystyle \underset{x=1}{\overset{p}{}}}d_N(x)\left({\displaystyle \frac{x}{p}}\right)+`$
$`+(\left(\begin{array}{c}N1\\ s1\end{array}\right)\left(\begin{array}{c}N2\\ s2\end{array}\right)){\displaystyle \underset{\left(\frac{x}{p}\right)=1}{}}1].`$
Finally, we have the following expression for the variance $`\sigma _s^2(m)`$ of the magnetization $`m`$ over the space $`\mathrm{\Omega }_s`$:
$$\sigma _s^2(m)=\frac{4}{N^2}\mathrm{E}_s\left(\left(\underset{x\tau }{}\left(\frac{x}{p}\right)\right)^2\right)\frac{4s^2(m^0)^2}{N^2},$$
from which one easily gets formula (4.33).
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# The Effect of Gasdynamics on the Structure of Dark Matter Halos
## 0.1 Introduction
N-body simulations of structure formation that have as initial conditions a spectrum of density fluctuations consistent with a cold dark matter (CDM) universe have revealed that dark matter halos possess a universal density profile that diverges as $`r^\gamma `$ near the center, with $`1\gamma 2`$ (e.g. Navarro, Frenk, & White 1997 “NFW”; Moore et al. 1998). Much effort has been made attempting to understand the apparent discrepancy between these singular density profiles found in N-body simulations and current observations of the rotation curves of nearby dwarf galaxies and of strong lensing of background galaxies by the galaxy cluster CL0024+1654 (?), which suggest that cosmological halos of all scales have flat density cores, instead. Several ways to reconcile the theory with observations have been hypothesized, among which are:
* Problems with the N-body gravity solvers. Upon further scrutiny, however, the simulations have been found to represent faithfully the actual gravitational processes at work in hierarchical clustering.
* Self-interacting dark matter. If the dark matter is not collisionless as previously assumed but rather can interact nongravitationally with itself, this may change the dynamics in such a way as to produce constant density cores (?).
* Gasdynamical effects due to the baryonic component. The inclusion of gas dynamics in the N-body simulations may lead to constant density cores.
It is the latter case which we address here. In particular, we analyze the density profiles of cosmological halos formed in numerical simulations of cosmological pancakes with and without a substantial baryonic component.
## 0.2 Halo Formation by Gravitational Instability During Pancake Collapse
Consider the growing mode of a single sinusoidal plane-wave density fluctuation of comoving wavelength $`\lambda _p`$ and dimensionless wavevector $`𝐤_𝐩=\widehat{𝐱}`$ (length unit = $`\lambda _p`$) in an Einstein-de Sitter universe dominated by cold, collisionless dark matter. Let the initial amplitude $`\delta _i`$ at scale factor $`a_i`$ be chosen so that a density caustic forms in the collisionless component at scale factor $`a=a_c=a_i/\delta _i`$.
Pancakes modelled in this way have been shown to be gravitationally unstable, leading to filamentation and fragmentation during the collapse (?). Such instability can be triggered if we perturb the 1D fluctuation described above by adding to the initial conditions two transverse, plane-wave density fluctuations with equal wavelength $`\lambda _s=\lambda _p`$, wavevectors $`𝐤_𝐬`$ pointing along the orthogonal vectors $`\widehat{𝐲}`$ and $`\widehat{𝐳}`$, and amplitudes $`ϵ_y\delta _i`$ and $`ϵ_z\delta _i`$, respectively, where $`ϵ_y1`$ and $`ϵ_z1`$. A pancake perturbed by two such density perturbations will be referred to as $`S_{1,ϵ_y,ϵ_z}`$. All results presented here refer to the case $`S_{1,0.2,0.2}`$ unless otherwise noted. Such a perturbation leads to the formation of a quasi-spherical mass concentration in the pancake plane at the intersection of two filaments (Figure 1). Halos formed from pancake collapse as modelled above have a density profile similar in shape to those found in simulations of hierarchical structure formation with realistic initial fluctuation spectra (??). As such, pancake collapse and fragmentation can be used as a test-bed model for halo formation which retains the realistic features of anisotropic collapse, continuous infall, and cosmological boundary conditions. Questions we seek to answer quantitatively are:
* Is the halo that results from pancake instability and fragmentation in a relaxed equilibrium state?
* How isotropic is the resulting velocity distribution for the dark matter, and how isothermal are the gas and dark matter?
* What effect does including gas in the simulations have on the dark matter density profile?
The pancake problem (without radiative cooling) is self-similar and scale-free, once distance is expressed in units of the pancake wavelength $`\lambda _p`$ and time is expressed in terms of the cosmic scale factor $`a`$ in units of the scale factor $`a_c`$ at which caustics form in the dark matter and shocks in the gas. In the currently-favored flat, cosmological-constant-dominated universe, however, this self-similarity is broken because $`\mathrm{\Omega }_M/\mathrm{\Omega }_\mathrm{\Lambda }`$ decreases with time, where $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ are the matter and vacuum energy density parameters, respectively. For objects which collapse at high redshift in such a universe (e.g. dwarf galaxies), the Einstein-de Sitter results are still applicable as long as we take $`(\mathrm{\Omega }_B/\mathrm{\Omega }_{DM})_{EdeS}=(\mathrm{\Omega }_B/\mathrm{\Omega }_{DM})_\mathrm{\Lambda }`$, where $`\mathrm{\Omega }_B`$ and $`\mathrm{\Omega }_{DM}`$ are the baryon and dark matter density parameters. If $`\mathrm{\Omega }_B=0.045`$, $`\mathrm{\Omega }_{DM}=0.255`$, and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ at present, then the E-deS results are applicable if we take $`\mathrm{\Omega }_B=0.15`$ and $`\mathrm{\Omega }_{DM}=0.85`$, instead.
## 0.3 Results
The numerical technique used here to simulate the collapse and fragmentation of a cosmological pancake couples the gasdynamical method Adaptive Smoothed Particle Hydrodynamics (ASPH) (??) with a particle-particle, particle-mesh (P<sup>3</sup>M) N-body gravity solver (?). Results presented here are for a simulation with $`64^3`$ particles each of dark matter and gas when hydrodynamics is included, and $`64^3`$ dark matter particles for the case with gravity only, on a P<sup>3</sup>M grid with $`128^3`$ cells. We consider an Einstein-de Sitter universe ($`\mathrm{\Omega }_M=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$) and two cases: $`(\mathrm{\Omega }_B,\mathrm{\Omega }_{DM})=(0,1)`$ or $`(0.15,0.85)`$. The pancake-filament-halo structure of our results is illustrated by Figure 1. By $`a/a_c=5`$, this halo contains approximately $`10^5`$ particles each of dark matter and gas.
### 0.3.1 The equilibrium halo structure and similarity to CDM $`N`$-body simulation halos
The left panel of Figure 2 shows the evolution of the density profile of the dark matter halo in the simulation with no gas. By $`a/a_c=5`$, the halo has asymptoted to an equilibrium profile. A $`\chi ^2`$-minimization fit was performed on this profile over the range between the gravitational softening length $`r_{soft}`$ and $`r_{200}`$, the radius within which the mean density is 200 times the cosmic mean density at that epoch, where $`r_{soft}/\lambda _p2.3\times 10^3`$ and $`r_{200}/\lambda _p7.5\times 10^2`$, for each component X of the matter density. The fitting function we use, in general, is of the form
$$\rho _X=\frac{\rho _s}{(\frac{r}{r_s})^\gamma (1+(\frac{r}{r_s})^{1/\alpha })^{\alpha (\beta \gamma )}}$$
(1)
(??), which approaches $`\rho _Xr^\gamma `$ at $`rr_s`$ and $`\rho _Xr^\beta `$ at $`rr_s`$, where $`\rho _X`$ and $`\rho _s`$ are in units of the mean cosmic density of component X at the same epoch. We found $`\{\alpha ,\beta ,\gamma ,r_s/\lambda _p,\rho _s\}=\{0.866,3.18,1.18,1.39\times 10^2,9.19\times 10^3\}`$. This profile is similar in shape to the universal density profile found for CDM halos by NFW based upon N-body simulations, for which $`\{\alpha ,\beta ,\gamma \}=\{1,3,1\}`$. In the notation of NFW, our best-fit halo concentration parameter, $`cr_{200}/r_s`$, is $`c5`$, while our $`\rho _s`$ corresponds to NFW’s $`\delta _c`$. According to NFW, $`\delta _c=(200/3)c^3/[ln(1+c)c/(1+c)]`$. For our value of $`c5`$, this yields $`\delta _c9\times 10^3`$, very close to our best-fit value of $`\rho _s`$. Such a density profile is thus a more general outcome of cosmological halo formation by gravitational condensation, not limited to hierarchical clustering.
A common assumption is that “virialized” objects resulting from cosmological perturbations have an approximately isothermal temperature distribution and a dark matter velocity distribution that is isotropic. Figure 3 shows the spherically-averaged specific thermal energy of the gas ($`\frac{3}{2}k_BT/\mu m_p`$) and kinetic energy of the dark matter ($`\frac{1}{2}v^2`$). The gas in the halo is nearly isothermal from the softening length to $`r_{200}`$, with $`T(r_{soft})/T(r_{200})2`$. The dark matter kinetic energy follows a similar profile with a slightly steeper drop closer to the center. Figure 4 shows the parameter $`\beta _{an}1v_t^2/2v_r^2`$, which is a measure of the anisotropy of the velocity distribution of the collisionless dark matter, as a function of distance from the center. In the limits of isotropic and radial motion, $`\beta _{an}=0`$ and $`1`$, respectively. The distribution is somewhat more anisotropic than usually found in hierarchical clustering simulations, due to the highly anisotropic infall which channels matter along the filaments that intersect at our halo, favoring radial motion. While this is true for both cases, the case with gas included seems to suppress radial motion at small radii.
### 0.3.2 Comparison of density profiles with and without gas
Figure 2 shows the profile fits for each component. The parameters $`\{\alpha ,\beta ,\gamma ,r_s/\lambda _p,\rho _s\}`$ for each component are:
* Dark matter for N-body only: $`\{0.866,3.18,1.18,1.39\times 10^2,9.19\times 10^3\}`$
* Dark matter for N-body/Gas: $`\{0.141,2.78,1.95,2.82\times 10^2,6.15\times 10^2\}`$
* Gas: $`\{3.07\times 10^3,2.84,1.98,2.47\times 10^2,8.06\times 10^2\}`$
The density profiles shown in the right panel of Figure 2 show that when gas is added to the simulation, the dark matter density profile steepens at the center. The dark matter parameter $`\gamma `$, the negative of the asymptotic logarithmic slope for inner radii, is $`1.18`$ for the case with no gas and $`1.95`$ for the case with. While the fitting function does not make a direct comparison of the fits at inner radii easy because of the difference in scale radius $`r_s`$ for each component, it does show that the halo in the case with gas has a steeper inner slope. As such, this effect can only worsen the disagreement between the predicted singular profiles and observations of constant density cores. More detailed analysis at higher resolution is needed before this conclusion can be made with certainty, since the part of the profile in question is at radii very close to the numerical softening length for the gravity calculation, even though there were more than $`10^5`$ particles each of gas and dark matter within $`r_{200}`$ at this epoch.
## 0.4 Summary and Discussion
We have analyzed the structure of the halo which forms by gravitational instability when a cosmological pancake is perturbed during its collapse, based upon ASPH/P<sup>3</sup>M simulations. Such halos resemble those formed in CDM simulations, with a density profile like that of NFW, when no gas is included. Our results strengthen the conclusion reached elsewhere that halos with a universal structure like that identified by NFW arise from gravitational instability under a wider range of circumstances than those involving hierarchical clustering (???). When gas is included, we find that the halo gas acquires a steeper inner profile than that of the pure N-body results, almost as steep as $`\rho r^2`$, as does the dark matter in this case. The halo appears to be in a quasi-isothermal equilibrium state, although matter continues to fall into the halo, mainly along the filaments. The persistence of filamentary substructure in the halo seems to manifest itself by sustaining the radial motion of the infalling particles even after “virialization”, in contrast to results found from simulations of hierarchical clustering (?). The inclusion of gas lessens this effect somewhat at the center of the halo.
Our result that the presence of a gasdynamical component steepens the inner profile of the dark matter halo relative to the pure N-body result is consistent with some other recent simulations of cluster formation in the CDM model (e.g Lewis et al. 1999). Tittley and Couchman (1999) report $`\rho _{DM}r^{1.8}`$ at small r in their simulations, in rough agreement with our result, although they note that the NFW profile is also a reasonable fit. However, for CDM-like initial conditions, they actually report a steeper inner profile for the gas than for the dark matter, while our gas and dark matter inner profiles are equally steep. CDM simulations of a single cluster by many different codes, summarized in Frenk et al. (2000), report somewhat different results from those of Tittley and Couchman (1999). In particular, while the dark matter profile was found in Frenk et al. (2000) to be well-fit by the NFW profile when gas was included, the gas density profile was even flatter than this at small radii. Lewis et al. (1999) also find a gas density inner profile which is flatter than that of the dark matter. It would be tempting to attribute the difference between our result that the density profiles of both gas and dark matter are steeper than NFW when gas is included, and the flatter inner profiles reported in Frenk et al. (2000), to the difference between CDM-like and pancake-instability initial conditions. However, this would not explain the difference between the results of Frenk et al. (2000) and Tittley and Couchman (1999), both for CDM initial conditions.
###### Acknowledgements.
This work was supported in part by grants NASA NAGC-7363 and NAG5-7821, Texas ARP No. 3658-0624-1999, and NSF SBR-9873326 and ACR-9982297.
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# Low Redshift QSO Ly𝛼 Absorption Line Systems Associated with Galaxies
## 1 Introduction
The origin and nature of low redshift Ly$`\alpha `$ line absorbers are still a matter of debate (cf. Chen et al. 1998, hereafter CLWB; Tripp, Lu, & Savage 1998). These absorbers are thought to arise either from gaseous galactic haloes/discs or from the underdense web-like regions of filaments and sheets in the large scale structure of the cosmic matter. The first suggestion seems reasonable because galaxies in general seem to possess extended gaseous haloes (Bahcall & Spitzer 1969) or huge gaseous discs (Maloney 1992; Hoffman et al. 1993). The second suggestion is drawn from the studies of high redshift Ly$`\alpha `$ absorption lines. At high redshift, it is widely believed that most Ly$`\alpha `$ absorption lines are tracers of intergalactic hydrogen as suggested by Sargent et al. (1980). A diffuse intergalactic medium (IGM) model for the Ly$`\alpha `$ forest was investigated by Bi, Börner, & Chu (1992). Several authors investigated the scenario of the Ly$`\alpha `$ forest produced by the IGM in the context of cosmological simulations (e.g., Cen et al. 1994; Petitjean, Mücket, & Kates 1995; Miralda-Escudé et al. 1996; Hernquist et al. 1996; Haehnelt, Steinmetz, & Rauch 1996; Cen & Simcoe 1997; Zhang et al. 1997, 1998; Theuns et al. 1998; Machacek et al. 1999). It is natural to extend this model to low redshift. Some insights have been provided by hydrodynamic simulations predicting absorber properties down to $`z=0`$. For instance, Theuns, Leonard, & Efstathiou (1998) find that the observed decrease in the rate of evolution of Ly$`\alpha `$ absorption lines at $`z2`$ can be explained by the steep decline in the photoionizing background resulting from the rapid decline in the quasar numbers at low redshift (see also Riediger, Petitjean, & Mücket 1998). Davé et al. (1999) find that shocked or radiatively cooled gas of higher overdensity can give rise to the majority of strong Ly$`\alpha `$ lines at every redshift. However on the observational side, Lanzetta et al. (1995, hereafter LBTW) & CLWB claim that luminous galaxies can account for at least about $`50`$ per cent (and even more) of the strong Ly$`\alpha `$ absorption lines ($`W_r0.3`$ Å) observed by the Hubble Space Telescope (HST) Quasar Absorption Line Key Project (Bahcall et al. 1993; Bahcall et al. 1996; Weymann et al. 1998). Although the fraction is still quite uncertain because of the unknown galaxy space density (or luminosity function), the unknown gas absorption cross section of galaxies and the uncertainties of the observed number density of Ly$`\alpha `$ absorption systems, there is no doubt in principle that some absorbers are physically associated with galaxies especially for those with column densities above $`10^{14}\mathrm{cm}^2`$. Thus the origin of Ly$`\alpha `$ lines at low redshift is still an unsolved problem.
In order to clarify this question, it is particularly important to identify galaxies giving rise to QSO absorption and to analyse their physical properties. So far, tremendous efforts have been made to locate responsible galaxies in QSO fields (e.g., Morris et al. 1993; Lanzetta et al. 1995; Bowen et al. 1996; Le Brun, Bergeron, & Boissé 1996; van Gorkom et al. 1996; CLWB; Bowen, Pettini, & Boyle 1998; Tripp et al. 1998; Impey, Petry, & Flint 1999). In general, these investigations of the physical link between an individual galaxy and an individual absorption system aim at answering (1) whether absorbers are physically associated with galaxies and what the percentage of absorption lines is arising from galaxies, (2) whether there is an anti-correlation between projected distance (i.e., impact parameter) from the line-of-sight (hereafter LOS) to the galaxy centre and the absorption rest frame equivalent width (hereafter REW). For example, Morris et al. (1993) carried out a redshift survey of galaxies in the field of 3C273 and found no galaxies within a projected distance of $`230h_{80}^1\mathrm{kpc}`$ of any of the 12 lines toward 3C273 with REW exceeding 50 mÅ. One of their conclusions is that the absorbers are not randomly distributed with respect to the galaxies, though the absorber-galaxy correlation is not as strong as the galaxy-galaxy correlation. In a similar program but with different LOS, Stocke et al. (1995) and Shull et al. (1996) suggested that most of the Ly$`\alpha `$ absorbers are located within large-scale galaxy structures. In contrast, using an imaging and spectroscopic survey (in the field of HST spectroscopic target QSOs), LBTW claimed that the fraction of absorbers arising from galaxies is quite high and there is a distinct anti-correlation of REW versus projected distance. CLWB confirmed these results with more LOSs and concluded that most galaxies are surrounded by extended gaseous envelopes of $`170h^1\mathrm{kpc}`$ in radius and that many or most Ly$`\alpha `$ absorption systems arise in galaxies. Tripp et al. (1998) reached similar conclusions, but cautioned that selection effects could artificially tighten the anti-correlation, and also that the galaxy survey may be incomplete. They pointed out that there could be fainter galaxies at smaller projected distance (also suggested by Linder 1998; see also Impey et al. 1999) which could be revealed in a deeper survey. Moreover, they found some missing lines from the CLWB samples and from their LOSs. These missing lines would weaken the anti-correlation, if included. Recently, Impey, Petry & Flint (1999) studied Ly$`\alpha `$ QSO absorbers in the nearby universe ($`0<z<0.22`$) based on the spectroscopy of ten quasars obtained with the Goddard High Resolution Spectrograph (GHRS) of the HST. At odds with the results of LBTW & CLWB, they concluded that nothing in their data would specifically lead to associate absorbers preferentially with haloes of luminous galaxies.
Another very useful tool for the analysis of physical properties of absorbers is provided by observations of the intervening absorption in multiple LOSs either from close quasar pairs or from multiple, gravitational lensed quasar images (see Rauch 1998 for the part of review on this subject). For example, double LOSs observations (e.g., Dinshaw et al. 1995; Dinshaw et al. 1994; Fang et al. 1996; Petitjean et al. 1998) have shown that the absorber size is about hundreds of kiloparsecs, Rauch et al. (1999) analysed spectra of images of a lensed quasar at redshift 3.628 and found some low-ionization lines arising from the ISM.
Given these observational results, it is important to build theoretical models to understand the origin of the Ly$`\alpha `$ absorbers at low redshift. Some theoretical efforts have been made to relate absorption systems with galaxies (e.g., Mo & Morris 1994; Mo 1994; Morris & van den Bergh 1994; Mo & Miralda-Escudé 1996; Linder 1998; Linder 1999). Unfortunately, even for those Ly$`\alpha `$ absorption lines genuinely arising from galaxy haloes, we do not know a priori which part of the galaxy gives rise to the absorption. In other words, it is unclear whether the absorbing clouds are located in the outer regions of the halo as infalling clouds or in the rotating disc as interstellar medium clouds or in the satellites. There are some competing models. Morris and van den Bergh (1994) estimated that tidal tails can explain $`20`$ per cent of the low redshift Ly$`\alpha `$ absorbers, but so far there is a lack of detailed models. Mo & Miralda-Escudé (1996) concluded that gaseous galaxy haloes can account for all absorbers with HI column density $`N_{\mathrm{HI}}10^{17}\mathrm{cm}^2`$ at redshift $`z2`$. Recently a model in which absorption is due to gas in an extended disc was proposed by Linder (1998, 1999), who argued that high surface brightness galaxies together with low surface brightness galaxies can account for the majority of Ly$`\alpha `$ absorption line systems. However this picture only incorporates spiral galaxies while there are some absorbers associated with E/S0 galaxies (cf. CLWB). This model requires a large number of low surface brightness galaxies and it is unknown if a spiral galaxy can possess a disc extending beyond $`100h^1\mathrm{kpc}`$.
In current models of galaxy formation, galaxies are considered to possess haloes, discs (for spiral galaxies) and satellites which can give rise to absorption. Motivated by these considerations, we perform Monte-Carlo simulations using semi-analytic models with plausible assumptions, given our current knowledge about the properties of these components. Our aim is to study: (1) which component is most important. (2) whether current observational results can be explained by the models. (3) whether selection effects (which should be applied when pairing an absorber with a luminous galaxy) can tighten the correlations between REW and projected distance. (4) what kind of future observations are needed to discriminate models and to examine the correlations.
The paper is organized as follows. In § 2, we describe our Monte-Carlo simulation methods and give results. We construct simulations with allowed parameters and compare the predicted line number density $`(\frac{dN}{dz})`$ with observational results for $`W_r0.3`$Å Ly$`\alpha `$ lines, Lyman-limit systems, and damped Ly$`\alpha `$ systems. Detailed properties of absorbers, correlations of equivalent width versus projected distance, galaxy luminosity and redshift, are studied. The average covering factor is estimated. In § 3, we study selection effects. After applying selection criteria, we compare our results with imaging and spectroscopic surveys. At last we make some predictions. A discussion is presented in § 4. A summary of the results is given in § 5. Throughout the paper, we adopt a dimensionless Hubble constant $`h=H_0/(100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1)`$. Our presentation is mainly based on the $`\mathrm{\Lambda }`$CDM cosmogony (with $`\mathrm{\Omega }_0=0.3,\mathrm{\Omega }_{\mathrm{\Lambda },0}=0.7,h=0.7`$), results based on the SCDM cosmogony (with $`\mathrm{\Omega }_0=1.0,\mathrm{\Omega }_{\mathrm{\Lambda },0}=0.0,h=0.5`$) are also discussed for comparison.
## 2 The Monte-Carlo Simulations
We start our Monte-Carlo simulations with galaxy samples (whose luminosity distribution is consistent with observational luminosity functions) along a QSO LOS. These galaxies are placed along the LOS randomly within a cylinder volume in co-moving coordinate (then we assign a redshift and a projected distance for each galaxy). For one particular galaxy whose halo is characterized by the circular velocity (derived from the Tully-Fisher relation or the Faber-Jackson law), we model its absorbing components (discs, halo gas clouds and satellites) in detail so as to determine its gas cross section and cloud properties, such as HI column density, temperature and LOS velocity. The total equivalent line width is then calculated assuming a Voigt profile for each cloud. This procedure produces a catalogue of absorber-galaxy pairs with information about the absorbing galaxies for further analysis.
### 2.1 Galaxy samples
There will be a number of galaxies intersecting a particular LOS with random projected distances to a given QSO at redshift $`z_\mathrm{q}`$. We consider those galaxies in a cylinder within a radius of $`R_{\mathrm{cy}}=400h^1\mathrm{kpc}`$ to the QSO LOS, since in our models there is no absorption cloud outside of $`400h^1\mathrm{kpc}`$ and the upper limit of the projected distance in imaging surveys of absorbers is less than this radius. The number of galaxies at redshift $`z`$ with interval $`\mathrm{\Delta }z`$ is
$$\mathrm{\Delta }N_\mathrm{g}=n_c(1+z)^3\pi R_{\mathrm{cy}}^2\frac{cdt}{dz}\mathrm{\Delta }z,0<z<z_\mathrm{q},$$
(1)
where
$$\frac{dt}{dz}=\frac{1}{(1+z)H(z)},$$
(2)
$$H(z)=H_0[\mathrm{\Omega }_{\mathrm{\Lambda },0}+(1\mathrm{\Omega }_{\mathrm{\Lambda },0}\mathrm{\Omega }_0)(1+z)^2+\mathrm{\Omega }_0(1+z)^3]^{1/2}.$$
(3)
The co-moving galaxy density $`n_c`$ is obtained by integration over the B-band Schechter luminosity function
$$n_c=_{L_{B\mathrm{min}}}^{\mathrm{}}\varphi (L_B)𝑑L_B$$
$$=_{L_{B\mathrm{min}}}^{\mathrm{}}\varphi ^{}(L_B/L_B)^\alpha e^{L_B/L_B}𝑑L_B/L_B,$$
(4)
where $`L_{B\mathrm{min}}`$ is the minimum B-band luminosity. The luminosities of these galaxies are selected in such a way that their distribution is consistent with the luminosity function $`\varphi (L_B)`$.
The luminosity functions for different morphological types over the range $`22M_B+5\mathrm{log}h14`$ are as follows (Marzke et al. 1998): (1) For late-type galaxies (Spiral), $`\varphi ^{}=8.0\pm 1.4\times 10^3h^3\mathrm{Mpc}^3`$, $`\alpha =1.11_{0.06}^{+0.07}`$, $`M_B=19.43_{0.08}^{+0.08}+5\mathrm{log}h`$, and $`0.0067L_BL_B10.7L_B`$. (2) For early-type galaxies (E/S0), $`\varphi ^{}=4.4\pm 0.8\times 10^3h^3\mathrm{Mpc}^3`$, $`\alpha =1.00_{0.09}^{+0.09}`$, $`M_B=19.37_{0.11}^{+0.10}+5\mathrm{log}h`$, and $`0.007L_BL_B11.3L_B`$. We do not consider Irr/Pec galaxies because they are rare ($`\varphi ^{}=0.2\pm 0.08\times 10^3h^3\mathrm{Mpc}^3`$) <sup>1</sup><sup>1</sup>1Only normal galaxies are considered in the models. In reality however, galaxies may contain HI tidal tails and could give rise to absorption. We will discuss the problem in § 4..
The luminosity functions above are valid only at very low redshift ($`z<0.05`$). They are derived from the recently enlarged Second Southern Sky Redshift Survey (SSRS2). Some other determinations give higher normalizations. For example, the galaxy luminosity functions from the ESO Slice Project (ESP) galaxy redshift survey (Zucca et al. 1997) is characterized by $`\varphi ^{}=2.0\pm 0.4\times 10^2h^3\mathrm{Mpc}^3`$, $`M_{}=19.61_{0.08}^{+0.06}`$ and $`\alpha =1.22_{0.07}^{+0.06}`$ over the redshift interval $`z<0.3`$. The ESP luminosity functions are in good agreement with those of the AUTOFIB redshift survey (Ellis et al. 1996) which are characterized by $`\varphi ^{}=2.45_{0.31}^{+0.37}\times 10^2h^3\mathrm{Mpc}^3`$, $`\alpha =1.16_{0.12}^{+0.15}`$, $`M_{}=19.30_{0.12}^{+0.15}`$ over the redshift interval $`0.02<z<0.15`$, $`\varphi ^{}=1.48_{0.19}^{+0.30}\times 10^2h^3\mathrm{Mpc}^3`$, $`\alpha =1.41_{0.07}^{+0.12}`$, $`M_{}=19.65_{0.10}^{+0.12}`$ over the redshift interval $`0.15<z<0.35`$, and $`\varphi ^{}=3.55_{2.00}^{+2.91}\times 10^2h^3\mathrm{Mpc}^3`$, $`\alpha =1.45_{0.18}^{+0.16}`$, $`M_{}=19.38_{0.25}^{+0.27}`$ over the redshift interval $`0.35<z<0.75`$. The luminosity function of Lilly et al. (1995) is characterized by $`\varphi ^{}=2.72\pm 0.4\times 10^2h^3\mathrm{Mpc}^3`$, which is about twice that of SSRS2 (Marzke et al. 1998), and $`\alpha =1.03`$ over the redshift interval $`0.2<z<0.5`$.
Since the faint-end slope $`\alpha `$ and the characteristic magnitude $`M_{}`$ of various luminosity functions are not much different, we apply the AUTOFIB luminosity function normalization over $`0.02<z<0.15`$ in our simulations and assign it to the two morphological types with the same ratio for spirals and E/S0 galaxies as in the SSRS2 luminosity function. Namely, the characteristic luminosity is $`\varphi ^{}1.58\times 10^2h^3\mathrm{Mpc}^3`$ for spiral galaxies and $`\varphi ^{}0.87\times 10^2h^3\mathrm{Mpc}^3`$ for E/S0 galaxies. We will discuss results for other luminosity functions.
As the luminosity of a galaxy in the cone along the LOS and its redshift $`z_\mathrm{g}`$ are known, we calculate the apparent magnitude of the galaxy applying the k-correction and cosmological dimming,
$$m_B=M_B+5\mathrm{log}(D_L)+25+k_B(z_\mathrm{g}),$$
(5)
where $`D_L=(1+z)D_M`$ is luminosity distance of the galaxy in $`\mathrm{Mpc}`$, and $`D_M`$, the proper motion distance is
$$D_M=c_0^{z_\mathrm{g}}H(z)^1𝑑z$$
(6)
for a flat universe ($`\mathrm{\Omega }_k=1\mathrm{\Omega }_0\mathrm{\Omega }_\mathrm{\Lambda }=0`$), which is the case in this paper. We adopt B-band $`k`$ corrections for galaxies of different morphological types as in Pence (1976).
At any given epoch, haloes can be parameterized by their circular velocity $`V_{\mathrm{cir}}`$, which is simply related with galaxy morphological type and luminosity. Empirical relations are known between B-band magnitude and LOS velocity dispersion, $`\sigma `$, of the matter in galaxies both for ellipticals and for spirals (Faber and Jackson 1976; Tully and Fisher 1977).
The Faber-Jackson relation is
$$M_B+5\mathrm{log}h=(19.39\pm 0.07)+10(\mathrm{log}\sigma 2.3)$$
(7)
for ellipticals, and
$$M_B+5\mathrm{log}h=(19.75\pm 0.07)+10(\mathrm{log}\sigma 2.3)$$
(8)
for S0’s (Fukugita & Turner 1991). The circular velocity of the halo for elliptical and S0 galaxies is $`V_{\mathrm{cir}}=\sqrt{2}\sigma `$. Using the E/S0 type luminosity function of Marzke et al. (1998), we get $`81.6\mathrm{km}\mathrm{s}^1V_{\mathrm{cir}}514.6\mathrm{km}\mathrm{s}^1`$ ($`V_{\mathrm{cir}}^{}280.1\mathrm{km}\mathrm{s}^1`$) for Elliptical galaxies and $`75.1\mathrm{km}\mathrm{s}^1V_{\mathrm{cir}}473.6\mathrm{km}\mathrm{s}^1`$ ($`V_{\mathrm{cir}}^{}258.5\mathrm{km}\mathrm{s}^1`$) for S0 galaxies.
For spirals the Tully-Fisher relation is used to derive the LOS velocity width $`\mathrm{\Delta }v`$. We take (Fukugita & Turner 1991)
$$M_B+5\mathrm{log}h=(19.18\pm 0.10)+(6.56\pm 0.48)(\mathrm{log}\mathrm{\Delta }v2.5).$$
(9)
The halo circular velocity for a spiral is $`V_{\mathrm{cir}}=\mathrm{\Delta }v/2`$. We get $`25.7\mathrm{km}\mathrm{s}^1V_{\mathrm{cir}}425.4\mathrm{km}\mathrm{s}^1`$ ($`V_{\mathrm{cir}}^{}172.6\mathrm{km}\mathrm{s}^1`$) using the spiral-type luminosity function of Marzke et al. (1998). Note that the upper limits of the circular velocity are very large, there are, however, no such large galaxies in the sample due to the exponential cutoff in the Schechter luminosity function.
### 2.2 Gaseous galactic haloes
We model the gaseous galactic haloes following the work by Mo & Miralda-Escudé (1996). In such semi-analytic models, it is assumed that the gas in a halo has a two-phase (a hot phase and a cold phase) structure which, in principle is described by the density profiles and the temperature profile. The density profiles are characterized by the so-called cooling radius and virial radius. And the temperature profile of the hot gas is characterized by the virial temperature. Our modeling is summarized as follows (see Mo & Miralda-Escudé 1996 for more details):
In cooling flow models, when the gravitational potential is important, the core radius of the hot gas profile is similar to the cooling radius (Waxman & Miralda-Escudé 1995). Thus, a self-similar density profile for hot gas is assumed as,
$$\rho _h(r)=\rho _h(r_c)\frac{2r_c^2}{r(r+r_c)},$$
(10)
where
$$\rho _h(r_c)=\frac{5\mu kT_v}{2\mathrm{\Lambda }(T_v)t_\mathrm{M}}.$$
(11)
We assume
$$\rho _h(r_c)=\frac{f_gV_{\mathrm{cir}}^2}{4\pi Gr_c^2},$$
(12)
so that the density of the hot gas at this radius is a fraction $`f_g`$ of the total density of the halo. $`\mathrm{\Lambda }(T_v)`$ is the cooling rate of the gas at the virial temperature. The gas mass fraction $`f_g`$ is assumed $`0.030.05`$. $`\mu `$ is the average mass per particle, which is $`0.6m_\mathrm{H}`$ with $`m_\mathrm{H}`$ being the mass of hydrogen nucleus. The cooling radius $`r_c`$ is determined by eq. (11) and eq. (12),
$$r_c=\sqrt{\frac{f_g\mathrm{\Lambda }(T_v)t_\mathrm{M}}{5\pi G\mu ^2}}124.6\mathrm{kpc}\sqrt{\mathrm{\Lambda }_{23}t_\mathrm{M}/10\mathrm{G}\mathrm{yrs}},$$
(13)
where $`\mathrm{\Lambda }_{23}`$ is the cooling rate in units of $`10^{23}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^3`$. The hot gas is taken to be isothermal, so that $`T_h(r)T_v=\mu V_{\mathrm{cir}}^2/2k`$. $`t_\mathrm{M}=t/(1+\mathrm{\Omega }_0)`$, is the time interval between major mergers, since the gas is then shock heated to a stage from which it starts cooling. The cooling function $`\mathrm{\Lambda }(T_v)`$ is adopted from Sutherland & Dopita (1993). The age of a halo at redshift $`z_1`$ is an integration of eq. (2)
$$t=_{z_1}^{\mathrm{}}\frac{1}{(1+z)H(z)}𝑑z.$$
(14)
The virial radius is
$$r_v=V_{\mathrm{cir}}/[10H(z)].$$
(15)
When the hot gas is shock heated and starts to cool, it will sink to the galaxy centre with velocity $`\stackrel{}{u}=\widehat{r}u(r)`$. The cooling flow can be described as,
$$\frac{\rho _c}{t}+(\rho _c\stackrel{}{u})=\frac{\mathrm{\Lambda }(T_h)}{\frac{5}{2}\mu kT_h}\rho _h^2(r).$$
(16)
We assume $`u(r)`$ to be a constant:
$$u(r)=v_c$$
(17)
where $`v_c`$ is the infall velocity which must be of the order of $`V_{\mathrm{cir}}`$. In eq. (13), $`r_c`$ is a function of $`t`$ (here $`t`$ stands for $`t_\mathrm{M}`$, and we omit the subscript hereafter), and we have $`\dot{r_c}=\frac{dr_c}{dt}=r_c/2t`$ if we assume a constant $`T_v`$. Let us set
$$\begin{array}{ccc}x\hfill & \hfill & \frac{r}{r_c}\hfill \\ \overline{\rho }\hfill & \hfill & \frac{\rho }{\rho _h(r_c)}\hfill \\ \overline{u}\hfill & \hfill & \frac{u}{\dot{r_c}}.\hfill \end{array}$$
Equation (16) has only one variable $`r`$ and can be simplified to
$$(x+\overline{u})\frac{d\overline{\rho _c}}{dx}+\overline{\rho _c}\frac{1}{x^2}\frac{d}{dx}(x^2\overline{u})=\frac{8}{x^2(1+x)^2}.$$
(18)
This equation can be solved analytically (see Appendix A for details).
When $`r_c>r_v`$, $`r_c`$ is simply a parameter rather than a physical cooling radius. Thus we assume that the residual hot gas at $`r_v`$ is still at the virial temperature and has a density such that its cooling time is equal to t, so that $`T_h(r_v)=T_v`$ and $`\rho _h(r_v)=(5\mu kT_v)/[2\mathrm{\Lambda }(T_v)t]`$. In such a case, we replace $`r_c`$ in eq. (10), (11) by $`r_v`$. And we can not use eq. (16) to describe the cooling flow. In this case, most of the accreted gas will have cooled. Since the total gas mass accreted in a halo with circular velocity $`V_{\mathrm{cir}}`$ is $`f_gV_{\mathrm{cir}}^2r_v/G`$, the total mass of gas that has been in the cold phase can be written as
$$M=\frac{f_gV_{\mathrm{cir}}^2r_v}{G}_0^{r_v}4\pi x^2\rho _h(x)𝑑x,$$
(19)
where $`\rho _h(x)`$ is the density of gas in the hot phase as discussed above. The cold gas should form clouds that will fall through the halo. Assume a mass flow rate $`\dot{M}=M/t`$, and assume that the clouds move to the halo centre with a constant velocity $`v_c`$ (which is of the same order as the halo circular velocity). Assuming also spherical symmetry for the gas distribution, we can write the density of the cold gas as (see Mo & Miralda-Escudé 1996 for more details),
$$\rho _c(r)=\frac{f_gV_{\mathrm{cir}}^2r_{\mathrm{min}}}{4\pi Gr^2v_ct}\left[1\frac{r_{\mathrm{min}}^2}{r_c^2}_0^1x^2\overline{\rho _h}(x)𝑑x\right],$$
(20)
where $`r_{\mathrm{min}}=\mathrm{min}(r_v,r_c)`$.
The cold gas is assumed to be, for simplicity, in spherical clouds with masses constrainted by various physical processes, such as gravitational instability, evaporation by hot gas, hydrodynamic instability, etc. Too large clouds will eventually collapse to form stars and small clouds will evaporate by heat conduction and also be disrupted by hydrodynamic instabilities. The net effects of these processes is to preferentially destroy low mass clouds, so it is possible to end-up with a log-normal mass function, like the mass function of star clusters observed in the galaxy (e.g., Gnedin & Ostriker 1997), if we begin with a power-law mass function. For this reason and for lack of knowledge of the mass distribution of cold clouds, we assume a log-normal distribution of cloud masses $`M_c`$,
$$p(M_c)dM_c=\frac{1}{\sqrt{2\pi }\sigma _\mathrm{M}}exp\left[\frac{ln^2(M_c/\overline{M})}{2\sigma _\mathrm{M}^2}\right]\frac{dM_c}{M_c}.$$
(21)
Here $`\overline{M}`$ is the cloud mean mass. As discussed by Mo & Miralda-Escudé (1996), the mean mass of the clouds is approximately $`10^510^6M_{}`$. In this paper, we choose $`\overline{M}=`$ several $`\times 10^5M_{}`$ and $`\sigma _\mathrm{M}=0.10.3`$. We also use a constant cloud mass in simulations, but the results do not change much.
We model the clouds as spheres of uniform, isothermal photoionized gas confined by the pressure of the hot medium, so that $`\rho _{\mathrm{cloud}}T_c=\rho _hT_h`$, where $`T_c`$, the temperature of the clouds, is about $`2\times 10^4K`$. The cloud radius is $`R_c=\left(3M_c/4\pi \rho _{\mathrm{cloud}}\right)^{1/3}`$ (typically $`R_c110\mathrm{kpc}`$ at a radius within $`100\mathrm{kpc}`$, cf. Mo 1994). The total hydrogen column density through the cloud centre will be $`N_0(\mathrm{H})=R_c\rho _{\mathrm{cloud}}/2.3\mu `$, and the H number density is $`n(\mathrm{H})=\rho _{\mathrm{cloud}}/2.3\mu `$.
We assume that the cloud is almost completely photoionized by a constant UV background, and in ionization equilibrium. The fraction of hydrogen in the neutral state (HI atom), is determined by the flux of the UV background ionization field $`J(\nu )`$ and $`n_\mathrm{H}`$. We take
$$J(\nu )=J_{21}(z)\times 10^{21}\left(\frac{\nu }{\nu _{\mathrm{HI}}}\right)^\alpha \mathrm{\Theta }(\nu )\mathrm{erg}\mathrm{cm}^2\mathrm{sr}^1\mathrm{Hz}^1\mathrm{s}^1,$$
(22)
where $`\nu _{\mathrm{HI}}`$ is the hydrogen Lyman limit frequency, $`J_{21}(z)=0.5`$ for $`z>2`$, and $`J_{21}=0.5\times [(1+z)/3]^2`$ for $`z<2`$. A break in the spectrum at $`\nu _44Ry`$ (due to continuum absorption by He II), with $`\mathrm{\Theta }(\nu <\nu _4)=1`$ and $`\mathrm{\Theta }_4\mathrm{\Theta }(\nu \nu _4)=0.1`$, is included (cf. Miralda-Escudé & Ostriker 1990; Madau 1992). We take $`\alpha =0.5`$.
For $`N(\mathrm{H})10^{19}\mathrm{cm}^2`$, which is the case for most clouds, the cloud is optically thin to the ionizing field with ionization parameter $`U=\frac{\mathrm{\Phi }(\mathrm{H})}{n(\mathrm{H})c}`$, where the ionizing photon flux $`\mathrm{\Phi }(\mathrm{H})=_{\nu _{\mathrm{HI}}}^{\mathrm{}}\frac{4\pi J(\nu )}{\mathrm{h}\nu }𝑑\nu `$, h is Planck’s constant. The neutral hydrogen column density $`N_{\mathrm{HI}}`$ can be derived from the code CLOUDY 90 (Ferland 1996). Then we obtain the HI column density of a cloud at a distance to the LOS $`l`$,
$$N_{\mathrm{HI}}=N_{\mathrm{HI}}(0)\sqrt{1l^2/R_c^2},$$
(23)
where $`N_{\mathrm{HI}}(0)`$ is the HI column density through the cloud centre.
There might be one or more absorbing clouds in the LOS with different velocities with respect to the galaxy centre. The velocity structure follows eq. (17) and the LOS velocity can be calculated easily.
### 2.3 Dark matter satellites
According to the hierarchical clustering scenario, galaxies are assembled by merging and accretion of numerous dark matter satellites of different sizes and masses. As pointed out by Klypin et al. (1999), this ongoing process does not destroy all the accreted satellites. Their paper gives results of satellite population around a big galaxy-size halo by high-resolution cosmological simulations. The VDF (velocity distribution function) of satellites within $`200h^1\mathrm{kpc}`$ and $`400h^1\mathrm{kpc}`$ is
$$n(>V_{\mathrm{cir},\mathrm{sat}})5000(V_{\mathrm{cir},\mathrm{sat}}/10\mathrm{km}\mathrm{s}^1)^{2.75}h^3\mathrm{Mpc}^3$$
(24)
and
$$n(>V_{\mathrm{cir},\mathrm{sat}})1200(V_{\mathrm{cir},\mathrm{sat}}/10\mathrm{km}\mathrm{s}^1)^{2.75}h^3\mathrm{Mpc}^3,$$
(25)
respectively, where $`V_{\mathrm{cir},\mathrm{sat}}=(1070)\mathrm{km}\mathrm{s}^1`$. This number of satellites is roughly proportional to $`(V_{\mathrm{cir}}/220\mathrm{km}\mathrm{s}^1)^3`$. The velocity dispersion of the satellites is of the order of the circular velocity of the central halo. The number of satellites in the models and in the Local Group agrees well for massive satellites with $`V_{\mathrm{cir}}>50\mathrm{km}\mathrm{s}^1`$, but disagrees by a factor of ten for low mass satellites with $`V_{\mathrm{cir}}`$ about $`10\mathrm{km}\mathrm{s}^130\mathrm{km}\mathrm{s}^1`$ (see Klypin et al. 1999 for discusion). Possibly, most of these low mass satellites are dark matter mini-haloes or analogy of high-velocity clouds (HVCs) at distance $`>100\mathrm{kpc}`$ in the halo of Milky Way galaxy. In addition, the distant HVCs are interpreted as gas contained within DM ‘minihalos’. (e.g., Blitz et al. 1999). It is possible that these satellites possess gas but have little or no star formation, so that they are faint and can not be found in optical surveys. The simulation of the survival of DM satellites has included dynamics friction and tidal stripping. Gas in DM satellite also suffers from ram-pressure stripping by hot gas in the central halo. If the surviving DM satellites can accrete gas and the gas is not stripped away, they can contribute to QSO absorption because there could be some fraction of gas in the neutral state with detectable HI column density (Mo & Morris 1994),
$$N_{\mathrm{HI}}6\times 10^{14}\mathrm{cm}^2\left(\frac{V_{\mathrm{cir},\mathrm{sat}}}{30\mathrm{km}\mathrm{s}^1}\right)^4$$
$$\times \left(\frac{10\mathrm{kpc}}{R}\right)^3\left(\frac{f_g}{0.05}\right)^2\frac{T_{4.5}^{3/4}}{J_{21}},$$
(26)
where $`R`$ is the galaxy projected distance (distance from satellite centre to LOS) and $`T_{4.5}=T/10^{4.5}K`$. We choose $`T_{4.5}=1.0`$ here. The spatial distribution of these satellites in the vicinity of the central galaxy is assumed to follow an inverse square law of distance to the centre.
If the population of DM satellites around big central haloes is not predicted correctly by N-body simulations or these satellites do not possess much gas (for example, due to tidal striping, ram-pressure striping, photoevaporation, or supernova-driven ejection), the absorption by satellites will be overestimated.
### 2.4 Galaxy discs
Observations of low redshift damped Ly$`\alpha `$ systems show that some of these systems are possibly not in normal disc galaxies and their host galaxies are ambiguous (e.g., they could be low surface brightness galaxies or faint dwarf galaxies, or failed galaxies which are not detected, see Steidel et al. 1994, Le Brun et al. 1997, Rao & Turnshek 1998). But in general, some damped Ly$`\alpha `$ systems must arise from galaxy discs (e.g, Prochaska & Wolfe 1997, 1998 and references therein). In this paper we only simulate damped systems arising from galaxy discs.
Unfortunately the HI extent of discs is uncertain so far. Several studies involving $`21\mathrm{cm}`$ mapping of galaxy discs have found ‘sharp edges’ where the HI column density falls off dramatically from a few times $`10^{19}\mathrm{cm}^2`$ to an undetectable level ($`4\times 10^{18}\mathrm{cm}^2`$). Such edges have been explained by models where the ionizing level increases rapidly from the inner optically thick to the outer optically thin regime (Maloney 1993; Corbelli & Salpeter 1993; Dove & Shull 1994a). Maloney (1993) assumed an exponential hydrogen disc and used a transition region model to calculate the HI column density of NGC 3198. His results are in good agreement with observations. Other authors have suggested an extended power-law disc in the highly ionized regime (Hoffman et al. 1993; Linder 1998; Linder 1999). However we will take the plausible assumption of an exponential disc extending from the centre to the outer part of a galaxy. For an exponential disc, we adopt the model of Mo, Mao, & White (1998, hereafter MMW). The galaxy disc is assumed to be thin, to be in centrifugal balance, and to have an exponential surface density profile,
$$\mathrm{\Sigma }(R)=\mathrm{\Sigma }_0exp(R/R_d).$$
(27)
Here $`R_d`$, $`\mathrm{\Sigma }_0`$ and $`R`$ are the disc scalelength, central surface density and distance to the centre respectively. Following MMW, we have
$$R_d8.8h^1\mathrm{kpc}\left(\frac{\lambda }{0.05}\right)\left(\frac{V_c}{250\mathrm{km}\mathrm{s}^1}\right)\left[\frac{H(z)}{H_0}\right]^1\left(\frac{j_d}{m_d}\right),$$
(28)
and
$$\mathrm{\Sigma }_04.8\times 10^{22}h\mathrm{cm}^2m_\mathrm{H}\left(\frac{m_d}{0.05}\right)\left(\frac{\lambda }{0.05}\right)^2$$
$$\times \left(\frac{V_c}{250\mathrm{km}\mathrm{s}^1}\right)\left[\frac{H(z)}{H_0}\right]\left(\frac{m_d}{j_d}\right)^2,$$
(29)
where $`m_d`$ and $`j_d`$ are the fixed ratios of disc mass to halo total mass and disk angular momentum of halo total angular momentum respectively. Generally, we choose $`m_dj_d0.05`$ throughout this work without considering the instability of galaxy discs and evolution of these two parameters. $`\lambda `$ is defined as the halo spin parameter, whose distribution is
$$p(\lambda )d\lambda =\frac{1}{\sqrt{2\pi }\sigma _\lambda }exp\left[\frac{ln^2(\lambda /\overline{\lambda })}{2\sigma _\lambda ^2}\right]\frac{d\lambda }{\lambda },$$
(30)
where $`\overline{\lambda }=0.05`$ and $`\sigma _\lambda =0.5`$ (MMW).
A galaxy disc is thought to have a vertical structure. It is a good assumption (expect for very flattened haloes) to ignore the change in the halo density with a height $`Z`$ above the midplane (cf. Maloney 1993). Then in the limit of negligible self-gravity the vertical profile of the gas will be a Gaussian,
$$n_\mathrm{H}(R,Z)=n_\mathrm{H}(R,0)e^{Z^2/2\sigma _h^2}$$
(31)
where the scale height is given by
$$\sigma _h(R)R\frac{\sigma _{zz}}{V_A}.$$
(32)
Here we assume the core radius of the halo to be much smaller than $`R`$ (cf. Maloney 1992). The asymptotic velocity $`V_A`$ is assumed to be of the same order as the halo circular velocity $`V_{\mathrm{cir}}`$. And we take the typical velocity dispersion $`\sigma _{zz}6\mathrm{km}\mathrm{s}^1`$. The midplane density is
$$n_\mathrm{H}(R,0)=\frac{N_\mathrm{H}^{\mathrm{tot}}(R)}{(2\pi )^{1/2}\sigma _h},$$
(33)
where the total hydrogen column density $`N_\mathrm{H}^{\mathrm{tot}}(R)=\mathrm{\Sigma }(R)/m_\mathrm{H}`$. The incident ionizing photons come from the top of the gas disc with a flux $`\varphi _{i,ex}`$ photons $`\mathrm{cm}^2\mathrm{s}^1`$. This photon flux will ionize the gas to a depth $`Z_i`$ at which the column recombination rate equals the ionizing photon flux, i.e.,
$$_{Z_i}^{\mathrm{}}\alpha _{\mathrm{rec}}n_\mathrm{H}^2(R,0)e^{Z^2/\sigma _h^2}𝑑Z=\frac{1}{2}\varphi _{i,ex}.$$
(34)
Here $`\alpha _{\mathrm{rec}}=4.18\times 10^{13}\mathrm{cm}^3\mathrm{s}^1T_{e,4}^{0.72}`$ is the recombination coefficient at a temperature of $`T_{e,4}=T_e/10,000K`$. We assume $`T_e=20,000K`$. In the optically thick regime, the UV ionizing field is incident from one side so that the ionizing photon flux is $`\frac{1}{2}\varphi _{i,ex}`$, while $`\varphi _{i,ex}=5.4\times 10^4I_{\mathrm{ly}}\mathrm{photons}\mathrm{cm}^2s^1`$ and $`I_{\mathrm{ly}}=J_{21}/(0.04\mathrm{erg}\mathrm{cm}^2sr^1\mathrm{Hz}^1\mathrm{s}^1)`$. We define
$$b\frac{\mathrm{\Phi }_{i,ex}}{\sqrt{\pi }\alpha _{\mathrm{rec}}n_\mathrm{H}^2(R,0)\sigma _h}1.36\times 10^3\left(\frac{\sigma _h}{\mathrm{kpc}}\right)T_{e,4}^{0.72}I_{\mathrm{ly}}^1/N_{18}^2,$$
where $`N_{18}=N_\mathrm{H}^{tot}/(10^{18}\mathrm{cm}^2)`$. Then the equation (34) becomes
$$1\mathrm{Erf}(Z_i/\sigma _h)=b,$$
(35)
where Erf is the error function.
When $`b=1`$, $`Z_i=0`$, one can get a critical column density below which the hydrogen will be highly ionized. This critical column density is $``$ a few $`\times 10^{19}\mathrm{cm}^2`$ (see Maloney 1993).
When $`b<1`$, the disc is optically thick and one can derive a $`Z_i`$ from equation (35). The HI in disc has a sandwich structure. If $`Z<Z_i`$, all the hydrogen is assumed to be neutral (the central HI layer) so that the HI fraction is
$$\chi _{\mathrm{HI}}1.$$
(36)
Above the central HI layer, which is the case for $`ZZ_i`$, the hydrogen is highly ionized. Thus assuming ionization equilibrium, the fraction of HI is
$$\chi _{\mathrm{HI}}2n_\mathrm{H}(\alpha +3)T_e^{0.72}I_{\mathrm{ly}}^1,$$
(37)
given the ionizing field is incident from one side of the disc.
When $`b>1`$, the whole disc becomes optically thin and thus the fraction of HI can be estimated as (given the ionizing field is incident from both sides of the disk)
$$\chi _{\mathrm{HI}}n_\mathrm{H}(\alpha +3)T_e^{0.72}I_{\mathrm{ly}}^1$$
(38)
assuming ionization equilibrium (cf. Maloney 1992). Here $`\alpha `$ is the spectral index of the UV background ionizing field, which is $`0.51.5`$.
Once $`\chi _{\mathrm{HI}}`$ is known, we can calculate the total HI column density along a LOS by
$$N_{\mathrm{HI}}=_{LOS}\chi _{\mathrm{HI}}n_\mathrm{H}(R,Z^{})𝑑Z^{},$$
(39)
where the LOS has mid-plane distance $`R_0`$, inclination angle $`\gamma `$ and orientation angle $`\alpha `$. For a point ($`R`$, $`Z^{}`$) along the LOS, we have
$$R(Z^{})=\sqrt{R_0^2+(Z^{}\mathrm{tan}\gamma )^2+2R_0Z^{}\mathrm{tan}\gamma \mathrm{cos}\alpha }.$$
(40)
Thus we can calculate $`N_{\mathrm{HI}}`$ numerically. We also do calculations using Cloudy 90 (Ferland 1996) and find the result is consistent with our calculation of the above sandwich structure of the disc using the analytic methods.
In the optically thick regime, our predicted HI column density is in good agreement with the calculation by Maloney (1993) and the fraction of hydrogen in the HI phase is about 2/3. However as we shall see in the next section, the gas cross section of this regime could be too large and the models predict too many damped Ly$`\alpha `$ systems (DLAs) at low redshift. Thus we assume the fraction of the hydrogen column density in the HI phase to be $`\kappa `$ in this regime and simulate for different values of $`\kappa `$ to compare with the observational number density of DLAs.
### 2.5 Rest frame equivalent width
We ‘observe’ haloes, discs, and satellites by random LOS and obtain HI clouds which contribute to Ly$`\alpha `$ absorption. Each absorbing cloud (halo cloud, disc, satellite) is assumed to have the Voigt profile. The optical depth of a line is
$$\tau _\nu 2.65\times 10^2f_{jk}N_j\varphi (\nu ;\nu _{jk})$$
(41)
for $`h\nu _{jk}kT`$, where $`f_{jk}`$ is the oscillator strength of the line and $`N_j`$ is column density, $`\varphi `$ is the Voigt profile, $`j,k`$ are the lower and upper energy level indexs respectively (for Ly$`\alpha `$ line, $`j=1,k=2`$). The REW, defined in frequency units is,
$$W=(1e^{\tau _\nu })𝑑\nu .$$
(42)
In accordance with observational usage, $`W`$ is defined in wavelength units, so the value must be multiplied by $`\lambda /\nu `$. For the HI Ly$`\alpha `$ line, $`\lambda =1215.670`$Å, and $`f_{jk}=0.4162`$.
Because there may be $`n`$ components in a LOS (a direct sum or a blend of two or more lines), we compute $`\tau (\nu )={\displaystyle \underset{i=1}{\overset{n}{}}}\tau _{\nu ,i}`$ and then calculate $`W`$ numerically. For $`N_{HI}>10^{19}\mathrm{cm}^2`$, which is the case when LOS intersects an optically thick disc, the REW is determined from the column density accurately as (Petitjean 1998)
$$W_r=\sqrt{N_{\mathrm{HI}}/(1.88\times 10^{18}\mathrm{cm}^2)}\mathrm{\AA }.$$
(43)
### 2.6 The predicted line number densities
There are two kinds of basic observational facts for Ly$`\alpha `$ absorbers at low redshift. From spectroscopic surveys in QSO spectra one can derive the line number densities per unit redshift interval, $`(\frac{dN}{dz})`$ (for absorbers with $`W_r0.3`$Å, for metal absorption systems, for LLSs, DLAs). Another kind of observation is the imaging and spectroscopic survey in the QSO fields from which one can relate absorbers with luminous galaxies. In this part of the paper, we compare the predicted $`(\frac{dN}{dz})`$ with that observed to set constraints on model parameters and absorbing components. The only selection criterion for $`(\frac{dN}{dz})`$ is the lower limit of the line width (or HI column densities). However, to compare our results with results of imaging and spectroscopic surveys, it is necessary to study selection effects (such as limitations on apparent magnitude, velocity separation, angular separation) when relating absorbers with luminous galaxies and studying the properties of absorber-galaxy pairs. This study will be presented in § 3.
Several models are simulated with different absorbing components and different parameters in plausible ranges. The parameters used to describe the absorbing components are as follows: $`f_g`$ (the baryon fraction), $`Z`$ (the metallicity, in unit of $`Z_{}`$), $`M_5`$ (the mean mass of halo clouds, in units of $`10^5M_{}`$), $`v_{\mathrm{inf}}`$ (the infall velocity of halo clouds, scaled by the halo circular velocity, $`V_{\mathrm{cir}}`$), $`\kappa `$ (the HI fraction of total hydrogen in the optically thick part of the galaxy disc), the velocity distribution function (VDF) chosen for satellites \[we denote the VDF of satellites following eq.(24) as VDF1 and VDF following eq.(25) as VDF2\], and the flux of the UV background ionizing field $`J_{21}(z)=J_{21}^0(1+z)^2`$ ($`J_{21}^0`$ is the flux at $`z=0`$).
Five types of models are considered: (1) Model A: discs only, (2) Model B: halo clouds only, (3) Model C: satellites only, (4) Model D: disks and halo clouds, (5) Model F: disks, halo clouds and satellites. Model A, B, C are used to predict the fraction of absorption by galaxy discs, halo clouds and satellites respectively. Model D and F are used to see what fraction the plausible combination of absorbing components can explain observational line number densities. The standard model is Model F3 which includes all absorption components and has standard parameters:
$`f_g=0.05,Z=0.3Z_{},M_5=5,v_{\mathrm{inf}}=1.0,J_{21}^0=0.056,\kappa =0.1`$, and VDF2 (the reasons of choosing these values are given below).
Model A has only one parameter, $`\kappa `$ (see below for its meaning) which is listed in Table 1. The detailed parameters of the other models are listed in the second column of Table 2. In the table, we only list those non-standard parameters and assume a $`\mathrm{\Lambda }`$CDM cosmogony. But we also discuss the SCDM cosmogony in some cases (see notation in the table).
To get statistically stable results, we ‘observe’ through sufficient numbers of LOSs (typically several hundreds). The redshift interval for each LOSs is from 0 to 1, and contains approximately 300 galaxies in the simulations with $`V_{\mathrm{cir}}30\mathrm{km}\mathrm{s}^1`$ within a column with a radius of $`400h^1\mathrm{kpc}`$ along a LOS. We simulate 500 LOSs for mode A and 200 LOSs for model B to F (i.e., with $`6\times 10^4`$ galaxies in total) to get stable results. The results are listed in Table 1 and Table 2. The details of the models are as follows:
1. Model A: As mentioned above, the HI fraction of the disc model in the inner part of the disk is $``$ 2/3. This predicts too many low redshift DLAs (see below and Table 1). However some gas in the disc may form stars or may be in molecular phase, only a fraction of the gas in the inner part of disk is in HI phase. So we define this fraction as a free parameter $`\kappa `$, and simulate for a set of $`\kappa `$ to get constraints on it by comparing the predicted DLA number densities with the observational results. Note that we adopt the traditional definition of DLA with $`N_{\mathrm{HI}}2\times 10^{20}\mathrm{cm}^2`$. The observational number density of DLAs derived by Rao et al. (1995) is $`n_{\mathrm{DLA}}(z)=dN_{\mathrm{DLA}}/dz=(0.015\pm 0.004)(1+z)^{2.27\pm 0.25}`$. For a mean redshift of $`z=0.5`$, $`n_{\mathrm{DLA}}(0.5)0.038\pm 0.014`$. The result of the HST QSO Absorption Line Key Project (Jannuzi et al. 1998) at $`z=0.58`$ is about 0.020. Our simulation results for 500 LOSs with total redshift interval of 500 are listed in Table 1. Our predicted number densities of DLAs are $`0.040.06`$, comparable to the observed counterparts, if $`\kappa =0.10.2`$. Models with larger $`\kappa `$ will predicted too many low redshift DLAs. Thus a reasonable value of $`\kappa `$ is 0.1 and we use this as the standard parameter.
2. Model B: Using model B, we perform simulations for different values of $`f_g`$, metallicity $`Z`$, mean cloud mass as well as cloud infall velocity. We only take into account cold gas inside $`r_{\mathrm{min}}`$ of haloes since we assume that there is no cold gas outside this radius. The changing pattern of total absorber number with different parameters is conceivable. For comparison (see Table 2), the results of models B1, B2 and B3 show, that the predicted number density will decline for larger mean cloud masses because there are fewer clouds and thus the covering factor is smaller. The mean cloud mass will be several $`\times 10^5M_{}`$. Cold clouds of such mass can survive from evaporation by heat conduction and gravitational instability, as well as hydrodynamic instability (Mo & Miralda-Escudé 1996). The results of models B3 and B4 show that the predicted number densities do not change significantly for $`Z=0.3Z_{}`$ and $`Z=0.1Z_{}`$. We use $`Z0.3Z_{}`$ hereafter. The results of models B2 and B5 show, that the predicted number density will increase when the gas fraction $`f_g`$ increases (thus increase the total gas cross section). Comparing B2 and B6, if the infall velocities of the clouds is only some fraction of the halo circular velocity $`V_{\mathrm{cir}}`$, the predicted number density increases because there should be more cold gas nowadays in the haloes, which is apparent from eq. (20). The results of model B7 shows, if the UV background ionizing field were lower at $`z<1`$ due to the declining quasar density, more absorbers can be expected.
3. Model C: Using model C, we compare results for different $`f_g`$ and VDF of satellites. The simulation results are as follows (see Table 2): model C1 (C3) predicts more absorbers than model C2 (C4) because of the larger gas fraction \[see eq. (26)\]; model C3 (C4) can produce more absorbers than model C1 (C2) because more satellites are included. Model C3 can predict $`40`$ per cent of the line number density of the observed strong Ly$`\alpha `$ lines. Our population of satellites is chosen from Klypin et al. (1999). If it is true that a big galaxy possesses a large number of satellites and they possess gas, theses satellites may play an important role for strong Ly$`\alpha `$ absorption lines (with $`W_r0.3`$Å).
4. Model D: The predicted number density of $`W_r0.3`$Å lines in model D is similar to those of model B since the contribution by galaxy discs is only a small fraction (see Table 1 & Table 2). Galaxy halo clouds together with galaxy discs cannot explain the majority of observational number densities for the strong lines.
5. Model F: All possible components are considered in this type of model. The predicted absorber number densities are larger than those in halo-only models and satellite-only models (see Table 2). For example, model F3 predicts that the number density of absorbers can be as high as $`55\pm 22`$ per cent of the observed number density of the HST QSO Absorption Line Key Project. Model F5 can predict $`59\pm 23`$ per cent of the observed strong absorption lines. This is in good agreement with the results of LBTW and CLWB. It is easy to understand that the predicted number density should be higher if the UV background is weaker: For example, in model F6, whose UV background is assumed to be five times lower, the predicted fraction can reach 92 per cent, because more gas in the galactic haloes is in the HI phase and the total absorption cross section of satellite-halos becomes larger. This model, however, predicts too many Lyman-limit systems (see below).
6. If we choose a SCDM cosmogony, for example in model B8, D4, F7, the total absorber number density will decrease for two reasons. One reason is, that for a fixed baryon fraction, the cosmic time becomes shorter than in a $`\mathrm{\Lambda }`$CDM model so that less gas can cool (the cooling radius decreases, as does absorbing gas cross section). Another reason is that in a SCDM universe there are fewer galaxies (about 221) along a LOS than in a $`\mathrm{\Lambda }`$CDM universe. In summary adopting a SCDM cosmogony reduces the fraction of absorber number density to only about 28 per cent in model F7, which is too small compared with the HST result. We conclude that models with a SCDM cosmogony can not predict sufficient absorbers at low redshift under the same assumption.
The observed $`(\frac{dN}{dz})`$ of the HST QSO Absorption Line Key Project (Bahcall et al. 1996) for strong Ly$`\alpha `$ lines ($`W_r0.3`$Å) is
$$(18.2\pm 5.0)(1+z)^\gamma ,\gamma =0.58\pm 0.50.$$
The observed $`(\frac{dN}{dz})`$ for Lyman-limit systems is
$$0.5\pm 0.3\mathrm{at}z=0.5$$
or
$$0.7\pm 0.2\mathrm{with}<z_{\mathrm{LLS}}>=0.7$$
(Stengler-Larrea et al. 1995; see also Mo & Miralda-Escudé 1996 for model predictions). The observed $`(\frac{dN}{dz})`$ for damped Ly$`\alpha `$ systems has been given above. The results of LBTW & CLWB show that at least 30 per cent of the absorbers with $`W_r0.3`$Å are related with luminous galaxies. With comparison to these results, we summarize: (1) model B can be ruled out because the predicted $`(\frac{dN}{dz})`$ of strong Ly$`\alpha `$ lines are small and no damped systems can be predicted in such models. (2) model C can be ruled out for lack of damped systems and for insufficient $`(\frac{dN}{dz})`$ of Lyman-limits systems. (3) models D, F2 and F7 can be excluded because they give insufficient $`(\frac{dN}{dz})`$ of strong Ly$`\alpha `$ lines. (4) models B7 and F6 predict too many Lyman-limit systems and should be ruled out also. Thus, we come to conclusion that models F1, F3, F4, F5 are reasonable models, because they predict plausible number densities for strong Ly$`\alpha `$ lines ($`W_r0.3`$Å), Lyman-limit systems, and DLAs.
#### 2.6.1 Velocity spread of absorption systems
The line number densities above are predicted by measuring the overall REW of all the subcomponents of a galaxy/absorber. There are some cases where more than one subcomponent in a galaxy is crossed by the LOS. Sometimes it happens that the velocity spread of a line system is very large. When the velocity spread of a line is less than $`100200\mathrm{km}\mathrm{s}^1`$, the HST key project would possibly have detected just one line, but otherwise this would have been counted as 2 lines. To study this effect, we plot the cumulative distribution of velocity spread of the simulated galaxy/absorber systems for model F3 in Fig.1. As we can see, about 78 (87) per cent of the systems have a velocity spread less than $`100(200)\mathrm{km}\mathrm{s}^1`$. So, if we treat those lines with velocity spread larger than $`200\mathrm{km}\mathrm{s}^1`$ as 2 lines, the predicted line number density will increase by about 15 per cent.
#### 2.6.2 Photoionization flux contributed by galaxies
Hot stars (O, B stars) in a galaxy can contribute some ionizing photons (at $`\lambda <912\AA `$, namely Lyman continuum, hereafter Lyc), which may reduce the neutral fraction in the galactic halo. Here we make a simple estimation about the effect of extra photoionization by star formation in the galaxy itself.
Stellar synthesis models suggest that the number of ionizing photons emitted from a galaxy in a unit time ($`Q(H^0)`$) is related to the star formation rate (SFR) by
$$\mathrm{SFR}=1.08\left(\frac{Q(H^0)}{10^{53}\mathrm{s}^1}\right)M_{}\mathrm{yr}^1$$
(see Kennicutt 1998 and references therein). Of the Lyc photons, only $`14`$ per cent can escape the OB association and enter the diffuse ionized medium (“Reynolds layer”) above the galaxy disk (Dove & Shull 1994b). About 65 per cent of the escaping photons are not absorbed in the HII layer and either escape from the galaxy or are absorbed by additional gas at high latitude (Dove & Shull 1994b). Thus approximately 9 per cent of the Lyc photons can reach the outer halo. However, Leitherer et al. (1995) observed a sample of 4 starburst galaxies and concluded that less than 3 per cent of the ionizing photons can escape. As an approximation, we use an escaping fraction $`f_{esc}=0.05`$. The rate of the escaping Lyc photons is then $`f_{esc}Q(H^0)`$. The ionizing photon flux at galactic distance $`R`$ is
$$\mathrm{\Phi }_{Lyc}10^8f_{esc}Q_{51}R_{\mathrm{kpc}}^2\mathrm{cm}^2\mathrm{s}^1,$$
where $`R_{\mathrm{kpc}}=R/\mathrm{kpc}`$, $`Q_{51}=Q(H^0)/(10^{51}\mathrm{s}^1)`$ and the energy flux <sup>2</sup><sup>2</sup>2If we take the shape of $`J_v`$ into account, the result should be multiplied by some factor. However because all the estimations here are quite uncertain, we omit this factor. is
$$J_\nu =\frac{\mathrm{h}}{4\pi }\mathrm{\Phi }_{Lyc}\left(\frac{\mathrm{\Phi }_{Lyc}}{10^7\mathrm{cm}^2\mathrm{s}^1}\right)\mathrm{erg}\mathrm{cm}^2\mathrm{sr}^1\mathrm{s}^1\mathrm{Hz}^1,$$
where h is Planck’s constant. In dimensionless form we have
$$J_{21}10f_{esc}Q_{51}R_{\mathrm{kpc}}^2.$$
For a normal spiral galaxy like the Milky Way Galaxy, with a SFR of $`12M_{}\mathrm{yr}^1`$, $`Q_{51}`$ is about 100, and $`J_{21}=50R_{\mathrm{kpc}}^2`$. At a distance of 10 kpc, $`J_{21}=0.5`$, which is larger than what we have used ($`J_{21}=0.05`$ at $`z=0`$) by an order of magnitude. At a distance of 30 kpc, $`J_{21}=0.05`$, which is comparable to what we have used. At a distance of 100 kpc $`J_{21}=0.005`$. For most of the halo clouds in our models, the typical distances to the galaxy center are larger than 30 kpc, so the additional photoionization flux by the galaxy itself may be negligible. This is also true for the satellites, because they are located at even larger distances. The typical absorption radius is about 100 kpc. If we make the extreme assumption that all clouds within 30 kpc are fully ionized, the number of lines would be reduced by about 10 per cent. For early type galaxies, the SFR is lower and so the local photoionization can be neglected. The situation may be different for starburst galaxies with SFR bigger than $`10M_{}\mathrm{yr}^1`$, but the number density of starburst galaxies is small and furthermore, the rate of emission of ionizing photons drops quickly in a short period of time (see Nulsen, Barcons, & Fabian 1998 for discussion). To estimate the effect we have made simulations of the model F3 but with an ionizing flux increased by a factor of 2. The value of $`\frac{dN}{dz}`$ for strong lines is reduced to $`8.8`$ from 11.9. We therefore expect that our results will not be affected significantly by the local ionizing sources.
### 2.7 Properties of absorbing galaxies
Our simulations provide additional information about the galaxy/absorber systems. For example, Fig.2 presents the distribution histograms of the projected distances, equivalent line widths, luminosities, absolute magnitudes, circular velocities and redshifts of absorbing galaxies for the standard model F3. The results are summarized as follows:
1. The REW distribution has a peak near 0.3Å, the majority of REWs are between 0.1Å and 1.0Å.
2. About 70 per cent of projected distances are between $`20h^1\mathrm{kpc}`$ and $`200h^1\mathrm{kpc}`$, only 10 per cent are at projected distances less then $`20h^1\mathrm{kpc}`$ and 20 per cent are at projected distances larger then $`200h^1\mathrm{kpc}`$.
3. About 70 per cent of the luminosities are between $`0.1L_B`$ and $`1.0L_B`$, which implies absorbing galaxies are relatively luminous galaxies in the model.
4. At least 80 per cent of the circular velocities are between $`100\mathrm{km}\mathrm{s}^1`$ to $`300\mathrm{km}\mathrm{s}^1`$, the fraction of absorbers with circular velocities less than $`100\mathrm{km}\mathrm{s}^1`$ is only about 15 per cent.
5. The number density for LLS is about 0.69 which is in good agreement with observations (Stengler-Larrea et al. 1995).
6. The distribution with redshift is flat and the average $`n(z)`$ for absorbers with $`W_r0.3\mathrm{\AA }`$ is about 11.9 which can account for about 55 per cent of the sources observed by the HST.
Note that in the lowest right panel of Fig.2, the predicted number density is almost independent of redshift. This result allows us to average the number density over the whole redshift interval in Table 2. In observation, the number density of strong Ly$`\alpha `$ line evolves slowly with redshift at low redshift ($`z<1`$). It is valuable to investigate the correlations of REW versus projected distance, galaxy luminosity, galaxy/absorber redshift and study the fractions of absorbers produced by different morphological types of galaxies.
#### 2.7.1 Projected distance
We investigate the anti-correlation between REW, $`W_r`$ and projected distance of LOS to galaxy centre, $`\rho `$. A power-law relation is adopted,
$$\mathrm{log}W_r=\alpha \mathrm{log}\rho +C,$$
(44)
where $`\alpha `$ is the slope and $`C`$ is constant. The results are summarized in Table 2. In the table, we list results for typical models B2, C1, C3, D1, F1, F3 (the standard model), and F5. In model Fs, we still get an apparent anti-correlation and a slope of $``$ -0.4. Note that in all models the majority of absorption lines by halo clouds and/or satellites definitely decides the character of the anti-correlation (for example, see Fig.3). The difference in the results found in $`r_s`$ and $`r_p`$ is reasonable because these three types of absorption components which do not necessarily have the same pattern of dependence. This difference has been noticed before (Le Brun et al. 1996).
We plot the distribution of $`W_r`$ and $`\rho `$ for some models. Panels in Fig.3 are for the results of model B2, C1, D1, F3 respectively.
Small galaxies may have almost no cold gas because gas could be heated by supernova explosions or all the gas may have cooled to form stars, and there is no more gas to accrete. Thus these galaxies cannot contribute to absorption. In practice, if we exclude those with $`V_{\mathrm{cir}}100\mathrm{km}\mathrm{s}^1`$ (as we can see in model F3, there are about 80 per cent of the absorbers with circular velocity of $`100\mathrm{km}\mathrm{s}^1`$ or more), there will be a stronger anti-correlation and a steeper slope of the linear fit at a highly significant level for almost all models. For model Fs, typically the Pearson co-efficient, $`r_p`$ is $``$ -0.5 and the slope is $``$ -0.5 at very high significance. This can be seen clearly from the results for ‘sample b’ in Table 3.
In summary, our predicted anti-correlation is significant but the dependence of $`W_r`$ on $`\rho `$ is not as strong as that of CLWB. However, as pointed out by Tripp et al. (1998), selection effects could artificially tighten this anti-correlation. We will discuss this problem in the next section.
#### 2.7.2 Galaxy luminosity
As suggested by CLWB, there is a power-law relationship between $`W`$ and $`\rho `$ and $`L_B`$,
$$\mathrm{log}W_r=\alpha \mathrm{log}\rho +\beta \mathrm{log}L_B+C.$$
(45)
We apply the analysis to model F1, F3 and F5. The results are listed in Table 4. For example, for model F1, the analysis for sample b yields, $`C=1.21`$, $`\alpha =0.79`$, $`\beta =0.19`$. We can determine the absorption radius of a galaxy $`rL_B^t`$ ($`t=\beta /\alpha `$) with $`t0.240.32`$. The value of $`t`$ is comparable but smaller than 0.37 which was derived by CLWB. However it is similar to that derived from MgII obsorbers (Bergeron & Boissé 1991; Bergeron et al. 1992; Le Brun et al. 1993; Steidel 1995). Again, selection effects could lead to misleading conclusions (see discussion in § 3.3).
#### 2.7.3 Absorber redshift
We also analyse the dependence of the line equivalent width on absorber redshift assuming
$$\mathrm{log}W_r=\alpha \mathrm{log}\rho \gamma \mathrm{log}(1+z)+C,$$
(46)
and
$$\mathrm{log}W_r=\alpha \mathrm{log}\rho +\beta \mathrm{log}L_B\gamma \mathrm{log}(1+z)+C.$$
(47)
We apply the analysis to model F1, F3 and F5. The results are listed in Table 5 and Table 6. In summary, the relationship between REW and projected distance together with absorber redshift is marginally superior (with larger $`|r_p|`$ or $`|r_s|`$) to the relationship between REW and the projected distance but marginally inferior (with smaller $`|r_p|`$ or $`|r_s|`$) to the relationship between REW and projected distance accounting for $`L_B`$, and the anti-correlations between REW and projected distance accounting for $`L_B`$ together with $`z`$ is superior to the relationship between REW and projected distance accounting for $`L_B`$. For the analysis of eq.(47), we have $`\alpha 0.480.80`$, $`\beta 0.140.21`$, $`\gamma 0.490.61`$. This means the absorption radius of a galaxy $`rL_B^t(1+z)^u`$ ($`t=\beta /\alpha `$, $`u=\gamma /\alpha `$) with $`t0.240.32`$ and $`u0.761.02`$. Our result of dependence on absorber redshift is different with the result of CLWB. CLWB concluded that REW do not depend on absorber redshift. However as we will discuss below, selection effects should be considered in the imaging and spectroscopic survey and the total redshift interval of LOS in observation may be not large enough to determine the relations.
#### 2.7.4 Covering factor
From eq. (47), and the results in Table 6, we can estimate the average absorption radius and covering factor from redshift of 0 to 1. For example, for ‘sample a’ (see notation in Table 3 for its meaning) of model F1, absorbers with REW larger than 0.3Å follow the relation
$$\frac{r}{r_{}}=\left(\frac{L_B}{L_B}\right)^t(1+z)^u,$$
(48)
where $`r_{}=0.23h^1\mathrm{Mpc}`$ and $`t=0.32`$, $`u=0.86`$. Thus we can calculate the total number by integration
$$N_{\mathrm{total}}=_0^1\pi r_{}^{}{}_{}{}^{2}\frac{c(1+z)^{22u}}{H(z)}C_l(z)𝑑z$$
$$\times \varphi ^{}\mathrm{\Gamma }(1+2ts,L_{B_{\mathrm{min}}}/L_B),$$
(49)
where $`C_l(z)`$ is the covering factor and $`\mathrm{\Gamma }`$ is the incomplete gamma function. Note that for these absorbers with $`W_r0.3`$ Å within $`r_{}`$, it is not necessary that the covering factor is always larger than one, that is, for a LOS with a large projected distance it is not always possible to find an absorber with $`W_r0.3`$ Å. If we choose $`s=1.1`$, $`L_{B_{\mathrm{min}}}/L_B=0.007`$, $`\varphi ^{}=0.027h^3\mathrm{Mpc}^3`$ then we get $`N_{\mathrm{total}}=20.4\times F(z)`$, where $`F(z)=_0^1\frac{(1+z)^{0.28}}{\sqrt{\mathrm{\Lambda }_0+\mathrm{\Omega }_0(1+z)^3}}C_l(z)𝑑z`$. Comparing with the total number listed in Table 2, $`N_{\mathrm{total}}`$ should be 7.3, the average value of $`F(z)`$ should be about 0.36. However the integration $`_0^1\frac{(1+z)^{0.28}}{\sqrt{\mathrm{\Lambda }_0+\mathrm{\Omega }_0(1+z)^3}}𝑑z`$ is about 0.85. Thus the average covering factor within $`230h^1\mathrm{kpc}`$ should be $``$ 0.42 if $`C_l(z)`$ can be treated roughly as a constant. The effective gas absorption radius should be $`\sqrt{0.42}\times r_{}150h^1\mathrm{kpc}`$. Therefore the covering factor within $`250h^1\mathrm{kpc}`$ is $`(150/250)^20.36`$.
Our predicted covering factor is in good agreement with the LBTW paper. In that paper, almost every galaxy with $`\rho <70h^1\mathrm{kpc}`$ gives rise to absorption, about five of 10 galaxies with $`70h^1\mathrm{kpc}<\rho <160h^1\mathrm{kpc}`$ give rise to absorption, and just one of 9 galaxies with $`\rho >160h^1\mathrm{kpc}`$ give rise to absorption. Thus the covering factor within $`250h^1\mathrm{kpc}`$ is about 0.31. However our predicted effective absorption radius ($`150h^1\mathrm{kpc}`$) is a bit smaller than the $`174h^1\mathrm{kpc}`$ derived by CLWB. The reason of this difference could be the large covering factor used in CLWB. Independently, Bowen et al. (1996) derive a covering factor of 0.50 within $`\rho <160h^1\mathrm{kpc}`$ (three of six galaxies give rise to absorption).
#### 2.7.5 Galaxy morphological type
CLWB conclude that galaxies which produce Ly$`\alpha `$ absorption systems span a wide range of morphological types from elliptical or S0 galaxies through late-type spiral galaxies. Consistently the models show every type of galaxy can produce Ly$`\alpha `$ absorption systems. The predicted fractions of absorbing galaxies in spiral, S0, elliptical galaxies are defined as f1, f1, f3 respectively and listed in Table 2. As we can see, the fractions vary from model to model. In model B and D, about half of the absorbers arise in spiral galaxies. In model C, this fraction is only about one third because E/S0 galaxies possess more satellites than spirals. In model F, about 35 per cent of the absorbers are produced by spirals. It is, however, not clear, if there are many clouds inside haloes of ellipticals/S0 galaxies and how many absorbing systems can arise from tidal tails related with spirals. If we exclude absorption by halo clouds of ellipticals and S0 galaxies, for instance, comparing models B2 and F1, the fraction by spirals is estimated to be $`(7.3\times 38\%)/(7.33.7\times 54\%)52`$ per cent. A crude estimate using Fig.2 of CLWB shows that this number is about 70 per cent, however the information on galaxy morphology in observations may be inadequate to tackle this problem. We suggest that these numbers should be investigated further to see whether the gas absorption sections of different types of galaxies are the same or not.
## 3 Selection effects: Mock Imaging And Spectroscopic Surveys
In section 2 we have studied $`(\frac{dN}{dz})`$ and the overall properties of absorbers limited only by the lowest line width of 0.1Å, but without considering selection criteria. In carrying out comparisons with results of imaging and spectroscopic surveys, it is absolutely necessary to construct absorber catalogues with observational selection criteria and investigate the possibility of mis-identification of absorbers (i.e., optically unseen absorbers mis-matched with a bright neighbour galaxy).
### 3.1 Selection criteria of absorber-galaxy pairs
Our galaxy-absorber pair selection criteria are similar to those of CLWB. We only select absorbers with a rest frame line width
$$W_r0.1\mathrm{\AA },$$
(50)
and with the B band luminosity satisfying
$$m_B24.3.$$
(51)
Similar to LBTW, we only select those galaxies within angular distances <sup>3</sup><sup>3</sup>3See LBTW’s Fig.2. for the variation of projected distance threshold with redshift for an angular distance threshold $`\theta =1^{}3`$. to the QSOs satisfying
$$\theta 1^{}3.$$
(52)
A small velocity separation between absorber and galaxy centre, $`v=|cz_{\mathrm{gal}}cz_{\mathrm{abs}}|<500\mathrm{km}\mathrm{s}^1`$, is required to relate an absorber with a luminous galaxy (Lanzetta et al. 1997; CLWB; Bowen et al. 1996). This small velocity separation excludes almost all random galaxy-absorber pairs.
### 3.2 Mis-identification
LBTW argued that it is unlikely that the absorbing gas arises in dwarf companions of the luminous galaxies because no such dwarf galaxy was found in the LOS toward QSOs at redshifts $`z<0.2`$ down to a luminosity of $`0.05L_{}`$ <sup>4</sup><sup>4</sup>4In our models, a spiral galaxy with $`L_B0.05L_B`$ may have $`V_{vir}55\mathrm{km}\mathrm{s}^1`$ and $`m_B23.2`$ at $`z0.2`$ in a $`\mathrm{\Lambda }`$CDM cosmology, for example.. But van Gorkom et al. (1996) and Hoffman et al. (1998) have located faint dwarf galaxies at the redshifts of a few low-$`z`$ Ly$`\alpha `$ lines. In addition, as we can see in §2.1, the lower limit in the luminosity function can be as low as $`0.007L_{}`$. Satellites are even fainter with typical circular velocity around $`30\mathrm{km}\mathrm{s}^1`$. Of course, at very low redshift, these faint satellites can be identified in imaging surveys with a large field of view. However, an angular threshold of $`1^{}3`$ at $`z=0.1`$ means in general a distance of $`100h^1\mathrm{kpc}`$ so that many of these satellites are excluded from pair samples because they typically have large distances to the central galaxy.
As suggested by Tripp, Lu, & Savage (1998), the Ly$`\alpha `$ lines could be due to undetected faint dwarf galaxies that are clustered with the observed luminous galaxies. This kind of Ly$`\alpha `$ lines arise at small projected distances, but the host galaxies are too faint to be identified in galaxy imaging surveys, especially at modest to high redshifts. These un-seen absorbers (optically uncatalogued galaxies) can cause erroneous identifications, since one could make a mistake to relate the corresponding line with a nearby luminous galaxy at a larger projected distance. We simulate nearby luminous galaxies around a central galaxy by the two-point correlation function of normal galaxies, which is
$$\xi (r)=\left(\frac{5h^1\mathrm{Mpc}}{r}\right)^\gamma ,$$
(53)
where $`r`$ is the separation of two galaxies and $`\gamma 1.8`$. Thus the galaxy number in a volume with radius of $`h^1\mathrm{Mpc}`$ is
$$N(r,z)=n_c(1+z)^3_0^14\pi r^2[1+\xi (r,z)]𝑑r,$$
(54)
where $`n_c`$ is the co-moving galaxy density described in §2.1, and $`z`$ is the redshift of the central galaxy. The result of the Canada-France Redshift Survey shows a strong redshift evelution of the galaxy two-points correlation function which is
$$\xi (r,z)=\xi (r)(1+z)^{(3+ϵ)},$$
(55)
where $`ϵ02`$ at $`z<1.3`$ (Le Fèvre et al. 1996). Because $`ϵ`$ is uncertain observationally, we use $`ϵ=1.5`$ for simplicity (Shepherd et al. 1997). Our results are not sensitive to the choice, because the redshift range covered is small. The luminosity distribution of these galaxies is consistent with the luminosity function and they are distributed around the central galaxy following $`\xi (r,z)`$. The apparent magnitude of these galaxies is calculated from eq. (5). The distribution of the pairwise velocity is a Gaussian with dispersion $`\sigma =400\mathrm{km}\mathrm{s}^1`$ (Efstathiou 1996; Mo, Jing & Börner 1993; Jing, Mo & Börner 1998).
We define a galaxy-absorber pair intimately linking the absorber with a bright galaxy (an absorbing galaxy with apparent magnitude brighter than the luminosity limit, i.e. $`m_B24.3`$) as a ‘physically associated pair’ or a ‘physical pair’ for simplicity. On the contrary a mis-matched galaxy-absorber pair is called ‘spurious pair’. The method to find a mis-matched pair is as follows: When there is a faint absorber (whose apparent magnitude is fainter than the luminosity limit), its neighbours will be simulated to see whether there is a nearby bright galaxy (brighter than the luminosity limit) with larger projected distance (however not larger than $`400h^1\mathrm{kpc}`$). For a positive result, this bright neighbour will be chosen to pair with the absorption line arising from the faint absorber and the new projected distance will be chosen. For a negative result in the search of a bright neighbour, we classify the galaxy-absorber pair as a ‘missing pair’. We define those bright ‘physical pairs’ and ‘spurious pairs’ as ‘bright pairs’. We also call a pair outside a certain angular distance threshold as a ‘missing pair’. For instance, at very low redshift some ‘bright pairs’ have large angular separations to a QSO LOS.
### 3.3 Effect on the $`W_r\rho L_B`$ relations
When comparing the $`W_r\rho L_B`$ relations for simulated galaxy/absorber pairs to the observations, the selection effects mentioned above may have impacts on the statistics. For example, there could be some ‘spurious’ galaxies at large impact parameters within the redshift window of the absorbers. In some of the surveys carrried out, the sky area surveyed is so large that there is always one galaxy (by chance) within the 500 $`\mathrm{km}\mathrm{s}^1`$ of the redshift. This may cause mis-identification of the absorbing galaxy and add noise to the correlations. On the other hand, faint absorbing galaxies without bright neighbours within the 500$`\mathrm{km}\mathrm{s}^1`$ window will not be listed in the catalogues, which may reduce the noise and strengthen the correlations. In observations, the correlations are the results of the balance of these two effects. In principle, both effects may lead to misleading conclusions about the average galaxy/absorber distance.
We simulate 200 sight lines as in section 2 using methods described above. For model F1, our result shows that, if all bright galaxy/absorber pairs are used , we get $`\mathrm{log}W_r=0.860.59\mathrm{log}\rho +0.20\mathrm{log}(L_B/L_B)`$ and $`\frac{r}{r_{}}=\left(\frac{L_B}{L_B}\right)^{0.34}`$, where $`r_{}=220.7h^1\mathrm{kpc}`$ ($`W_r0.3`$Å). In contrast, for all physical pairs, we have $`\mathrm{log}W_r=0.890.63\mathrm{log}\rho +0.20\mathrm{log}(L_B/L_B)`$ and $`\frac{r}{r_{}}=\left(\frac{L_B}{L_B}\right)^{0.32}`$, where $`r_{}=174.8h^1\mathrm{kpc}`$ ($`W_r0.3`$Å). As we can see, a larger $`r_{}`$ is derived if spurious pairs at large distances are used.
### 3.4 Results for mock observations of known QSO LOSs
We simulate observations for 10 QSOs at the redshifts listed in the CLWB paper. The redshifts of these QSOs are 0.200, 0.264, 0.329, 0.371, 0.513, 0.534, 0.574, 0.616, 0.719, 0.927. The total redshift interval is about 5. Because of the proximity effect, galaxies within $`3000\mathrm{km}\mathrm{s}^1`$ of the quasar redshift are excluded. With 100 mock observations, we get the statistical number of galaxy-absorber pairs and the statistical properties of galaxy-absorber pairs.
The total number of ‘physical pairs’ with $`W_r0.3`$Å for the ten QSOs is
$$36.0\pm 6.3,\mathrm{\hspace{0.17em}57.8}\pm 7.5,\mathrm{\hspace{0.17em}64.2}\pm 7.0$$
for model F1, F3, F5 respectively. After applying the selection criteria, the total number of ‘bright pairs’ with $`W_r0.3`$Å is
$$21.0\pm 4.8,\mathrm{\hspace{0.17em}26.1}\pm 4.8,\mathrm{\hspace{0.17em}29.9}\pm 5.3$$
for model F1, F3, F5 respectively. The predicted ‘bright pair’ numbers are in good agreement with that of CLWB paper in which there are 26 galaxies giving rise to absorption.
If the model prediction is correct, we argue that the galaxy imaging survey at the faint end could be incomplete. This can be seen from the lower-left panel of Fig.4, in which the dotted line shows our predicted distribution of apparent magnitudes of ‘bright pairs’ and the solid line represents the distribution for all ‘physical pairs’, while the dashed line represents the distribution of apparent magnitudes of pairs from the CLWB paper. As we can see, considerable numbers of the predicted $`m_B`$ of pairs are fainter than 25, and in the range of 23 to 25, our predicted ‘bright pairs’ are more numerous than the observed ones. In the lowest-right panel of Fig.4, the predicted number of ‘physical pairs’ in redshift bins at $`z<0.5`$ is larger than that of CLWB, and comparable with that of CLWB at $`z>0.5`$, while for ‘bright pairs’, the number is comparable with that of CLWB at $`z<0.5`$ but is lower than that of CLWB at $`z>0.5`$. However the observed number density at high redshift is uncertain because the redshift interval there is quite small. Of course, if the difference is real, there are some implications. For example, the luminosity function at $`z>0.5`$ would be higher than at low redshift and galaxies could be brighter in the blue band because of intense star formation. In conclusion, from the number of pairs and distribution of pair redshifts, the model predictions are consistent with current observations.
In Fig.4, we also plot the distributions of $`\rho `$, $`W_r`$, $`L_B/L_B`$, $`M_B`$ as well as $`V_{\mathrm{cir}}`$ and compare the results with those of CLWB except for circular velocity. We apply the Chi-Square test to see if our predicted distributions are similar to those of CLWB (the data are from their table 4). The null hypothesis that the data sets are similar has probabilities of 97.0 per cent, 72.4 per cent, 65.7 per cent, 97.9 per cent, 31.1 per cent for $`W_r`$, $`\rho `$, $`L_B/L_B`$, $`M_B`$, $`z`$ respectively.
We give statistical results of the anti-correlation of $`W_r`$ versus $`\rho `$ for all mock runs. Fig.5 contains the histograms of the Spearman rank-order coefficient $`r_s`$, the Pearson coefficiet $`r_p`$, and the zero point C, the slope $`\alpha `$ of the linear fit as in eq. (44). In the figure, the solid lines are for ‘bright pairs’ satisfying the selection criteria, the dashed lines are for ‘bright pairs’ with $`V_{\mathrm{cir}}100\mathrm{km}\mathrm{s}^1`$, and the dotted lines are for all ‘physical pairs’. Statistically the distributions of the slope $`\alpha `$ shift from small to large values for ‘physical pairs’, ‘bright pairs’, ‘bright pairs’ with $`V_{\mathrm{cir}}100\mathrm{km}\mathrm{s}^1`$ and so on. Similarly, there are also shifts of lines in other panels. These shifts mean that selection effects can statistically strengthen the anti-correlation of $`W_r`$ and $`\rho `$. We have to point out that this strengthening only has a statistical meaning and does not necessarily occur for every specific run in our simulations, because in some cases selection effects may also weaken the anti-correlations. Available results from observations (CLWB; Tripp et al. 1998) are also shown with $`1\sigma `$ error bars. Obviously some simulation runs can give consistent results compared with the observations. Note that the results of mock runs have considerable scatter, which may imply, as will be discussed later in §4, that in the models the same total redshift interval as in present observations is not adequate to predict the real correlation.
As mentioned above, the outcome of each run could differ case by case. Some examples of mock observations are given in Fig.6, Fig.7, Fig.8.
The results for the anti-correlation of REW versus projected distance are shown in Fig.6 (for run No.20). All the real galaxy-absorber pairs (‘physical pairs’) are drawn in the left panel. For the $`108`$ ‘physical pairs’ (of which $`54`$ pairs have $`W_r0.3`$Å), we get the Spearman rank-order coefficient, $`r_s=0.41`$ (with significance level $`4.6\sigma `$) and the Pearson coefficient $`r_p=0.62`$ (with significance level $`8.0\sigma `$) and best fit line $`\mathrm{log}W_r=(0.85\pm 0.16)(0.60\pm 0.07)\mathrm{log}\rho `$. The galaxy-absorber pairs after applying selection criteria are drawn in the right panel. For $`49`$ ‘bright pairs’ (of which $`30`$ pairs have $`W_r0.3`$Å), we get $`r_s=0.58`$ (with significance level $`4.9\sigma `$), $`r_p=0.75`$ (with significance level $`7.7\sigma `$) and $`\mathrm{log}W_r=(1.63\pm 0.26)(0.95\pm 0.12)\mathrm{log}\rho `$.
In comparison, CLWB give
$$\mathrm{log}W_r=(1.34\pm 0.22)(0.93\pm 0.13)\mathrm{log}\rho $$
and Tripp et al. (1998) give
$$\mathrm{log}W_r=(1.32\pm 0.20)(0.80\pm 0.10)\mathrm{log}\rho $$
by adding more LOSs. The results of ‘bright pairs’ for this run are in good agreement with those of CLWB and Tripp et al. (1998). As we can see, selection effects do strengthen the anti-correlation in this specific run. However for some runs, selection effects do not strengthen the anti-correlation at all, because they only have a statistical meaning in the simulations.
We analyse the relation between REW and galaxy luminosity using eq. (45) for the run. The results are shown in Fig.7. In the left panel, $`C`$, $`\alpha `$, $`\beta `$ are $`0.96\pm 0.16`$, $`0.63\pm 0.07`$, $`0.13\pm 0.04`$ respectively. In the right panel, they are $`1.80\pm 0.26`$ of $`1.02\pm 0.12`$, $`0.18\pm 0.09`$ respectively. Again, in comparison, CLWB give
$$\mathrm{log}W_r=(1.78\pm 0.20)(1.02\pm 0.12)\mathrm{log}\rho $$
$$+(0.37\pm 0.10)\mathrm{log}(L_B/L_B).$$
As we can see, the zero point $`C`$ and $`\alpha `$ of our prediction for ‘bright pairs’ in run No.20 are in good agreement with those of CLTW. However our results for $`\beta `$ in this run are less than that of CLWB (still within the $`2\sigma `$ standard deviation), but in agreement with the value of $`0.10.2`$ suggested by Bowen et al. (1996).
There is a substantial number of missing pairs in our mock observations. Fig. 8 gives the results for the run. As we can see, there are a number of bright ‘physical pairs’ with angular distances to the LOS outside the threshold of $`1^{}3`$ and some faint ‘physical pairs’ without bright neighbours located within $`400h^1\mathrm{kpc}`$. These pairs are defined as missing pairs and could not be listed in optical catalogues of absorbers. Note that there are also a few spurious pairs with angular distance outside the threshold.
## 4 Discussion
As we remarked in the last section, if the total redshift interval is small, the results of correlations can have large statistical deviations. In order to investigate this, we randomly choose 40 groups of LOS with the same number of LOS in each group from the simulation. We calculate the correlation coeffecient, the slope and zero point of the linear fit for every group. At last the average values over all the groups as well as the standard deviations can be calculated. Then we increase the number of LOSs in each group. These values will change until the total redshift interval is large enough to get stable values with small deviations. We make plots of the dependence, from which we can determine how large a redshift interval is adequate for an observational sample. The results for model F3 for ‘bright pairs’ are presented in Fig.9. Clearly for ‘bright pairs’ a total redshift interval $`10`$ is necessary to get statistically accurate value of $`r_p`$, the Spearman coefficient, and $`\alpha `$, the slope of the linear fit in eq. (44), with standard deviation less than 0.1. If we want to determine the zero point (i.e. with standard deviation less than 0.2), the total redshift interval should be about 20 for model F3. This implies that results of the anti-correlation of $`W_r`$ versus $`\rho `$ by different surveys may differ from each other if the total redshift interval in the survey is not large enough.
In the models, only one third of the absorbers reside inside galactic haloes and two thirds of them are satelltes around central halos. This picture may reconcile different conclusions by various authors (Morris et al. 1993; LBTW; CLWB; Bowen, Blades, & Pettini 1996; Le Brun et al. 1996; Tripp et al. 1998; Impey at al. 1999). Firstly, the models predict a large absorption radius and a reasonable covering factor. Secondly, the predicted absorbers are still closely related with galaxies. Furthermore, it is possible that there are still a number of satelltes around central haloes even at distance $`>400h^1\mathrm{kpc}`$. Thus if a catalogue of galaxy-absorber pairs includes those pairs with large projected distances, the anti-correlation of $`W_r`$ versus $`\rho `$ could be weaken.
As listed in Table 2, our prefered models can explain up to $``$ 55 per cent (and even more if an absorption system with large velocity spread is counted as 2 lines as would be the case with the HST spectral resolution) of the HST observed counterpart for lines wider than 0.3Å. Can our models predict more absorbers? Actually the evolution of the galaxy luminosity functions with redshift can change our predicted absorber number density. Because the numbers of aborbers in our simulations are directly proportional to the galaxy number density, our predicted number of absorbers can be higher at higher redshift if the galaxy number density there is higher. For example, if we chose the AUTOFIB luminosity function at $`0.35<z<0.75`$ whose $`\varphi ^{}`$ is about 1.5 times the $`\varphi ^{}`$ at $`0.02<z<0.15`$, the predicted number density could be as high as 17.8, which may account for 80 per cent of the observed counterparts. Furthermore, the higher galaxy density at higher redshift may also increase the absorber density at higher redshift in Fig.4 and solve the possible discrepancy in high-redshift absorber number between our prediction and that of CLWB.
Caution should be taken with our results, not only because there could be some alternative Ly$`\alpha `$ absorption arising in other parts of the galaxies, but also because it is unclear under present circumstances whether there are many satellites in the vicinity of big central galaxies and whether they possess gas (see Klypin et al. 1999 for further discusion; see also Bullock, Kravtsov, & Weinberg 2000; Charlton, Churchill, & Rigby 2000). After all we did not include all possible absorption components related with galaxies. For example, Morris & van den Bergh (1994) suggested that a significant fraction of weak lines could arise from pressure-confined tidal debris in the enviroment of small groups and clusters of galaxies (cf. Mo 1994). Tidal debris can increase the total gas cross section so as to increase the absorber number density and the corresponding absorption line can have a large projected distance of $`>100h^1\mathrm{kpc}`$. However the absorption line number arising from tidal tails depends on the unknown gas cross section and the generally unknown lifetime of the gas in tidal tails. Thus it is not possible to detemine exactly what fraction of absorbers arises in tidal debris. Another possibility is that some low surface brightness galaxies could possess huge gaseous haloes or discs which can also give rise to absorption. In addition, galactic wind is also another possible source (Wang 1995). Of course, the absorption by the IGM still could play an important role.
In observations, it is difficult to assign a galaxy to an absorber counterpart, because an imaging survey of galaxies is never quite complete down to the faint end. Also absorbing galaxies may be outside of the angular extent of the survey. More LOSs to QSOs with higher resolution of the UV spectroscopy and more complete imaging surveys are necessary to investigate the physical origin and environment of the absorbers. At very low redshifts, it is possible to identify satellites with $`V_{\mathrm{cir}}30\mathrm{km}\mathrm{s}^1`$ in optical imaging surveys, and then we can examine whether these satellites can give rise to Ly$`\alpha `$ absorption lines. To discriminate models, it is also important to get more physical information about the absorbing components, such as size, temperature, metallicity, ionizing parameter, rather than only informations about $`W_r`$ and $`\rho `$. Observations of multi-LOS (QSO pairs or lensed images) can provide useful tools to get more insight into absorbers (see Rauch 1998 and references therein). Furthermore, the possible detection of line emission from extended gas in galactic haloes may also help to determine the properties of the gaseous haloes \[Ćirković, Bland-Hwathorn, & Samurović, 1999\].
## 5 Summary and Conclusions
In this paper, we present results of Monte-Carlo simulations of Ly$`\alpha `$ absorption line systems at redshift $`z<1`$. To get constraints on the parameters, we simulate a set of models with different absorption components and various parameters. We compare the predicted absorption line densities for strong Ly$`\alpha `$ lines, Lyman-limit systems as well as damped Ly$`\alpha `$ systems with observations. From these comparisons, some models can be excluded. In summary:
1. Models with a single absorption component (galaxy disc, or cold clouds in a galactic halo, or satellites) or models with disk and clouds in a galactic halo, cannot explain the observed number densites for strong Ly$`\alpha `$ lines, for Lyman-limit systems and for damped Ly$`\alpha `$ systems at low redshift.
2. Models with all three components (galaxy disc, cold clouds in a galactic halo, and satellites) can explain up to 60 per cent of the observed number density for strong Ly$`\alpha `$ lines at low redshift. These models can also predict reasonable number densities of Lyman-limit systems and damped Ly$`\alpha `$ systems at low redshift.
3. The fraction of the line number density for strong Ly$`\alpha `$ lines due to satellites is $`40`$ per cent more than that due to clouds in galactic haloes (which is $`20`$ per cent) by about a factor of 2. The exponential galaxy discs can only account for a small amount of strong Ly$`\alpha `$ lines. If indeed there are large numbers of satellites surrounding big central galaxies and they possess gas, these satellites may play an important role for strong Ly$`\alpha `$ lines at low redshift.
4. The predicted $`(\frac{dN}{dz})`$ for Lyman-limit systems due to cold clouds in galactic haloes is $`0.4`$, which can account for most of the observed Lyman-limit systems. The line number density of Lyman-limit systems due to satellites is only $``$ 0.1, which is four times smaller than that due to clouds in galactic haloes.
The properties of the predicted absorbers, such as REW, projected distance, galaxy luminosity, circular velocity and absorber redshift have been analysed. The predicted dependence of line width on projected distance is $`W_r\rho ^\alpha `$ with $`\alpha 0.40.6`$, rather than $`\alpha 0.80.9`$ (cf. CLWB; Tripp et al. 1998). This predicted anti-correlation is weaker than the observed one because we include all faint absorbers (with apparent magnitude fainter than magnitude limit in optical surveys) which have small impact parameters and we include absorbers with impact parameters larger than $`200h^1\mathrm{kpc}`$. Other correlations of REW versus luminosity and/or absorber redshift have also been investigated. In general, if we assume $`W_r\rho ^\alpha L_B^\beta (1+z)^\gamma `$, the analysis gives $`\alpha 0.5,\beta 0.15,\gamma 0.5`$. This means the average absorption radius of a galaxy $`rL_B^t(1+z)^u`$ with $`t0.3`$ and $`u1.0`$. The average covering factor within $`250h^1\mathrm{kpc}`$ is estimated as $`0.36`$ which is in good agreement with previous results (LBTW). The effective absorption radius is estimated to be $`150h^1\mathrm{kpc}`$, which is consistent with the observational result $`170h^1\mathrm{kpc}`$ derived by CLWB.
To compare with results of imaging and spectroscopic surveys, it is necessary to study selection effects. Selection effects have impacts on the statistics of the galaxy/absorber properties. The present of ‘spurious’ galaxies at larger impact parameters within the redshift window of the absorber and the ‘missing’ of faint absorber at small impact parameter may lead to misleading conclusions that the average galaxy/absorber separation is very large. Our simulations show that this is indeed the case. We construct mock observations with the same known QSO LOS as CLWB applying selection criteria which are similar. By an adequate number of mock runs, the total number of galaxy-absorber pairs can be predicted and the correlations mentioned above can be analysed. The predicted number of galaxy-absorber pairs with $`W_r0.3`$Å is $`26\pm 5`$, in good agreement with CLWB ($`26`$). The analysis of the anti-correlation between $`W_r`$ and $`\rho `$ shows that selection effects can statistically strengthen the anti-correlation. Some results for mock runs can produce anti-correlations consistent with CLWB. We also predict some ‘missing galaxy-absorber pairs’ which are excluded by the selection criteria.
We estimate the redshift interval adequate to predict accurate anti-correlations of $`W_r`$ versus $`\rho `$. To get results with a small scatter, it is found that in the standard model a total redshift interval of $`10`$ is required. This redshift interval is twice that of the LOSs in CLWB. This may imply that the total redshift interval in present surveys is not large enough to reveal the real anti-correlation.
## Acknowledgments
The authors thank Shude Mao, Tom Theuns for helpful discussions, careful reading of the manuscript and useful comments. We thank the referee, Xavier Barcons for helpful suggestions. WPL thanks Zhenglong Zou, Zugan Deng, Xiaoyang Xia for help. WPL acknowledges the Max-Planck-Institut für Astrophysik for hospitality and, gratefully acknowledges the predoctoral fellowship under the MPG-CAS exchange program. This work was also partly supported by SFB375.
## Appendix A Self-similar solution for cooling flow
From eq. (18), we write,
$$p(x)=\frac{1}{(x+\overline{u})x^2}\frac{d}{dx}(x^2\overline{u})$$
(56)
$$q(x)=\frac{8}{(x+\overline{u})x^2(1+x)^2},$$
(57)
and get the general solution as
$$\overline{\rho _c}=e^{{\scriptscriptstyle p(x)𝑑x}}\left[q(x)e^{{\scriptscriptstyle p(x)𝑑x}}𝑑x+c\right].$$
(58)
It is easy to solve this equation for the case that $`\overline{u}`$ does not depend on $`x`$. Then
$$\overline{\rho _c}=\left(\frac{x+b}{x}\right)^2\times \{\frac{8}{(1+x)(b1)^3}$$
$$+\frac{16}{(b+x)(b1)^3}+\frac{4}{(b+x)^2(b1)^2}+\frac{24ln\frac{1+x}{b+x}}{(b1)^4}+c\}$$
(59)
Here we let $`\overline{u}b`$. For $`x\mathrm{}`$, $`\rho _c(x)0`$, and have $`c=0`$.
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# Vortex states in binary mixture of Bose-Einstein condensates
## Abstract
The vortex configurations in the Bose-Einstein condensate of the mixture of two different spin states $`|F=1,m_f=1>`$ and $`|2,1>`$ of $`{}_{}{}^{87}Rb`$ atoms corresponding to the recent experiments by Matthews et. al. (Phys. Rev. Lett. 83, 2498 (1999)) are considered in the framework of the Thomas-Fermi approximation as functions of $`N_2/N_1`$, where $`N_1`$ is the number of atoms in the state $`|1,1>`$ and $`N_2`$ \- in the state $`|2,1>`$. It is shown that for nonrotating condensates the configuration with the $`|1,1>`$ fluid forming the shell about the $`|2,1>`$ fluid (configuration ”a”) has lower energy than the opposite configuration (configuration ”b”) for all values of $`N_2/N_1`$. When the $`|1,1>`$ fluid has net angular momentum and forms an equatorial ring around the resting central condensate $`|2,1>`$, the total energy of the system is higher than the ground energy, but the configuration ”a” has lower energy than the configuration ”b” for all $`N_2/N_1`$. On the other hand, when the $`|2>`$ fluid has the net angular momentum, for the lowest value of the angular momentum $`\mathrm{}l`$ ($`l=1`$) there is the range of the ratio $`N_2/N_1`$ where the configuration ”b” has lower energy than the configuration ”a”. For higher values of the angular momentum the configuration ”b” is stable for all values of $`N_2/N_1`$.
The realization of Bose-Einstein Condensation (BEC) in dilute atomic gases offers new opportunities for studying quantum degenerate fluids . These condensates which contain thousands of atoms confined to microscale clouds have similarities to superfluidity and laser, and provide new testing ground for many body physics.
Bose-Einstein condensates are quantum fluids in which many particles have the same quantum state. Intrinsic property of the interacting Bose-Einstein condensate is superfluidity. Bulk superfluids are distinguished from normal fluids by their ability to support dissipationless flow. Superfluidity of $`{}_{}{}^{4}He`$ atoms has been widely studied, however, only recently evidence for critical velocity - a key characteristic of superfluids - was observed in a Bose condensate of sodium atoms .
Superfluidity is closely related to the existence of stable quantized vortices which are well known for superfluid $`{}_{}{}^{4}He`$. Vortices can be created in superfluid $`{}_{}{}^{4}He`$ by cooling a rotating container with helium through the superfluid transition. Quite recently the vortices in one-component Bose- Einstein condensate were generated by simple ”stirring” of $`{}_{}{}^{87}Rb`$ atoms condensate . Vortex states in binary mixture of dilute Bose-Einstein condensates have been created using the method proposed by Williams and Holland . Their new scheme exploits the possibility of simultaneously trapping otherwise identical atoms of $`{}_{}{}^{87}Rb`$ in two different ”hyperfine” spin states $`|1>`$ and $`|2>`$. ($`|1>`$ and $`|2>`$ denote $`|F=1,m_f=1>`$ and $`|2,1>`$ spin states of $`{}_{}{}^{87}Rb`$ atoms respectively). Matthews et al. applied a microwave to the condensate and focused a laser beam at various points around its circumference, splitting the atoms into two hyperfine states. The result is a double condensate with one condensate at rest in the center of the cloud and the other in a unit (or multiple) vortex state. Controlling the temporal and spatial dependence of the microwave-induced conversion of $`|1>`$ into $`|2>`$ (and vice versa), one can directly create a $`|2>`$ (or $`|1>`$) state wave function having a wide variety of shapes out of a $`|1>`$ (or $`|2>`$) ground-state wave function or different numbers of atoms in each state.
In two types of vortices have been formed. Matthews et al. put the initial condensate into either the $`|1>`$ or $`|2>`$ state and then made a vortex in the $`|2>`$ or $`|1>`$ state. They considered the time evolution of the vortices to study their dynamics and stability and showed that the dynamics of the $`|1>`$ state vortices is different from the dynamics of the $`|2>`$ state vortices due to different scattering lengths.
The simultaneously trapped resting condensates consisting of the $`{}_{}{}^{87}Rb`$ atoms in the $`|1,1>`$ and $`|2,1>`$ spin states have been considered experimentally in . In this case the intraspecies and interspecies scattering lengths denoted correspondingly as $`a_{11},a_{12},,a_{22}`$ are in the proportion $`a_{11}:a_{12}:a_{22}=1.03:1:0.97`$ with the average of the three being $`55(3)\AA `$. It has been observed that the atoms with the larger scattering length $`a_{11}`$ in the state $`|1>`$ form a lower-density shell about the atoms with the smaller scattering length $`a_{22}`$ . In it has been shown that when the $`|1>`$ fluid forms an equatorial ring around the resting central fluid this configuration is rather stable. Conversely, the $`|2>`$ vortex ring sinks in towards the trap center and breaks up.
The aim of this article is to study the relative stability of different vortex configurations as a function of the ratio $`N_2/N_1`$, where $`N_1`$ is the number of atoms in the state $`|1>`$ and $`N_2`$ \- in the state $`|2>`$. It is shown that for nonrotating phase-separated condensates the configuration with the $`|1>`$ fluid forming the shell about the $`|2>`$ fluid (configuration ”a”) has lower energy than the opposite configuration (configuration ”b”) for all values of $`N_2/N_1`$. When the $`|1>`$ fluid has net angular momentum, the total energy of the system in nonrotating container is higher than the ground energy, but configuration ”a” has lower energy than the configuration ”b” for all $`N_2/N_1`$. On the other hand, when the $`|2>`$ fluid has the net angular momentum, for the lowest value of the angular momentum $`\mathrm{}l`$ ($`l=1`$) there is the range of the ratio $`N_2/N_1`$ where the configuration ”b” has lower energy than the configuration ”a”. For higher values of the angular momentum the configuration ”b” is stable for all values of $`N_2/N_1`$. It should be noted that the experiment of Matthews et al. corresponds to the case where configuration ”b” is unstable ($`N_2/N_11,l=1`$). In our calculations we use the parameters corresponding to the experiments by Matthews et al. : approximately $`8\times 10^5`$ atom in a spherically symmetric potential with oscillation frequencies of $`7.8Hz`$ in the radial and axial directions for both spin states. Let us describe our results in detail.
The modern theoretical description of dilute BEC was originated by the seminal Bogoliubov’s 1947 paper where he showed that weak repulsive interaction qualitatively changes the excitation spectra from quadratic free particle form to a linear phonon-like structure. To describe the trapped condensates at $`T=0`$ one can use the Gross-Pitaevskii (GP) (nonlinear Schrodinger) equation for the condensate wave function . This equation appears as the generalization of the Bogoliubov theory for the inhomogeneous phase. It was widely used to discuss the ground state properties and collective excitations in one-component BEC .
In order to derive analytic results, some approximations must be used. A commonly used one is the Thomas-Fermi Approximation (TFA), which ignores the kinetic energy terms. It has been shown that in the case of one component condensates the TFA results agree well with the numerical calculations for large particle numbers, except for a small region near the boundary of the condensate . In fact, even for small numbers of particles TFA still usually gives qualitatively correct results. In some situations the TFA cannot be used to predict quantitative features of the binary mixture of BEC. For example, the TFA solution considerably underestimate the degree of overlap between the condensates, however, it provides an excellent starting point of study.
Let us first consider the phase separation in binary mixture without rotation.
In the case of a two-species condensate, letting $`\psi _i(𝐫)`$ $`(i=1,2)`$ be the wave function of species $`i`$ with particle number $`N_i`$, we can write two coupled nonlinear Schrodinger (Gross- Pitaevskii) equations as:
$``$ $`{\displaystyle \frac{\mathrm{}^2}{2m_1}}^2\psi _1(𝐫)+{\displaystyle \frac{1}{2}}m_1\omega _1^2(x^2+y^2+\lambda ^2z^2)\psi _1(𝐫)`$ (1)
$``$ $`\mu _1\psi _1(𝐫)+G_{11}|\psi _1(𝐫)|^2\psi _1(𝐫)+G_{12}|\psi _2(𝐫)|^2\psi _1(𝐫)=0`$ (2)
$``$ $`{\displaystyle \frac{\mathrm{}^2}{2m_2}}^2\psi _2(𝐫)+{\displaystyle \frac{1}{2}}m_2\omega _2^2(x^2+y^2+\lambda ^2z^2)\psi _2(𝐫)`$ (3)
$``$ $`\mu _2\psi _2(𝐫)+G_{22}|\psi _2(𝐫)|^2\psi _2(𝐫)+G_{12}|\psi _1(𝐫)|^2\psi _2(𝐫)=0`$ (4)
Equations (2) and (4) were obtained by minimization of the energy functional of the trapped bosons of masses $`m_1`$ and $`m_2`$ given by:
$`E(\psi _1,\psi _2)=={\displaystyle }d^3r[{\displaystyle \frac{\mathrm{}^2}{2m_1}}|\psi _1(𝐫)|^2+`$ (5)
$`+`$ $`{\displaystyle \frac{1}{2}}m_1\omega _1^2(x^2+y^2+\lambda ^2z^2)|\psi _1(𝐫)|^2+`$ (6)
$`+`$ $`{\displaystyle \frac{\mathrm{}^2}{2m_2}}|\psi _2(𝐫)|^2+{\displaystyle \frac{1}{2}}m_2\omega _2^2(x^2+y^2+\lambda ^2z^2)|\psi _2(𝐫)|^2+`$ (7)
$`+`$ $`{\displaystyle \frac{G_{11}}{2}}|\psi _1(𝐫)|^4+{\displaystyle \frac{G_{22}}{2}}|\psi _2(𝐫)|^4+G_{12}|\psi _1(𝐫)|^2|\psi _2(𝐫)|^2].`$ (8)
The chemical potentials $`\mu _1`$ and $`\mu _2`$ are determined by the relations $`d^3r|\psi _i|^2=N_i`$. The trap potential is approximated by an effective three-dimensional harmonic-oscillator potential well, which is cylindrically symmetric about $`z`$ axis, $`\lambda `$ being the ratio of angular frequencies in the axial direction $`\omega _{zi}`$ to that in the transverse direction $`\lambda =\omega _{zi}/\omega _i`$. The interaction strengths, $`G_{11},G_{22},G_{12}`$ are determined by the $`s`$-wave scattering lengths for binary collisions of like and unlike bosons: $`G_{ii}=4\pi \mathrm{}^2a_{ii}/m_i;G_{12}=2\pi \mathrm{}^2a_{12}/m`$, where $`m^1=m_1^1+m_2^1`$.
Let us consider now the phase separation due to interaction between two condensates. In our case of $`|1,1>`$ and $`|2,1>`$ we have $`\frac{1}{2}m_1\omega _1^2=\frac{1}{2}m_2\omega _2^2.`$ We simplify the equations by using dimensionless variables. Let us define the length scale $`a_{}=\left(\frac{\mathrm{}}{m_1\omega _1}\right)^{1/2}`$, and define the dimensionless variables $`𝐫=a_{}𝐫^{},E=\mathrm{}\omega _1E^{},\psi _i(𝐫)=\sqrt{N_i/a_{}^3}\psi _i^{}(𝐫^{}).`$ The wave function $`\psi _i^{}(𝐫^{})`$ is normalized to $`1`$. In terms of these variables the Gross-Pitaevskii energy functional takes the form:
$`E^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }d^3r^{}[N_1|^{}\psi _1^{}|^2+N_1(x^2+y^2+\lambda ^2z^2)|\psi _1^{}|^2+`$ (9)
$`+`$ $`N_2\beta ^2|^{}\psi _2^{}|^2+N_2(x^2+y^2+\lambda ^2z^2)|\psi _2^{}|^2+`$ (10)
$`+`$ $`{\displaystyle \frac{1}{2}}N_1u_1|\psi _1^{}|^4+{\displaystyle \frac{1}{2}}N_2u_2\underset{¯}{\text{e}}ta^2|\psi _2^{}|^4+`$ (11)
$`+`$ $`{\displaystyle \frac{2\pi a_{12}}{a_{}}}{\displaystyle \frac{m_1}{m}}N_1N_2|\psi _1^{}|^2|\psi _2^{}|^2].`$ (12)
Here $`\beta ^2=m_1/m_2=\omega _2^2/\omega _1^2`$ and $`u_i=8\pi a_{ii}N_i/a_{}`$. Eqs. (2) and (4) are rewritten as:
$``$ $`^2\psi _1^{}+(x^2+y^2+\lambda ^2z^2)\psi _1^{}`$ (13)
$``$ $`\mu _1^{}\psi _1^{}+u_1|\psi _1^{}|^2\psi _1^{}++{\displaystyle \frac{4\pi a_{12}N_2}{a_{}}}{\displaystyle \frac{m_1}{m}}|\psi _2^{}|^2\psi _1^{}=0;`$ (14)
$``$ $`\beta ^2^2\psi _2^{}+(x^2+y^2+\lambda ^2z^2)\psi _2^{}\mu _2^{}\psi _2^{}+`$ (15)
$`+`$ $`u_2\beta ^2|\psi _2^{}|^2\psi _2^{}++{\displaystyle \frac{4\pi a_{12}N_1}{a_{}}}{\displaystyle \frac{m_1}{m}}|\psi _1^{}|^2\psi _2^{}=0;`$ (16)
where $`\mu _i^{}=2\mu _i/\mathrm{}\omega _1`$.
In the TFA, Eqs. (12), (14) and (16) can be further simplified by omitting the kinetic energy. For TFA the phase segregated condensates do not overlap, so we can neglect the last terms in Eqs. (12), (14) and (16), obtaining from (14) and (16), in different regions (to be determined later), the simple algebraic equations:
$`|\psi _1^{}(𝐫^{})|^2`$ $`=`$ $`{\displaystyle \frac{1}{u_1}}\left(\mu _1^{}(\rho ^2+\lambda ^2z^2)\right);`$ (17)
$`|\psi _2^{}(𝐫^{})|^2`$ $`=`$ $`{\displaystyle \frac{1}{u_2\beta ^2}}\left(\mu _2^{}(\rho ^2+\lambda ^2z^2)\right).`$ (18)
Here $`\rho ^2=x^2+y^2`$. From Eqs. (17) and (18) one can see that the condensate density has the ellipsoidal form.
In the case of phase separation, the energy of the system can be written in the form
$$E=E_1+E_2,$$
(19)
where
$`E_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega _1N_1\left[\mu _1^{}{\displaystyle \frac{1}{2}}u_1{\displaystyle d^3r^{}|\psi _1^{}|^4}\right],`$ (20)
$`E_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega _1N_2\left[\mu _2^{}{\displaystyle \frac{1}{2}}u_2\beta ^2{\displaystyle d^3r^{}|\psi _2^{}|^4}\right].`$ (21)
In order to obtain Eqs. (20)-(21), Eqs. (17)-(18) have been used.
To investigate the phase separation in the mixture we first suppose that the condensate 1 atoms form an ellipsoidal shell about the condensate 2 atoms (configuration ”a”). To determine the position of the boundary between the condensates, we use the condition of thermodynamic equilibrium : the pressures exerted by both condensates must be equal:
$$P_1=P_2.$$
(22)
Pressure is given by :
$$P_i=\frac{G_{ii}}{2}|\psi _i|^4.$$
(23)
Condensate 2 has the form of the ellipsoid with long semiaxis $`q`$:
$$\rho ^2+\lambda ^2z^2=q^2.$$
(24)
From Eqs. (17)-(18) and (22)-(24) one has the equation for $`q`$:
$$\mu _1^{}q^2=\kappa \mu _2^{}\kappa q^2,$$
(25)
where $`\kappa =\sqrt{(a_{11}m_2)/(a_{22}m_1)}`$.
Chemical potentials $`\mu _1^{}`$ and $`\mu _2^{}`$ can be obtained using the normalization conditions $`d^3r^{}|\psi _1^{}|^2=d^3r^{}|\psi _2^{}|^2=1`$ and are given by:
$$\mu _1^{}=\frac{\mu _1^0}{\left(1\frac{5}{2}q^3+\frac{3}{2}q^5\right)^{2/5}},$$
(26)
$$\mu _2^{}=\frac{3}{(\mu _1^{})^{3/2}q^3}\left(\frac{2\beta ^2(\mu _2^0)^{5/2}}{15}+\frac{(\mu _1^{})^{5/2}q^5}{5}\right),$$
(27)
where $`q=\sqrt{\mu _1^{}}q^{}`$ and
$$\mu _i^0=\left(\frac{15\lambda u_i}{8\pi }\right)^{2/5}.$$
(28)
From equations (26)-(28) one can determine the chemical potentials $`\mu _1^{}`$ and $`\mu _2^{}`$ and the semiaxis of the phase boundary ellipsoid $`q`$ as a function of $`N_1`$ and $`N_2`$. The energy of the configuration ”a” $`E_a=E_{a1}+E_{a2}`$ is given by:
$`E_{a1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega _1N_1\{\mu _1^{}{\displaystyle \frac{15}{4}}{\displaystyle \frac{(\mu _1^{})^{7/2}}{(\mu _1^0)^{5/2}}}\times `$ (29)
$`\times `$ $`[{\displaystyle \frac{8}{105}}({\displaystyle \frac{q^3}{3}}{\displaystyle \frac{2}{5}}q^5+{\displaystyle \frac{q^7}{7}})]\},`$ (30)
$`E_{a2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega _1N_2\{\mu _2^{}{\displaystyle \frac{15}{4}}{\displaystyle \frac{(\mu _1^{})^{3/2}}{\beta ^2(\mu _2^0)^{5/2}}}\times `$ (31)
$`\times `$ $`(\mu _{2}^{}{}_{}{}^{2}{\displaystyle \frac{q^3}{3}}2\mu _2^{}\mu _1^{}{\displaystyle \frac{q^5}{5}}+\mu _1^2{\displaystyle \frac{q^7}{7}})\}.`$ (32)
Let us now consider the opposite case when the condensate $`|2>`$ atoms form an ellipsoidal shell about the condensate $`|1>`$ atoms (configuration ”b”). In this case Eqs. (25)- (32) can be rewritten in the form:
$$\mu _1^{\prime \prime }q_1^2=\kappa (\mu _2^{\prime \prime }q_1^2),$$
(33)
$$(\mu _2^{\prime \prime })^{5/2}=\frac{\beta ^2(\mu _2^0)^{5/2}}{1\frac{5}{2}q_1^3+\frac{3}{2}q_1^5},$$
(34)
$$\frac{15}{2}\frac{(\mu _2^{\prime \prime })^{3/2}}{(\mu _1^0)^{5/2}}\left(\frac{\mu _1^{\prime \prime }q_1^3}{3}\frac{\mu _2^{\prime \prime }q_1^5}{5}\right)=1,$$
(35)
$`E_b`$ $`=`$ $`E_{b1}+E_{b2},`$ (36)
$`E_{b1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega _1N_1\{\mu _1^{\prime \prime }{\displaystyle \frac{15}{4}}{\displaystyle \frac{(\mu _{2}^{}{}_{}{}^{\prime \prime })^{3/2}}{(\mu _1^0)^{5/2}}}\times `$ (37)
$`\times `$ $`(\mu _{1}^{}{}_{}{}^{\prime \prime 2}{\displaystyle \frac{q_1^3}{3}}2\mu _1^{\prime \prime }\mu _2^{\prime \prime }{\displaystyle \frac{q_1^5}{5}}+\mu _2^{\prime \prime 2}{\displaystyle \frac{q_1^7}{7}})\},`$ (38)
$`E_{b2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega _1N_2\{\mu _2^{\prime \prime }{\displaystyle \frac{15}{4}}{\displaystyle \frac{(\mu _2^{\prime \prime })^{7/2}}{\beta ^2(\mu _2^0)^{5/2}}}\times `$ (39)
$`\times `$ $`[{\displaystyle \frac{8}{105}}({\displaystyle \frac{q_1^3}{3}}{\displaystyle \frac{2}{5}}q_1^5+{\displaystyle \frac{q_1^7}{7}})]\}.`$ (40)
Here $`\mu _1^{\prime \prime }`$ and $`\mu _2^{\prime \prime }`$ are the chemical potentials in the configuration ”b”, $`q_1=\sqrt{\mu _2^{\prime \prime }}q_1^{}`$ is the long semiaxis of the boundary ellipsoid, $`E_b`$ is the energy of the configuration ”b”.
To estimate which configuration is stable, one has to compare $`E_a`$ and $`E_b`$.
To evaluate $`\mathrm{\Delta }E=E_aE_b`$ in general case it is worth first to estimate the energy of the phase boundary which arises due to gradient terms omitted in TFA. The surface energy per unit area, the surface tension, is defined as $`\sigma =E_s/S`$, where $`E_s`$ is the surface energy, and $`S`$ is the interface area. $`\sigma `$ may be written in the form :
$`\sigma `$ $`=`$ $`{\displaystyle \frac{\mathrm{}\omega _1}{2\sqrt{2}a_{}^2}}({\displaystyle \frac{a_{12}}{\sqrt{a_{11}a_{22}}}}1)^{1/2}(u_1u_2N_1N_2)^{1/4}\times `$ (41)
$`\times `$ $`|\psi _1^{}||\psi _2^{}|(N_1|\psi _1^{}|^2+N_2|\psi _2^{}|^2)^{1/2}.`$ (42)
Taking into account that the surface area of the ellipsoid with the semiaxis $`a_{}q`$ has the form:
$$S=2\pi a_{}^2q^2\left(1+\frac{1}{\lambda \sqrt{\lambda ^21}}\mathrm{log}\frac{1}{\lambda \sqrt{\lambda ^21}}\right),$$
(43)
one can estimate the contribution of the surface energy $`E_s=\sigma S`$ to the total energy of each configuration. To be specific, we will use the parameters corresponding to the experiments of Matthews et. al. on $`{}_{}{}^{87}Rb`$ atoms. In this case $`m_1=m_2`$, $`a_{}=3.88\times 10^4cm`$, $`\lambda =1`$, $`N=N_1+N_2=0.8\times 10^6`$ atoms. As mentioned in the introduction, $`a_{11}:a_{12}:a_{22}=1.03:1:0.97`$ with the average of the three being $`55(3)\AA `$.
In Fig. 1 we show the energies of configurations ”a” and ”b” (including the surface energy) $`E_a/(\mathrm{}\omega _1)`$ (solid line) and $`E_b/(\mathrm{}\omega _1)`$ (dashed line) as functions of $`\mathrm{log}_{10}(N_2/N_1)`$. One can see that $`E_a`$ is always lower than $`E_b`$. This is consistent with the qualitative assertion and experimental observation that it is energetically favorable for the atoms with the larger scattering length to form a lower-density shell about the atoms with the smaller scattering length . It should be noted that the surface energy is much smaller than the interaction energy because the scattering lengths $`a_{ij}`$ have very close values (see Eq. (42)).
Let us now consider the condensates with nonzero net momentum. For vortex excitation with angular momentum $`\mathrm{}l_j`$ , the condensate wave function is given by
$$\psi _{l_j}(𝐫)=|\psi _{l_j}(𝐫)|e^{il_j\varphi }.$$
(44)
After substituting the wave function for the vortex excitation (44) in Eq. (19), the equation describes the ground state of bosons in an effective confinement potential $`l_1^2\mathrm{}^2/2m_1\rho ^2+l_2^2\mathrm{}^2/2m_2\rho ^2+V_1+V_2`$, where $`V_i=m_i\omega _i(\rho ^2+\lambda ^2z^2)/2`$ and $`\rho ^2=x^2+y^2`$. So within the TFA the density of the vortex state has the form:
$`|\psi _1^{}(𝐫^{})|^2`$ $`=`$ $`{\displaystyle \frac{1}{u_1}}\left(\mu _1^{}(l_1)(\rho ^2+\lambda ^2z^2){\displaystyle \frac{l_1^2}{\rho ^2}}\right);`$ (45)
$`|\psi _2^{}(𝐫^{})|^2`$ $`=`$ $`{\displaystyle \frac{1}{u_2\beta ^2}}\left(\mu _2^{}(l_2)(\rho ^2+\lambda ^2z^2){\displaystyle \frac{\beta ^2l_2^2}{\rho ^2}}\right).`$ (46)
The important new qualitative feature of a vortex in the TFA is the appearance of a small hole of radius $`\xi `$, $`\xi _i^2l_i^2/\mu _i(l_i)`$, but the remainder of the condensate density is essentially unchanged. The fractional change in the chemical potentials caused by the vortex $`(\mu _i^{}(l_i)\mu _i^{})/\mu _i^{}`$ can be shown to be small , of the order of $`1/N^{4/5}`$. In calculation of physical quantities containing condensate density it is sufficient to retain the no- vortex density and simply cut off any divergent radial integrals at the appropriate core sizes $`\xi _1^2=l_1^2/\mu _1^{}`$ or $`\xi _2^2=\beta ^2l_2^2/\mu _2^{}`$. Note that using the unperturbed density for calculation of the vortex properties corresponds to the hydrodynamic limit.
In the case of the phase segregated condensate, one finds from Eqs. (45-46) and (20-21) that the energy change due to the presence of the vortices $`\mathrm{\Delta }E(l_1,l_2)`$ has the form:
$`\mathrm{\Delta }E`$ $`=`$ $`\mathrm{\Delta }E_{N_1}+\mathrm{\Delta }E_{N_2}={\displaystyle \frac{1}{2}}\mathrm{}\omega _1N_1{\displaystyle d_3r^{}\frac{l_1^2}{\rho ^2}|\psi _1^{}|^2}+`$ (47)
$`+`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega _1N_2{\displaystyle d_3r^{}\frac{l_2^2\beta ^2}{\rho ^2}|\psi _2^{}|^2}.`$ (48)
In the hydrodynamic limit $`\psi _i^{}`$ is given by Eqs. (17) and (18).
Let us consider the configuration ”a”. In the hydrodynamic limit the location of the phase boundary is given by Eq. (24). From (48) one has:
$`{\displaystyle \frac{\mathrm{\Delta }E_{N_1}^a}{\frac{1}{2}\mathrm{}\omega _1N_1}}`$ $`=`$ $`{\displaystyle \frac{5l_1^2(\mu _1^{})^{3/2}}{(\mu _1^0)^{5/2}}}\{[\mathrm{ln}{\displaystyle \frac{2\mu _1^{}}{l_1}}{\displaystyle \frac{4}{3}}]{\displaystyle \frac{3}{2}}q^{}\times `$ (49)
$`\times `$ $`[(1{\displaystyle \frac{1}{3}}q^2)\mathrm{ln}{\displaystyle \frac{2\mu _1^{}q^{}}{l_1}}(1{\displaystyle \frac{q^2}{9}})]\}`$ (50)
$`{\displaystyle \frac{\mathrm{\Delta }E_{N_2}^a}{\frac{1}{2}\mathrm{}\omega _1N_2}}`$ $`=`$ $`{\displaystyle \frac{15l_2^2(\mu _1^{})^{1/2}q^{}}{2(\mu _2^0)^{5/2}}}[(\mu _2^{}{\displaystyle \frac{1}{3}}\mu _1^{}q^2)\times `$ (51)
$`\times `$ $`\mathrm{ln}{\displaystyle \frac{2\sqrt{\mu _1^{}\mu _2^{}}q^{}}{l_2\beta }}(\mu _2^{}{\displaystyle \frac{\mu _1^{}q^2}{9}})]`$ (52)
The energy for the configuration ”b” has the form (the location of the phase boundary is given by Eqs. (33-35)):
$`{\displaystyle \frac{\mathrm{\Delta }E_{N_2}^b}{\frac{1}{2}\mathrm{}\omega _1N_2}}`$ $`=`$ $`{\displaystyle \frac{5l_3^2(\mu _2^{\prime \prime })^{3/2}}{(\mu _2^0)^{5/2}}}\{[\mathrm{ln}{\displaystyle \frac{2mu_2^{\prime \prime }}{l_3\beta }}{\displaystyle \frac{4}{3}}]{\displaystyle \frac{3}{2}}q_1^{}\times `$ (53)
$`\times `$ $`[(1{\displaystyle \frac{1}{3}}q_1^2)\mathrm{ln}{\displaystyle \frac{2\mu _2^{}q_2^{}}{l_3\beta }}(1{\displaystyle \frac{q_1^2}{9}})]\}`$ (54)
$`{\displaystyle \frac{\mathrm{\Delta }E_{N_1}^b}{\frac{1}{2}\mathrm{}\omega _1N_1}}`$ $`=`$ $`{\displaystyle \frac{15l_4^2(\mu _2^{\prime \prime })^{1/2}q_1^{}}{2(\mu _1^0)^{5/2}}}[(\mu _1^{\prime \prime }{\displaystyle \frac{1}{3}}\mu _2^{\prime \prime }q_1^2)\times `$ (55)
$`\times `$ $`\mathrm{ln}{\displaystyle \frac{2\sqrt{\mu _1^{\prime \prime }\mu _2^{\prime \prime }}q_1^{}}{l_4}}(\mu _1^{\prime \prime }{\displaystyle \frac{\mu _2^{\prime \prime }q_1^2}{9}})]`$ (56)
Using Eqs.(50-56) one can compare the energies of the configurations for different values of the vortex excitation net angular momenta and different number of state $`|1>`$ and state $`|2>`$ atoms. The results of calculations for the case of nonzero net momenta are presented on the figures 2-7. The parameters corresponding to the experiments by Matthews et al. were used.
The figures 2-4 show the situation with nonrotating condensate $`|2>`$ atoms and a vortex of condensate $`|1>`$ atoms. In Fig.2 we show the rotational parts (50-56) of the energies of configurations ”a” and ”b” $`\mathrm{\Delta }E_a/(\mathrm{}\omega _1)`$ (solid line) and $`\mathrm{\Delta }E_b/(\mathrm{}\omega _1)`$ (dashed line) as functions of $`\mathrm{log}_{10}(N_2/N_1)`$ in the case when the condensate of $`|1>`$ atoms has net angular momentum $`l=1`$ (lower curves) and $`l=2`$ (upper curves). For the configuration ”a” this means the rotation of the outer shell and for the configuration ”b” – of the inner part.
In Fig. 3 the corresponding total energy is shown, and Fig. 4 shows the difference between $`E_a`$ and $`E_b`$ for this case. One can see that $`E_a`$ is always lower than $`E_b`$.
Let us consider now the case when the condensate $`|2>`$ atoms has the net angular momentum $`l=1`$ or $`l=2`$. This case corresponds to the rotation of the outer shell for the configuration ”b” and of the interior part of the condensate for the configuration ”a”. The rotational part is shown on the Fig. 5, the total energy – on Fig. 6. The difference between the total energies $`E_a`$ and $`E_b`$ shown in Fig. 7 presents the most interesting of our results. In the $`l=2`$ case the ”b” configuration with vortex ring of condensate $`|2>`$ atoms in the outer shell is stable in the whole range of the relative concentrations $`N_1`$ and $`N_2`$ while in the case of $`l=1`$ vortex of condensate $`|2>`$ atoms this configuration is stable only in the regions of low $`N_2/N_1`$ or $`N_1/N_2`$ concentrations. When the values $`N_1`$ and $`N_2`$ approach one another the condensate $`|2>`$ atoms ring sinks in toward the trap center.
These results are in agreement with the experiment at approximately equal values of $`N_i`$: when the condensate $`|1>`$ has a net angular momentum, it forms an equatorial ring around the central condensate $`|2>`$. Conversely, a condensate $`|2>`$ vortex forms a ring that tends to contract down into the condensate $`|1>`$.
It should be noted that the similar conclusions have been made in Ref. . Using the dynamic stability analysis of the Gross-Pitaevskii equations, the authors of Ref. have investigated the instability mechanism of the configuration ”b” and concluded that stabilization of this configuration can be attained by the controlling the relative population of both species.
To summarize, we have shown that for nonrotating phase-separated condensates the configuration with the $`|1>`$ fluid forming the shell about the $`|2>`$ fluid (configuration ”a”) has lower energy than the opposite configuration (configuration ”b”) for all values of $`N_2/N_1`$. When the $`|1>`$ fluid has net angular momentum, the total energy of the system is higher than the ground energy, but the configuration ”a” has lower energy than the configuration ”b” for all $`N_2/N_1`$. On the other hand, when the $`|2>`$ fluid has the net angular momentum, for the lowest value of the angular momentum $`\mathrm{}l`$ ($`l=1`$) there is the range of the ratio $`N_2/N_1`$ where the configuration ”b” has lower energy than the configuration ”a”. For higher values of the angular momentum the configuration ”b” is stable for all values of $`N_2/N_1`$.
This work was supported in part by NATO Grant No.PST.CLG.976038. V.N.R and E.E.T acknowledge the financial support from the Russian Science Foundation through the Grant No.98-02-16805. STC was partly supported by NASA under grant No. CRG8-1427. The authors are grateful to V.M. Perez-Garcia and J.J. Garcia-Ripoll for drawing their attention to Ref. .
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# Propagators for scalar bound states at finite temperature in a NJL model This work was partially supported by the National Natural Science Foundation of China and by Grant No.LWTZ-1298 of the Chinese Academy of Sciences.
## 1 Introduction
Finite temperature field theory has attracted much research interest of people due to its application to evolution of early universe and phase transition of the hadron matter \[1-5\]. However, the complete equivalence between its two formalisms - the imaginary-time and the real-time formalism - has been a subtle issue. It is usually assumed that the two formalisms should give the same results. Nevertheless, in some actual problems where only amputated Green functions are involved, the calculations in the two formalisms often show different results . Much work has been contributed to seeking correspondence between the two formalisms for some amputated functions \[7-10\]. The general conclusion is that the difference between the amputated functions obtained in the two formalisms could originate from that one actually deal with different Green functions in the two cases thus they should be used for different physical purposes . In a recent research on the Nambu-Goldstone mechanism of dynamical electroweak symmetry breaking at finite temperature based on a Nambu-Jona-Lasinio (NJL) model , we calculate the propagators for scalar bound states which show different imaginary parts in their denominators in the two formalisms. It was naturally supposed that we had calculated different Green functions in the two cases. However, this inference is open to question. The reason is that the analytic continuation used there of the Matsubara frequency to real energy, though not in the most general form, was essentially made as the way leading to a causal Green function which should just be that obtained in the real-time formalism. In addition, one notes that the propagators for scalar bound states at finite temperature in a NJL model correspond some four-point amputated functions, and the calculations of four-point amputated functions can be effectively reduced to the ones of some two-point functions. It is accepted that a two-point function should be equivalent in the two formalisms of thermal field theory . Therefore, it is necessary for us to reexamine the whole calculations of a NJL model. In this paper, we will do that by means of a chiral $`U_L(1)\times U_R(1)`$ NJL model. After rigorous and careful calculations, we will finally prove that the propagators for scalar bound states including the imaginary parts in their denominators are in fact identical in the two formalisms. We also find a remarkable result that, in the real-time formalism, the thermal transformation matrix of the matrix propagators for scalar bound states is precisely the one of the matrix propagator for an elementary scalar particle. The key-points to reach the above conclusions which were ignored in Ref. are that, except keeping the most general form of the analytic continuation in the imaginary-time formalism, we must carefully consider and separate the imaginary part of the zero-temperature loop integral from relevant expressions. This is very similar to the case of a complete calculation of a two-point function , except that in the case of a NJL model we must additionally use the gap equation so as to eliminate some parts of the Green functions. We will discuss the propagators for scalar bound states in the one-loop approximation respectively in the imaginary-time and real-time formalism , then compare the derived results.
## 2 The imaginary-time formalism
In the imaginary-time formalism, the Lagrangian of the four-fermion interactions will be
$$_{4\mathrm{F}}^I=\frac{G}{4}[(\overline{\psi }\psi )^2(\overline{\psi }\gamma _5\psi )^2],$$
where $`\psi `$ is the fermion field with single flavor and $`N`$ colors, $`G`$ is the coupling constant. Assume the scalar interactions in $`_{4\mathrm{F}}^I`$ may lead to the formation of the fermion condensate $`\overline{\psi }\psi _T`$ at temperature $`T`$, then we will obtain the gap equation which determines the dynamical fermion mass $`m(T)m`$
$$1=GI,I=2N\frac{d^3l}{(2\pi )^3}T\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{(\omega _n+i\mu )^2+\stackrel{}{l}^2+m^2},$$
(1)
where $`\omega _n=(2n+1)\pi /\beta (n=0,\pm 1,\pm 2,\mathrm{};\beta =1/T)`$ is the Matrubara frequency of the fermions and $`\mu `$ is the chemical potential of the fermions. A scalar bound state $`\varphi _S=(\overline{\psi }\psi )`$ could be formed only through the four-fermion interactions, so we must define the propagators of $`\varphi _S`$ as the four-point amputated function $`\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_m,\stackrel{}{p})`$ for the transition from $`(\overline{\psi }\psi )`$ to $`(\overline{\psi }\psi )`$, where $`\mathrm{\Omega }_m=2\pi m/\beta (m=0,\pm 1,\pm 2,\mathrm{})`$ is the Matsubara frequency corresponding to the external energy. $`\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_m,\stackrel{}{p})`$ submits to the the Schwinger-Dyson equation
$$\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_m,\stackrel{}{p})=\frac{G}{2}[1+2N(i\mathrm{\Omega }_m,\stackrel{}{p})\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_m,\stackrel{}{p})],$$
(2)
where the fermion loop contribution
$`N(i\mathrm{\Omega }_m,\stackrel{}{p})`$ $`=`$ $`2N{\displaystyle \frac{d^3l}{(2\pi )^3}T\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{(\omega _n+i\mu )(\omega _n+i\mu +\mathrm{\Omega }_m)+\stackrel{}{l}(\stackrel{}{l}+\stackrel{}{p})m^2}{[(\omega _n+i\mu )^2+\omega _l^2][(\omega _n+i\mu +\mathrm{\Omega }_m)^2+\omega _{l+p}^2]}}`$
$`=`$ $`I+2N{\displaystyle \frac{d^3l}{(2\pi )^3}T\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{(\mathrm{\Omega }_m^2+\stackrel{}{p}^2)2m^2(\omega _n+i\mu )\mathrm{\Omega }_m\stackrel{}{l}\stackrel{}{p}}{[(\omega _n+i\mu )^2+\omega _l^2][(\omega _n+\mathrm{\Omega }_m+i\mu )^2+\omega _{l+p}^2]}}`$
with the denotations $`\omega _l^2=\stackrel{}{l}^2+m^2`$ and $`\omega _{l+p}^2=(\stackrel{}{l}+\stackrel{}{p})^2+m^2`$. By means of the gap equation (1), we can write the solution of (2) as
$$\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_m,\stackrel{}{p})=1/2N[(\mathrm{\Omega }_m^2+\stackrel{}{p}^2)4m^2]\frac{d^3l}{(2\pi )^3}A(i\mathrm{\Omega }_m,\stackrel{}{p},\stackrel{}{l}),$$
$$A(i\mathrm{\Omega }_m,\stackrel{}{p},\stackrel{}{l})=T\underset{n}{}\frac{1}{(\omega _n+i\mu )^2+\omega _l^2}\frac{1}{(\omega _n+\mathrm{\Omega }_m+i\mu )^2+\omega _{l+p}^2},$$
(3)
where we have used the property that, owing to the Lorentz invariance, $`\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_m,\stackrel{}{p})`$ must be a function of $`(\mathrm{\Omega }_m^2+\stackrel{}{p}^2)`$ , thus
$$\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_m,\stackrel{}{p})=\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_m,\stackrel{}{p})=\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_m,\stackrel{}{p}).$$
Doing the sum of the Matsubara frequency $`\omega _n`$ in (3) by the standard procedure then making the analytic continuation of the external energy $`i\mathrm{\Omega }_mp^0+i\epsilon p^0(\epsilon =0_+)`$ and keeping the general form of the replacement, we may write $`A(i\mathrm{\Omega }_mp^0+i\epsilon p^0,\stackrel{}{p},\stackrel{}{l})`$ by
$`A(p,\stackrel{}{l})`$ $`=`$ $`{\displaystyle \frac{1}{4\omega _l\omega _{l+p}}}\{{\displaystyle \frac{1n(\omega _l+\mu )n(\omega _{l+p}\mu )}{p^0+\omega _l+\omega _{l+p}i\epsilon \eta (p^0)}}+{\displaystyle \frac{n(\omega _l+\mu )n(\omega _{l+p}+\mu )}{p^0+\omega _l\omega _{l+p}i\epsilon \eta (p^0)}}`$
$`{\displaystyle \frac{n(\omega _l\mu )n(\omega _{l+p}\mu )}{p^0\omega _l+\omega _{l+p}i\epsilon \eta (p^0)}}{\displaystyle \frac{1n(\omega _l\mu )n(\omega _{l+p}+\mu )}{p^0\omega _l\omega _{l+p}i\epsilon \eta (p^0)}}\},`$
where $`n(\omega _l\pm \mu )=1/[e^{\beta (\omega _l\pm \mu )}+1]`$ and $`\eta (p^0)=p^0/|p^0|`$. For the convenience of making a comparison with the following results in the real-time formalism we express $`A(p,\stackrel{}{l})`$ as a integral of $`l^0`$ by using
$$\mathrm{sin}^2\theta (l^0,\mu )=\frac{\theta (l^0)}{\mathrm{exp}[\beta (l^0\mu )]+1}+\frac{\theta (l^0)}{\mathrm{exp}[\beta (l^0+\mu )]+1}$$
(4)
and the formula $`1/(X+i\epsilon )=X/(X^2+\epsilon ^2)i\pi \delta (X)`$. Eventually we obtain the physical causal propagator for $`\varphi _S`$ in the imaginary-time formalism
$`\mathrm{\Gamma }_{IF}^{\varphi _S}(p)`$ $``$ $`i\mathrm{\Gamma }_I^{\varphi _S}(i\mathrm{\Omega }_mp^0+i\epsilon p^0,\stackrel{}{p})`$ (5)
$`=`$ $`i/(p^24m^2+i\epsilon )[K(p)+H(p)iS^I(p)],`$
where
$`K(p)`$ $`=`$ $`2N{\displaystyle \frac{id^4l}{(2\pi )^4}\frac{1}{(l^2m^2+i\epsilon )[(l+p)^2m^2+i\epsilon ]}}`$ (6)
$`=`$ $`{\displaystyle \frac{N}{8\pi ^2}}{\displaystyle \underset{0}{\overset{1}{}}}𝑑x\left(\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2+M^2}{M^2}}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2+M^2}}\right),M^2=m^2p^2x(1x)`$
is the contribution from the zero temperature fermion loop with the four-dimension Euclidean momentum cut-off $`\mathrm{\Lambda }`$,
$$H(p)=4\pi N\frac{d^4l}{(2\pi )^4}\left\{\frac{(l+p)^2m^2}{[(l+p)^2m^2]^2+\epsilon ^2}+(pp)\right\}\delta (l^2m^2)\mathrm{sin}^2\theta (l^0,\mu )$$
(7)
and
$`S^I(p)`$ $`=`$ $`\eta (p^0)4\pi ^2N{\displaystyle \frac{d^4l}{(2\pi )^4}\delta (l^2m^2)\delta [(l+p)^2m^2]}`$ (8)
$`\times [\mathrm{sin}^2\theta (l^0,\mu )\eta (l^0+p^0)+\mathrm{sin}^2\theta (l^0+p^0,\mu )\eta (l^0)].`$
It is emphasized that $`S^I(p)`$ does not contain any pinch singularity due to the factors $`\eta (l^0+p^0)`$ and $`\eta (l^0)`$ in its integrand and since
$$\delta (l^2m^2)\delta [(l+p)^2m^2]=0,\mathrm{when}\mathrm{\hspace{0.33em}\hspace{0.33em}0}p^2<4m^2,$$
we must have $`S^I(p)=0`$, when $`0p^2<4m^2`$. In addition, from (6), $`K(p)`$ is real when $`p^2<4m^2`$ and from (7), $`H(p)`$ is always real. . Similarly, we may find out the physical causal propagator for the pseudoscalar bound state $`\varphi _P=(\overline{\psi }i\gamma _5\psi )`$
$$\mathrm{\Gamma }_{IF}^{\varphi _P}(p)=i/(p^2+i\epsilon )[K(p)+H(p)iS^I(p)].$$
(9)
The equations (5) and (9) show that $`\varphi _S`$ and $`\varphi _P`$ each have the mass $`2m`$ and $`0`$ thus can be respectively identified with the massive ”Higgs” scalar particle and the massless Nambu-Goldstone boson for the spontaneous symmetry breaking $`U_L(1)\times U_R(1)U_{L+R}(1)`$. This represents the Nambu-Goldstone theorem at finite temperature in the model.
## 3 The real-time formalism
In the real-time formalism, the Lagrangian of the four-fermion interactions will be
$$_{4F}^R=\frac{G}{4}\underset{a=1}{\overset{2}{}}\{[(\overline{\psi }\psi )^{(a)}]^2[(\overline{\psi }\gamma _5\psi )^{(a)}]^2\}(1)^{a+1},$$
(10)
where $`a=1`$ denotes physical fields and $`a=2`$ ghost fields. As a result of the thermal condensates $`(\overline{\psi }\psi )^{(a)}_T0,(a=1,2)`$, the gap equation becomes
$$1=GI,I=2N\frac{d^4l}{(2\pi )^4}\left[\frac{i}{l^2m^2+i\epsilon }2\pi \delta (l^2m^2)\mathrm{sin}^2\theta (l^0,\mu )\right],$$
(11)
which is actually identical to (1) obtained in the imaginary-time formalism .
The propagator for the scalar bound state $`\varphi _S`$ is now a $`2\times 2`$ matrix whose elements correspond to the four-point amputated functions $`\mathrm{\Gamma }_S^{ba}(p)`$ for the transition from $`(\overline{\psi }\psi )^{(a)}`$ to $`(\overline{\psi }\psi )^{(b)}(a,b=1,2)`$. $`\mathrm{\Gamma }_S^{ba}(p)`$ obey the Schwinger-Dyson equations
$$\mathrm{\Gamma }_S^{bc}(p)[\delta ^{ca}GN^{ca}(1)^{a+1}]=i\frac{G}{2}\delta ^{ba}(1)^{a+1},$$
$$N^{ca}(p)=\frac{i}{2}N\frac{d^4l}{(2\pi )^4}\mathrm{tr}\left[iS^{ca}(l,m)iS^{ac}(l+p,m)\right],$$
(12)
where $`iS^{ca}(l,m)`$ is the elements of the thermal matrix propagator for the fermions . After using the gap equation (11), the solution of (12) can be expressed by the matrix
$`\left(\begin{array}{cc}\mathrm{\Gamma }_S^{11}(p)& \mathrm{\Gamma }_S^{12}(p)\\ \mathrm{\Gamma }_S^{21}(p)& \mathrm{\Gamma }_S^{22}(p)\end{array}\right)`$ $`=`$ $`{\displaystyle \frac{1}{|K(p)+H(p)iS(p)|^2R^2(p)}}`$ (13)
$`\times \left(\begin{array}{cc}\frac{i\left[K^{}\left(p\right)+H\left(p\right)+iS\left(p\right)\right]}{p^24m^2+i\epsilon }& \frac{\left(p^24m^2\right)R\left(p\right)}{\left(p^24m^2\right)^2+\epsilon ^2}\\ \frac{\left(p^24m^2\right)R\left(p\right)}{\left(p^24m^2\right)^2+\epsilon ^2}& \frac{i\left[K\left(p\right)+H\left(p\right)iS\left(p\right)\right]}{p^24m^2i\epsilon }\end{array}\right),`$
where
$`S(p)`$ $`=`$ $`4\pi ^2N{\displaystyle \frac{d^4l}{(2\pi )^4}\delta (l^2m^2)\delta [(l+p)^2m^2]}`$ (14)
$`\times [\mathrm{sin}^2\theta (l^0+p^0,\mu )\mathrm{cos}^2\theta (l^0,\mu )+\mathrm{cos}^2\theta (l^0+p^0,\mu )\mathrm{sin}^2\theta (l^0,\mu )]`$
and
$$R(p)=2\pi ^2N\frac{d^4l}{(2\pi )^4}\delta (l^2m^2)\delta [(l+p)^2m^2]\mathrm{sin}2\theta (l^0,\mu )\mathrm{sin}2\theta (l^0+p^0,\mu )$$
(15)
are the terms containing the pinch singularities. $`\mathrm{\Gamma }_S^{11}(p)`$ in (13) is of the form of a causal propagator , hence if it is identified with the physical propagator for the scalar bound state $`\varphi _S`$, as made in Refs. , then the main feature of $`\varphi _S`$ including its mass $`2m`$ could be shown. However, the expression of $`\mathrm{\Gamma }_S^{11}(p)`$ has considerable difference from $`\mathrm{\Gamma }_{IF}^{\varphi _S}(p)`$ in (5) in the imaginary-time formalism. For finding a closer correspondence between the physical propagators for $`\varphi _S`$ in the two formalisms, we will seek a thermal transformation matrix $`𝖬`$ which can diagonalize the matrix propagator (13) so that
$$\left(\begin{array}{cc}\mathrm{\Gamma }_S^{11}(p)& \mathrm{\Gamma }_S^{12}(p)\\ \mathrm{\Gamma }_S^{21}(p)& \mathrm{\Gamma }_S^{22}(p)\end{array}\right)=𝖬^1\left(\begin{array}{cc}\mathrm{\Gamma }_{RF}^{\varphi _S}(p)& 0\\ 0& \mathrm{\Gamma }_{RF}^{\varphi _S}{}_{}{}^{}(p)\end{array}\right)𝖬^1,$$
(16)
where
$$𝖬=\left(\begin{array}{cc}\mathrm{cosh}\mathrm{\Theta }& \mathrm{sinh}\mathrm{\Theta }\\ \mathrm{sinh}\mathrm{\Theta }& \mathrm{cosh}\mathrm{\Theta }\end{array}\right)$$
(17)
and $`\mathrm{\Gamma }_{RF}^{\varphi _S}(p)`$ is now defined as the physical causal propagator for $`\varphi _S`$. It is seen from (13) that $`\mathrm{\Gamma }_S^{22}(p)=[\mathrm{\Gamma }_S^{11}(p)]^{}`$ and $`\mathrm{\Gamma }_S^{21}(p)=\mathrm{\Gamma }_S^{12}(p)=\mathrm{\Gamma }_S^{12}(p)^{}`$, thus (16) may be reduced to three independent algebraic equations for $`\mathrm{Re}\mathrm{\Gamma }_{RF}^{\varphi _S}(p)`$, $`\mathrm{Im}\mathrm{\Gamma }_{RF}^{\varphi _S}(p)`$ and $`\mathrm{sinh}\mathrm{\Theta }`$ (or $`\mathrm{cosh}\mathrm{\Theta }`$ )
$`\mathrm{Re}\mathrm{\Gamma }_S^{11}(p)`$ $`=`$ $`(\mathrm{cosh}^2\mathrm{\Theta }+\mathrm{sinh}^2\mathrm{\Theta })\mathrm{Re}\mathrm{\Gamma }_{RF}^{\varphi _S}(p),\mathrm{Im}\mathrm{\Gamma }_S^{11}(p)=\mathrm{Im}\mathrm{\Gamma }_{RF}^{\varphi _S}(p),`$
$`\mathrm{\Gamma }_S^{12}(p)`$ $`=`$ $`2\mathrm{sinh}\mathrm{\Theta }\mathrm{cosh}\mathrm{\Theta }\mathrm{Re}\mathrm{\Gamma }_{RF}^{\varphi _S}(p).`$ (18)
Considering (13), we obtain from (18)
$$\frac{\mathrm{cosh}^2\mathrm{\Theta }+\mathrm{sinh}^2\mathrm{\Theta }}{2\mathrm{sinh}\mathrm{\Theta }\mathrm{cosh}\mathrm{\Theta }}=\frac{S^{}(p)}{R(p)},S^{}(p)=S(p)\mathrm{Im}K(p),$$
(19)
and furthermore,
$$\mathrm{cosh}\mathrm{\Theta }=\frac{1}{\sqrt{2}}\left[\frac{S^{}(p)}{\sqrt{S_{}^{}{}_{}{}^{2}(p)R^2(p)}}+1\right]^{1/2},\mathrm{sinh}\mathrm{\Theta }=\frac{1}{\sqrt{2}}\left[\frac{S^{}(p)}{\sqrt{S_{}^{}{}_{}{}^{2}(p)R^2(p)}}1\right]^{1/2}.$$
(20)
Then (18) will lead to the physical causal propagator
$$\mathrm{\Gamma }_{RF}^{\varphi _S}(p)=i/(p^24m^2+i\epsilon )\left[\mathrm{Re}K(p)+H(p)i\sqrt{S_{}^{}{}_{}{}^{2}(p)R^2(p)}\right].$$
(21)
Based on the interactions (10) and parallel derivation, we can give the pseudoscalar matrix propagator for the transition from $`(\overline{\psi }i\gamma _5\psi )^{(a)}`$ to $`(\overline{\psi }i\gamma _5\psi )^{(b)}(a,b=1,2)`$ with the elements $`\mathrm{\Gamma }_P^{ba}(p)=\mathrm{\Gamma }_S^{ba}(p)|_{m=0}`$, then diagonalize $`\mathrm{\Gamma }_P^{ba}(p)`$ by the same thermal matrix as $`𝖬`$ in (17) and obtain physical causal propagator $`\mathrm{\Gamma }_{RF}^{\varphi _P}(p)`$ for the pseudoscalar bound state $`\varphi _P=(\overline{\psi }i\gamma _5\psi )`$
$$\mathrm{\Gamma }_{RF}^{\varphi _P}(p)=i/(p^2+i\epsilon )\left[\mathrm{Re}K(p)+H(p)i\sqrt{S_{}^{}{}_{}{}^{2}(p)R^2(p)}\right].$$
(22)
The elements (20) of the matrix $`𝖬`$ depend on $`S(p)`$, $`R(p)`$ and the imaginary part $`\mathrm{Im}K(p)`$ of the zero-temperature loop integral and seem to have rather complicated expressions. However, the final result is remarkable, i.e. the matrix $`𝖬`$ is identical to the thermal transformation matrix of the matrix propagator for an elementary scalar particle. In fact, from the expressions (15) and (14) and the definition (4), we can write
$$R(p)=2\pi ^2N\frac{d^4l}{(2\pi )^4}\delta (l^2m^2)\delta [(l+p)^2m^2]\frac{\eta (l^0)\eta (l^0+p^0)}{\mathrm{cosh}[\beta (l^0\mu )/2]\mathrm{cosh}[\beta (l^0+p^0\mu )/2)]}$$
(23)
and
$$S(p)=\mathrm{cosh}(\beta p^0/2)R(p)+D(p)$$
(24)
$$D(p)=\frac{N}{16\pi ^2}\frac{d^3l}{\omega _l\omega _{l+p}}[\delta (p^0+\omega _l+\omega _{l+p})+\delta (p^0\omega _l\omega _{l+p})].$$
(25)
On the other hand, we may calculate $`K(p)`$ expressed by (6) in terms of the residue theorem and obtain
$$K(p)=\frac{N}{16\pi ^3}\frac{d^3l}{\omega _l\omega _{l+p}}\left[\frac{1}{p^0+\omega _l+\omega _{l+p}i\epsilon }\frac{1}{p^0\omega _l\omega _{l+p}+i\epsilon }\right].$$
Hence the imaginary part of $`K(p)`$
$$\mathrm{Im}K(p)=D(p).$$
(26)
It is seen from (26) and (25) that, $`\mathrm{Im}K(p)0`$ only if $`\delta [p_{}^{0}{}_{}{}^{2}(\omega _l+\omega _{l+p})^2]0`$ or $`p^2=(\omega _l+\omega _{l+p})^2`$. But the latter can be satisfied only if $`p^24m^2`$. This reproduces the former conclusion that $`K(p)`$ is real when $`p^2<4m^2`$. The equations (24) and (26) show that we may separate the imaginary part $`\mathrm{Im}K(p)`$ of the zero-temperature loop integral from $`S(p)`$ and this fact is essential for the following results. Substituting (26) into (24) and considering (19), we will have
$$S^{}(p)=S(p)\mathrm{Im}K(p)=\mathrm{cosh}(\beta p^0/2)R(p),$$
(27)
then from (20) obtain
$$\mathrm{cosh}\mathrm{\Theta }=[1+n(p^0)]^{1/2},\mathrm{sinh}\mathrm{\Theta }=[n(p^0)]^{1/2},n(p^0)=1/(e^{\beta |p^0|}1),$$
(28)
which are precisely the elements of the thermal transformation matrix of the matrix propagator for an elementary scalar particle with zero chemical potential, though now we are dealing with the scalar and pseudoscalar bound states $`\varphi _S`$ and $`\varphi _P`$ composed of fermions and antifermions in the NJL model. This implies that scalar particles, whether elementary or composite, seem always to have the same thermodynamic property.
## 4 Equivalence of the two formalisms and causal, retarded and advanced propagators
We will prove that $`\mathrm{\Gamma }_{IF}^{\varphi _S}(p)`$ and $`\mathrm{\Gamma }_{IF}^{\varphi _P}(p)`$ expressed by (5) and (9) in the imaginary-time formalism are respectively the same as $`\mathrm{\Gamma }_{RF}^{\varphi _S}(p)`$ and $`\mathrm{\Gamma }_{RF}^{\varphi _P}(p)`$ expressed by (21) and (22) in the real-time formalism. It is easy to see that, for this purpose, we only need prove that $`K(p)+H(p)iS^I(p)=\mathrm{Re}K(p)+H(p)i\sqrt{S_{}^{}{}_{}{}^{2}(p)R^2(p)}`$ or
$`S^I(p)`$ $`=`$ $`\sqrt{S_{}^{}{}_{}{}^{2}(p)R^2(p)}+\mathrm{Im}K(p)`$ (29)
$`=`$ $`\mathrm{sinh}(\beta |p^0|/2)R(p)+\mathrm{Im}K(p),`$
where (27) has been used. In fact, from (8) together with (4), (23) (25) and (26) we can obtain
$`S^I(p)`$ $`=`$ $`\eta (p^0)\mathrm{sinh}(\beta p^0/2)2\pi ^2N{\displaystyle \frac{d^4l}{(2\pi )^4}\delta (l^2m^2)\delta [(l+p)^2m^2]}`$
$`\times {\displaystyle \frac{\eta (l^0)\eta (l^0+p^0)}{\mathrm{cosh}[\beta (l^0\mu )/2]\mathrm{cosh}[\beta (l^0+p^0\mu )/2)]}}`$
$`+\eta (p^0)2\pi ^2N{\displaystyle \frac{d^4l}{(2\pi )^4}\delta (l^2m^2)\delta [(l+p)^2m^2][\eta (l^0+p^0)\eta (l^0)]}`$
$`=`$ $`\eta (p^0)\mathrm{sinh}(\beta p^0/2)R(p)`$
$`+\eta (p^0){\displaystyle \frac{N}{16\pi ^2}}{\displaystyle \frac{d^3l}{\omega _l\omega _{l+p}}[\delta (p^0+\omega _l+\omega _{l+p})+\delta (p^0\omega _l\omega _{l+p})]}`$
$`=`$ $`\mathrm{sinh}(\beta |p^0|/2)R(p)+{\displaystyle \frac{N}{16\pi ^2}}{\displaystyle \frac{d^3l}{\omega _l\omega _{l+p}}[\delta (p^0+\omega _l+\omega _{l+p})+\delta (p^0\omega _l\omega _{l+p})]}`$
$`=`$ $`\mathrm{sinh}(\beta |p^0|/2)R(p)+\mathrm{Im}K(p),`$
i.e. (29) is valid indeed. Here separation of $`\mathrm{Im}K(p)`$ from $`S^I(p)`$ is also essential. Hence we can reach the conclusion that, in a NJL model, the four-point amputated functions corresponding to a scalar or pseudoscalar bound state are identical in the imaginary-time and the real-time formalism of thermal field theory. This is not surprising because the calculations of four-point amputated functions in a NJL model can be effectively reduced to the ones of usual two-point functions, but with a new feature that the propagators now discussed for the bound states, are only subtracted four-point amputated functions, rather than the whole of them, because some parts of them have been subtracted through the use of the gap equation. It is emphasized that the key-point of proving such equivalence and deriving the matrix elements (28) of $`𝖬`$ lies in that one must carefully consider and separate the imaginary part $`\mathrm{Im}K(p)`$ of the zero-temperature loop integral from relevant expressions e.g. $`K(p)`$, $`S(p)`$ and $`S^I(p)`$, which is often possibly ignored in usual calculations .
Since the causal propagators for scalar or pseudoscalar bound state are identical in the two formalisms, we can omit the subscript $`\mathrm{"}I\mathrm{"}`$ and $`\mathrm{"}R\mathrm{"}`$ and uniquely express them respectively by
$`\mathrm{\Gamma }_F^{\varphi _S}(p)`$ $`=`$ $`\mathrm{\Gamma }_{IF}^{\varphi _S}(p)=\mathrm{\Gamma }_{RF}^{\varphi _S}(p)`$ (30)
$`=`$ $`i/(p^24m^2+i\epsilon )[\mathrm{Re}K(p)+H(p)i\mathrm{sinh}(\beta |p^0|/2)R(p)]`$
and
$`\mathrm{\Gamma }_F^{\varphi _P}(p)`$ $`=`$ $`\mathrm{\Gamma }_{IF}^{\varphi _P}(p)=\mathrm{\Gamma }_{RF}^{\varphi _P}(p)`$ (31)
$`=`$ $`i/(p^2+i\epsilon )[\mathrm{Re}K(p)+H(p)i\mathrm{sinh}(\beta |p^0|/2)R(p)].`$
Based on the identity of the causal propagators in the two formalisms, it is easy to obtain the retarded and the advanced propagator $`\mathrm{\Gamma }_r^\varphi (p)`$ and $`\mathrm{\Gamma }_a^\varphi (p)`$ (where $`\varphi `$ for $`\varphi _S`$ or $`\varphi _P`$) which are also the same respectively in the two formalisms. If first in the imaginary-time formalism, then we may define
$`\mathrm{\Gamma }_{IF}^\varphi (p)`$ $``$ $`i\mathrm{\Gamma }_I^\varphi (i\mathrm{\Omega }_mp^0+i\epsilon p^0),`$
$`\mathrm{\Gamma }_{Ir}^\varphi (p)`$ $``$ $`i\mathrm{\Gamma }_I^\varphi (i\mathrm{\Omega }_mp^0+i\epsilon ),`$
$`\mathrm{\Gamma }_{Ia}^\varphi (p)`$ $``$ $`i\mathrm{\Gamma }_I^\varphi (i\mathrm{\Omega }_mp^0i\epsilon ).`$ (32)
From (32), we can derive
$$\mathrm{\Gamma }_{IF}^\varphi (p)=\theta (p^0)\mathrm{\Gamma }_{Ir}^\varphi (p)+\theta (p^0)\mathrm{\Gamma }_{Ia}^\varphi (p),$$
$$[\mathrm{\Gamma }_{Ir}^\varphi (p)]^{}=\mathrm{\Gamma }_{Ia}^\varphi (p),$$
and furthermore,
$$\mathrm{\Gamma }_{Ir}^\varphi (p)=\theta (p^0)\mathrm{\Gamma }_{IF}^\varphi (p)\theta (p^0)[\mathrm{\Gamma }_{IF}^\varphi (p)]^{}.$$
(33)
Identifying $`\mathrm{\Gamma }_{IF}^\varphi (p)`$ in (33) with the common $`\mathrm{\Gamma }_F^\varphi (p)`$ in the two formalisms expressed by (30) and (31), we will have
$`\mathrm{\Gamma }_r^{\varphi _S}(p)`$ $`=`$ $`i/(p^24m^2+i\epsilon p^0)[\mathrm{Re}K(p)+H(p)i\mathrm{sinh}(\beta p^0/2)R(p)],`$
$`\mathrm{\Gamma }_r^{\varphi _P}(p)`$ $`=`$ $`i/(p^2+i\epsilon p^0)[\mathrm{Re}K(p)+H(p)i\mathrm{sinh}(\beta p^0/2)R(p)],`$
$`\mathrm{\Gamma }_a^\varphi (p)`$ $`=`$ $`[\mathrm{\Gamma }_r^\varphi (p)]^{},\varphi =\varphi _S\mathrm{or}\varphi _P,`$ (34)
which represent the retarded and the advanced propagators for the bound state $`\varphi `$. The result (34) can also be obtained from the matrix propagator (13) in the real-time formalism by so called the transformation in the RA basis , which will be discussed elsewhere.
## 5 Conclusions
We have proven that the physical causal (as well as retarded and advanced) propagators for scalar bound states in a chiral $`U_L(1)\times U_R(1)`$ NJL model, defined by the four-point amputated functions subtracted through the use of the gap equation, are identical in the imaginary-time and real-time-time formalism. This result convincingly shows equivalence of the two formalisms of thermal field theory in the NJL model. The key-points to complete the above proof lie in keeping the general form of the analytic continuation of the Matsubara frequency in the imaginary-time formalism and separating carefully the imaginary part of the zero temperature loop integral from relevant expressions, and these are also certainly of crucial importance for general explicit demonstration of equivalence of the two formalisms, for instance, in the calculations to many-loop order and/or of n-point Green functions. In addition, we have found that, in the real-time formalism, the thermal transformation matrices of the matrix propagators for scalar bound states are precisely the one for an elementary scalar particle and this fact strongly indicates similarity of thermodynamic property between a composite and an elementary scalar particle.
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# Stability of the monoclinic phase in the ferroelectric perovskite PbZr1-xTixO3.
## I Introduction
Exceptionally striking dielectric and piezoelectric properties are found in PbZr<sub>1-x</sub>Ti<sub>x</sub>O<sub>3</sub> (PZT), the perovskite-type oxide system which is the basis of practically all transducers and other piezoelectric devices. This solid solution is cubic at high temperatures but becomes slightly distorted at lower temperatures, where it is ferroelectric. Except for a narrow region close to PbZrO<sub>3</sub>, the ferroelectric phase is divided in two regions of different symmetry, rhombohedral for Zr-rich compositions and tetragonal for Ti-rich compositions. The highest piezoelectric response in this system is found at the nearly vertical boundary between these two phases, at x$``$ 0.47; the so-called morphotropic phase boundary(MPB), as defined by Jaffe et al. . The PZT phase diagram for compositions around the MPB is shown in Fig.1, where the open circles represent the data of Jaffe et al. , which define the MPB above room temperature. The sharpness of this line is such that a composition fluctuation of $`\mathrm{\Delta }`$x= 0.01 corresponds to a temperature uncertainty of $`\mathrm{\Delta }`$T$``$ 90 K. Recently, high resolution x-ray diffraction measurements on extremely homogenous samples by Noheda et al. showed that an intermediate monoclinic phase exists between the rhombohedral and tetragonal PZT phases . The observation of this monoclinic phase in two different compositions, x= 0.48 and x=0.50 , has allowed a preliminary modification of the phase diagram, as shown in Fig. 1. Furthermore, the discovery of this phase around the MPB in PZT answers many of the questions raised by previous investigators about the
nature of the MPB and the underlying basis for the special physical properties of PZT in this region of the phase diagram, especially in the context of the coexistence of rhombohedral and tetragonal phases.
The monoclinic unit cell is doubled with respect to the tetragonal one and has $`b`$ as the unique axis. $`a_m`$ and $`b_m`$ are directed along the pseudo-cubic \[$`\overline{1}\overline{1}`$0\] and \[1$`\overline{1}`$0\] directions, respectively, while c<sub>m</sub> is close to the tetragonal $`c`$ axis, along , but tilted away form it such that the angle, $`\beta `$, between a<sub>m</sub> and c<sub>m</sub> is slightly larger than 90. This monoclinic phase has unique characteristics in comparison to all other ferroelectric perovskite phases. The polar axis is not determined by symmetry and can be directed anywhere within the monoclinic ac plane; that is, the polar axis is allowed to rotate within this plane. In the case of PZT, the pseudo-cubic and directions are contained within the monoclinic plane and the monoclinic polar axis is tilted away from the polar axis of the tetragonal phase, , towards that of the rhombohedral phase, . As has already been pointed out by other authors , the diffraction data show clearly that the local structure of PZT differs from that of the average one. A structure analysis of rhombohedral PZT by Corker et al. indicated that the Zr/Ti cations in the Zr-rich compositions are distributed among three locally-disordered sites with monoclinic symmetry (see the gray circles in the top-left plot of Fig. 1), resulting in average rhombohedral symmetry (black circle in the top-left plot of Fig. 1.). In a similar structure analysis of tetragonal PZT close to the MPB, the diffraction data were shown to be consistent with Zr/Ti cations distributed among four locally-disordered cation sites with monoclinic symmetry, resulting in average tetragonal symmetry.
In recent years, the development of first-principles calculations applied to the study of ferroelectric perovskites has contributed greatly to the understanding of the physical properties of these materials (see, e.g.,refs. ). The incorporation of a compositional degree of freedom to allow for the study of solid solutions has been an important advance which has opened up the possibility of investigating more complex ferroelectric materials such as PZT and related systems . Very recently, Bellaiche et al. have succeeded to derive the monoclinic phase of PZT from first-principles calculations. These authors also show that the value of the piezoelectric coefficient calculated taking into account rotation of the polarization vector in the monoclinic plane is in good agreement with the high values observed in PZT.
In the present work, the stability region of the monoclinic phase in PZT is characterized by means of synchrotron x-ray powder diffraction measurements made on PZT compositions at closely-spaced intervals in the range x= 0.42-0.52. The monoclinic phase is observed at 20K for 0.46 $``$x$``$ 0.51 and this composition range narrows as the temperature increases. The transition temperature between the tetragonal and monoclinic phases is very steep as a function of composition and coincides with the previously mentioned MPB of Jaffe et al. above ambient temperature.
## II Experimental <br>
PZT samples with x= 0.42, 0.45, 0.46, 0.47, 0.50, 0.51 and 0.52 were prepared by conventional solid-state reaction techniques similar to those used previously for PZT with x= 0.48 . During the calcination, two steps were used. First, the desired solid solution was formed at 900 C using the appropriate amounts of reagent-grade powders of lead carbonate, zirconium oxide and titanium oxide with chemical purity better than 99.9 %. Second, the formed product was pulverized and allowed to reach homogeneity by heating for six hours at 850C (lower than the temperature at which PbO evaporates). Pellets were then pressed using an organic binder and, after burn out of the binder, heated to 1250 C at a ramp rate of 10 C/min, held at this temperature in a covered crucible for 2 hours, and cooled down to room temperature. During sintering, PbZrO<sub>3</sub> was used as a lead source to maintain a PbO-rich atmosphere.
Several sets of high-resolution synchrotron x-ray powder diffraction measurements were made on different occasions at beam line X7A at the Brookhaven National Synchrotron Light Source. Data were collected from the ceramic disks in symmetric flat-plate reflection geometry using $`\theta `$-$`2\theta `$ scans over selected angular regions in the temperature range 20-750 K. The sample was rocked 1-2 during data collection to improve powder averaging. In all these experiments a Ge(111) double-crystal monochromator was use to provide an incident beam with a wavelength close to 0.7 Å . A Ge(220) crystal analyzer and scintillation detector were mounted in the diffracted beam, giving an instrumental resolution of about 0.01 on the 2$`\theta `$ scale. As described in ref. , measurements above room temperature were performed with the disk mounted inside a wire-wound BN tube furnace. The accuracy of the temperature in the furnace is estimated to be about 10 K. For low-temperature measurements, the pellet was loaded in a closed-cycle He cryostat, which has an estimated temperature accuracy of 2 K. With this type of diffraction geometry it is not always possible to eliminate preferred orientation and texture effects, but the peak positions, on which the present results are based, are not affected.
In many cases the peak profiles were quite complex, necessitating a very detailed and careful peak-fitting analysis. The peak positions were determined from least-squares fits to the profile recorded for each of the selected regions. A pseudo-Voigt peak shape function with an asymmetry correction was used , and factors such as anisotropic peak widths, coexisting phases and diffuse scattering between peaks were taken into account. The lattice parameter of individual phases were obtained from fits to the observed peak positions for several reflections. Because of the complicated peak shapes, we found that the above procedure gave more consistent results than standard profile-fitting programs.
Examples of selected regions of the diffraction patterns for the three PZT phases, tetragonal (top), monoclinic (center) and rhombohedral (bottom), around the morphotropic phase boundary are shown in Figure 2. The narrow width of the peaks demonstrates the excellent quality of the ceramic samples and allows the specific characteristics of each phase to be clearly distinguished. In particular, the monoclinic phase exhibit unique features that cannot be accounted for either of the other phase or a mixture of them. In the monoclinic phase the unit cell is doubled in volume with respect to the tetragonal one, with $`a_m`$ and $`b_m`$ lying along the tetragonal \[$`\overline{1}\overline{1}`$0\] and \[1$`\overline{1}`$0\] directions, and $`c_m`$ tilted slightly away from the direction. The monoclinic phase illustrated in this figure corresponds to the composition x= 0.48 at 20K, described in detail in ref., with $`a_m`$= 5.721 Å, $`b_m`$= 5.708 Å, $`c_m`$= 4.138 Åand $`\beta `$= 90.50. The $`c/a`$ value in Fig. 2 (center) is defined as $`2\sqrt{2c_m}/(a_m+b_m)`$, in order to correspond to the $`c_t/a_t`$ ratio in the tetragonal case (top).
## III Phase transitions
The evolution of the different structures as a function of temperature has been determined for all the PZT samples in the present study (x= 0.42, 0.45, 0.46, 0.47, 0.51 and 0.52), and combined with previous data obtained for x= 0.48 and x= 0.50 . These results give a complete picture of the phase transitions occurring around the morphotropic phase boundary from 20 to 700K. Three different low temperature phases, with rhombohedral, monoclinic and tetragonal symmetry, are observed. An important result is that the MPB defined by Jaffe et al. is shown to correspond to the limit of the monoclinic phase rather than that of the rhombohedral phase, and is a very robust line that is reproduced for all the samples under investigation. Both the tetragonal-monoclinic and the tetragonal-rhombohedral phase transitions will be described in this section, as well as the rhombohedral-rhombohedral phase transitions observed for x= 0.42 at lower temperatures.
### A <br>Tetragonal-monoclinic phase transition
The tetragonal phase in PZT is very similar to that of pure PbTiO<sub>3</sub> . The effects of Zr substitution on the structure of the tetragonal phase are basically two: first, as the Zr content increases, the tetragonal strain, $`c_t`$/$`a_t`$, decreases, and, second, the cubic-to-tetragonal phase transition evolves from first-order to second-order. Figure 3 (top) shows the lattice parameters of PZT with x= 0.51 as a function of temperature. The paraelectric-ferroelectric phase transition at T$``$ 660 K is of second order, as expected , and the ferroelectric phase is purely tetragonal down to 100 K. Below this temperature structural changes can be noticed; in particular,the tetragonal $`(h0h)_t`$ and $`(hhh)_t`$ reflections broaden markedly. This is apparent in the lower part of Fig. 3, where the pseudo-cubic $`(220)_{pc}`$ reflections are shown at high temperature (right), at an intermediate temperature (center), and at low temperature (left). Based on a careful peak-fitting analysis, the broadening at low temperatures of these reflections can be attributed to two separated peaks, consistent with the monoclinic symmetry observed in PZT with x=0.48 , also illustrated in Fig. 2. However, the monoclinic distortion is quite small being $`a_m`$$``$$`b_m`$ and $`\beta 90.2^{}`$. Similar behavior was observed for a sample of PZT with x= 0.50 prepared under slightly different conditions, to be discussed later. As seen from Fig. 3, the monoclinic angle, $`\beta `$, is small and the tetragonal-to-monoclinic transition temperature can only be approximately defined at $`T_{TM}50`$ K. On the other hand, data collected from PZT with x= 0.52 show a well-defined tetragonal phase down to 20K.
The evolution of the lattice parameters with temperature for PZT with x= 0.46 is shown in Figure 4 (top). The features displayed by this composition are similar to those of PZT with x= 0.48 . A comparison with the latter data at low temperatures shows that the monoclinic angle, $`\beta `$, is larger for x= 0.46 than for x= 0.48. With decreasing x (Ti content), the differences between $`a_m`$ and $`b_m`$ also increase, while the difference between $`a_m`$ and $`c_m`$ decreases, corresponding to the evolution to a rhombohedral phase in which $`\mathrm{"}a_m=b_m=c_m\mathrm{"}`$ . The monoclinic phase is very well-defined at low temperatures, as shown by the pseudo-cubic $`(110)_{pc}`$ reflections plotted at the bottom left of Fig 4. The evolution of $`\beta `$-$`90^{}`$ shows a transition to a tetragonal phase at $`T_{TM}450K`$, in agreement with the MPB of Jaffe et al.. However, the characteristic features of the tetragonal phase also appear well below this temperature, and there is a wide region of phase coexistence between the tetragonal and monoclinic phases, as shown in the central plot at the bottom of the figure. In this plot the peak positions for the pseudo-cubic $`(110)_{pc}`$ reflections corresponding to the monoclinic and tetragonal phases are shown together with the experimental data. From the observed data, a reliable peak fitting analysis can be carried out and the lattice parameters determined for both phases in this region, as plotted as a function of temperature at the top of Fig. 4. The measurements on PZT with x= 0.47 show similar behavior, but with a narrower coexistence region (300 K $`<`$ T$`<`$ 400 K). For this composition the evolution of the order parameter, $`\beta 90^{}`$, suggests a tetragonal-to-monoclinic phase transition at $`T_{TM}`$ 310 K, very close to that observed for x= 0.48, but the sample is not fully tetragonal until the temperature is larger than 400 K, corresponding to the MPB of Jaffe el al. for this composition.
### B <br>Tetragonal-rhombohedral phase transition
A similar analysis for PZT with x= 0.45 yields completely different results, as shown by the evolution of the $`(200)_{pc}`$ reflection, in the lower part of Fig. 5. At low temperature the sample is rhombohedral (left) and remains rhombohedral up to T $``$ 500 K, while for T $`>`$ 550 K, this composition is tetragonal. Some diffuse scattering is observed between the tetragonal peaks, as shown in the bottom right of Fig. 5. This feature is present in all compositions in the study, as previously noted in ref. , and is associated with the existence of twin boundaries in the tetragonal ferroelectric phase . From the evolution of the order parameter (90 \- $`\alpha _r`$), it is possible to determine that the tetragonal-rhombohedral phase transition is complete at $`T_{TR}580K`$. A coexistence region is observed in the interval 500 K $`<`$ T$`<`$ 580 K. In the central plot at the bottom of the figure, the $`(200)_{pc}`$ reflection in this region is depicted, together with the calculated peak positions for both phases.
### C Low temperature rhombohedral phase
The data obtained for the PZT sample with x= 0.42 show that this composition has rhombohedral symmetry all the way down to 20 K from the Curie point at $`T_c`$$``$ 650 K. At 20 K the rhombohedral lattice parameters are $`a`$= 4.0921 Å and $`\alpha _r`$= 89.61. With increasing temperature, the rhombohedral angle, $`\alpha _r`$, increases gradually until the cubic phase is reached, while $`a`$ remains practically constant.
Two different rhombohedral phases have been observed in PZT: a high-temperature phase ($`F_{R(HT)}`$) and a low temperature rhombohedral phase ($`F_{R(LT)}`$) , which have space groups R3m and R3c, respectively. In the latter phase, adjacent oxygen octahedra along the polar axis are rotated about this axis in opposite directions, so that the unit cell is doubled with respect to the high-temperature phase . The corresponding phase boundary was also determined by Jaffe et al. in the region of the phase diagram above room temperature. An extension of this boundary below room temperature was reported in a neutron powder diffraction study by Amin et al. , who investigated the superlattice peaks from a sample with x= 0.40 and found the transition temperature to occur at about 250 K. In the present work, we were also able to observe very weak superlattice peaks from a composition with x= 0.42 below room temperature in the synchrotron x-ray patterns. The phase boundary in this case was found to lie at approximately 175 K (see fig. 6).
We have also observed one very weak superlattice peak in a recent neutron diffraction study of a sample with x= 0.48 at 20 K. This peak can be indexed in terms of a monoclinic cell with a doubled c-axis, but the nature of the distortion and any possible relationship with that in the low-temperature rhombohedral phase has not yet been determined.
## IV Discussion
The results presented above are summarized and compared with previous data from x= 0.48 and x= 0.50 in Figure 6, which represents the new PZT phase diagram around the MPB. The data obtained for the tetragonal-(monoclinic/rhombohedral) transition temperatures for 0.45$`<`$ x$`<`$0.51 are very consistent and lie on a well-defined line, which reproduces the MPB of Jaffe et al. above ambient temperature. The boundary between the rhombohedral and monoclinic regions is shown as a vertical line between 0.45$`<`$ x$`<`$ 0.46, since no evidence of a monoclinic-rhombohedral phase transition has been observed. The lattice parameters at 20 and 300 K for the compositions under study are listed in Table 1, which also shows clearly the widening of the monoclinic region at lower temperatures. Figure 7 shows the evolution of the lattice parameters of the different phases as a function of composition at 300 K, from the rhombohedral to the tetragonal PZT phases via the monoclinic phase. At
the top of the figure, the unit cell volume shows an essentially linear behavior with composition in the range studied. The monoclinic angle, $`\beta `$, and lattice strain, $`c/a`$, at 300 K are also plotted as a function of composition in Fig. 7, where the rhombohedral cell with lattice parameters $`a`$ and $`\alpha _r`$ (see Table 1) has been expressed in terms of the monoclinic cell . $`c/a`$ corresponds to $`c_t/a_t`$, $`2\sqrt{2c_m}/(a_m+b_m)`$ and $`1`$, in the tetragonal, monoclinic and rhombohedral cases, respectively. $`\beta `$ is $`90^{}`$ for a tetragonal cell and is the monoclinic angle for the monoclinic cell. The role of the monoclinic phase as a ”bridge” between the tetragonal and rhombohedral phases in PZT is clearly demonstrated in these plots. The monoclinic phase has also been observed by Raman scattering in a very recent paper by Souza Filho et al. .
The structural studies reported here, together with those in refs. \[2-4\], comprise data from ten samples from two different origins spanning the composition range 0.42$``$x$``$0.52. As mentioned earlier, only one of these samples was inconsistent with the picture shown in Fig. 6, namely PZT with x= 0.47 described in ref. . For this composition it was found that the tetragonal phase transformed to a rhombohedral phase at low temperatures, while, at intermediate temperatures, a poorly-defined region of coexisting phases was observed. On the other hand, the data for the x=0.47 sample studied in the present work shows, as described above, characteristics similar to those of x= 0.46 or x= 0.48; in particular, the existence of a monoclinic phase at low temperatures and no traces of a rhombohedral phase. It is noteworthy that an analysis of the peak-widths in the cubic phase shows clear differences in the microstructure of the two sets of samples. The microstrain, $`\mathrm{\Delta }`$d/d, and crystallite-size of the samples used in the present study are estimated to be about 3x10<sup>-4</sup> and $`1\mu `$m, respectively . A similar analysis for the x= 0.47 sample described in ref. yields values of 11x$`10^4`$ and 0.2 $`\mu `$m, respectively. One possible explanation is that because of the smaller crystallite size
in the ceramic samples in ref., the inhomogeneous internal stress ”prematurely” induces the tetragonal-rhombohedral phase transitions and inhibits the formation of the intermediate monoclinic phase. With a larger crystallite size, the internal strain is more easily relieved , presumably through the formation of non-180 domains , and the monoclinic phase transition is stabilized.
It is interesting to address the question of why the monoclinic phase was not observed in any of the previous studies. One important factor is the much superior resolution of synchrotron powder diffraction equipment compared to that of laboratory equipment; a second is the presence of wide regions of rhombohedral-tetragonal phase coexistence in many of these studies, due to compositional fluctuations and/or small grain sizes , which would obscure the evidence for monoclinic symmetry. For samples prepared by conventional ”dry” solid-state techniques, a much narrower range of compositional fluctuations and large grain size can be achieved by the use of a final heat treatment at 1250 C, as in the present case, or by the use of ”semi-wet” methods of synthesis and lower firing temperatures . However, perhaps the key element for clarifying the phase diagram is to carry out the structural studies at low temperatures, as clearly demonstrated in Fig. 6.
Very recently, experiments on poled samples by Guo et al. have further underlined the crucial role of the monoclinic phase in PZT. These experiments have revealed that poling induces the monoclinic distortion. The application of an electric field causes the rotation of the polar axis and an associated monoclinic distortion, which is retained after the field is removed. These features are shown to be the origin of the high piezoelectric response in PZT. It is observed that for rhombohedral PZT close to the MPB the region of stability of the monoclinic phase increases after an electric field is applied. The field-induced monoclinic phase is found to be considerably wider on the Zr side of the phase diagram, at room temperature, extending at least to a Ti content of x= 0.42. These experiments validate the microscopic model for the MPB proposed in ref.; i.e. the application of a field would favor one of the local sites, which corresponds exactly to the observed monoclinic distortion (see the dashed arrow in the top-left plot of Fig. 1), and induces the monoclinic phase (see the top-right plot in Fig. 1). Further studies of the poled samples are in progress, and will be reported in a subsequent publication.
Acknowledgments
The authors are especially grateful to J.A. Gonzalo and S-E. Park, who were collaborators in the initial stages of this investigation, for their advice and encouragement. We also wish to thank L. Bellaiche, T. Egami, E. Salje, B.A Tuttle, D. Vanderbilt and T. Vogt for very helpful discussions. Financial support by the U.S. Department of Energy, Division of Materials Sciences (contract No. DE-AC02-98CH10886) and ONR (MURI project N00014-96-1-1173) is also acknowledged.
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# Novikov - Shubin signatures, II
## 1. Introduction
In the classical topology of manifolds it is well known that the torsion subgroup $`T`$ of one-dimensional integral homology of a 3-dimensional oriented manifold supports a symmetric nondegenerate linking form
$$T\times T𝐐/𝐙.$$
Equivariant analogs of the linking form (for example, the Blanchfield form) play an important role in the knot theory.
The purpose of the present paper is to construct a linking form of a new kind, which lives on the torsion part of the extended $`L^2`$-homology. Namely, let $`M`$ be a closed oriented manifold of dimension $`2q+1`$. For any flat Hermitian bundle $``$ over $`M`$ with fiber and monodromy in a finite von Neumann category $`𝒞`$, we defined in , , , extended $`L^2`$ homology $`_i(M;)`$, which is an object of an abelian category $`(𝒞)`$ determined by $`𝒞`$. The homology $`_i(M;)`$ splits canonically as a direct sum of its projective and torsion parts. We show in this paper that the torsion part of the middle dimensional homology $`T(_q(M;))`$ supports a canonical linking form which carries a topological information.
Classification of the torsion Hermitian forms of this sort was performed in the previous paper .
In the present paper we study numerical invariants of such linking forms, the Novikov - Shubin signature and the torsion signature. These invariants use an additional structure, the trace on category $`𝒞`$, and for the torsion signature we need this trace to be not normal, i.e. Dixmier type. The Novikov - Shubin signature is based on the Novikov - Shubin invariant; the torsion signature is based on the concept of the tosion dimension which was introduced by the author in .
We compute the Novikov - Shubin signature and the torsion signature in the (commutative) case of manifolds with fundamental group $`𝐙`$; we show that in this case they can be expressed completely in terms of the Blanchfield form.
The importance of the torsion signature and the Novikov - Shubin signature follows from the fact that these invariants are defined in a general non-commutative situation. We should also emphasize existence of a very large variety of von Neumann flat Hermitian bundles over a given manifold; we refer to , where such bundles are associated with growth processes.
I hope that the algebraic invariants of linking forms, in the spirit of the ones constructed in the present paper, will provide useful topological invariants of manifolds. For instance, these invariants can be applied to study knots and links. I plan to describe applications to the knot concordance problem in a separate article.
## 2. Push forward Hermitian forms
In this section we describe a technique, which allows to apply functors to Hermitian forms. Our major motivation is as follows. Starting from a closed PL manifold, we have the Poincaré duality, which may be viewed as a Hermitian form in the homotopy category of chain complexes (cf. section 2); the push forward forms of the Poincaré duality will be different intersection and linking forms. One of these (the linking form on the torsion part of the extended $`L^2`$ cohomology) is the main subject of this paper.
### 2.1.
We will use the formalism of Hermitian forms in categories with duality as described in , §1.
Suppose that $`𝒞`$ and $`𝒞^{}`$ are two categories with dualities $`(D,s)`$ and $`(D^{},s^{})`$, correspondingly. Let $`F:𝒞𝒞^{}`$ be a covariant functor. We will suppose that $`F`$ has the following property. For any symmetry $`ϵG(𝒞)`$ (cf. , §1.2) there exists a unique symmetry $`ϵ^{}G(𝒞^{})`$, such that for any $`M\mathrm{Ob}(𝒞)`$ holds $`F(ϵ_M)=ϵ_{F(M)}^{}.`$ In this situation $`F`$ determines a homomorphism
$$F_{}:G(𝒞)G(𝒞^{}),F_{}(ϵ)=ϵ^{}$$
between the groups of symmetries.
The functor $`F`$ determines also two contravariant functors $`FD,D^{}F:𝒞𝒞^{}`$.
Let $`\gamma :FDD^{}F`$ be a natural transformation.
Given a form $`\varphi :MD(M)`$ in category with duality $`𝒞`$ (cf. , §1.3), consider the following form in $`𝒞^{}`$:
$`\psi :F(M)D^{}F(M),\text{where}\psi =\gamma (M)F(\varphi ).`$ (2.1)
We will say that form $`\psi `$ is derived from $`\varphi `$ by means of the pair $`(F,\gamma )`$, or that it is a push forward of $`\varphi `$.
If $`\gamma `$ is a natural isomorphism of functors, then the push forward form $`\psi `$ of any non-degenerate form $`\varphi `$ is non-degenerate.
We wish to understand under which conditions the push forward form $`\psi `$ is Hermitian.
###### Definition 2.1.
We will say that the pair $`(F,\gamma )`$ is $`\eta `$-Hermitian, where $`\eta G(𝒞^{})`$, if the following diagram of natural transformations
$`\begin{array}{ccc}F& \stackrel{Fs}{}& FDD\\ \eta s^{}F& & \gamma D& & \\ D^{}D^{}F& \underset{D^{}\gamma }{}& D^{}FD\end{array}`$ (2.2)
commutes.
Here $`Fs`$ denotes the natural transformation between the functors $`F`$ and $`FDD`$ which to any object $`M\mathrm{Ob}(𝒞^{})`$ associates the morphism $`F(s(M)):F(M)`$ $`FDD(M)`$. Similarly, the natural transformation $`\gamma D:FDDD^{}FD`$ associates to any $`M\mathrm{Ob}(𝒞^{})`$ the morphism $`\gamma (D(M)):FDD(M)D^{}FD(M)`$. This explains our notations.
###### Lemma 2.2.
Suppose that $`(F,\gamma )`$ is $`\eta `$-Hermitian. If the initial form $`\varphi :MD(M)`$ is $`ϵ`$-Hermitian, then the push forward form $`\psi :F(M)D^{}F(M)`$ in $`𝒞^{}`$ is $`ϵ^{}`$-Hermitian, where $`ϵ^{}=\eta ^1F_{}(ϵ)G(𝒞^{})`$.
###### Proof.
Consider the following diagram
$`\begin{array}{ccccc}F(M)& \stackrel{F(s(M))}{}& FDD(M)& \stackrel{FD(\varphi )}{}& FD(M)\\ \eta s^{}\left(F\left(M\right)\right)& & \gamma \left(D\left(M\right)\right)& & \gamma \left(M\right)& & \\ D^{}D^{}F(M)& \underset{D^{}(\gamma (M))}{}& D^{}FD(M)& \underset{D^{}F(\varphi )}{}& D^{}F(M).\end{array}`$ (2.3)
The square on the right is commutative because $`\gamma `$ is a natural transformation; the square on the left is commutative since (2.2) holds. One of the compositions from the left upper corner of (2.3) to the right lower corner is
$$\eta D^{}F(\varphi )D^{}(\gamma (M))s^{}(F(M))=\eta \psi ^{}.$$
The other composition equals
$$\gamma (M)F(\varphi ^{})=\gamma (M)F(ϵ\varphi )=F_{}(ϵ)\gamma (M)F(\varphi )=F_{}(ϵ)\psi .$$
## 3. Poincaé duality as a Hermitian form
### 3.1. Duality in the category of chain complexes
Here we will recall the construction of a duality in the category of chain complexes, which geometrically corresponds to the Poincaré duality for manifolds.
Let $`𝒞`$ be an additive category with duality $`(D,s)`$. Fix a non-negative integer $`n`$ and consider the category $`𝒞_n(𝒞)`$, which has as its objects chain complexes of fixed length $`n+1`$
$$C=(0C_n\mathrm{}C_{i+1}\stackrel{d}{}C_i\stackrel{d}{}C_{i1}\mathrm{}C_00)$$
over $`𝒞`$. Morphisms of this category are homotopy classes of $`𝒞`$-chain morphisms.
We will describe a duality in $`𝒞_n(𝒞)`$. For any chain complex $`C\mathrm{Ob}(𝒞_n(𝒞))`$ denote by $`𝔇(C)`$ the following chain complex
$$𝔇(C)=(0D_n\stackrel{(1)^nd^{}}{}D_{n1}\mathrm{}D_i\stackrel{(1)^id^{}}{}D_{i1}\mathrm{}D_00),$$
where $`D_i=D(C_{ni})=C_{ni}^{}`$ is the dual of $`C_{ni}`$ with respect to the duality $`(D,s)`$ in $`𝒞`$, and $`d^{}:D_iD_{i1}`$ is $`d^{}=D(d):D(C_{ni})D(C_{ni+1})`$, the dual of the boundary homomorphisms $`d:C_{ni+1}C_{ni}`$ (with respect to $`(D,s)`$).
If $`f:CC^{}`$ is a chain map between chain complexes in $`𝒞_n(𝒞)`$, which is given by $`𝒞`$-morphisms $`f_i:C_iC_i^{}`$, for $`i=0,1,\mathrm{},n`$, then the dual morphism $`𝔇(f):𝔇(C^{})𝔇(C)`$ is given by the sequence $`D(f_i):D(C_i^{})D(C_i)`$. Note also that the homotopy class of $`𝔇(f)`$ is determined by the homotopy class of $`f`$. Thus we obtain a contravariant functor $`𝔇:𝒞_n(𝒞)𝒞_n(𝒞)`$.
The natural isomorphism $`𝔰:\mathrm{id}_{𝒞_n(𝒞)}𝔇𝔇,`$ which we have to specify in order to define a duality in the category $`𝒞_n(𝒞)`$ (cf. , §1.1) assigns to any chain complex $`C`$ the chain map $`𝔰(C):C𝔇𝔇(C)`$, which is given by the isomorphisms $`ϵ_is_i:C_iDD(C_i)`$, where $`i=0,1,2,\mathrm{},n`$, where $`s_i=s(C_i)`$, and $`ϵ_i=(1)^{i(n+1)}.`$
The above construction applies to the category $`\mathrm{\Lambda }`$-mod of finitely generated projective left $`\mathrm{\Lambda }`$-modules, where $`\mathrm{\Lambda }`$ is the ring with involution, with its canonical duality, cf. , §1.4. The corresponding homotopy category of chain complexes $`𝒞_n(\mathrm{\Lambda }\text{-mod})`$ will be briefly denoted $`𝒞_n(\mathrm{\Lambda })`$.
### 3.2. Poincaré duality
Here we recall the classical construction, which associates to any $`n`$-dimensional manifold a Hermitian form in the homotopy category of chain complexes $`𝒞_n(\mathrm{\Lambda })`$.
Let $`K`$ be a closed piecewise linear manifold of dimension $`n`$. We will denote by $`\pi `$ its fundamental group $`\pi =\pi _1(K)`$. Let $`w:\pi \{1,1\}`$ be the orientation character of $`K`$ (the first Stiefel - Whitney class). It determines an involution on the group ring $`\mathrm{\Lambda }=𝐐[\pi ]`$ given by $`g\overline{g}=w(g)g^1`$ for $`g\pi `$.
Fix two mutually dual triangulations on $`K`$ and denote by $`C(\stackrel{~}{K})`$ and $`C^{}(\stackrel{~}{K})`$ the simplicial chain complexes of the universal covering $`\stackrel{~}{K}`$ with respect to these triangulations. $`C(\stackrel{~}{K})`$ and $`C^{}(\stackrel{~}{K})`$ consist of free $`\mathrm{\Lambda }`$-modules and have length $`n+1`$ and so they can be viewed as objects of category $`𝒞_n(\mathrm{\Lambda })`$. Any choice of orientation of the universal covering $`\stackrel{~}{K}`$ determines a non-degenerate intersection pairing (cf. , )
$`[,]:C(\stackrel{~}{K})\times C^{}(\stackrel{~}{K})\mathrm{\Lambda },`$ (3.1)
which is $`\mathrm{\Lambda }`$-linear with respect to the first variable and anti-linear with respect to the second variable. It has the property
$`[dx,y]=(1)^p[x,dy]`$ (3.2)
for $`xC(\stackrel{~}{K})`$ and $`yC^{}(\stackrel{~}{K})`$, where $`p=|x|`$ denotes the dimension of $`x`$. Note that $`[x,y]`$ is 0 unless $`|x|+|y|=n`$. Recall that $`[,]`$ is defined by the formula
$$[x,y]=\underset{g\pi }{}g^1(gx),y,$$
where $`x,y`$ denotes the usual geometric intersection number of chains $`x`$ and $`y`$ with respect to a given orientation of $`\stackrel{~}{K}`$ (adding signs at each intersection point).
Since $`C(\stackrel{~}{K})`$ and $`C^{}(\stackrel{~}{K})`$ are constructed with respect to two different triangulation of the same space, there is a chain homotopy equivalence $`h:C(\stackrel{~}{K})C^{}(\stackrel{~}{K})`$, which is determined uniquely up to homotopy. In fact (since $`\frac{1}{2}\mathrm{\Lambda })`$) $`h`$ can be chosen so that
$`[x,h(y)]=(1)^{pq}\overline{[y,h(x)]},\text{where}p=|x|,q=|y|,`$ (3.3)
for all $`x,yC(\stackrel{~}{K})`$.
The map $`x(y[x,h(y)]\mathrm{\Lambda })`$ determines a chain map
$`\varphi :C(\stackrel{~}{K})𝔇(C(\stackrel{~}{K})),\varphi (x)(y)=[x,h(y)]`$ (3.4)
(this statement is equivalent to (3.2), which we will view as a non-degenerate form in category $`𝒞(\mathrm{\Lambda })`$ supplied with duality $`(𝔇,𝔰)`$. Note that the canonical isomorphism $`𝔰`$ in $`𝒞_n(\mathrm{\Lambda })`$ is constructed as explained in section 2.1. An easy check using the definitions shows that
$`\varphi ^{}=\varphi ,`$
i.e. $`\varphi `$ is Hermitian.
###### Corollary 3.1.
Any closed piecewise linear manifold $`K`$ of dimension $`n`$ with oriented universal cover $`\stackrel{~}{K}`$ determines canonically (via (3.4) a nondegenerate Hermitian form in category $`𝒞_n(\mathrm{\Lambda })`$ with respect to duality $`(𝔇,𝔰)`$.
## 4. Linking form on extended $`L^2`$-homology
In this section we introduce the linking form on the extended $`L^2`$-homology and cohomology.
### 4.1. Categories $`\mathrm{\Lambda }`$-mod-$`𝒞`$ and $`𝒞`$-mod-$`\mathrm{\Lambda }`$
Let $`𝒞`$ be a Hilbertian von Neumann category (cf. , §2.2). If $`\mathrm{\Lambda }`$ is a ring, we may consider the category of left $`\mathrm{\Lambda }`$-modules in $`𝒞`$, cf. , §6. An object of $`\mathrm{\Lambda }`$-mod-$`𝒞`$ is defined as $`𝒳\mathrm{Ob}(𝒞)`$ with a given ring homomorphism $`\mathrm{\Lambda }\mathrm{hom}_𝒞(𝒳,𝒳)`$. We may think of $`\mathrm{\Lambda }`$ as acting on $`𝒳`$ from the left. Morphisms of $`\mathrm{\Lambda }`$-mod-$`𝒞`$ are the morphisms of $`𝒞`$ which commute with the action of $`\mathrm{\Lambda }`$. It is clear that $`\mathrm{\Lambda }`$-mod-$`𝒞`$ is an additive category.
Similarly one may consider the category $`𝒞`$-mod-$`\mathrm{\Lambda }`$ of right $`\mathrm{\Lambda }`$-modules in $`𝒞`$. A right $`\mathrm{\Lambda }`$-module in $`𝒞`$ is an object $`𝒳`$ of $`𝒞`$ with a ring homomorphism $`\mathrm{\Lambda }^{op}\mathrm{hom}_𝒞(𝒳,𝒳)`$, where $`\mathrm{\Lambda }^{op}`$ is the opposite ring of $`\mathrm{\Lambda }`$.
If $``$ is a left $`\mathrm{\Lambda }`$-module in $`𝒞`$ then the dual $`^{}`$ in $`𝒞`$ (cf. , §2) is naturally defined as a right $`\mathrm{\Lambda }`$-module:
$`(\varphi \lambda )(m)=\varphi (\lambda m)\text{for}\lambda \mathrm{\Lambda },m.`$
Conversely, the dual in $`𝒞`$ of a right $`\mathrm{\Lambda }`$-module is a left $`\mathrm{\Lambda }`$-module.
Suppose that we are given an involution on $`\mathrm{\Lambda }`$. Then we may canonically construct a duality $`(D_\mathrm{\Lambda },s_\mathrm{\Lambda })`$ in $`\mathrm{\Lambda }`$-mod-$`𝒞`$. Namely, given a $`\mathrm{\Lambda }`$-module $``$ in $`𝒞`$, we will define $`D_\mathrm{\Lambda }()`$ to be the dual of $``$ in $`𝒞`$ (cf. 2.2 in , i.e. $`D_\mathrm{\Lambda }()`$ is the set of all anti-linear functionals on $``$) with the following action of $`\mathrm{\Lambda }`$. If $`\varphi D_\mathrm{\Lambda }()`$ and $`\lambda \mathrm{\Lambda }`$ we set
$`(\lambda \varphi )(m)=\varphi (\overline{\lambda }m),\text{for}m.`$
One checks that the canonical isomorphism $`s_\mathrm{\Lambda }=s:D_\mathrm{\Lambda }D_\mathrm{\Lambda }()`$ in $`𝒞`$ is now an isomorphism in $`\mathrm{\Lambda }`$-mod-$`𝒞`$, i.e. it commutes with the $`\mathrm{\Lambda }`$-action. Let us emphasize that the duality in $`\mathrm{\Lambda }`$-mod-$`𝒞`$ depends on the involution of $`\mathrm{\Lambda }`$ in an essential way.
Similarly, one introduces a duality in $`𝒞`$-mod-$`\mathrm{\Lambda }`$ using an involution of the ring $`\mathrm{\Lambda }`$.
### 4.2.
Let $`K`$ be a closed PL manifold of dimension $`n`$. Let $`C(\stackrel{~}{K})`$ denote the simplicial chain complex of the universal covering $`\stackrel{~}{K}`$. From §2 we know that $`C(\stackrel{~}{K})`$ carries a nondegenerate Hermitian form $`\varphi :C(\stackrel{~}{K})𝔇(C(\stackrel{~}{K}))`$ in $`𝒞_n(\mathrm{\Lambda })`$, whose sign depends on the choice of orientation of $`\stackrel{~}{K}`$. Here $`\mathrm{\Lambda }=𝐐[\pi ]`$, where $`\pi =\pi _1(K)`$, considered together with the involution determined by the first Stiefel-Whitney class $`w:\pi \{1,1\}`$.
Let $`𝒞`$ be a Hilbertian von Neumann category. Let $`\mathrm{Ob}(𝒞\text{-mod-}\mathrm{\Lambda })`$ be a right $`\mathrm{\Lambda }`$-module in $`𝒞`$, and let $`\psi :^{}`$ be an Hermitian form in $`\mathrm{Ob}(𝒞\text{-mod-}\mathrm{\Lambda })`$. $``$ is simply a unitary representation of $`\pi `$ with respect to the scalar product determined by $`\psi `$. More precisely, the scalar product $`x,y=\psi (x)(y)`$ for $`x,y`$ and the right action of $`\pi `$ on $``$ satisfy
$$xg,yg=w(g)x,y,g\pi ,\text{and}x,y=\overline{y,x}.$$
The following composite
$`\sigma :\stackrel{~}{}_\pi C(\stackrel{~}{K})\stackrel{\psi \varphi }{}^{}\stackrel{~}{}_\pi 𝔇(C(\stackrel{~}{K}))\stackrel{}{}𝔇(\stackrel{~}{}_\pi C(\stackrel{~}{K}))`$ (4.1)
can be viewed as a Hermitian form in category $`𝒞_n(𝒞)`$. We will use the following Lemma and the push forward construction from section 2 in order to define derived Hermitian forms of (4.1).
###### Lemma 4.1.
Let $`𝒞`$ be a Hilbertian von Neumann category and let $`𝒯`$ denote the torsion subcategory of the extended category $`(𝒞)`$. Let $`F:𝒞_{2q+1}(𝒞)𝒯`$ be the functor which assigns to a chain complex $`C\mathrm{Ob}(𝒞_{2q+1}(𝒞))`$ the torsion part of its $`q`$-dimensional extended homology
$`F(C)=T(_q(C))\mathrm{Ob}(𝒯).`$
Consider the duality $`(𝔇,𝔰)`$ in $`𝒞_{2q+1}(𝒞)`$ (cf. 3.1) and the canonical duality $`(𝔢,s)`$ in $`𝒯`$, cf. , §3.3. Then there exists a natural equivalence
$`\gamma :F𝔇𝔢F,`$
so that the pair $`(F,\gamma )`$ is $`(1)^{q+1}`$-Hermitian (in the sense of §1).
Proof will be given later in subsection 4.3.
###### Theorem 4.2.
Let $`K`$ be a closed piecewise linear manifold of odd dimension $`2q+1`$. Let $`\mathrm{Ob}(𝒞\text{-mod-}\mathrm{\Lambda })`$ be a representation of $`\pi =\pi _1(K)`$, where $`𝒞`$ is a Hilbertian von Neumann category. Let $`T_q`$ denote the torsion part of the extended $`L^2`$-homology $`T_q=T(_q(K;))`$. Then a choice of orientation of the universal covering $`\stackrel{~}{K}`$ and of a Hermitian form $`\psi :^{}`$ in $`𝒞\text{-mod-}\mathrm{\Lambda }`$ determine canonically a non-degenerate $`(1)^{q+1}`$-Hermitian form
$`:T_q𝔢(T_q).`$ (4.2)
This Theorem follows immediately from Lemma 2.2, Corollary 3.1 and Lemma 4.1.
Form $``$ above will be called the linking form.
### 4.3. Proof of Lemma 4.1
First, note that the functor $`F`$ applied to a chain complex $`C`$ in $`𝒞_{2q+1}(𝒞)`$ gives the following torsion object
$`F(C)=(d:C_{q+1}/Z_{q+1}\overline{B}_q),`$
where $`Z_{q+1}`$ denotes the space of cycles in $`C_{q+1}`$ and $`B_q`$ denote the subspace of boundaries in $`C_q`$; cf. , . Applying the duality functor $`𝔢`$ (cf. , §3.3) we obtain
$`𝔢(F(C))=(d^{}:(\overline{B}_q)^{}(C_{q+1}/Z_{q+1})^{}),`$
where the star denotes the duality functor in $`𝒞`$, cf. , §2.2.
Now we want to compute $`F(𝔇(C))`$. The chain complex $`𝔇(C)`$ in dimensions $`q+1`$ and $`q`$ looks as follows
$`\mathrm{}C_q^{}\stackrel{(1)^{q+1}d^{}}{}C_{q+1}^{}\mathrm{}.`$
Clearly $`\mathrm{cl}(\mathrm{im}(d^{}))=\{\varphi C_{q+1}^{};\varphi |_{Z_{q+1}}=0\}=(C_{q+1}/Z_{q+1})^{}`$ and similarly
$$\mathrm{ker}(d^{})=\{\varphi C_q^{};\varphi |_{\overline{B}_q}=0\},\text{and thus}C_q^{}/\mathrm{ker}(d^{})(\overline{B}_q)^{}.$$
Hence we obtain that $`F(𝔇(C))`$ can be identified with
$`F(𝔇(C))=((1)^{q+1}d^{}:(\overline{B}_q)^{}(C_{q+1}/Z_{q+1})^{}).`$
Therefore, we may define the natural transformation $`\gamma (C):F𝔇(C)𝔢F(C)`$ as the morphism in the extended category $`(𝒞)`$ represented by the following commutative diagram
$`\begin{array}{cccc}((1)^{q+1}d^{}:& (\overline{B}_q)^{}& & (C_{q+1}/Z_{q+1})^{})\\ & \left(1\right)^{q+1}\mathrm{id}& & \mathrm{id}& & \\ (d^{}:& (\overline{B}_q)^{}& & (C_{q+1}/Z_{q+1})^{}),\end{array}`$
which is clearly an isomorphism in $`(𝒞)`$.
Now examining diagram (2.2) we find that it is commutative with $`\eta =(1)^{q+1}`$ and hence the pair $`(F,\gamma )`$ is $`(1)^{q+1}`$-Hermitian. ∎
### 4.4. Computing the linking form
Here we will show how one may practically compute the linking form (4.2).
We will assume that $`K`$ is a odd-dimensional manifold, $`n=dimK=2q+1`$. $`\pi `$ will denote the fundamental group of $`K`$, $`w:\pi 𝐙_2`$ the orientation character, and the group ring $`\mathrm{\Lambda }=𝐐[\pi ]`$ will be considered with the involution determined by $`w`$.
Fix two mutually dual triangulations of $`K`$ and consider the chain complexes of the universal cover $`\stackrel{~}{K}`$ which we will denote by $`C_{}(\stackrel{~}{K})`$ and $`C_{}^{}(\stackrel{~}{K})`$. Now, we have the following commutative diagram
$$\begin{array}{ccc}C_{q+1}(\stackrel{~}{K})& \stackrel{d}{}& C_q(\stackrel{~}{K})\\ h_{q+1}& & h_q\\ C_{q+1}^{}(\stackrel{~}{K})& \stackrel{d^{}}{}& C_q^{}(\stackrel{~}{K})\end{array}$$
(4.3)
where $`h_{}`$ denotes a symmetric (i.e. satisfying (3.3) lift to the universal cover of a simplicial approximation of the identity map $`KK`$. Also, we have the following two non-degenerate pairings (cf. (3.1))
$`C_{q+1}(\stackrel{~}{K})\times C_q^{}(\stackrel{~}{K})\mathrm{\Lambda }`$
and
$`C_q(\stackrel{~}{K})\times C_{q+1}^{}(\stackrel{~}{K})\mathrm{\Lambda }`$
which are linear with respect to the first variable and anti-linear with respect to the second variable. They determine isomorphisms
$`C_{q+1}^{}(\stackrel{~}{K})\stackrel{}{}D(C_q(\stackrel{~}{K})),`$
$`C_q^{}(\stackrel{~}{K})\stackrel{}{}D(C_{q+1}(\stackrel{~}{K})),`$
where $`D`$ denotes the duality functor in the category of projective $`\mathrm{\Lambda }`$-modules (cf. , §1.4). We have
$`d^{}=(1)^{q+1}D(d)\text{and, using (}\text{3.3}\text{),}h_{q+1}=(1)^{q(q+1)}D(h_q)=D(h_q)`$
and thus we may rewrite diagram (4.3) as follows
$`\begin{array}{ccc}C_{q+1}(\stackrel{~}{K})& \stackrel{d}{}& C_q(\stackrel{~}{K})\\ D\left(h\right)& & h& & \\ D(C_q(\stackrel{~}{K}))& \stackrel{(1)^{q+1}D(d)}{}& D(C_{q+1}(\stackrel{~}{K})).\end{array}`$ (4.4)
Now, if we are given a right $`\mathrm{\Lambda }`$-module $``$ in a Hilbertian category $`𝒞`$ with an Hermitian non-degenerate form $`\psi :D_\mathrm{\Lambda }()`$, applying the functor $`\stackrel{~}{}_\mathrm{\Lambda }`$ (cf. , §6.4) to diagram (4.4, we obtain the following commutative diagram
$`\begin{array}{ccc}\stackrel{~}{}_\mathrm{\Lambda }C_{q+1}(\stackrel{~}{K})& \stackrel{d}{}& \stackrel{~}{}_\mathrm{\Lambda }C_q(\stackrel{~}{K})\\ \left(1\right)^{q+1}D\left(h^{}\right)& & h^{}& & \\ (\stackrel{~}{}_\mathrm{\Lambda }C_q(\stackrel{~}{K}))^{}& \stackrel{d^{}}{}& (\stackrel{~}{}_\mathrm{\Lambda }C_{q+1}(\stackrel{~}{K}))^{},\end{array}`$ (4.5)
where $`h^{}`$ denotes the composition
$$\stackrel{~}{}_\mathrm{\Lambda }C_q(\stackrel{~}{K})\stackrel{\psi h}{}D_\mathrm{\Lambda }()\stackrel{~}{}_\mathrm{\Lambda }D(C_{q+1}(\stackrel{~}{K}))(\stackrel{~}{}_\mathrm{\Lambda }C_{q+1}(\stackrel{~}{K}))^{}.$$
The last diagram (4.5 clearly represents a $`(1)^{q+1}`$-Hermitian form on the torsion part $`T(\mathrm{coker}d)=T(_q(K,))`$ (compare , §4), which coincides with the linking form (4.2).
### 4.5. Example: Linking form of the circle
Here we will compute explicitly the linking form of the circle $`S^1`$.
As the initial data we have to specify an object $``$ of a von Neumann category $`𝒞`$ with a nondegenerate Hermitian form $`\psi :^{}`$, $`\psi ^{}=\psi `$ and a right unitary action of the fundamental group $`\pi _1(S^1)`$ on $``$. Fix a generator $`t\pi =\pi _1(S^1)`$. Then a unitary action of $`\pi `$ on $``$ is given by a $`𝒞`$-morphism $`t:`$ such that $`t^{}\psi t=\psi `$. We will assume that $`t1:`$ is injective and has a dense image. We will view $`\psi `$ simply as a scalar product on $``$.
Fix a cell decomposition of $`K=S^1`$, consisting of one zero-dimensional cell $`e^0`$ and one one-dimensional cell $`e^1`$. The corresponding cell decomposition of the universal covering $`\stackrel{~}{K}`$ consists of the lifts $`t^ie^0`$ and $`t^je^1`$ of the cells $`e^0`$ and $`e^1`$, where $`i,j𝐙`$, and the boundary homomorphism $`d:C_1(\stackrel{~}{K})C_0(\stackrel{~}{K})`$ is given by
$$d(e^1)=(t1)e^0.$$
The dual cell decomposition also has only two cells: one zero-dimensional cell $`f^0`$ and one one-dimensional cell $`f^1`$. The corresponding dual cell decomposition of $`\stackrel{~}{K}`$ is a shifted up version of the initial cell decomposition of $`\stackrel{~}{K}`$ given by $`t^ie^0`$ and $`t^je^1`$.
As above, the boundary homomorphism of dual chain complex $`C_{}^{}(\stackrel{~}{K})`$ is given by
$$d(f^1)=(t1)f^0.$$
Intersection pairing (3.1) is given by
$$[e^1,f^0]=1,[e^0,f^1]=t^1.$$
Now we have to find a symmetric (i.e. satisfying (3.3)) simplicial approximation $`h:C_{}(\stackrel{~}{K})C_{}^{}(\stackrel{~}{K})`$ of the identity map $`\stackrel{~}{K}\stackrel{~}{K}`$. Such symmetric simplicial approximation is given by
$$h(e^0)=\frac{1}{2}(1+t^1)f^0,h(e^1)=\frac{1}{2}(1+t^1)f^1.$$
One checks that diagram (4.4 in this case has the form
$$\begin{array}{ccc}C_1(\stackrel{~}{K})& \stackrel{t1}{}& C_0(\stackrel{~}{K})\\ {\scriptscriptstyle \frac{1}{2}}\left(t+1\right)& & {\scriptscriptstyle \frac{1}{2}}\left(t^1+1\right)& & \\ (C_0(\stackrel{~}{K}))^{}& \stackrel{1t^1}{}& (C_1(\stackrel{~}{K}))^{}.\end{array}$$
Therefore, diagram (4.5 becomes
$`\begin{array}{ccc}& \stackrel{t1}{}& \\ {\scriptscriptstyle \frac{1}{2}}\left(t+1\right)& & {\scriptscriptstyle \frac{1}{2}}\left(t^1+1\right)& & \\ & \stackrel{t^11}{}& .\end{array}`$ (4.6)
This diagram represents the linking form of the circle $`S^1`$ (cf. , §4.1). The zero-dimensional extended $`L^2`$ homology of the circle is $`_0(S^1;)=((t1):)`$, which is all torsion (since we assume that $`t1:`$ has dense image) and (4.6) represents a (-1)-Hermitian linking form on $`_0(S^1;)`$.
As the next step we want to express the linking form (4.6) as a discriminant form of a Hermitian form on a projective object, cf. , §4. In , §4.9 we have described a general algorithm for this. Applying it in this situation we arrive at the following. Consider the spectral decomposition
$$2tt^1=_0^{\mathrm{}}\lambda 𝑑E_\lambda $$
of the self-adjoint non-negative operator $`(t1)(t1)^{}=2tt^1.`$ Pick a small $`ϵ>0`$ and split $``$ as the direct sum
$$=_ϵQ,\text{where}_ϵ=E_ϵ,\text{and}Q=(1E_ϵ).$$
In other words, $`t1`$ is small on $`_ϵ`$ and it is large on $`Q`$; therefore, we may make an excision (cf. , §4.6) and cut out $`Q`$. As the result (cf. diagrams (4-14) and (4-15) in ) we obtain that the linking form of the circle equals the discriminant form of the following $`(1)`$-Hermitian form
$`(tt^1)/2:_ϵ_ϵ.`$ (4.7)
In order to make the problem more specific, let $`=L^2(Z,\mu )`$ where $`Z`$ is a locally compact Hausdorff space with a positive Radon measure $`\mu `$, cf. , §2, Example 7 and also , §7.3. Assume that the unitary action of $`\pi _1(S^1)`$ on $`=L^2(Z,\mu )`$ is given as the multiplication on a function $`t(z)=\mathrm{exp}(if(z))L_𝐂^{\mathrm{}}(Z,\mu )`$, where $`f:Z[\pi ,\pi ]`$ is a real valued function $`fL_𝐑^{\mathrm{}}(Z,\mu )`$. Then $`_ϵ`$ is $`L^2(Z_ϵ,\mu _ϵ)`$, where
$$Z_ϵ=\{zZ;|\mathrm{exp}(if(z))1|^2<ϵ\},\mu _ϵ=\mu |_{Z_ϵ}.$$
The splitting of Theorem 7.7 from produces two torsion objects $`𝒳_+`$ and $`𝒳_{}`$ out of the linking form, $`𝒳_+𝒳_{}=_0(S^1;)`$. Their spectral density functions are given by
$$F_+(\lambda )=\mu \{zZ_ϵ;f(z)\mathrm{sin}^1(0,\lambda ]\},F_{}(\lambda )=\mu \{zZ_ϵ;f(z)\mathrm{sin}^1([\lambda ,0)\},$$
compare , formula (4-4). Clearly, the above formulae allow to compute (in terms of the function $`f`$ and the measure $`\mu `$) the Novikov-Shubin numbers of the positive and negative parts $`𝒳_+`$, $`𝒳_{}`$. The result depends only on the behavior of function $`f`$ near points $`zZ`$ where $`f(z)`$ is ”small”, but in contract with the usual Novikov-Shubin invariants, the points where $`f`$ is small and positive are treated differently compared to points where $`f`$ is small and negative.
The following is a cohomological version of Theorem 4.2.
###### Theorem 4.3.
Let $`K`$ be a closed piecewise linear $`2q+1`$-dimensional manifold. Denote by $`\pi `$ the fundamental group of $`K`$ and let $`w:\pi 𝐙_2`$ be the first Stiefel - Whitney class of $`K`$. Consider the group ring $`𝐂[\pi ]`$ with the involution determined by $`w`$, i.e. $`\overline{g}=w(g)g^1`$ for $`g\pi `$. Then any non-degenerate Hermitian form $`\psi :D_\mathrm{\Lambda }()`$ in $`\mathrm{\Lambda }`$-mod-$`𝒞`$ and a choice of orientation of $`\stackrel{~}{K}`$ determine canonically a $`(1)^{q+1}`$-Hermitian form
$`:T^{q+1}𝔢(T^{q+1}),\text{where}T^{q+1}=T(^{q+1}(K,)).`$ (4.8)
We call (4.8) the cohomological linking form.
### 4.6.
Given a closed PL manifold $`K`$, the data for the linking form (4.2) consist of specifying an orientation of $`\stackrel{~}{K}`$ and an object $``$ of a Hilbertian von Neumann category $`𝒞`$, a right action of $`\pi =\pi _1(K)`$ on $``$ and also a Hermitian form $`\psi :^{}`$. Geometrically, all these data can be understood as a flat bundle $`K`$ with fiber $``$, whose monodromy coincides with the given action of $`\pi `$ on $``$, together with a flat bundle map $`,:𝐂_w`$, where $`𝐂_w`$ denotes a flat complex line bundle with monodromy given by the first Stiefel-Whitney class $`w:\pi \{1,1\}`$.
In case when the manifold $`K`$ is orientable, the bundle $``$ is simply a flat unitary bundle over $`K`$ with fiber and monodromy in category $`𝒞`$.
The same data determine the cohomological linking form (4.8).
## 5. Manifolds with boundary
In this section we show how to compute the linking form assuming that the given odd-dimensional manifold is a boundary $`K`$, where $`dimK=2q`$.
### 5.1. The intersection form
Consider a compact oriented $`(2q)`$-dimensional manifold $`K`$ with boundary $`K`$ and a flat unitary Hilbertian bundle $``$ over $`K`$. The fiber and the monodromy of $``$ belong to a Hilbertian von Neumann category $`𝒞`$. We have the extended homology $`_i(K;)`$ and $`_i(K;)`$. We will describe now the intersection form
$`:P_q(K)(P_q(K))^{},`$ (5.1)
where we denote
$`P_q(K)=P(_q(K;)/\mathrm{ker}[j:P(_q(K;))P(_q(K,K;))].`$ (5.2)
Recall that $`P(_q(K;))`$ denotes the projective part of $`_q(K;)`$, the extended $`L^2`$-homology with coefficients in $``$. The intersection form (5.1) is defined as follows. We have the following commutative diagram
$$\begin{array}{ccc}_q(K;)& \stackrel{j}{}& _q(K,K;)\\ & & & & \\ ^q(K,K;)& \stackrel{j^{}}{}& ^q(K;)\end{array}$$
in the extended abelian category $`(𝒞)`$, where the vertical maps are Poincaré duality isomorphisms and the horizontal maps are induced by the inclusion $`j`$. Applying the functor of projective part we obtain the commutative diagram
$$\begin{array}{ccc}P(_q(K;))& \stackrel{j}{}& P(_q(K,K;))\\ & & & & \\ P(^q(K,K;))& \stackrel{j^{}}{}& P(^q(K;)).\end{array}$$
By the Universal Coefficients Theorem, $`P(^q(K,K;))P(_q(K,K;))^{}`$, and also $`P(^q(K;))P(_q(K;))^{}`$; moreover, the homomorphism $`j^{}`$ (downstairs) is dual to $`j`$ (upstairs). The morphism $`P(_q(K;))P(^q(K;))(P(_q(K;))^{}`$ (acting from the left upper corner to the right lower corner) vanishes on $`\mathrm{ker}j`$ and takes the values in $`((P(_q(K;))/\mathrm{ker}j)^{}(P(_q(K;))^{}`$. Recall that by the definition, $`P(_q(K;))/\mathrm{ker}j=P_q(K)`$. Hence the above diagram determines
$$:P_q(K)(P_q(K))^{},$$
which by definition is the intersection form.
The intersection form $``$ is $`(1)^q`$-Hermitian and is weakly nondegenerate, i.e. the morphism (5.1) has zero kernel. Note that in general the intersection form $``$ fails to be non-degenerate. In fact, our goal in this section is to compute the discriminant form of $``$ (cf. , §4.4), which describes the way $``$ degenerates.
The intersection form $``$ is nondegenerate if $`K`$ is closed $`K=\mathrm{}`$.
As we know from Theorem 3.4, in the situation above there also exists the linking form
$`:T_{q1}(K)𝔢(T_{q1}(K)),`$ (5.3)
where $`T_{q1}(K)`$ denotes $`T(_{q1}(K;))`$, the torsion part of the extended homology. The linking form $``$ is non-degenerate and $`(1)^q`$-Hermitian.
The following is the main result of this section.
###### Theorem 5.1.
Let $`𝒳T_{q1}(K)`$ denote the image of the boundary homomorphism
$`𝒳=\mathrm{im}[:T(_q(K,K;))T_{q1}(K)].`$ (5.4)
Then $`𝒳`$ is contained in its annihilator $`𝒳^{}`$ with respect to the linking form (5.3) and the induced by $``$ the non-degenerate form on $`𝒳^{}/𝒳`$ is congruent to the discriminant form (cf. , §4.4) of the intersection form $``$, cf (5.1).
Let’s explain the terms used in the statement of Theorem 5.1. Suppose that $`:T𝔢(T)`$ is a torsion Hermitian form. If $`i:𝒳T`$ is an embedding (in the sense of the abelian category $`(𝒞)`$) then the annihilator $`𝒳^{}T`$ is defined as the kernel of $`𝔢(i):T𝔢(𝒳)`$. The assumption that $`𝒳`$ is contained in $`𝒳^{}`$ is equivalent to vanishing of the composite $`𝒳\mathrm{@}>i>>T\mathrm{@}>𝔢(i)>>𝔢(𝒳)`$. The induced torsion Hermitian form on $`𝒳^{}/𝒳`$ is defined as follows. Consider the commutative diagram
$$\begin{array}{ccccc}𝒳& \stackrel{i}{}& T& \stackrel{𝔢(i)}{}& 𝔢(𝒳)\\ 1& & & & 1& & \\ 𝒳& \underset{i}{}& 𝔢(𝒯)& \underset{𝔢(i)}{}& 𝔢(𝒳).\end{array}$$
Both rows are chain complexes and the homology of the upper row is $`𝒳^{}/𝒳`$. Homology of the lower row is canonically isomorphic to $`𝔢(𝒳^{}/𝒳)`$ (use Lemma in , §1.1). Since the vertical arrows in the above diagram are isomorphisms, it determines an isomorphism $`𝒳^{}/𝒳𝔢(𝒳^{}/𝒳)`$, i.e. a nondegenerate form on $`𝒳^{}/𝒳`$.
###### Corollary 5.2.
The intersection form $``$ is non-degenerate if and only if the torsion submodule $`𝒳T_{q1}(K)`$, given by (5.4), is a metabolizer of the linking form $``$ (i.e. $`𝒳=𝒳^{}`$). ∎
###### Corollary 5.3.
Suppose that under the conditions of Theorem 5.1 it is additionally known that $`T(_{q1}(K;))=0`$. Then the linking form $`:T_{q1}(K)𝔢(T_{q1}(K))`$ is congruent to the discriminant of the intersection form $``$.
###### Proof.
If $`T(_{q1}(K;))=0`$ then by the Poincaré duality $`T(_q(K,K;))=0`$, and hence $`𝒳=0`$ (cf. (4.4)). Therefore in this case $`𝒳^{}`$ coincides with $`T_{q1}(K)`$. The result now follows from Theorem 5.1. ∎
In the proof of Theorem 5.1 we will use the following Lemma.
###### Lemma 5.4.
Suppose that $`_{}=(\mathrm{}_{i+1}_i_{i1}\mathrm{})`$ is an exact sequence in the extended abelian category $`(𝒞)`$. Form two chain complexes
$$P_{}=(\mathrm{}P(_{i+1})P(_i)P(_{i1})\mathrm{})$$
and
$$T_{}=(\mathrm{}T(_{i+1})T(_i)T(_{i1})\mathrm{})$$
by taking the projective and torsion parts of $`_{}`$ respectively. Consider homology of $`P_{}`$ and $`T_{}`$ in $`(𝒞)`$, which we denote $`_i(P_{})`$ and $`_i(T_{})`$ correspondingly. Then for any $`i`$ there is a natural isomorphism
$$_i(P_{})_{i1}(T_{}).$$
###### Proof.
We have a short exact sequence $`0T_{}_{}P_{}0`$ of chain complexes in $`(𝒞)`$. The corresponding long homological sequence provides the required isomorphism, since $`_{}`$ is exact. ∎
### 5.2. Proof of Theorem 5.1
Consider the following commutative diagram with exact rows
$`\begin{array}{ccccc}\mathrm{}_q(K,K;)& \stackrel{}{}& _{q1}(K;)& \stackrel{i_{}}{}& _{q1}(K;)\mathrm{}\\ & & & & & & \\ \mathrm{}^q(K;)& \stackrel{i^{}}{}& ^q(K;)& \stackrel{\delta }{}& ^{q+1}(K,K;)\mathrm{}\end{array}`$ (5.5)
The vertical arrows denote the isomorphisms of Poincare duality. We will apply to this diagram the functors of torsion and projective parts and then use the natural isomorphisms of Lemma 5.4. Applying the functor of the torsion part we obtain the commutative diagram
$`\begin{array}{ccccc}T_q(K,K)& \stackrel{}{}& T_{q1}(K)& \stackrel{i_{}}{}& T_{q1}(K)\\ & & & & & & \\ 𝔢(T_{q1}(K))& \stackrel{𝔢(i_{})}{}& 𝔢(T_{q1}(K))& \stackrel{𝔢()}{}& 𝔢(T_q(K,K)).\end{array}`$ (5.6)
The vertical morphisms here can be viewed as three different linking forms induced by the Poincare duality. The homology of the upper horizontal sequence in (5.6) can be easily identified with $`𝒳^{}/𝒳`$. The lower horizontal sequence is dual to the upper one.The middle vertical isomorphism of (5.6) induces an isomorphism from the homology of the upper horizontal sequence to the homology of the lower one; this isomorphism is by the definition the induced form $`𝒳^{}/𝒳𝔢(𝒳^{}/𝒳)`$.
Now we apply the functor of projective part $`P`$ to diagram (5.5) (at the place shifted by one) to get the commutative diagram
$`\begin{array}{ccccc}P(_q(K;))& \stackrel{j}{}& P(_q(K,K;))& \stackrel{}{}& P(_{q1}(K;))\\ & & & & & & \\ P(_q(K,K;))^{}& \stackrel{j^{}}{}& P(_q(K;))^{}& & P(_q(K;))^{}.\end{array}`$ (5.7)
The left square of this diagram appeared above in subsection 5.1 in the definition of the intersection form $``$. We want to compute the homology of the two horizontal sequences in the middle and to identify the results with the corresponding homology of the rows of (5.6) using Lemma 5.4. We find that the diagram
$`\begin{array}{ccccc}P_q(M)& \stackrel{}{}& P_q(M)^{}& & 0\\ 1& & 1& & & & \\ P_q(M)& \stackrel{}{}& P_q(M)^{}& & 0,\end{array}`$ (5.8)
which is built entirely out of the intersection pairing, has homology of the horizontal sequences in the middle term identical (in a canonical way) to the homology of (5.7). Thus we obtain that (5.8) can be mapped onto (5.7) inducing isomorphisms of the extended homology in the middle places. Comparing with the definition of discriminant form in , §4.4, we see that this isomorphism between $`𝒳^{}/𝒳`$ and $`(:P_q(M)P_q(M)^{})`$ gives a congruence between the induced form on $`𝒳^{}/𝒳`$ and the discriminant form of the intersection form. ∎
## 6. Hyperbolicity of the linking form
###### Theorem 6.1.
Suppose that $`M`$ is a closed orientable $`(2q)`$-dimensional manifold and $`M`$ is a flat unitary Hilbertian bundle having the fiber and the monodromy in a finite von Neumann category $`𝒞`$. Suppose that $`WM`$ is a closed codimension one submanifold, which separates $`M`$ into two parts $`M_+`$ and $`M_{}`$. If the middle dimensional extended homology $`_q(M;)=0`$ vanishes, then the linking form of $`W`$
$$:T(_{q1}(W;))𝔢(T(_{q1}(W;)))$$
is hyperbolic (cf. , §5).
###### Proof.
Applying Theorem 5.1 to two manifolds $`M_+`$ and $`M_{}`$, we obtain two submodules $`𝒳_\pm T(_{q1}(W;))`$, described in Theorem 5.1.
First, we observe that each of $`𝒳_\pm `$ is a metabolizer, i.e. $`𝒳_\pm ^{}=𝒳_\pm `$. Indeed, by Theorem 5.1, the factor $`𝒳_\pm ^{}/𝒳_\pm `$ is isomorphic to the cokernel of the intersection pairing $`_\pm :P_q(M_\pm )(P_q(M_\pm ))^{}`$ (we use the notations introduced in section 4.1). However, our assumption $`_q(M;)=0`$ implies that $`P_q(M_\pm )=0`$.
Now we show that the metabolizers $`𝒳_+`$ and $`𝒳_{}`$ are ”disjoint”, i.e. $`𝒳_+𝒳_{}=0`$ (here the intersection is understood as intersection of two subobjects of an object of an abelian category). Indeed, according to Theorem 5.1 and the previous arguments, we may identify $`𝒳_\pm `$ with the kernel of the morphism $`T(_{q1}(W;))T(_{q1}(M_\pm ;))`$ induced by the inclusion. Using the Mayer - Vietoris sequence
$`0_{q1}(W;)\stackrel{i_{+}^{}{}_{}{}^{}+i_{}^{}{}_{}{}^{}}{}_{q1}(M_+;)_{q1}(M_{};)_{q1}(M;)`$ (6.1)
and our assumption $`_q(M;)=0`$, we obtain $`𝒳_+𝒳_{}=0`$.
Now we want to show that $`𝒳_++𝒳_{}=T(_{q1}(W;))`$. We examine the Mayer - Vietoris sequence (6.1) again. Using the Poincaré duality we find that our assumption $`_q(M;)=0`$ implies $`T(_{q1}(M;))=0`$. Therefore we conclude that the direct sum $`T(_{q1}(M_+;))T(_{q1}(M_{};))`$ belongs to the image of $`i_{+}^{}{}_{}{}^{}+i_{}^{}{}_{}{}^{}`$. Thus we get an isomorphism
$$i_{+}^{}{}_{}{}^{}+i_{}^{}{}_{}{}^{}:T(_{q1}(W;))\stackrel{}{}T(_{q1}(M_+;))T(_{q1}(M_{};)),$$
which implies $`𝒳_++𝒳_{}=T(_{q1}(W;))`$.
As a result, $`𝒳_\pm `$ are mutually complementary metabolizers in $`T(_{q1}(W;)))`$, and this implies hyperbolicity of the linking form of $`W`$. ∎
## 7. Novikov - Shubin signatures
In this section we define a new numerical invariant of torsion Hermitian forms which we call the Novikov - Shubin signature.
### 7.1.
Let $`K`$ be a closed oriented $`2q+1`$-dimensional manifold. Let $`K`$ be a flat unitary bundle with fiber and monodromy in a superfinite von Neumann category $`𝒞`$, cf. , §7. Then we have the $`(1)^{q+1}`$-Hermitian linking form
$`:T_q𝔢(T_q).`$ (7.1)
Here $`T_q`$ denotes the torsion part of the extended $`L^2`$-homology, $`T_q=T(_q(K;))`$. The form
$`\{\begin{array}{cc},\text{if}q\text{is odd},\hfill & \\ \sqrt{1},\text{if}q\text{is even}\hfill & \end{array}`$ (7.2)
is Hermitian, and by Theorem 7.7 from , $`T_q`$ splits canonically as a direct sum
$`T_q=(T_q)_+(T_q)_{}`$
of a positive and a negative definite parts of (7.2). This splitting contains all the information about the linking form $``$, cf. , §7.
Hence numerical invariants of the linking form $``$ can be obtained by computing numerical invariants of the positive $`(T_q)_+`$ and negative $`(T_q)_{}`$ parts of $`T_q`$.
### 7.2.
The best known invariant of torsion objects is the Novikov-Shubin number, introduced by S.P. Novikov and M.A. Shubin in . In the context of von Neumann categories the Novikov - Shubin number was studied in , . In order to define the Novikov - Shubin number one should have specified (additionally to the above mentioned data) a trace on the category $`𝒞`$, cf. , §3.9. The Novikov-Shubin number of a torsion object $`𝒳`$ is denoted $`\mathrm{ns}(𝒳)`$, or $`\mathrm{ns}_{\mathrm{tr}}(𝒳)`$, if we want to emphasize the used trace $`\mathrm{tr}`$.
It was observed in , that a more convenient invariant (which is equivalent to $`\mathrm{ns}(𝒳)`$) is given by
$`𝔠(𝒳)=\mathrm{ns}(𝒳)^1[0,\mathrm{}].`$
It is called the Novikov - Shubin capacity of $`𝒳`$. The advantage of capacity against the Novikov-Shubin number consists in the fact that it adequately describes the size of a torsion object: larger torsion objects have larger capacity, trivial object has zero capacity, there may also exist torsion objects having infinite capacity.
### 7.3.
We my use the Novikov-Shubin capacity to obtain invariants of linking forms by setting
$`𝔠_+()=𝔠((T_q)_+),𝔠_{}()=𝔠((T_q)_{}).`$ (7.3)
The pair of numbers $`(𝔠_+(),𝔠_{}())`$ will be called the Novikov - Shubin signature of $``$.
In fact,
$`\mathrm{max}\{𝔠_+(),𝔠_{}()\}=𝔠(T_q)=𝔠(T(_q(K;)))`$ (7.4)
(cf. , (3-22)). Therefore, only one of the numbers $`𝔠_+()`$ or $`𝔠_{}()`$ may actually carry a new information (in case it is smaller than the Novikov - Shubin capacity $`𝔠(T_q)`$).
Invariants (7.3) may be useful in order to show that the linking form $``$ is not hyperbolic.
###### Proposition 7.1.
(a) If linking form (7.1) is hyperbolic then
$$𝔠_+()=𝔠_{}().$$
(b) If $`=_1_2`$ is an orthogonal sum of forms $`_1`$ and $`_2`$, then
$$𝔠_\pm ()=\mathrm{max}\{𝔠_\pm (_1),𝔠_\pm (_2)\}.$$
###### Proof.
(a) follows from , Theorem 7.10. (b) follows from Proposition 3.10 in .∎
### 7.4. Example
We will compute now the invariants (7.3) in the case of manifolds with $`\pi _1(K)=𝐙`$. We will assume here that the unitary flat bundle $``$ has the fiber $`\mathrm{}^2(𝐙)`$ and the monodromy is given by the standard action of $`𝐙`$ on $`\mathrm{}^2(𝐙)`$. It is well known (using Fourier transform) that we may equivalently think of $`\mathrm{}^2(𝐙)`$ as of $`\mathrm{}^2(S^1)`$, such that the action of the generator $`t𝐙`$ is given as the multiplication $`\mathrm{}^2(S^1)\mathrm{}^2(S^1)`$ by the identity function $`t=\mathrm{id}:S^1S^1𝐂`$. The von Neumann category $`𝒞`$ generated by the $`𝐂[𝐙]`$-module $`\mathrm{}^2(𝐙)`$ is isomorphic to the category of finite dimensional measurable fields of Hilbert spaces over $`S^1`$, which is superfinite, according to , §7.3. A trace on this category can be specified by choosing a measure on the circle $`S^1`$. (Note that there are traces on this category which do not come from measures on $`S^1`$).
We will fix this measure as follows. Let $`JS^1`$ be a closed interval and let $`\mu _J`$ be the measure on $`S^1`$, which is the usual Lesbegue measure restricted to $`J`$ and which is zero outside $`J`$. We will assume that $`J`$ is small enough, cf. below.
Let $`K`$ be a $`(2q+1)`$-dimensional closed oriented manifold with $`\pi _1(𝐙)`$. Consider the universal covering $`\stackrel{~}{K}K`$. It has $`𝐙`$ as the group of covering transformations. The homology $`H_q(\stackrel{~}{K};𝐂)`$ is a direct sum of a free $`𝐂[𝐙]`$-module and of a torsion $`𝐂[𝐙]`$-module. The torsion submodule also splits as a finite direct sum of modules of the form
$$M_{c,m}=𝐂[𝐙]/(tc)^m𝐂[Z]$$
where $`c𝐂`$ and $`m0`$, $`m𝐙`$. We will denote by $`𝒮_q𝐂`$ the set of those $`c`$, which lie on the unit circle $`S^1`$. $`𝒮_q(K)`$ is a finite set, which we will call the support of $`K`$.
We will assume that interval $`JS^1`$ is so small that it contains at most one point $`c𝒮_q(K)`$.
For $`c𝒮_q(K)J`$ consider the $`(tc)`$-adic part $`(H_q(\stackrel{~}{K};𝐂))_c`$ of $`H_q(\stackrel{~}{K};𝐂)`$, i.e. the set of all homology classes, which can be killed by a power $`(tc)^m`$ with some large $`m`$. There is a $`(1)^{q+1}`$-Hermitian Blanchfield pairing
$`\{,\}:(H_q(\stackrel{~}{K};𝐂))_c\times (H_q(\stackrel{~}{K};𝐂))_c/𝐂[𝐙]_c,`$ (7.5)
where $``$ denotes the field of fractions of $`𝐂[𝐙]`$ and $`𝐂[𝐙]_c`$ denotes the ring rational function in $`t`$ such that their denominators are prime to $`tc`$, cf. , §9. The involution in $`/𝐂[𝐙]_c`$ is given by the complex conjugation and by $`\overline{t}=t^1`$.
We will explain below (cf. 7.5) how the Blanchfield pairing (7.5) determines nonnegative integers $`n_j^+(c)`$ and $`n_j^{}(c)`$, where $`j=1,2,\mathrm{}`$ such that only finitely many of them are nonzero.
Now we define odd and even height numbers as follows:
$`h_{\text{odd}}(c)=\mathrm{max}\{j\text{odd, such that}n_j^+(c)>0\text{or}n_j^{}(c)>0\},`$ (7.6)
$`h_{\text{ev}}^\pm (c)=\mathrm{max}\{j\text{even, such that}n_j^\pm (c)>0\}.`$ (7.7)
###### Theorem 7.2.
In the situation described in 6.5, assume that $`J𝒮_q(K)`$ consists of a single interior point $`cJ`$. Then the positive and negative Novikov-Shubin capacities (7.3) of the linking form
$`:T_q𝔢(T_q),T_q=T(_q(K;\mathrm{}^2(𝐙)))`$ (7.8)
with respect to the trace determined by the measure $`\mu _J`$ are given by
$`𝔠_+()=\mathrm{max}\{h_{\text{odd}}(c),h_{\text{ev}}^+(c)\},𝔠_{}()=\mathrm{max}\{h_{\text{odd}}(c),h_{\text{ev}}^{}(c)\}.`$ (7.9)
Note that
$$\mathrm{max}\{h_{\text{odd}}(c),h_{\text{ev}}^+(c),h_{\text{ev}}^{}(c)\}$$
coincides with the height of $`(H_q(\stackrel{~}{K};𝐂))_c`$, and hence by , Theorem 4.13, it coincides with the Novikov-Shubin capacity of $`_q(K;\mathrm{}^2(𝐙))`$ with respect to $`\mu _J`$.
Theorem 7.2 demonstrates that the numbers $`𝔠_+()`$ and $`𝔠_{}()`$ can be arbitrary except that their maximum is fixed by relation (7.4).
Let us consider now situation when $`𝒮_q(K)J`$ is a single point, which is an end point of the interval $`J`$. We assume that the circle $`S^1`$ is canonically oriented (anticlockwise) and hence we may speak about the initial and the terminal points of $`J`$.
###### Theorem 7.3.
In the situation described in 7.4, assume that the intersection $`J𝒮_q(K)`$ consists of a single point $`c`$, which is the terminal point of the interval $`J`$. Then the positive and negative Novikov-Shubin capacities (7.3) of the linking form
$$:T_q𝔢(T_q),T_q=T(_q(K;\mathrm{}^2(𝐙)))$$
with respect to the trace determined by the measure $`\mu _J`$ are given by
$`𝔠_+()=\mathrm{max}\{h_{\text{ev}}^+(c),h_{\text{odd}}^{}(c)\},𝔠_{}()=\mathrm{max}\{h_{\text{ev}}^{}(c),h_{\text{odd}}^+(c)\}.`$ (7.10)
### 7.5. Numbers $`n_j^\pm (c)`$
Here we will explain how the Blanchfield form (7.5) determines the numbers $`n_j^\pm `$, where $`j=1,2,\mathrm{}`$, which appear in Theorems (7.6), (7.7) and (7.10). Here we essentially follow , §2.
Any element $`f/𝐂[𝐙]_c`$ can be uniquely represented in the form
$$f=\alpha _1g+\alpha _2g^2+\alpha _3g^3+\mathrm{},\text{where}g=ic(tc)^1,$$
with $`\alpha _j𝐂`$; only finitely many $`\alpha _j`$’s are nonzero. Note that $`g/𝐂[𝐙]_c`$ is ”real”, i.e. $`\overline{g}=g`$ assuming that $`cS^1`$.
Given $`x,y(H_q(\stackrel{~}{K};𝐂))_c`$, the value $`\{x,y\}/𝐂[𝐙]_c`$ of the Blanchfield form (7.5) can be uniquely expressed as
$`\{x,y\}=\alpha _1(x,y)g+\alpha _2(x,y)g^2+\alpha _3(x,y)g^3+\mathrm{},`$
and this defines a sequence of $`(1)^{q+1}`$-Hermitian forms
$`\alpha _j:(H_q(\stackrel{~}{K};𝐂))_c\times (H_q(\stackrel{~}{K};𝐂))_c𝐂,j=1,2,\mathrm{}.`$
Let $`T_j(H_q(\stackrel{~}{K};𝐂))_c`$ denote the subspace, consisting of cycles $`z(H_q(\stackrel{~}{K};𝐂))_c`$, such that $`(tc)^jz=0`$. Now, we define the numbers $`n_j^+(c)`$ and $`n_j^{}(c)`$ as the numbers of positive and negative squares (correspondingly) in the diagonal representation of the Hermitian form
$`\{\begin{array}{cc}\alpha _j:T_j\times T_j𝐂,\text{if }q\text{ is odd},\hfill & \\ \sqrt{1}\alpha _j:T_j\times T_j𝐂,\text{if }q\text{ is even}.\hfill & \end{array}`$
### 7.6. Beginning of the proof of Theorems 7.2 and 7.3
As in , let $`𝒪`$ denote the ring of germs at the origin of complex valued holomorphic functions $`f:(ϵ,ϵ)𝐂`$. An element of $`𝒪`$ can also be represented by a power series
$$f(\tau )=\underset{n0}{}a_n\tau ^n,a_n𝐂$$
having a nonzero radius of convergence. The ring operations are given by pointwise addition and multiplication. The involution in $`𝒪`$ is given by $`\overline{\tau }=\tau `$ and by the complex conjugation.
The algebraic structure of $`𝒪`$ is extremely simple since it is a discrete valuation ring. Its maximal ideal $`𝔪𝒪`$ is given by $`𝔪=\{f𝒪;f(0)=0\}`$. Any finitely generated $`𝒪`$-module is a direct sum of a free module and a torsion submodule. The torsion submodule can be represented as a direct sum of finitely many modules of the form $`𝒪/𝔪^k`$, where $`k𝐙`$.
Let $`𝒞^{}`$ denote the category of finitely generated $`𝒪`$-modules and let $`𝒯^{}𝒞^{}`$ be the full subcategory generated by torsion modules.
We will describe a duality $`(𝔢^{},s^{})`$ in $`𝒯^{}`$, cf. , §1. The functor $`𝔢^{}:𝒯^{}𝒯^{}`$ is given by $`𝔢^{}(X)=\overline{\mathrm{Hom}}_𝒪(X,/𝒪)`$, where $`X\mathrm{Ob}(𝒯^{})`$, and $``$ denotes the ring of fractions of $`𝒪`$. The bar over $`\mathrm{Hom}`$ means that we consider the set of all anti-homomorphisms, compare , §1.4. Note that $``$ can be identified with the ring of germs of meromorphic function $`f:(ϵ,ϵ)𝐂`$. The canonical isomorphism $`s^{}:X𝔢^{}𝔢^{}(X)`$ is given by formula (1-13) in (evaluation and conjugation).
As in 7.4, $`𝒞`$ denotes the von Neumann category of measurable fields of finite dimensional Hilbert spaces over $`S^1`$ and $`𝒯(𝒞)`$ denotes its torsion subcategory. Fix a point $`cS^1`$. We will describe now a covariant functor
$`F_c:𝒯^{}𝒯,`$
compatible with the dualities $`𝔢`$ in $`𝒯`$ (cf. , §3) and $`𝔢^{}`$ in $`𝒯^{}`$. Given a torsion object $`X\mathrm{Ob}(𝒯^{})`$, we may find a free resolution
$`0𝒪^n\stackrel{\alpha }{}𝒪^nX0,`$ (7.11)
where $`\alpha =(a_{ij}(\tau ))`$ is an $`n\times n`$-matrix over $`𝒪`$, such that $`\alpha `$ is invertible over the ring of fractions $``$. It follows that for some small $`ϵ>0`$ all functions $`a_{ij}(\tau )`$ are defined and analytic for $`\tau [ϵ,ϵ]`$ and the matrix $`(a_{ij}(\tau ))`$ is invertible for $`\tau [ϵ,ϵ]`$, $`\tau 0`$. Let $`I_ϵS^1`$ be the interval $`I_ϵ=\{c\mathrm{exp}(i\varphi );\varphi (ϵ,ϵ)\}`$ on the unit circle. We define
$`F_c(X)=((a_{ij}):\mathrm{}^2(I_ϵ,\mu )^n\mathrm{}^2(I_ϵ,\mu )^n)\mathrm{Ob}(𝒯).`$ (7.12)
Here $`\mu `$ denotes the Lesbegue measure on $`S^1`$ and $`(a_{ij})`$ acts by the matrix multiplication. $`\mathrm{}^2(I_ϵ,\mu )^n`$ is the space of $`L^2`$ section of a trivial bundle of rank $`n`$ over $`I_ϵ`$. It is clear that the obtained object $`F_c(X)`$ is torsion and it is independent of the choice of $`ϵ`$. To show that $`F_c`$ is functorial, suppose that
$`0𝒪^n^{}\stackrel{\alpha ^{}}{}𝒪^n^{}X^{}0`$
is another torsion object of $`𝒯^{}`$. Any $`𝒪`$-homomorphism $`h:XX^{}`$ leads to a commutative diagram
$`\begin{array}{ccccccccc}0& & 𝒪^n& \stackrel{\alpha }{}& 𝒪^n& & X& & 0,\\ & & \varphi & & \psi & & h& & \\ 0& & 𝒪^n^{}& \stackrel{\alpha ^{}}{}& 𝒪^n^{}& & X^{}& & 0,\end{array}`$
and hence we obtain the following morphism of $`𝒯^{}`$
$`\begin{array}{cccc}(\alpha :& (\mathrm{}^2(I_\delta ,\mu ))^n& & (\mathrm{}^2(I_\delta ,\mu ))^n)\\ & \varphi & & \psi & & \\ (\alpha ^{}:& (\mathrm{}^2(I_\delta ,\mu ))^n^{}& & (\mathrm{}^2(I_\delta ,\mu ))^n^{}),\end{array}`$
which, viewed as morphism of the extended category $`(𝒞)`$ (cf. , §1.3), clearly depends only on $`h`$ and does not depend on the choice of $`\varphi `$ and $`\psi `$; we will denote it $`F_c(h):F_c(X)F_c(X^{})`$. Here $`\delta >0`$ is some small number.
Let us show that there is a canonical isomorphism $`\gamma :F_c𝔢^{}𝔢F_c`$. If $`𝒳\mathrm{Ob}(𝒯^{})`$ is given by resolution (7.11), where $`\alpha =(a_{ij})`$, then it is clear that $`𝔢^{}(X)`$ has resolution
$`0𝒪^n\stackrel{\alpha ^{}}{}𝒪^n𝔢^{}(X)0,`$
where $`\alpha ^{}=(\overline{a_{ji}})`$ (transposition and conjugation) and hence $`F_c(𝔢^{}(X))`$ is given by
$`((\overline{a_{ji}}):\mathrm{}^2(I_ϵ,\mu )^n\mathrm{}^2(I_ϵ,\mu )^n)\mathrm{Ob}(𝒯).`$ (7.13)
The same object (7.13) is obtained if we first apply $`F_c(X)`$ (cf. (7.12)) and then the duality $`𝔢(F_c(X))`$.
Using $`F_c`$ and $`\gamma `$ we may apply the construction of push forward Hermitian forms, cf. §2.1. Namely, if $`\psi :X𝔢^{}(X)`$ is a Hermitian form in $`𝒯^{}`$ then the push forward form $`\varphi :F_c(X)𝔢(F_c(X))`$ is a Hermitian form in $`𝒯`$.
###### Theorem 7.4.
Let $`cS^1`$ and let $`\pi _1(K)𝒪`$ be given by $`t(\tau c\mathrm{exp}(i\tau ))𝒪`$. This determines a local system $`^𝒪`$ of free $`𝒪`$-modules over $`K`$ of rank one. The homology $`H_q(K;^𝒪)`$ is a direct sum of a free module and a torsion submodule; we will denote the latter by $`T_q^𝒪`$. As in , we have the linking form $`^𝒪:T_q^𝒪𝔢^{}(T_q^𝒪)`$. On the other hand, consider an interval $`I_ϵS^1`$, containing $`c`$ in its interior, and let $`𝒞_ϵ`$ be the von Neumann category of measurable fields of finite dimensional Hilbert spaces over $`I_ϵ`$. The multiplication operator $`\mathrm{}^2(I_ϵ,\mu )\mathrm{}^2(I_ϵ,\mu )`$ by the function $`\mathrm{id}:I_ϵS^1𝐂`$ determines a flat unitary bundle $``$ over $`K`$ with fiber and monodromy in category $`𝒞_ϵ`$. Let $`T_q=T(_q(K;)=T(_q(K;\mathrm{}^2(I_ϵ,\mu )))`$ denote the torsion part of the extended $`L^2`$ homology in dimension $`q`$. Then for small enough $`ϵ`$ the linking form $`^ϵ:T_q𝔢(T_q)`$ is congruent to the push forward, with respect to the pair $`(F_c,\gamma )`$, cf. §1.1, of the form $`^𝒪`$ .
Note that $`𝒞_ϵ`$ can be naturally viewed as a full subcategory of $`𝒞`$; hence the linking form $`^ϵ`$ can be viewed as a Hermitian form in $`𝒯`$.
$`^ϵ`$ is a localized version (at point $`c`$) of the form (7.8).
Theorem 7.4 allows to recover the form $`^ϵ`$ in terms of the linking form of the deformation, which was studied in .
### 7.7. Proof of Theorem 7.4
Consider two mutually dual triangulations of $`K`$. Consider diagram (4.4. Tensoring it by $`𝒪`$ we obtain a commutative diagram of the form
$`\begin{array}{ccc}F_1& \stackrel{\alpha }{}& F_0\\ \left(1\right)^{q+1}h^{}& & h& & \\ F_0^{}& \stackrel{\alpha ^{}}{}& F_1^{}\end{array}`$ (7.14)
consisting of free finitely generated $`𝒪`$-modules such that:
* the torsion submodule of $`\mathrm{coker}(\alpha )`$ is $`T_q^𝒪`$;
* the torsion submodule of $`\mathrm{coker}(\alpha ^{})`$ is canonically isomorphic to $`𝔢^{}(T_q^𝒪)`$;
* the vertical maps of diagram (7.14) determine an isomorphism $`T_q^𝒪𝔢^{}(T_q^𝒪)`$, which coincides with the linking form $`^𝒪`$.
The star $``$ in diagram (7.14) denotes the duality for finitely generated projective modules, cf. , §1.4. For example $`F_0^{}`$ denotes the set of all anti-linear homomorphisms $`F_0𝒪`$. Note that in , §1, we defined a cohomological version of the linking form $`^𝒪`$; in statement (c) above we refer to its homological analogue.
Morphism $`\alpha `$ and $`h`$ in (7.14) can be represented by matrices with entries in $`𝒪`$. Hence, we may find some $`ϵ>0`$ such that all matrix elements of $`\alpha `$ and $`h`$ are analytic on $`[ϵ,ϵ]`$. Hence, ”tensoring” (7.14) by $`\mathrm{}^2(I_ϵ,\mu )`$, we obtain the following commutative diagram
$`\begin{array}{ccc}(\mathrm{}^2(I_ϵ,\mu ))^{n_1}& \stackrel{\stackrel{~}{\alpha }}{}& (\mathrm{}^2(I_ϵ,\mu ))^{n_0}\\ \left(1\right)^{q+1}\stackrel{~}{h}^{}& & \stackrel{~}{h}& & \\ (\mathrm{}^2(I_ϵ,\mu ))^{n_0}& \stackrel{\stackrel{~}{\alpha }^{}}{}& (\mathrm{}^2(I_ϵ,\mu ))^{n_1}\end{array}`$ (7.15)
where $`n_i`$ denotes the rank of $`F_i`$, $`i=0,1`$. Here $`\stackrel{~}{\alpha }`$ and $`\stackrel{~}{h}`$ are given by the same matrices as $`\alpha `$ and $`h`$, where instead of germs we have chosen analytic functions $`I_ϵ𝐂`$. More precisely, we use correspondence between functions $`f:(ϵ,ϵ)𝐂`$ and functions $`\stackrel{~}{f}:I_ϵ𝐂`$, where $`\stackrel{~}{f}(c\mathrm{exp}(i\tau ))=f(\tau )`$, $`\tau [ϵ,ϵ]`$.
We claim that for $`ϵ>0`$ small enough:
* the torsion submodule of $`\mathrm{coker}(\stackrel{~}{\alpha })`$ is $`T_q`$, the torsion part of extended homology $`_q(K;\mathrm{}^2(I_ϵ,\mu ))`$;
* the torsion submodule of $`\mathrm{coker}(\stackrel{~}{\alpha }^{})`$ is canonically isomorphic to $`𝔢(T_q)`$;
* the vertical maps of diagram (7.15) determine an isomorphism $`T_q𝔢(T_q)`$, which coincides with the linking form $``$.
Here we view $`\stackrel{~}{\alpha }`$ and $`\stackrel{~}{\alpha }^{}`$ as morphisms of the extended category, and their cokernels are understood in this sense.
The statements (a’), (b’), (c’) follow from the discussion of section 4.4 (passage from diagram (4.4) to diagram (4.5).
Let us show that for small enough $`ϵ`$ there is a canonical isomorphism $`F_c(T_q^𝒪)T_q`$ and via this isomorphism the form $``$ turns into the push forward form of $`^𝒪`$.
To show this, first note that without loss of generality we may assume that in diagram (7.14) $`𝒪`$-homomorphism $`\alpha `$ is injective. Indeed, if it is not injective we may represent $`F_1`$ as a direct sum $`KF_1^{}`$, where $`K`$ is the kernel of $`\alpha `$ and $`\alpha |_{F_1^{}}`$ is injective. Then $`\alpha ^{}`$ takes values in $`(F_1^{})^{}`$. Therefore, replacing $`F_1`$ by $`F_1^{}`$ we arrive at a similar diagram with injective $`\alpha `$.
Let $`F_0^{}F_0`$ be the smallest $`𝒪`$-submodule containing $`\mathrm{im}(\alpha )`$ such that $`F_0/F_0^{}`$ has no torsion. Then $`F_0F_0^{}(F_0/F_0^{})`$ and $`\alpha ^{}`$ vanishes on $`(F_0/F_0^{})^{}`$. Hence the diagram
$`\begin{array}{ccc}F_1& \stackrel{\alpha }{}& F_0^{}\\ \left(1\right)^{q+1}h^{}& & h& & \\ F_{0}^{}{}_{}{}^{}& \stackrel{\alpha ^{}}{}& F_1^{}\end{array}`$
will replace (7.14) and here $`\alpha `$ is injective and $`\mathrm{coker}(\alpha )`$ is torsion.
Both above changes (cutting off the kernel and the free part of the cokernel of $`\alpha `$) on diagram (7.14) can be performed simultaneously on the corresponding diagram (7.15) assuming that $`ϵ`$ is small enough. Both changes do not influence the resulting torsion forms. Hence we may assume that the initial diagram (7.14) has injective $`\alpha `$ and $`\mathrm{coker}(\alpha )`$ is torsion. But then our italicized above statement follows directly from the definitions.
This completes the proof of Theorem 7.4. ∎
### 7.8. End of proof of Theorems 7.2 and 7.3
Consider first the case of $`q`$ odd, when $``$ and $`^𝒪`$ are Hermitian. We will assume that the interval $`J`$ contains $`c`$ as an internal point as assumed in Theorem 7.2.
Because of Theorem 7.4 and Remark 2.12 in , the form $``$ can be represented as an orthogonal sum of forms $`_{k,\pm }`$, which are defined as the discriminant form of
$`(\pm t^k:\mathrm{}^2([ϵ,ϵ],\mu )\mathrm{}^2([ϵ,ϵ],\mu )).`$ (7.16)
Here $`\mu `$ is the Lesbegue measure on the interval $`[ϵ,ϵ]`$ and $`\pm t^k`$ stands for the multiplication operator by the function $`\pm t^k`$. It is clear that
$`𝔠_+(_{k,+})=k,𝔠_{}(_{k,+})=\{\begin{array}{cc}k,\text{if }k\text{ is odd}\hfill & \\ 0,\text{if }k\text{ is even}\hfill & \end{array}`$ (7.17)
and similarly
$`𝔠_{}(_{k,})=k,𝔠_+(_{k,})=\{\begin{array}{cc}k,\text{if }k\text{ is odd}\hfill & \\ 0,\text{if }k\text{ is even}\hfill & \end{array}`$ (7.18)
Statements (7.17) and (7.18) are based on the calculation similar to that made in Theorem 4.13 in . If the form $``$ contains $`n_k^\pm (c)`$ copies of the form $`_{k,\pm }`$ we obtain that $`𝔠_+()`$ equals the maximum of numbers $`k`$, such that $`n_k^+>0(c)`$ or $`k`$ is odd and $`n_k^{}>0(c)`$ (by Proposition 7.1).
This proves Theorem 7.2 in case $`q`$ odd; here we use Proposition 9.8 from , expressing the linking form of deformation (defined in §1 of , cf. formula (6)) in terms of the Blanchfield form.
Assuming that the interval $`J`$ contains $`c`$ as its terminal point (as in Theorem 7.3) we will have instead of (7.17), (7.18)
$`𝔠_+(_{k,+})=𝔠_{}(_{k,})=\{\begin{array}{cc}k,\text{if }k\text{ is even}\hfill & \\ 0,\text{if }k\text{ is odd}\hfill & \end{array}`$
and
$`𝔠_+(_{k,})=𝔠_{}(_{k,+})=\{\begin{array}{cc}0,\text{if }k\text{ is even}\hfill & \\ k,\text{if }k\text{ is odd}\hfill & \end{array}`$
Taking maximum and using Proposition 7.1 leads to (7.10).
The case when $`q`$ is even is similar. ∎
## 8. The torsion signature
In this section we introduce another numerical invariant of torsion Hermitian forms - the torsion signature.
### 8.1.
Let $`𝒞`$ be a superfinite von Neumann category and let $`\mathrm{tr}`$ be a fixed not normal (i.e. Dixmier type) trace on $`𝒞`$. We refer to , , where these notions are described in detail.
In we introduced the notion of torsion dimension
$$𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)𝐑$$
for any torsion object $`𝒳\mathrm{Ob}((𝒞))`$ with respect to the trace $`\mathrm{tr}`$. This concept of dimension for torsion objects behaves similarly to the well known von Neumann dimension for projective objects. In particular,
$`𝔱𝔬𝔯𝔡𝔦𝔪(𝒳𝒴)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)+𝔱𝔬𝔯𝔡𝔦𝔪(𝒴).`$ (8.1)
Compare the property
$`𝔠(𝒳𝒴)=\mathrm{max}\{𝔠(𝒳),𝔠(𝒴)\}`$ (8.2)
for the Novikov-Shubin capacity.
Recall that $`𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)`$ is defined as the infimum of $`dimP`$ (with respect to the trace $`\mathrm{tr}`$), such that $`P`$ is projective and can be mapped onto $`𝒳`$. The torsion dimension $`𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)`$ can also be characterized and computed in terms of the spectral density function, cf. .
Using the notion of torsion dimension, we may define a new numerical invariants of torsion Hermitian forms $`:𝒳𝔢(𝒳)`$. Namely, by Theorem 7.7 of we can uniquely split $`𝒳=𝒳_+𝒳_{}`$ as the orthogonal sum of positively and negatively definite forms. Now we set
$`𝔱𝔬𝔯𝔰𝔦𝔤𝔫()=𝔱𝔬𝔯𝔰𝔦𝔤𝔫_{\mathrm{tr}}()=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳_+)𝔱𝔬𝔯𝔡𝔦𝔪(𝒳_{}).`$ (8.3)
We will call number (8.3) the torsion signature of form $``$.
It is clear, that the torsion signature of a hyperbolic form is zero.
It is certainly not true that the torsion signature of a metabolic form vanishes (since, as we know from ,§5, all torsion forms are metabolic).
Note also that the torsion signature of an orthogonal sum $`_1_2`$ equals the sum of their torsion signatures
$`𝔱𝔬𝔯𝔰𝔦𝔤𝔫(_1_2)=𝔱𝔬𝔯𝔰𝔦𝔤𝔫(_1)+𝔱𝔬𝔯𝔰𝔦𝔤𝔫(_2).`$ (8.4)
### 8.2. Torsion signatures of odd-dimensional manifolds
Let $`K`$ be a closed $`(2q+1)`$-dimensional closed oriented manifold and let $`K`$ be a unitary flat bundle with fiber and monodromy in a superfinite von Neumann category $`𝒞`$ supplied with a choice of a Dixmier type trace. By Theorem 4.2, there is a well-defined $`(1)^{q+1}`$-Hermitian non-degenerate linking form
$`:T(_q(K;))𝔢(T(_q(K;))).`$
We will define the torsion signature of $`K`$ as follows
$`𝔱𝔬𝔯𝔰𝔦𝔤𝔫_{\mathrm{tr}}(K;)=\{\begin{array}{cc}𝔱𝔬𝔯𝔰𝔦𝔤𝔫_{\mathrm{tr}}(),\text{if}q\text{is odd,}\hfill & \\ 𝔱𝔬𝔯𝔰𝔦𝔤𝔫_{\mathrm{tr}}(\sqrt{1}),\text{if}q\text{is even.}\hfill & \end{array}`$
The torsion signature is a homotopy invariant of $`K`$.
Changing the orientation of $`K`$ reverses the sign of the torsion signature:
$`𝔱𝔬𝔯𝔰𝔦𝔤𝔫_{\mathrm{tr}}(K;)=𝔱𝔬𝔯𝔰𝔦𝔤𝔫_{\mathrm{tr}}(K;).`$
There is always a canonical choice for the flat bundle $``$; namely, we may take for $``$ the bundle with fiber $`\mathrm{}^2(\pi _1(K))`$ with the standard action of the fundamental group $`\pi _1(K)`$.
Theorem 5.1 gives a necessary condition for vanishing of the torsion signature: it happens if our odd-dimensional manifold can be obtained as a slice of an even-dimensional manifold with vanishing middle-dimensional extended homology. In section 9 we will give a different criterion for vanishing of the torsion signature.
### 8.3. The torsion signature and the Blanchfield form
Here we will compute explicitly the torsion signature for closed oriented manifolds $`K`$, with $`dimK=2q+1`$, $`\pi _1(K)=𝐙`$; we will take for the flat bundle $``$ the natural flat bundle with fiber $`\mathrm{}^2(𝐙)=\mathrm{}^(S^1)`$, cf. section 7.4.
As in section 7.4, our von Neumann category $`𝒞`$ is the category of finite dimensional Hilbert spaces over the circle $`S^1`$. In Theorems 7.2 and 7.3 we considered two different traces on $`𝒞`$ and both these traces were normal. Now we will consider non-normal traces on $`𝒞`$, which we will construct as follows. Let $`c𝒮_q(K)`$ be a point of the unit circle, which belongs to the support of the torsion submodule of $`H_q(\stackrel{~}{K};𝐂)`$, cf. 7.4. Let $`\xi (\xi )`$ be a measurable field of finite dimensional Hilbert spaces over $`S^1`$ and let $`f:`$ be a decomposable linear map, cf. , part II, chapter 2. For almost all $`\xi S^1`$ we have a linear map $`f(\xi ):(\xi )(\xi )`$ and we will denote by $`\text{Tr}(f(\xi ))`$ its usual trace (the sum of diagonal entries in a matrix representation). We will define now two traces $`\mathrm{tr}_\pm `$ on category $`𝒞`$ by setting
$`\mathrm{tr}_+(f)=\text{Lim}_\omega [n{\displaystyle _0^{1/n}}\text{Tr}(f(c\mathrm{exp}(i\varphi )))𝑑\varphi ]`$ (8.5)
and
$`\mathrm{tr}_{}(f)=\text{Lim}_\omega [n{\displaystyle _{1/n}^0}\text{Tr}(f(c\mathrm{exp}(i\varphi )))𝑑\varphi ]`$ (8.6)
Note that the function $`\xi \text{Tr}(f(\xi )))`$, where $`\xi S^1`$, is essentially bounded; hence
$$n[n_0^{1/n}\text{Tr}(f(c\mathrm{exp}(i\varphi ))d\varphi ]$$
is a bounded sequence on complex numbers, and therefore we may apply in (8.5) the Banach limit $`\text{Lim}_\omega `$ as $`n\mathrm{}`$. The same explanation applies to (8.6). Formulae (8.5) and (8.6) define two different non-normal traces on $`𝒞`$. Both these traces are ”supported” at point $`cS^1`$, but ”they live on different sides of $`c`$”.
###### Theorem 8.1.
Let $`K`$ be a closed oriented $`(2q+1)`$-dimensional manifold with $`\pi _1(K)=𝐙`$, and let $`:T_q𝔢(T_q)`$, where $`T_q=T(_q(K;\mathrm{}^2(𝐙)))`$, be the linking form. Then the torsion signature of $``$ with respect to the trace $`\mathrm{tr}_+`$ equals
$`𝔱𝔬𝔯𝔰𝔦𝔤𝔫_{\mathrm{tr}_+}()=\sigma _{\text{ev}}(c)+\sigma _{\text{odd}}(c),`$ (8.7)
and the torsion signature of $``$ with respect to the trace $`\mathrm{tr}_{}`$ equals
$`𝔱𝔬𝔯𝔰𝔦𝔤𝔫_{\mathrm{tr}_{}}()=\sigma _{\text{ev}}(c)\sigma _{\text{odd}}(c).`$ (8.8)
Here
$`\sigma _{\text{ev}}(c)={\displaystyle \underset{k\text{even}}{}}[n_k^+(c)n_k^{}(c)]\text{and}\sigma _{\text{odd}}(c)={\displaystyle \underset{k\text{odd}}{}}[n_k^+(c)n_k^{}(c)],`$ (8.9)
where $`n_k^\pm (c)`$ are the numbers determined by the Blanchfield form, cf. section 7.5.
###### Proof.
The proof is similar to the proof of Theorems 7.2 and 7.3 described above in full detail. It is based on Theorem 7.4, which explicitly computes the linking form $``$. The distinction happens only at the very last stage of the proof, when one computes the torsion signatures of the elementary forms (7.16) and uses additivity (8.4). ∎
Comparing Theorems 7.2, 7.3 and 8.1 we conclude, that the torsion signature is more informative than the Novikov - Shubin signature.
The numbers $`\sigma _{\text{odd}}(c)`$ and $`\sigma _{\text{ev}}(c)`$ play an important role in , where they describe the behavior of the eta-invariant.
## 9. Excess of extension
We know that the torsion signature of a hyperbolic form vanishes. In this section we will give a different criterion for vanishing of the torsion signature.
Let $`𝒞`$ be a von Neumann category with a non-normal trace $`\mathrm{tr}`$. Everywhere in this section we will assume that the trace $`\mathrm{tr}`$ is non-negative, cf. .
For any short exact sequence
$`0𝒳^{}𝒳𝒳^{\prime \prime }0,`$ (9.1)
consisting of torsion objects of the extended abelian category $`(𝒞)`$, holds
$`\mathrm{max}\{𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{}),𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{\prime \prime })\}𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{})+𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{\prime \prime }),`$ (9.2)
cf. . Moreover, if sequence (9.1) splits, then $`𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{})+𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{\prime \prime }).`$ We may use inequality (9.2) in order to measure the complexity of extension (9.1).
###### Definition 9.1.
Given extension (9.1), define its excess (with respect to the chosen trace) as the following real number
$`\mathrm{xc}(𝒳^{}𝒳𝒳^{\prime \prime })=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{})+𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{\prime \prime })𝔱𝔬𝔯𝔡𝔦𝔪(𝒳).`$ (9.3)
Non triviality of the excess (with respect to at least one trace $`\mathrm{tr}`$ on $`𝒞`$) implies, that the extension does not split. Note that
$`0\mathrm{xc}(𝒳^{}𝒳𝒳^{\prime \prime })\mathrm{min}\{𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{}),𝔱𝔬𝔯𝔡𝔦𝔪(𝒳^{\prime \prime })\}.`$ (9.4)
Lemma 9.4 below provides many examples showing that the excess may be nonzero.
###### Theorem 9.2.
If a non-degenerate torsion Hermitian form $`:𝒳𝔢(𝒳)`$ in $`(𝒞)`$ admits a metabolizer $`𝒴𝒳`$, $`𝒴^{}=𝒴`$, with vanishing excess
$`\mathrm{xc}(𝒴𝒳𝒳/𝒴)=0,`$ (9.5)
then the torsion signature of $``$ vanishes:
$`𝔱𝔬𝔯𝔰𝔦𝔤𝔫()=0.`$ (9.6)
The proof of Theorem 9.2 will be given later at the end of this section. The proof will use the following Lemmas.
###### Lemma 9.3.
If $`𝒴𝒳`$ is a metabolizer, then vanishing of excess (9.5) is equivalent to
$`2𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳).`$ (9.7)
###### Proof.
For any subobject $`𝒴𝒳`$ we have
$`𝔱𝔬𝔯𝔡𝔦𝔪(𝒳/𝒴)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒴^{})`$ (9.8)
because of the canonical isomorphism $`𝒴^{}𝔢(𝒳/𝒴)`$, and $`𝔢(𝒳/𝒴)`$ is isomorphic to $`𝒳/𝒴`$ (not canonically). If $`𝒴`$ is the metabolizer, then $`𝔱𝔬𝔯𝔡𝔦𝔪(𝒳/𝒴)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)`$, (from (9.8)) and hence vanishing of excess (9.5) is equivalent to (9.7). ∎
###### Lemma 9.4.
Suppose that $`:𝒳𝔢(𝒳)`$ is a non-degenerate positive definite torsion Hermitian form. Then for the metabolizer $`𝒴𝒳`$ (which is unique by Proposition 6.2 of ) holds
$`\mathrm{xc}(𝒴𝒳𝒳/𝒴)=𝔱𝔬𝔯𝔡𝔦𝔪𝒳=𝔱𝔬𝔯𝔡𝔦𝔪𝒴=dim𝒳/𝒴,`$ (9.9)
i.e. this extension has the maximal possible excess.
Any subobject $`𝒵𝒳`$, such that $`𝒵^{}𝒵`$, has torsion dimension equal to torsion dimension of $`𝒳`$, $`𝔱𝔬𝔯𝔡𝔦𝔪𝒵=𝔱𝔬𝔯𝔡𝔦𝔪𝒳`$.
###### Proof.
From the proof of Proposition 6.2 in we know that without loss of generality we may assume that $`𝒳`$ is represented as $`(\alpha :AA)`$, where $`A`$ is an object of $`𝒞`$ and $`\alpha `$ is a self-adjoint positive operator, $`\alpha \mathrm{hom}_𝒞(A,A)`$, and the metabolizer $`𝒴𝒳`$ is represented as $`(\beta :AA)`$ with $`\beta `$ being the positive square root of $`\alpha `$. If $`F(\lambda )`$ is the spectral density function of $`𝒳`$ then $`F(\lambda ^2)`$ is the spectral density function of $`𝒴`$. Therefore we obtain
$$𝔱𝔬𝔯𝔡𝔦𝔪𝒳=\underset{\lambda +0}{lim}F(\lambda )=\underset{\lambda +0}{lim}F(\lambda ^2)=𝔱𝔬𝔯𝔡𝔦𝔪𝒴.$$
Now, formula (9.8) shows that
$$\mathrm{xc}(𝒴𝒳𝒳/𝒴)=\mathrm{\hspace{0.17em}2}𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)𝔱𝔬𝔯𝔡𝔦𝔪(𝒳),$$
which proves (9.7).
If $`𝒵^{}𝒵`$ then $`𝒵`$ contains the metabolizer of $`𝒳`$ (by the arguments in the proof of Propositions 6.1 and 6.2 in ) and therefore $`𝔱𝔬𝔯𝔡𝔦𝔪(𝒵)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)`$. This completes the proof. ∎
The following Corollary is a partial result in the direction of Theorem 9.2; it takes care of the case of reduction of a positively defined form with respect to a submodule with zero excess.
###### Corollary 9.5.
Let $`:𝒳𝔢(𝒳)`$ be a non-degenerate positive definite torsion Hermitian form. Suppose that $`𝒴𝒳`$ is a subobject such that $`𝒴𝒴^{}`$. Then the following conditions are equivalent:
* the excess $`\mathrm{xc}(𝒴𝒳𝒳/𝒴)=0`$ vanishes;
* $`𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)=0.`$
Either of these conditions implies
$`𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒴^{}/𝒴).`$ (9.10)
###### Proof.
If $`𝒴𝒴^{}`$ then by Lemma 9.4
$$𝔱𝔬𝔯𝔡𝔦𝔪(𝒴^{})=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳).$$
Thus, if the excess (i) vanishes, we get
$$\begin{array}{c}𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)+𝔱𝔬𝔯𝔡𝔦𝔪(𝒳/𝒴)\hfill \\ =𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)+𝔱𝔬𝔯𝔡𝔦𝔪(𝒴^{})\hfill \\ =𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)+𝔱𝔬𝔯𝔡𝔦𝔪(𝒳),\hfill \end{array}$$
which implies $`𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)=0`$.
Conversely, if $`𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)=0`$, then inequality (9.4) shows that $`\mathrm{xc}(𝒴𝒳𝒳/𝒴)=0`$.
If (ii) holds, then $`\mathrm{xc}(𝒴𝒴^{}𝒴^{}/𝒴)=0`$ (again, because of (9.4)) and
$$𝔱𝔬𝔯𝔡𝔦𝔪(𝒴^{}/𝒴)=dim(𝒴^{})=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳).$$
This completes the proof. ∎
### Proof of Theorem 9.2
Let $`𝒳=𝒳_+𝒳_{}`$ be the canonical decomposition determined by $``$ (cf. , Theorem 7.7). Denote by $`𝒴_\pm =i_\pm (𝒴)𝒳_\pm `$ the image of $`𝒴`$ under the projections $`j_\pm :𝒳𝒳_\pm `$. Then clearly $`𝒴_\pm ^{}𝒴_\pm `$ (otherwise we will get a contradiction to $`𝒴^{}=𝒴`$). Applying the second statement of Lemma 9.4 we find
$`𝔱𝔬𝔯𝔡𝔦𝔪(𝒴_\pm )=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳_\pm ).`$ (9.11)
On the other hand,
$`𝔱𝔬𝔯𝔡𝔦𝔪(𝒴_\pm )𝔱𝔬𝔯𝔡𝔦𝔪(𝒴).`$ (9.12)
Therefore we obtain
$`𝔱𝔬𝔯𝔡𝔦𝔪(𝒳_+)𝔱𝔬𝔯𝔡𝔦𝔪(𝒴),`$ (9.13)
$`𝔱𝔬𝔯𝔡𝔦𝔪(𝒳_{})𝔱𝔬𝔯𝔡𝔦𝔪(𝒴).`$ (9.14)
Now we will use $`\mathrm{xc}(𝒴𝒳𝒳/𝒴)=0`$ to conclude
$$2𝔱𝔬𝔯𝔡𝔦𝔪(𝒴)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳_+)+𝔱𝔬𝔯𝔡𝔦𝔪(𝒳_{})$$
This shows that inequalities (9.13) and (9.14) are in fact equalities, and therefore
$$𝔱𝔬𝔯𝔡𝔦𝔪(𝒳_+)=𝔱𝔬𝔯𝔡𝔦𝔪(𝒳_{})\text{and}𝔱𝔬𝔯𝔰𝔦𝔤𝔫()=0.$$
This completes the proof. ∎
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# SN 1998bw @ late phases1footnote 11footnote 1Based on observations collected at the European Southern Observatory, La Silla and Paranal, Chile.
## 1. Introduction
Supernova (SN) 1998bw was born famous because of its positional and temporal coincidence with the gamma-ray burst GRB 980425 (Galama et al. 1998a). The unique nature of the associated radio source, in terms of the high luminosity immediately after the explosion and the inferred relativistic expansion rate (Kulkarni et al. 1998), argues that SN 1998bw and GRB 980425 are related. Independent of the connection to the GRB, SN 1998bw is an interesting supernova. It was optically luminous and displayed very high expansion velocities. SN 1998bw was classified as a Type Ic (Patat & Piemonte 1998), and late time spectra reported here are consistent with this classification.
Modeling of the early light curve and spectra suggested an extremely energetic explosion of a massive carbon-oxygen star (Iwamoto et al. 1998; Woosley, Eastman, & Schmidt 1999). As the energy was more than ten times that of a canonical core-collapse supernova, the term ’hypernova’ was suggested for this event (Iwamoto et al. 1998). The mass of <sup>56</sup>Ni needed to power the early light curve in these models, $`0.50.7\mathrm{M}_{}`$, is much larger than the $`\stackrel{}{\stackrel{<}{}}\mathrm{\hspace{0.17em}\hspace{0.17em}0.1}\mathrm{M}_{}`$, typical for ’normal’ core collapse SNe (Patat et al. 1994; Schmidt et al. 1994). This large mass of <sup>56</sup>Ni for SN 1998bw was disputed by Höflich, Wheeler, & Wang (1999), who found that the early light curve can also be reproduced by an asymmetric SN Ic ejecting only $`0.2\mathrm{M}_{}`$ of <sup>56</sup>Ni, given the right viewing angle and degree of asymmetry.
The models for SN 1998bw differ considerably in terms of progenitor mass, nickel mass and explosion energy. We therefore focus on the late time evolution, when the nucleosynthesis and density distribution can be studied directly. The total emission in the nebular phase is also less sensitive to asymmetries, compared to the early light curve.
We present photometry and spectroscopy of SN 1998bw up to 500 days past explosion. We also compare these with results from spectral synthesis based on various progenitor models. A more detailed analysis will be given elsewhere.
## 2. Observations, Reductions and Results
### 2.1. Photometry
Photometry was obtained at eight epochs between 33 and 504 days past explosion, using the ESO-3.6m telescope on La Silla and the VLT/UT1 on Paranal. The images were bias subtracted and flat fielded using IRAF.
To obtain the supernova magnitudes we used DAOPHOT. The instrumental magnitudes were converted using local standards from Galama et al. (1998b)<sup>2</sup><sup>2</sup>2http://www.astro.uva.nl/$``$titus. At the earliest epochs, the errors are estimated as the standard deviations in the magnitudes of the local standards. They also encapsulate the neglect of color transformations.
At later phases, the main uncertainty lies in the background subtraction, as the supernova is superposed on an H II region. We constructed template PSFs using one or several stars, and the best PSF was chosen from judgments of the subtracted images. This PSF was subtracted from the SN position with a range of magnitudes, and the supernova magnitude was taken as the one which gave the smoothest background after subtraction. Upper and lower limits were estimated by noting when the subtracted image showed a hole, or a clear point source, at the supernova position. The supernova magnitudes are given in Table 1 and plotted in Figure 1.
### 2.2. Spectroscopy
Spectroscopic observations were also obtained to late phases (Fig. 2). At ESO-3.6m we used EFOSC2 with grisms 11 and 12, and at VLT/UT1 we used FORS1 with the 300-grisms. All spectra were bias subtracted, flat fielded and wavelength calibrated. Flux calibration was done with a spectrophotometric standard star, and the absolute fluxing was performed using the simultaneous $`V`$-band photometry. All spectra were taken close to the parallactic angle to reduce differential refraction.
The spectral evolution of SN 1998bw from day 33 to day 504 is shown in Figure 2. The evolution of the spectra in the nebular phase is fairly slow. The \[O I\] $`\lambda \lambda 6300,6364`$ are the strongest lines, and the 7300 and 8600 Å features, likely to be due to Ca II, decrease in strength relative to the \[O I\] lines. Other prominent features are seen at $`4570`$ and 5890 Å, and are probably due to Mg I\] (or \[Fe III\]) and Na I D (and possibly He I), respectively.
## 3. Discussion
### 3.1. General considerations
A lower limit of $`0.22\mathrm{M}_{}`$ of <sup>56</sup>Ni for SN 1998bw was achieved by McKenzie & Schaefer (1999) by comparing the $`V`$-magnitude at 170 days to that of SN 1987A. While the light curve of SN 1987A closely followed the decay rate of <sup>56</sup>Co, suggesting full trapping of the $`\gamma `$-rays, the light curve of SN 1998bw fell substantially faster. The assumption of full trapping at 170 days in SN 1998bw used by McKenzie & Schaefer is clearly too conservative. In fact, the light curve of SN 1998bw is similar to that of many other SNe Ib/c (Clocchiatti & Wheeler 1997; Sollerman, Leibundgut, & Spyromilio 1998), and can be reproduced by a simple model of a radioactively powered ejecta leaking $`\gamma `$-rays due to homologous expansion (Fig. 1). If the decay would follow full trapping from day 64 and onward, the flux in the $`V`$-band at day 170 would be 2.4 times brighter than actually observed (Fig. 1). Assuming full trapping at day 64, at the beginning of the fast decay, then provides an upper limit of $`41\%`$ on the trapping at day 170, increasing the lower limit of the Ni-mass to $`0.5\mathrm{M}_{}`$.
The above estimate depends, however, on the bolometric correction and its evolution for these supernovae. That the bolometric decay rate is similar to that of the $`V`$-band at this epoch is indicated by the simple estimates of the bolometric luminosity by McKenzie & Schaefer (1999). This is also confirmed by our detailed modeling (see below) and supported by the fit to the $`V`$-curve by the simple $`\gamma `$-leakage model. However, the bolometric correction could still be different for SNe 1987A and 1998bw. If SN 1998bw emitted a smaller fraction of its bolometric luminosity in the optical range than SN 1987A, its Ni-mass would be overestimated.
Adopting a distance of 35 Mpc ($`H_0=72)`$, a reddening of $`E(BV)=0.06`$ (Schlegel, Finkbeiner, & Davis 1998), and integrating our spectrum at day 141, gives a luminosity ($`L_{\mathrm{opt}}=3.5\times 10^{41}`$ ergs s<sup>-1</sup>) that requires $`0.1\mathrm{M}_{}`$ of <sup>56</sup>Ni. This limit clearly has to be corrected upwards as it assumes full $`\gamma `$-ray trapping, and that our spectrum contains all the flux emitted by the supernova. With the same argument as above, no more than $`48\%`$ of the energy is deposited at 141 days. Furthermore, detailed models (see below) show that the observed spectrum contains only $`71\%`$ of the total flux. These corrections indicate that the lower limit must be increased to $`0.3\mathrm{M}_{}`$. The dominant errors are the uncertainties in bolometric correction ($`10\%`$), distance, and reddening for SN 1998bw. Allowing for a 15$`\%`$ uncertainty in the distance scale, and $`\pm 0.1`$ mag in $`A_V`$, make the absolute fluxes (and the Ni-mass) uncertain by $`35\%`$.
### 3.2. Detailed modeling
To obtain more quantitative estimates of the Ni-mass, detailed modeling is required. Our analysis uses an updated version of the spectral code in Kozma & Fransson (1998a, 1998b). At the epochs discussed here, the ejecta emission is powered by the radioactive decay of <sup>56</sup>Co. The deposition of the resulting $`\gamma `$-rays and positrons is modeled in detail according to Kozma & Fransson (1992). The ionization and heating is balanced by thermal and radiative processes, producing the observed optical and infrared emission. This time-dependent code gives good agreement with light curves and late time spectra for SN 1987A (Kozma & Fransson 1998a, 1998b; Kozma 2000).
As input models we use different combinations of unmixed and artificially mixed models of SN 1998bw, based on the CO138-model calculated by Iwamoto et al. (1998) and the CO6-model from Woosley et al. (1999). Both models are spherically symmetric explosions of massive carbon-oxygen stars. The CO138-model has a mass $`M_{\mathrm{CO}}`$ = 13.8$`\mathrm{M}_{}`$, explosion energy $`E_{\mathrm{exp}}=3\times 10^{52}`$ ergs, and mass of <sup>56</sup>Ni, $`M`$(<sup>56</sup>Ni) = 0.7$`\mathrm{M}_{}`$. Such a carbon-oxygen core corresponds to a zero age main sequence (ZAMS) star of $`M_{\mathrm{ZAMS}}40\mathrm{M}_{}`$. The CO6C-model with $`M_{\mathrm{CO}}`$ = 6.55$`\mathrm{M}_{}`$, $`E_{\mathrm{exp}}=2.2\times 10^{52}`$ ergs, and $`M`$(<sup>56</sup>Ni) = 0.47$`\mathrm{M}_{}`$ corresponds to a $`M_{\mathrm{ZAMS}}25\mathrm{M}_{}`$ star.
In Figure 3 we show our modeled spectra together with the observations at 141 days. From the observed line profiles we conclude that mixing is important. Calculations based on non-mixed models give very broad and flat-topped line profiles which are not seen in the observations. In the mixed models, we have also decreased the expansion velocities to get a better agreement with the observed line widths. A mixed model with the original velocities results in nicely peaked, but too broad, line profiles. A decrease in velocity results in an increase of the densities given by the original explosion models. We have increased the densities by factors of 10 – 100, except in the Fe-rich regions where it is kept unchanged (see below). These changes also affect the kinetic energy in the modified models. In the layers below $`\mathrm{12\hspace{0.17em}000}\mathrm{km}\mathrm{s}^1`$, which is the only material which contributes significantly to the emission at late times, the kinetic energies are $`2\times 10^{51}`$ ergs (CO6) and $`5\times 10^{51}`$ ergs (CO138), respectively. Note that our analysis is not sensitive to the high velocity gas, which might also contain a significant amount of energy. Early time analysis is therefore more reliable for determining the total kinetic energy of the explosion. As seen in Figure 3, the synthetic line profiles are still somewhat broader than the observed. We can overcome this by decreasing the velocities even more, but the high density in the oxygen-rich regions then results in strong O I lines at $`\lambda \lambda 5577,7774,9265`$, which are not seen in the observations.
Most features in the observed spectrum have a counterpart in the modeled spectra, and in Figure 3 we include some of the line identifications. The $`40005500`$ Å regime is problematic, and preliminary results from improved models suggest that most of these features are a mixture of Fe II and Fe III lines. The models presented here have too low a degree of ionization, with mostly Fe I and Fe II emission. A higher density in the Fe-emitting material would result in even lower ionization, which is why we do not increase the density in the Fe-rich regions. The ionization can instead be increased by decreasing the density in these regions. Scattering and fluorescence may also play a significant role in understanding these features. We therefore regard the Fe I emission as dubious.
### 3.3. The Light Curve
In Figure 1 we show the light curves from our calculations of the mixed CO6- and CO138-models (solid lines), and from the original, non-mixed models (dotted lines). The CO138-model is brighter than CO6. This is not only due to the higher Ni-mass, but also a result of the higher ejecta-mass ($`M_{\mathrm{ej}}`$), and lower expansion velocity ($`v_{\mathrm{max}}`$), giving a larger optical depth to the $`\gamma `$-rays ($`\tau _\gamma `$); $`\tau _\gamma M_{\mathrm{ej}}/v_{\mathrm{max}}^2`$.
The effect of changing the expansion velocities can also be seen in Figure 1. While the two CO6-models have the same Ni-mass and ejecta-mass, they differ in expansion velocity. The mixed models are somewhat brighter than the non-mixed models, due to the higher $`\gamma `$-ray optical depth for the models with lower velocities. In the original CO6-model we find that at 200 (400, 600) days only $`7.2(4.4,3.9)\%`$ of the energy released in the <sup>56</sup>Co-decays is deposited in the ejecta. Out of this energy $`47(78,87)\%`$ is from the kinetic energy of the positrons. For the mixed CO6-model with decreased velocities, $`11(5.7,4.5)\%`$ of the energy is deposited.
From the light curves, and within the framework of our modeling, we conclude that the nickel mass needed to power the late emission is indeed very high. For the two models that best reproduce the late time observations, we calculate the deposition of the radioactive energy, and thus the required mass of nickel. The CO6-model requires $`0.9\mathrm{M}_{}`$ of <sup>56</sup>Ni to maintain the luminosity at $`340410`$ days. The CO138-models, which traps $`10\%`$ of the energy at 400 days, requires $`0.5\mathrm{M}_{}`$. The errors on these estimates are dominated by the uncertainty in distance and reddening stated in §3.1.
Höflich et al. (1999) found that the early light curve of SN 1998bw could be explained in terms of a more ’normal’ SN Ic, which ejected 0.2$`\mathrm{M}_{}`$ of <sup>56</sup>Ni, given the right degree of asymmetry. However, at the nebular phase studied here the ejecta are optically thin, and for a given density structure the luminosity is not as sensitive to asymmetries as is the early light curve. By modeling the late time line profiles, providing a measure of the energy input with velocity, we get a handle on the distribution of the ejecta, and can therefore model the gamma-deposition self-consistently.
The Ni-mass determined for the CO138-model is close to the lower limit from §3.1. This is because this massive model closely represents the maximum trapping allowed by the observed light curve. Also, as the mixed CO6-model is already dominated by positrons at 400 days, the estimated Ni-mass for this model must be close to the upper limit.
While SN 1998bw ejected a large mass of <sup>56</sup>Ni, a few other supernovae have shown very low amounts of ejected <sup>56</sup>Ni (SN 1994W \[$`0.015\mathrm{M}_{}`$, Sollerman, Cumming, & Lundqvist 1998\] and SN 1997D \[$`0.002\mathrm{M}_{}`$, Turatto et al. 1998\]). The iron yield from core-collapse supernovae thus seems to vary substantially.
At very late phases ($`>400`$ days), the low energy input means low temperatures for the ejecta, and most of the emission in our models comes out in the infrared. The optical light curves no longer directly follow the bolometric light curve. After 400 days the model light curves drop faster than the observations. A reduced density in the Fe-rich material (§3.2) can increase the temperature and boost in particular the $`V`$-curve at late phases.
Another possible source for the powering at late phases is interaction with a circumstellar medium (CSM). The absence of H and He in the spectra means that substantial amounts of such gas should be present. Models of the radio emission indicate that some interaction is already taking place (Li & Chevalier 1999), although the CSM density in these models is probably too low to affect the optical light curve, as the wind velocity is high ($`2000\mathrm{km}\mathrm{s}^1`$). We see no significant spectroscopic signature that indicates CSM interaction, in our late spectra.
## 4. Summary
We present photometry and spectroscopy of SN 1998bw up to 500 days past explosion. Starting with two different spherically symmetric explosion models, we calculate late time spectra and light curves using a detailed spectral synthesis code to model our observations. We conclude that mixing is important, and we also have to alter the density distribution macroscopically to obtain reasonable agreement with the line profiles in the observed spectra.
Representing the supernova density profile with the distribution which best reproduces the late time observations, we calculate the deposition of the radioactive energy, and determine how much <sup>56</sup>Ni is required to produce the late time luminosity. We find that a large amount, roughly half a solar mass, is needed to power the light curve one year after the explosion, although the exact number depends on the chosen explosion model.
Evidence for asymmetries in SN 1998bw comes from early polarization measurements (Kay et al. 1998). The modeling presented here is restricted to the available symmetric models, although the importance of mixing may point to an asymmetric explosion. Whether asymmetric models can reproduce the observations is of interest for future studies.
We thank the Swedish Natural Science Research Council, the Swedish National Space Board and the Knut and Alice Wallenberg Foundation for support. J. S. is grateful to the grant from ESO Director’s Discretionary fund, and from Holmbergs and Hiertas funds.
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# Integrable highest weight modules over affine superalgebras and Appell’s function
## 0 Introduction
In this paper we continue the study of integrable irreducible highest weight modules over affine superalgebras that we began in \[KW\].
First, let us recall the definition of an integrable module over an ordinary affine Kac-Moody algebra $`\widehat{𝔤}`$ \[K3\]. Let $`𝔤`$ be a finite-dimensional simple or abelian Lie algebra over $``$ with a symmetric invariant bilinear form $`(.|.)`$. Recall that the associated affine algebra is
(0.1)
$$\widehat{𝔤}=([t,t^1]_{}𝔤)Kd$$
with the following commutation relations ($`a,b𝔤;m,n`$ and $`a(m)`$ stands for $`t^ma`$) :
(0.2) $`[a(m),b(n)]=[a,b](m+n)+m\delta _{m,n}(a|b)K,[d,a(m)]=ma(m),[K,\widehat{𝔤}]=0.`$
We identify $`𝔤`$ with the subalgebra $`1𝔤`$. The bilinear form $`(.|.)`$ extends from $`𝔤`$ to a symmetric invariant bilinear form on $`\widehat{𝔤}`$ by:
$`(a(m)|b(n))`$ $`=`$ $`\delta _{m,n}(a|b),([t,t^1]𝔤|K+d)=0,`$
(0.3) $`(K|K)`$ $`=`$ $`(d|d)=0,(K|d)=1.`$
Choose a Cartan subalgebra $`𝔥`$ of $`𝔤`$ and let $`𝔤=𝔥(_{\alpha \mathrm{\Delta }}𝔤_\alpha )`$ be the root space decomposition, where $`𝔤_\alpha `$ denotes the root space attached to a root $`\alpha \mathrm{\Delta }𝔥^{}`$. Let
(0.4)
$$\widehat{𝔥}=𝔥+K+d$$
be the Cartan subalgebra of $`\widehat{𝔤}`$, and, as before, let $`𝔤_\alpha (m)=t^m𝔤_\alpha `$.
A $`\widehat{𝔤}`$-module $`V`$ is called *integrable* if the following two properties hold \[K3\]:
(0.5) $`\widehat{𝔥}\text{ is diagonalizable on }V,`$
(0.6) $`\text{all }𝔤_\alpha (m)(\alpha \text{ a root of }𝔤,m)\text{ are locally finite on }V.`$
(Property (0.6) means that $`dimU(𝔤_\alpha (m))v<\mathrm{}`$ for any $`vV`$.)
It is easy to show that these two properties imply
(0.7) $`𝔤\text{ is locally finite on }V(\text{ i.e., }dimU(𝔤)v<\mathrm{}\text{ for any }vV).`$
Here and further $`U(𝔞)`$ denotes the universal enveloping algebra of a Lie (super)algebra $`𝔞`$. Note also that condition (0.6) is vacuous if $`𝔤`$ is abelian.
Let now $`𝔤=𝔤_{\overline{0}}+𝔤_{\overline{1}}`$ be a finite-dimensional Lie superalgebra over $``$ with an even symmetric invariant bilinear for $`(.|.)`$ (for a background on Lie superalgebras see \[K1\]). Recall that “even” means that $`(𝔤_{\overline{0}}|𝔤_{\overline{1}})=0`$, “symmetric” means that $`(.|.)`$ is symmetric on $`𝔤_{\overline{0}}`$ and skewsymmetric on $`𝔤_{\overline{1}}`$, and “invariant” means that $`([a,b]|c)=(a|[b,c]),a,b,c𝔤`$. We shall assume, in addition, that $`𝔤_{\overline{0}}`$ is reductive:
(0.8)
$$𝔤_{\overline{0}}=_{j=0}^N𝔤_{\overline{0}j},$$
where $`𝔤_{\overline{0}0}`$ is abelian and $`𝔤_{\overline{0}j}`$ with $`j1`$ are simple Lie algebras.
The affine superalgebra $`\widehat{𝔤}`$ associated to the Lie superalgebra $`𝔤`$ and the bilinear form $`(.|.)`$ is defined in exactly the same way as in the Lie algebra case by relations (0.2). Likewise, the invariant even symmetric bilinear form $`(.|.)`$ on $`\widehat{𝔤}`$ is defined by (0.3), and the Cartan subalgebra $`\widehat{𝔥}`$ is defined by (0.4) after a choice of a Cartan subalgebra $`𝔥`$ of $`𝔤_{\overline{0}}`$. Note that for each$`j\{0,1,\mathrm{},N\}`$, the superalgebra $`\widehat{𝔤}`$ contains an affine Kac-Moody algebra $`\widehat{𝔤}_{\overline{0}j}`$ associated to $`𝔤_{\overline{0}j}`$.
We shall see that condition (0.6) of integrability is too strong in the superalgebra case, as for most of the affine superalgebras it allows only trivial highest weight modules. This forces us to consider weaker conditions (cf. \[KW\]):
###### Definition 0.1.
Given a subset $`J\{1,\mathrm{},N\}`$, a $`\widehat{𝔤}`$-module $`V`$ is called $`J`$-*integrable* if it satisfies conditions (0.5) and (0.7) and if it is integrable as $`\widehat{𝔤}_{\overline{0}j}`$-module for all $`jJ`$.
Let $`𝔤=𝔥(_{\alpha \mathrm{\Delta }}𝔤_\alpha )`$ be a root space decomposition of the Lie superalgebra $`𝔤`$ with respect to a Cartan subalgebra $`𝔥`$ of $`𝔤_{\overline{0}}`$. Choose a set of positive toots $`\mathrm{\Delta }_+`$ in $`\mathrm{\Delta }`$ and let $`𝔫_+=_{\alpha \mathrm{\Delta }_+}𝔤_\alpha `$. For each $`\mathrm{\Lambda }\widehat{𝔥}^{}`$ one defines an *irreducible highest weight module* $`L(\mathrm{\Lambda })`$ over $`\widehat{𝔤}`$ as the (unique) irreducible $`\widehat{𝔤}`$-module for which there exists a non-zero vector $`v_\mathrm{\Lambda }`$ such that
(0.9)
$$hv_\mathrm{\Lambda }=\mathrm{\Lambda }(h)v_\mathrm{\Lambda }\text{ for }h\widehat{𝔥},𝔫_+v_\mathrm{\Lambda }=0,𝔤(m)v_\mathrm{\Lambda }=0\text{ for }m>0,$$
where, as before $`𝔤(m)=t^m𝔤`$. The number $`k=\mathrm{\Lambda }(K)`$ is called the *level* of $`L(\mathrm{\Lambda })`$ and of $`\mathrm{\Lambda }`$. Note that $`K=kI`$ on $`L(\mathrm{\Lambda })`$ and that $`\overline{L}(\mathrm{\Lambda }):=U(𝔤)v_\mathrm{\Lambda }`$ is an irreducible highest weight module over $`𝔤`$.
In § 1 we describe a general approach to the classification of irreducible integrable highest weight modules over arbitrary Kac-Moody superalgebras, and in § 2 and § 6 give their complete classification in the affine (non-twisted) case, using Serganova’s odd reflections.
In § 3 we give a free field realization of all level $`1`$ integrable highest weight modules over $`g\mathrm{}(m|n)^\widehat{}`$, which leads to a “quasiparticle” character formula for these modules and to a “theta function” type character formula. This construction may be viewed as a generalization of the classical boson-fermion correspondence based on the oscillator algebra $`g\mathrm{}(1)^\widehat{}`$ and of the super boson-fermion correspondence based on $`g\mathrm{}(1|1)^\widehat{}`$ \[KL\]. The former produces the classical vertex operators and relates representation theory of $`g\mathrm{}(1)^\widehat{}`$ to the denominator identity for $`s\mathrm{}(2)^\widehat{}`$, while the latter produces vertex operators for the symplectic bosons and relates representation theory of $`g\mathrm{}(1|1)^\widehat{}`$ to the denominator identity for $`s\mathrm{}(2|1)^\widehat{}`$ (see \[K4\]).
In § 4 we show that the “theta function” type character formula for $`s\mathrm{}(m|1)^\widehat{}(m2)`$ is a product of a theta function, a power of the eta function, and a more “exotic” function, called a multivariable Appell function. The classical Appell function appeared in the 1880’s in the papers by Appell \[A\] and by Hermite in their study of elliptic functions. Most recently this function has been discussed in \[P\]. The study of asymptotics of Appell’s functions gives the high temperature asymptotics of integrable level 1 $`s\mathrm{}(m|1)^\widehat{}`$-characters. We also derive here formulas for branching functions for integrable level $`1`$ $`s\mathrm{}(m|1)^\widehat{}`$-modules restricted to the even subalgebra. They turn out to be certain “half” modular functions.
In § 5 we relate integrable level $`1`$ modules over $`g\mathrm{}(m|n)^\widehat{}`$ to the denominator identity for $`s\mathrm{}(m+1|n)^\widehat{}`$, and as a result, we derive for these modules yet another, a Weyl type, character formula.
In § 7 we give a free field realization of the two level $`1`$ integrable highest weight modules over $`osp(m|n)^\widehat{}`$, which generalizes the constructions for $`so(m)^\widehat{}`$ and $`sp(n)^\widehat{}`$ from \[KP1\], \[F\] and \[FF\]. These lead to character formulas and high temperature asymptotics of the characters.
In § 8 we show that integrability is a necessary condition for an irreducible highest weight $`\widehat{𝔤}`$-module to be a module over the associated vertex algebra, and that in the level $`1`$ case this condition is sufficient. We thus get examples of rational vertex algebras for which the $``$-span of normalized (super)characters is not $`SL(2,)`$-invariant. The latter property was proved in \[Z\] under certain additional assumptions, and it was generally believed that these assumptions were superfluous.
In § 9 we discuss some open problems.
It is interesting to note that in the “super” case a number of new interesting phenomena occur. The level gets quantized by the integrability condition, but in almost all cases the number of integrable modules is infinite. This is the case for the lowest, level $`1`$, integrable $`s\mathrm{}(m|n)^\widehat{}`$-modules which apparently causes the specialized characters and branching functions to lose their customary modularity properties, which are so ubiquitous in the affine Lie algebra case \[KP2\],\[K3\]. However, in the cases when the number of characters of given level is finite, like, for example, $`k=1osp(m|n)^\widehat{}`$ case, the specialized normalized characters are still modular, though their $``$-span is no longer $`SL(2,)`$-invariant as in the affine Lie algebra case.
It is also interesting to note that while the characters of affine Lie algebras are global sections of line bundles on abelian varieties, the characters of affine Lie superalgebras are related to global sections of rank $`2`$ vector bundles on abelian varieties, as the work of Polishchuk \[P\] on Appell’s function apparently indicates.
We are grateful to A. Polishchuk for giving us the idea of Lemma 4.1 and other very useful remarks. The first named author wishes to thank ENS, IHES and MSRI for their hospitality.
## 1 Integrability of highest weight modules over Kac-Moody superalgebras
Consider the following data:
$$𝒟=\{𝔥,I,I_1,\mathrm{\Pi }^{},\mathrm{\Pi }\},$$
where $`𝔥`$ is a vector space, $`I`$ is an index set, $`I_1`$ is a subset of $`I`$, $`\mathrm{\Pi }^{}=\{\alpha _i^{}\}_{iI}`$ and $`\mathrm{\Pi }=\{\alpha _i\}_{iI}`$ are linearly independent sets of vectors in $`𝔥`$ and $`𝔥^{}`$ respectively indexed by $`I`$. One associates to these data a Lie superalgebra $`𝔤(𝒟)`$ defined as the quotient of the Lie algebra on generators $`e_i,f_i(iI)`$ and $`𝔥`$, the generators $`e_i`$ and $`f_i`$ for $`iI_1`$ being odd and all other generators being even, and the *standard relations* $`(i,jI,h𝔥)`$:
$$[𝔥,𝔥]=0,[e_i,f_j]=\delta _{ij}\alpha _i^{},[h,e_i]=\alpha _i,he_i,[h,f_i]=\alpha _i,hf_i,$$
by the maximal graded with respect to the root space decomposition intersecting $`𝔥`$ trivially (cf. \[K1\], \[K3\] ).
The commutative $`ad`$-diagonizable subalgebra $`𝔥`$ of $`𝔤(𝒟)`$ is called the Cartan subalgebra, $`\mathrm{\Pi }`$ and $`\mathrm{\Pi }^{}`$ are called the sets of simple roots and coroots respectively, elements $`e_i`$ and $`f_i`$ $`(iI)`$ are called Chevalley generators, etc. One defines the notions of roots and root spaces in the usual way (cf. \[K1\], \[K3\]). Let $`𝔫_+`$ (resp. $`𝔫_{}`$) denote the subalgebra of $`𝔤`$ generated by the $`e_i`$’s (resp. $`f_i`$’s). Then, as usual, one has the triangular decomposition:
$$𝔤=𝔫_{}+𝔥+𝔫_+.$$
Let $`a_{ij}=\alpha _j,\alpha _i^{}`$. The matrix $`A=(a_{ij})_{i,jI}`$ is called the Cartan matrix of the data $`𝒟`$ (and of $`𝔤(𝒟)`$).
A root of $`𝔤(𝒟)`$ is called *even* (resp. *odd*) if the attached root space is even (resp. odd). For example a simple root $`\alpha _s`$ is called odd iff $`sI_1`$. An odd simple root $`\alpha _s`$ (and the coroot $`\alpha _s^{}`$) is called *isotropic* if $`a_{ss}=0`$. In what follows we let
$`p_{ij}=\{\begin{array}{cc}\hfill 1& \text{if both }\alpha _i\text{ and }\alpha _j\text{ are odd,}\hfill \\ \hfill 1& \text{othewise.}\hfill \end{array}`$
Note that $`𝔤(𝒟)`$ has an anti-involution $`\omega `$ defined by $`\omega (e_i)=f_i,\omega (f_i)=e_i,\omega |_𝔥=I`$. For that reason properties of the $`e_i`$’s automatically hold for the $`f_i`$’s.
###### Lemma 1.1.
1. An odd simple root $`\alpha _i`$ is isotropic iff $`[e_i,e_i]=0`$.
2. If $`ij`$, then $`[e_i,e_j]=0`$ iff $`a_{ij}=a_{ji}=0`$.
###### Proof.
It is clear that $`[f_j,[e_i,e_i]]=0`$ if $`ji`$, and one has: $`[f_i,[e_i,e_i]]=2a_{ii}e_i`$, which proves (a). The proof of (b) is similar.
It is straightforward to check the following relation $`(i,jI,ij)`$:
(1.2)
$$[[e_i,e_j],[f_i,f_j]]=p_{ij}(a_{ij}\alpha _j^{}p_{ii}a_{ji}\alpha _i^{}).$$
Further on we shall always assume the following property of the Cartan matrix $`A`$:
(1.3)
$$a_{ij}=0\text{ iff }a_{ji}=0.$$
Given $`sI_1`$ such that $`a_{ss}=0`$ (i.e., $`\alpha _s`$ is an odd isotropic simple root), define a new data
$$r_s(𝒟)=\{𝔥,I,r_s(I_1),r_s(\mathrm{\Pi }^{}),r_s(\mathrm{\Pi })\}$$
and new Chevalley generators $`r_s(e_i),r_s(f_i)`$ of $`𝔤(𝒟)`$ as follows (cf. \[S\], \[PS\], \[KW\]):
$`ir_s(I_1)\text{ iff }i`$ $``$ $`I_1\text{ in case }a_{si}0,ir_s(I_1)\text{ iff }iI_1\text{ otherwise;}`$
$`r_s(\alpha _s^{})`$ $`=`$ $`\alpha _s^{},r_s(\alpha _s)=\alpha _s,`$
$`r_s(\alpha _i^{})`$ $`=`$ $`\alpha _i^{}+{\displaystyle \frac{a_{is}}{a_{si}}}\alpha _s^{}\text{ and }r_s(\alpha _i)=\alpha _i+\alpha _s\text{ if }a_{si}0,`$
$`r_s(\alpha _i^{})`$ $`=`$ $`\alpha _i^{}\text{ and }r_s(\alpha _i)=\alpha _i\text{ in all other cases;}`$
$`r_s(e_s)`$ $`=`$ $`f_s,r_s(f_s)=e_s,`$
$`r_s(e_i)`$ $`=`$ $`[e_s,e_i]\text{ and }r_s(f_i)={\displaystyle \frac{1}{p_{si}a_{si}}}[f_s,f_i]\text{ if }a_{si}0,`$
$`r_s(e_i)`$ $`=`$ $`e_i,r_s(f_i)=f_i\text{ in all other cases.}`$
Denote by $`r_s(𝔫_+)`$ (resp. $`r_s(𝔫_{})`$) the subalgebra of $`𝔤(𝒟)`$ generated by the $`r_s(e_i)`$’s (resp. $`r_s(f_i)`$’s). The transformation $`r_s`$ is called an *odd reflection* (with respect to $`\alpha _s`$).
###### Lemma 1.2.
1. The data $`r_s(𝒟)`$ satisfy (1.3).
2. The new Chevalley generators satisfy the standard relations and together with $`𝔥`$ generate $`𝔤(𝒟)`$, so that $`𝔤(r_s(𝒟))𝔤(𝒟)`$.
3. One has the new triangular decomposition:
$$𝔤(𝒟)=r_s(𝔫_{})+𝔥+r_s(𝔫_+).$$
4. The data $`r_s(r_s(𝒟))`$ coincide with $`𝒟`$, and the Chevalley generators $`r_s(r_s(e_i))`$ (resp. $`r_s(r_s(f_i))`$) coincide, up to a non-zero factor, with $`e_i`$ (resp. $`f_i`$).
###### Proof.
is straightforward using (1.2) and the relation
$$[a,[a,b]]=\frac{1}{2}[[a,a],b]\text{ if }a\text{ ia an odd element.}$$
An element $`\rho 𝔥^{}`$ such that
$$\rho ,\alpha _i^{}=\frac{1}{2}a_{ii}\text{ for all }iI$$
is called a *Weyl vector* for $`\mathrm{\Pi }^{}`$.
###### Lemma 1.3.
If $`\rho `$ is a Weyl vector for $`\mathrm{\Pi }^{}`$, then $`\rho +\alpha _s`$ is a Weyl vector for $`r_s(\mathrm{\Pi }^{})`$.
###### Proof.
It suffices to check that in the case $`a_{is}0`$ one has: $`\rho +\alpha _s,a_{si}\alpha _i^{}+a_{is}\alpha _s^{}=`$$`\frac{1}{2}\alpha _i+\alpha _s,a_{si}\alpha _i^{}+a_{is}\alpha _s^{}`$, which is immediate.
Recall that for each $`\mathrm{\Lambda }𝔥^{}`$ one defines an irreducible highest weight module $`L(\mathrm{\Lambda })`$ over $`𝔤(𝒟)`$ as the (unique) irreducible $`𝔤(𝒟)`$-module for which there exists a non-zero vector $`v_\mathrm{\Lambda }`$ such that
(1.4)
$$hv_\mathrm{\Lambda }=\mathrm{\Lambda }(h)v_\mathrm{\Lambda }\text{ for }h𝔥,𝔫_+v_\mathrm{\Lambda }=0.$$
The vector $`v_\mathrm{\Lambda }`$, called a highest weight vector (with respect to $`𝔫_+`$) is determined uniquely up to a (non-zero) constant factor by the condition $`𝔫_+v_\mathrm{\Lambda }=0`$ (cf. \[K3\]). The linear function $`\mathrm{\Lambda }`$ is called the *highest weight* (with respect to $`𝔫_+`$) of $`L(\mathrm{\Lambda })`$.
###### Lemma 1.4.
Let $`\alpha _s`$ be an odd isotropic simple root and let $`𝔫_+^{}=r_s(𝔫_+)`$.
1. If $`\mathrm{\Lambda },\alpha _s^{}=0`$, then $`v_\mathrm{\Lambda }^{}=v_\mathrm{\Lambda }`$ is a highest weight vector with respect to $`𝔫_+^{}`$, so that the highest weight remains the same: $`\mathrm{\Lambda }^{}=\mathrm{\Lambda }`$.
2. If $`\mathrm{\Lambda },\alpha _s^{}0`$, then $`v_\mathrm{\Lambda }^{}=f_sv_\mathrm{\Lambda }`$ is a highest weight vector with respect to $`𝔫_+^{}`$ so that the highest weight vector becomes $`\mathrm{\Lambda }^{}=\mathrm{\Lambda }\alpha _s`$.
###### Proof.
is straightforward using the facts that $`f_s^2v_\mathrm{\Lambda }=\frac{1}{2}[f_s,f_s]v_\mathrm{\Lambda }=0`$, and $`f_sv_\mathrm{\Lambda }=0`$ iff$`\mathrm{\Lambda },\alpha _s^{}=0`$. ∎
As an immediate corollary of Lemmas 1.3 and 1.4 we obtain the following very useful formulas (cf. \[KW\]):
$`\mathrm{\Lambda }^{}+\rho ^{}`$ $`=`$ $`\mathrm{\Lambda }+\rho \text{ if }\mathrm{\Lambda },\alpha _s^{}(=\mathrm{\Lambda }+\rho ,\alpha _s^{})0,`$
(1.5) $`\mathrm{\Lambda }^{}+\rho ^{}`$ $`=`$ $`\mathrm{\Lambda }+\rho +\alpha _s\text{ if }\mathrm{\Lambda },\alpha _s^{}(=\mathrm{\Lambda }+\rho ,\alpha _s^{})=0.`$
Let $`\alpha 𝔥^{}`$ be a positive even root of $`𝔤(𝒟)`$ such that there exist root vectors $`e`$ attached to $`\alpha `$ and $`f`$ attached to $`\alpha `$ satisfying the following conditions:
1. $`adf`$ is locally nilpotent on $`𝔤(𝒟)`$,
2. $`[e,f]=\alpha ^{}𝔥,[\alpha ^{},e]=2e,[\alpha ^{},f]=2f`$.
Then we call $`f`$ an *integrable element* of $`𝔤(𝒟)`$. The following lemma is well-known (cf. \[K3\]).
###### Lemma 1.5.
Let $`f`$ be an integrable element attached to a negative root $`\alpha `$.
1. If $`f`$ is locally nilpotent on $`L(\mathrm{\Lambda })`$ then $`\mathrm{\Lambda },\alpha ^{}_+`$.
2. Provided that $`\alpha `$ is a simple root, $`f`$ is locally nilpotent on $`L(\mathrm{\Lambda })`$ iff $`\mathrm{\Lambda },\alpha ^{}_+`$.
Let $`\beta =\alpha _s`$ be an odd isotropic simple root. It will be convenient to use notation $`r_\beta `$ in place of $`r_s`$. Consider a sequence of roots $`\beta _0,\beta _1,\mathrm{},\beta _k`$ such that $`\beta _0`$ is an odd isotropic simple root from $`\mathrm{\Pi }^{(0)}:=\mathrm{\Pi }`$, $`\beta _1`$ is an odd isotropic simple root from $`\mathrm{\Pi }^{(1)}=r_{\beta _0}(\mathrm{\Pi }^{(0)}),\mathrm{}`$ , $`\beta _k`$ is an odd isotropic simple root from $`\mathrm{\Pi }^{(k)}=r_{\beta _{k1}}`$ $`(\mathrm{\Pi }^{(k1)})`$. Given $`\mathrm{\Lambda }𝔥^{}`$, let $`\mathrm{\Lambda }^{(0)}=\mathrm{\Lambda }`$ be the highest weight of $`L(\mathrm{\Lambda })`$ with respect to $`𝔫_+^{(0)}:=𝔫_+`$, $`\mathrm{\Lambda }^{(1)}`$ be the highest weight of $`L(\mathrm{\Lambda })`$ with respect to $`𝔫_+^{(1)}:=r_{\beta _0}(𝔫_+),\mathrm{},\mathrm{\Lambda }^{(k)}`$ be the highest weight of $`L(\mathrm{\Lambda })`$ with respect to $`𝔫_+^{(k)}=r_{\beta _{k1}}(𝔫_+^{(k1)})`$. Let $`\rho ^{(k)}`$ be a Weyl vector for $`\mathrm{\Pi }^{(k)}`$.
###### Proposition 1.1.
Let $`\alpha `$ be a positive root of $`𝔤(𝒟)`$ and let $`f`$ be an integrable root element attached to $`\alpha `$. Given $`\mathrm{\Lambda }𝔥^{}`$, let
$$S=\{i[0,1,\mathrm{},k1]|\mathrm{\Lambda }^{(i)},\beta _i^{}=0\}.$$
Suppose that $`\alpha \mathrm{\Pi }^{(k)}`$. Then the element $`f`$ is locally nilpotent on $`L(\mathrm{\Lambda })`$ if and only if
$$\mathrm{\Lambda }+\rho +\underset{iS}{}\beta _i,\alpha ^{}=\{1,2,\mathrm{}\}.$$
###### Proof.
It follows from (1) that
$$\mathrm{\Lambda }^{(k)}+\rho ^{(k)}=\mathrm{\Lambda }+\rho +\underset{iS}{}\beta _i.$$
Since $`\mathrm{\Lambda }^{(k)}+\rho ^{(k)},\alpha ^{}=\mathrm{\Lambda }^{(k)},\alpha ^{}+1`$, the proposition follows from Lemma 1.5b.
###### Proposition 1.2.
If, under the assumptions of Proposition 1.1, one has:
$$\mathrm{\Lambda }+\rho ,\alpha ^{},$$
then $`f`$ is integrable on $`L(\mathrm{\Lambda })`$.
###### Proof.
Due to Proposition 1.1, Proposition 1.2 holds if $`S=\mathrm{}`$. Let $`N=\mathrm{\Lambda },\alpha ^{}`$. It is well-known (cf. \[K3\]) that $`f`$ is integrable on $`L(\mathrm{\Lambda })`$ iff
(1.6)
$$f^{N+1}v_\mathrm{\Lambda }\text{ lies in a maximal submodule of the Verma module }M(\mathrm{\Lambda }).$$
But we have just shown that (1.6) holds for a Zariski open set of $`\lambda `$ on the hyperplane $`\lambda ,\alpha ^{}=N`$. Since (1.6) is a polynomial condition, we conclude that it holds for all $`\lambda `$ on this hyperplane. ∎
###### Proposition 1.3.
If, under the assumptions of Proposition 1.1, $`f`$ is integrable on $`L(\mathrm{\Lambda })`$ and
$$\mathrm{\Lambda }+\rho ,\beta _i^{}0\text{ for }i=0,1,\mathrm{},s(k),$$
then $`\mathrm{\Lambda }_{i=0}^s\beta _i,\alpha ^{}_+`$.
###### Proof.
We have: $`\mathrm{\Lambda },\beta _0^{}=\mathrm{\Lambda }+\rho ,\beta _0^{}0`$, hence, by (1) we have: $`\mathrm{\Lambda }+\rho =\mathrm{\Lambda }^{(1)}+\rho ^{(1)}`$, etc. Thus, $`\mathrm{\Lambda }^{(i)}+\rho ^{(i)}=\mathrm{\Lambda }+\rho `$ for $`i=1,\mathrm{},s`$. Therefore, by Lemma 1.4b , we have:
$$\mathrm{\Lambda }^{(s)}=\mathrm{\Lambda }\underset{i=0}{\overset{s}{}}\beta _i.$$
Now the proposition follows from Lemma 1.5a.
The calculation of coroots is facilitated by the following simple fact.
###### Proposition 1.4.
1. There exists a non-degenerate symmetric bilinear form $`(.|.)`$ on $`𝔥`$ such that, identifying $`𝔥`$ and $`𝔥^{}`$ via this form, we have:
(1.7)
$$\alpha _i^{}=\nu _i\alpha _i,\text{ where }\nu _i^\times ,$$
if and only if
(1.8)
$$A=diag(\nu _i)_{iI}B,\text{ where }B=(b_{ij})\text{ is a symmetric matrix.}$$
One then has: $`(\alpha _i|\alpha _j)=b_{ij}`$.
2. Let $`\mathrm{\Pi }^{}=\{\alpha _i^{}\}=r_s(\mathrm{\Pi }^{})`$ and $`\mathrm{\Pi }=\{\alpha _i^{}\}=r_s(\mathrm{\Pi })`$ where $`r_s`$ is an odd reflection, and suppose that (1.7) holds. Then
$$\alpha _i^{}=\nu _i\alpha _i^{}.$$
3. Provided that (1.8) holds and $`a_{ii}=2`$ or $`0`$ for all $`iI`$, one has for any non-isotropic root $`\alpha `$ which is obtained from a simple root by a sequence of odd reflections: $`\alpha ^{}=2\alpha /(\alpha |\alpha )`$.
###### Proof.
(a) is proved in \[K3\], (b) and (c) are easily checked. ∎
###### Remark 1.1.
A natural question is which of the Lie superalgebras $`𝔤(𝒟)`$ are of “Kac-Moody” type? The most natural answer, in our opinion, is that they should satisfy the following conditions:
1. $`𝔤(𝒟)_{\overline{0}}`$ is a (generalized) Kac-Moody algebra,
2. the $`𝔤(𝒟)_{\overline{0}}`$-module $`𝔤(𝒟)_{\overline{1}}`$ is integrable.
This definition covers the basic classical finite-dimensional Lie superalgebras and the associated affine superalgebras (including the twisted ones). Unfortunately, a well developed theory of generalized Kac-Moody superalgebras (see \[B\], \[R\] and references there) does not cover most of the latter superalgebras (because of the crucial assumption on the Cartan matrix that its off diagonal entries are non-positive).
## 2 Classification of integrable irreducible highest weight modules over $`g\mathrm{}(m|n)^\widehat{}`$
Consider the Lie superalgebra $`g\mathrm{}(m|n)`$, where $`m,n1`$ (see \[K1\]). Let $`e_{ij}`$ $`(1i,jm+n)`$ denote its standard basis. Denote by $`𝔥`$ the Cartan subalgebra of $`g\mathrm{}(m|n)`$ consisting of all diagonal matrices. Let $`ϵ_i`$ $`(1im+n)`$ be the basis of $`𝔥^{}`$ dual to the basis $`u_i:=e_{ii}`$ of $`𝔥`$. Then $`g\mathrm{}(m|n)=𝔤(𝒟)`$ for the following data $`𝒟=\{𝔥,I,I_1,\mathrm{\Pi }^{},\mathrm{\Pi }\}`$ (cf. \[K1\]). We let $`I=\{1,2,\mathrm{},m+n1\}`$, $`I_1=\{m\}`$; $`\alpha _i^{}=u_iu_{i+1}`$ for $`iI\backslash I_1`$, $`\alpha _m^{}=u_m+u_{m+1}`$, $`\alpha _i=ϵ_iϵ_{i+1}`$ for all $`iI`$. Its Cartan matrix is the following $`(m+n1)\times (m+n1)`$ matrix:
$`A=\left(\begin{array}{ccccccccc}\hfill 2& \hfill 1& \hfill 0& & & & & & \\ \hfill 1& \hfill 2& \hfill 1& & & & & & \\ & & \hfill \mathrm{}& & & & & & \\ & & \hfill 1& \hfill 2& \hfill 1& & & & \\ & & & \hfill 1& \hfill 0& \hfill 1& & & \\ & & & & \hfill 1& \hfill 2& \hfill 1& & \\ & & & & & \hfill \mathrm{}& & & \\ & & & & & & \hfill 1& \hfill 2& \end{array}\right)\text{m}\text{-th row}\text{.}`$
The Chevalley generators are as follows:
$$e_i=e_{i,i+1},f_i=e_{i+1,i}(i=1,\mathrm{},m+n1).$$
Note that $`\alpha _m`$ is the only odd simple root, and it is isotropic.
Consider the supertrace form on $`g\mathrm{}(m|n)`$:
$$(a|b)=strab.$$
This is a non-degenerate invariant supersymmetric bilinear form on $`g\mathrm{}(m|n)`$ whose restriction to $`𝔥`$ is non-degenerate and symmetric. Identifying $`𝔥`$ and $`𝔥^{}`$ via this bilinear form, we have:
$$ϵ_i=u_i\text{ for }i=1,\mathrm{},m;ϵ_i=u_i\text{ for }i=m+1,\mathrm{},m+n.$$
Hence we have:
(2.2)
$$\alpha _i^{}=\alpha _i\text{ for }i=1,\mathrm{},m,\alpha _i^{}=\alpha _i\text{ for }i=m+1,\mathrm{},m+n1,$$
and we may use Proposition 1.4. In particular,
$$((\alpha _i|\alpha _i))_{i,jI}=diag(\underset{m}{\underset{}{1,\mathrm{},1}},1,\mathrm{},1)A.$$
Likewise, the affine superalgebra $`g\mathrm{}(m|n)^\widehat{}`$ is isomorphic to $`𝔤(\widehat{𝒟})`$, where the data $`\widehat{𝒟}=\{\widehat{𝔥},\widehat{I},\widehat{I}_1,\widehat{\mathrm{\Pi }}^{},\widehat{\mathrm{\Pi }}\}`$ is an extension of the data $`𝒟`$ for $`g\mathrm{}(m|n)`$ defined as follows (cf. \[K3\]). The space $`\widehat{𝔥}`$ is defined by (0.4), $`\widehat{I}=I\{0\}`$, $`\widehat{I}_1=\{m,0\}`$, $`\widehat{\mathrm{\Pi }}^{}=\mathrm{\Pi }^{}\{\alpha _0^{}\}`$, $`\widehat{\mathrm{\Pi }}=\mathrm{\Pi }\{\alpha _0\}`$. Here the $`\alpha _i`$ for $`i\mathrm{\Pi }`$ are extended from $`𝔥`$ to $`\widehat{𝔥}`$ by letting $`\alpha _i(K)=\alpha _i(d)=0`$, $`\alpha _0=\delta \theta ,\alpha _0^{}=K\theta ^{}`$, where: $`\delta |_{𝔥+K}=0,\delta ,d=1,\theta =ϵ_1ϵ_{m+n}`$ is the highest root of $`g\mathrm{}(m|n),\theta ^{}=u_1+u_{m+n}`$. We extend the bilinear form $`(.|.)`$ from $`g\mathrm{}(m|n)`$ to $`g\mathrm{}(m|n)^\widehat{}`$ by (0.3). Identifying $`\widehat{𝔥}`$ with $`\widehat{𝔥}^{}`$ via this symmetric bilinear form, we get:
(2.3)
$$K=\delta ,\theta =\theta ^{},\alpha _0=\alpha _0^{}.$$
We have the following expression of $`\delta =K`$ in terms of simple roots and coroots:
(2.4)
$$\delta =K=\underset{i=0}{\overset{m+n1}{}}\alpha _i=\underset{i=0}{\overset{m}{}}\alpha _i^{}\underset{j=m+1}{\overset{m+n1}{}}\alpha _j^{}.$$
The Cartan matrix for $`\widehat{𝒟}`$ is
$`\widehat{A}=\left(\begin{array}{ccccccccccccc}\hfill 0& \hfill 1& \hfill 0& \mathrm{}& \hfill 0& \hfill 1& & & & & & & \\ \hfill 1& & & & & & & & & & & & \\ \hfill 0& & & & & & & & & & & & \\ \hfill \mathrm{}& & & A& & & & & & & & & \\ \hfill 0& & & & & & & & & & & & \\ \hfill 1& & & & & & & & & & & & \end{array}\right).`$
As above, we have:
$$((\alpha _i|\alpha _j))_{i,j\widehat{I}}=diag(\underset{m+1}{\underset{}{1,\mathrm{},1}},1,\mathrm{},1)\widehat{A}.$$
The even part of $`g\mathrm{}(m|n)`$ is $`g\mathrm{}(m)g\mathrm{}(n)`$, hence the even part of $`g\mathrm{}(m|n)^\widehat{}`$ is the sum $`g\mathrm{}(m)^\widehat{}+g\mathrm{}(n)^\widehat{}`$ with a common central element $`K`$ and a common scaling element $`d`$. Note that the restriction of the supertrace form to $`g\mathrm{}(m)`$ (resp. $`g\mathrm{}(n)`$) is the normalized (resp. negative of the normalized) invariant form, i.e., $`(\alpha |\alpha )=2`$ (resp. $`(\alpha |\alpha )=2`$) for any root $`\alpha `$.
The set of simple roots for $`g\mathrm{}(m)^\widehat{}`$ (resp. $`g\mathrm{}(n)^\widehat{}`$) is empty if $`m=1`$ (resp. $`n=1`$), and for $`m2`$ (resp. $`n2`$) it is as follows:
$$\widehat{\mathrm{\Pi }}^{}=\{\alpha _0^{}=\delta \theta ^{},\alpha _1,\mathrm{},\alpha _{m1}\}$$
(resp. $`\widehat{\mathrm{\Pi }}^{\prime \prime }=\{\alpha _0^{\prime \prime }=\delta \theta ^{\prime \prime },\alpha _{m+1},\mathrm{},\alpha _{m+n1}`$), where $`\theta ^{}={\displaystyle \underset{i=1}{\overset{m1}{}}}\alpha _i,\theta ^{\prime \prime }={\displaystyle \underset{i=m+1}{\overset{m+n1}{}}}\alpha _i`$. Assuming that $`m2`$, we have: $`(\theta ^{}|\theta ^{})=2`$, hence $`\theta ^{}=\theta ^{{}_{}{}^{}}`$, and we have:
(2.6)
$$\alpha _0^{}=\alpha _0^{{}_{}{}^{}}=\alpha _0+\underset{i=m}{\overset{m+n1}{}}\alpha _i=\alpha _0^{}+\alpha _m^{}\underset{i=m+n}{\overset{m+n1}{}}\alpha _i^{}.$$
A $`g\mathrm{}(m|n)^\widehat{}`$-module $`L(\mathrm{\Lambda })`$ is called *integrable*, if its restriction to $`g\mathrm{}(m)^\widehat{}`$ is integrable and its restriction to $`g\mathrm{}(m|n)`$ is locally finite. In this section we shall classify all such modules.
As usual, define fundamental weights $`\mathrm{\Lambda }_i\widehat{𝔥}(i=0,1,\mathrm{},m+n1)`$ by
$$\mathrm{\Lambda }_i,\alpha _j^{}=\delta _{ij},j=0,\mathrm{},m+n1,\mathrm{\Lambda }_i,d=0,$$
and *labels* of a weight $`\mathrm{\Lambda }`$ by:
$$k_i=\mathrm{\Lambda },\alpha _i^{}.$$
The following necessary conditions of integrability of $`L(\mathrm{\Lambda })`$ follow from Lemma 1.5a:
(2.7) $`k_i_+\text{ for }i=1,\mathrm{},m1,m+1,\mathrm{},m+n1,`$
(2.8) $`k^{}:=k_0+k_m{\displaystyle \underset{i=m+1}{\overset{m+n1}{}}}k_i_+.`$
We assume in (2.8) that $`m2`$ and use (2.6).
We call $`k^{}`$ the *partial level* of $`\mathrm{\Lambda }`$ since, using (2.4), we see that the *level* $`k:=\mathrm{\Lambda },K`$ is given by
(2.9)
$$k=\underset{i=1}{\overset{m1}{}}k_i+k^{}.$$
Hence, provided that $`m2`$, the level of an integrable $`g\mathrm{}(m|n)^\widehat{}`$-module is a non-negative integer.
###### Lemma 2.1.
Assume that $`m2`$. Then conditions (2.7) and (2.8) along with the condition
(2.10)
$$k^{}n$$
are sufficient for integrability of the $`g\mathrm{}(m|n)^\widehat{}`$-module $`L(\mathrm{\Lambda })`$.
###### Proof.
The lemma follows from Lemma 1.5 applied to the simple roots $`\alpha _i,i=1,\mathrm{},m1`$, and Proposition 1.2 applied to $`\alpha ^{}=\alpha _0^{{}_{}{}^{}}`$, since, due to (2.6) we have:
(2.11)
$$\rho ,\alpha _0^{{}_{}{}^{}}=n+1.$$
###### Lemma 2.2.
Let $`L(\mathrm{\Lambda })`$ be an integrable $`g\mathrm{}(m|n)^\widehat{}`$-module such that $`k^{}<n`$, and let $`m2`$. Then the following complementary condition holds:
1. there exist $`r,s_+`$ such that
1. $`k^{}=r+s`$,
2. $`k_0k_{m+n1}k_{m+n2}\mathrm{}k_{m+nr}r=0`$,
3. $`k_mk_{m+1}k_{m+2}\mathrm{}k_{m+s}s=0`$.
###### Proof.
Consider the following two sequences of roots of $`g\mathrm{}(m|n)^\widehat{}`$:
$`\beta _0`$ $`=`$ $`\alpha _0,\beta _1=\alpha _0+\alpha _{m+n1},\beta _2=\alpha _0+\alpha _{m+n1}+\alpha _{m+n2},\mathrm{},\beta _{n1}=\alpha _0+\alpha _{m+n1}+\mathrm{}+\alpha _{m+1};`$
$`\beta _0^{}`$ $`=`$ $`\alpha _m,\beta _1^{}=\alpha _m+\alpha _{m+1},\beta _2^{}=\alpha _m+\alpha _{m+1}+\alpha _{m+2},\mathrm{},\beta _{n1}^{}=\alpha _m+\mathrm{}+\alpha _{m+n1}.`$
It is clear by Proposition 1.4 that $`\beta _i^{}=\beta _i`$ and $`\beta ^{{}_{}{}^{}}=\beta _i^{}`$. Note that $`\mathrm{\Lambda }+\rho ,\beta _r^{}`$ (resp. $`\mathrm{\Lambda }+\rho ,\beta _s^{{}_{}{}^{}}`$) is equal to the left-hand side of (ii) (resp. (iii)). Note that
(2.12)
$$\beta _i,\alpha _0^{{}_{}{}^{}}=1=\beta _i^{},\alpha _0^{{}_{}{}^{}},i=0,\mathrm{},n1.$$
If $`\mathrm{\Lambda }+\rho ,\beta _i^{}0`$ for all $`i`$, using (2.8) and (2.12) we would conclude, by Proposition 1.3, that $`k^{}n0`$, in contradiction with the assumption of the lemma. Hence (ii) holds for some non-negative integer $`r`$ $`(<n)`$. Similarly, (iii) holds for some non-negative integer $`s`$ $`(<n)`$. Similarly, applying Proposition 1.3 to the union of sequences $`\beta _i`$ and $`\beta _i^{}`$ , we conclude that
(2.13)
$$r+sk^{}.$$
Hence, adding up (ii) and (iii) we get
(2.14)
$$k^{}+\underset{i=m+s+1}{\overset{m+nr1}{}}k_i=r+s.$$
Now (i) follows from (2.7), (2.13) and (2.14).
###### Remark 2.1.
Condition (\*) on $`\mathrm{\Lambda }`$ is equivalent to the following condition: there exists a non-negative integer $`sk^{}`$ $`(n1)`$ such that:
$$k_m=k_{m+1}+\mathrm{}+k_{m+s}+s\text{ and }k_{m+s+1}=\mathrm{}=k_{m+s+nk^{}1}=0.$$
This condition implies that $`\mathrm{\Lambda }`$ lies in a union of $`k^{}+1`$ hyperplanes of dimension $`k^{}+m1`$. Equivalently, there exists a non-negative integer $`rk^{}`$ $`(n1)`$ such that
$`k_0=k_{m+n1}+k_{m+n2}+\mathrm{}+k_{m+nr}+r\text{ and }k_{m+nr1}=\mathrm{}=k_{m+k^{}r+1}=0.`$
###### Theorem 2.1.
1. A $`g\mathrm{}(1|n)^\widehat{}`$-module $`L(\mathrm{\Lambda })`$ is integrable iff $`k_2,\mathrm{},k_n_+`$.
2. Provided that $`m2`$, a $`g\mathrm{}(m|n)^\widehat{}`$-module $`L(\mathrm{\Lambda })`$ is integrable iff conditions (2.7), (2.8) hold and, in the case $`k^{}<n`$, the complementary condition (\*) holds.
###### Proof.
In the case $`m=1`$, the only condition of integrability is local finiteness of $`g\mathrm{}(1,n)`$ on $`L(\mathrm{\Lambda })`$ which is equivalent to $`k_2,\mathrm{},k_n_+`$ due to Lemma 1.5b. It follows from Lemma 2.2 that in the case $`m2`$, the conditions listed by Theorem 2.1b are necessary. In view of Lemma 1.5b, it remains to show that these conditions are sufficient for local nilpotency of $`e_{\alpha _0^{}}`$. Due to Lemma 2.1, we may assume that
(2.15)
$$k^{}n1.$$
Consider the sequence of odd roots $`\beta _0,\mathrm{},\beta _{n1}`$ introduced in the proof of Lemma 2.2 and let $`\mathrm{\Pi }^{(0)}=\mathrm{\Pi }`$, $`\mathrm{\Pi }^{(1)}=r_{\beta _0}(\mathrm{\Pi }^{(0)}),\mathrm{},\mathrm{\Pi }^{(n)}=r_{\beta _{n1}}(\mathrm{\Pi }^{(n1)})`$, and notice that
$$\alpha _0^{}\mathrm{\Pi }^{(n)}.$$
Let $`\mathrm{\Lambda }^{(n)}`$ be the highest weight vector of $`L(\mathrm{\Lambda })`$ with respect to $`𝔫_+^{(n)}=r_{\beta _{n1}}\mathrm{}r_{\beta _0}(𝔫)`$. Due to Lemma 1.5b, it remains to show that conditions listed by Theorem 2.1b imply that
(2.16)
$$\mathrm{\Lambda }^{(n)},\alpha _0^{{}_{}{}^{}}_+.$$
Recall that by (1) we have:
(2.17)
$$\mathrm{\Lambda }^{(n)}+\rho ^{(n)}=\mathrm{\Lambda }+\rho +\underset{iS}{}\beta _i,$$
where $`S=\{i[0,\mathrm{},n1]|\mathrm{\Lambda }^{(i)},\beta _i^{}=0\}`$. Let $`t_i=\mathrm{\Lambda }^{(i)},\beta _i^{}`$ for short. Then condition (\*) gives for some $`r_+`$, $`r<n`$, that $`t_r=0`$. In view of Remark 2.1, we have:
(2.18)
$$t_r=t_{r+1}=\mathrm{}=t_{ns1}=0.$$
Hence, due to (2.17), (2.11), (2.12) and (2.18) we get:
$$\mathrm{\Lambda }^{(n)}+\rho ^{(n)},\alpha _0^{{}_{}{}^{}}=k^{}+(1n)+|S|k^{}+(1n)+(nsr)1,$$
proving (2.16), since $`\rho ^{(n)},\alpha _0^{{}_{}{}^{}}=1`$. ∎
###### Remark 2.2.
It follows from Theorem 2.1 that when $`m2`$, the only integrable $`g\mathrm{}(m|n)^\widehat{}`$-modules $`L(\mathrm{\Lambda })`$ of level $`k=0`$ are those for which all labels are $`0`$, in which case $`dimL(\mathrm{\Lambda })=1`$.
###### Remark 2.3.
If $`m2`$ and $`n2`$, then the only $`L(\mathrm{\Lambda })`$ which are integrable with respect to the whole even subalgebra are $`1`$-dimensional. (It is because the $`g\mathrm{}(m)^\widehat{}`$-integrability implies $`k0`$ and $`g\mathrm{}(n)^\widehat{}`$-integrability implies $`k0`$.)
###### Remark 2.4.
Define $`ϵ\widehat{𝔥}^{}`$ by letting $`ϵ|_𝔥=`$ supertrace, $`ϵ(K)=ϵ(d)=0`$. It follows from Theorem 2.1 that when $`m2`$, the complete list of highest weights of integrable $`g\mathrm{}(m|n)^\widehat{}`$-modules of level $`1`$, up to adding an arbitrary linear combination of $`ϵ`$ and $`\delta `$, is as follows:
$`\mathrm{\Lambda }_s(1sm1),(a+1)\mathrm{\Lambda }_m+a\mathrm{\Lambda }_{m+1}(a_+)(a+1)\mathrm{\Lambda }_0+a\mathrm{\Lambda }_{m+n1}(a_+).`$
###### Remark 2.5.
Consider the sequence of the sets of simple roots $`\mathrm{\Pi }^{(0)}=\mathrm{\Pi },\mathrm{},\mathrm{\Pi }^{(n)}=`$$`\{\alpha _0^{},\mathrm{},\alpha _{m+n1}^{}\}`$, introduced in the proof of Theorem 2.1. One has:
$`\alpha _0^{}`$ $`=`$ $`\alpha _0+\alpha _1,\alpha _1^{}=\alpha _2,\mathrm{},\alpha _{m2}^{}=\alpha _{m1},\alpha _{m1}^{}=\alpha _m+\alpha _{m+1}+\mathrm{}+\alpha _{m+n1}+\alpha _0,`$
$`\alpha _m^{}`$ $`=`$ $`(\alpha _{m+n}+\mathrm{}+\alpha _{m+n1}+\alpha _0),\alpha _j^{}=\alpha _j\text{ for }m+1jm+n1.`$
Let $`\mathrm{\Lambda }_j^{}`$ be the fundamental weights with respect to $`\mathrm{\Pi }^{(n)}`$. Given a weight $`\mathrm{\Lambda }`$, denote by $`\mathrm{\Lambda }^{(n)}`$ the highest weight of $`L(\mathrm{\Lambda })`$ with respect to $`\mathrm{\Pi }^{(n)}`$ (or rather $`𝔫_+^{(n)}`$). Using Lemma 1.4, it is easy to see that the weights listed in Remark 2.4 get changed under the map $`\mathrm{\Lambda }\mathrm{\Lambda }^{(n)}`$ as follows:
$`\mathrm{\Lambda }_j^{(n)}`$ $`=`$ $`\mathrm{\Lambda }_j(1jm),((a+1)\mathrm{\Lambda }_0+a\mathrm{\Lambda }_{m+n1})^{(n)}=(a+1)\mathrm{\Lambda }_0+a\mathrm{\Lambda }_{m+n1}\alpha _0,`$
$`((a+1)\mathrm{\Lambda }_m+a\mathrm{\Lambda }_{m+1})^{(n)}=(a+1)\mathrm{\Lambda }_m+a\mathrm{\Lambda }_{m+1}+\alpha _m^{}(a>0).`$
In terms of the fundamental weights $`\mathrm{\Lambda }_j^{}`$ the map $`\mathrm{\Lambda }\mathrm{\Lambda }^{(n)}`$ looks as follows:
$`\mathrm{\Lambda }_j\mathrm{\Lambda }_{j1}^{}(1jm),`$
$`(a+1)\mathrm{\Lambda }_0+\mathrm{\Lambda }_{m+n1}(a+2)\mathrm{\Lambda }_0^{}+(a+1)\mathrm{\Lambda }_{m+n1}^{}(a_+),`$
$`(a+1)\mathrm{\Lambda }_m+a\mathrm{\Lambda }_{m+1}a\mathrm{\Lambda }_m^{}+(a1)\mathrm{\Lambda }_{m+1}^{}(a).`$
It follows that all weights of level 1 listed by Remark 2.4 are conjugate to each other by odd reflections.
## 3 Free field realization of level $`1`$ integrable modules over $`g\mathrm{}(m|n)^\widehat{}`$.
Fix non-negative integers $`m`$ and $`n`$ such that $`m+n1`$ and denote by $`F`$ the vertex algebra generated by $`m`$ pairs of odd fields $`\psi ^i(z),\psi ^i(z)`$, $`(i=1,\mathrm{},m)`$ and $`n`$ pairs of even fields $`\phi ^j(z),\phi ^j(z)`$ $`(j=1,\mathrm{},n)`$, all pairwise local, subject to the following operator product expansions (as usual, we list only the non-trivial OPE):
$`\begin{array}{cc}\psi ^i(z)\psi ^j(w)\frac{\delta _{ij}}{zw},\hfill & \psi ^i(z)\psi ^j(w)\frac{\delta _{ij}}{zw},\hfill \\ \phi ^i(z)\phi ^j(w)\frac{\delta _{ij}}{zw},\hfill & \phi ^i(z)\phi ^j(w)\frac{\delta _{ij}}{zw}.\hfill \end{array}`$
This is called a free fermionic vertex algebra in the book \[K4\] to which we refer for foundations of the vertex algebra theory.
This vertex algebra has a family of Virasoro fields \[K4\], from which it is convenient to choose the following one:
$`L(z){\displaystyle \underset{k}{}}L_kz^{k2}`$ $`=`$ $`\frac{1}{2}{\displaystyle \underset{i=1}{\overset{m}{}}}(:\psi ^i(z)\psi ^i(z):+:\psi ^i(z)\psi ^i(z):)`$
$`+`$ $`\frac{1}{2}{\displaystyle \underset{j=1}{\overset{n}{}}}(:\phi ^j(z)\phi ^j(z)::\phi ^j(z)\phi ^j(z):).`$
With respect to $`L(z)`$ the fields $`\psi ^i(z)`$, $`\psi ^i(z)`$, $`\phi ^j(z)`$ and $`\phi ^j(z)`$ are primary of conformal weight $`1/2`$. We therefore write all these fields in the form $`x^i(z)=_{k\frac{1}{2}+}x_k^{(i)}z^{k1/2}`$ where $`x=\psi `$ , $`\psi ^{}`$, $`\phi `$ or $`\phi ^{}`$, and we have the following conditions on the vacuum $`|0`$:
$$\psi _k^{(i)}|0=0,\psi _k^{(i)}|0=0,\phi _k^{(i)}|0=0,\phi _k^{(i)}|0=0\text{ for }k>0.$$
The operator $`L_0`$ is called the *energy operator* or (Hamiltonian) and its eigenvalues are called the *energies* of the corresponding eigenvectors. The energy can be calculated from the following relations:
(3.3)
$$energy|0=0,energy(\psi _k^{(i)},\psi _k^{(i)},\phi _k^{(j)},\phi _k^{(j)})=k.$$
The second relation means that $`\psi _k^{(i)}`$, etc., changes the energy by $`k`$, i.e.,
$$energy(\psi _k^{(i)}v)=energy(v)k,\text{ etc.}$$
Next, for each pair $`i,j`$ that may occur introduce the following fields of conformal weight $`1`$:
$`\begin{array}{cc}a^{ij+}(z)=:\psi ^i(z)\psi ^j(z):,\hfill & a^{ij}(z)=:\phi ^i(z)\phi ^j(z):\hfill \\ E^{ij+}(z)=:\psi ^i(z)\phi ^j(z):,\hfill & E^{ij}(z)=:\phi ^i(z)\psi ^j(z):.\hfill \end{array}`$
###### Proposition 3.1.
1. Consider the affine superalgebra $`g\mathrm{}(m|n)^\widehat{}`$ and let $`A(z)=`$$`_k(t^kA)z^{k1}`$ for $`Ag\mathrm{}(m|n)`$. Then the linear map $`\sigma `$ given by
$`\begin{array}{cc}e_{ij}(z)a^{ij+}(z),\hfill & e_{i+m,j+m}(z)a^{ij}(z),\hfill \\ e_{i,j+m}(z)E^{ij+}(z),\hfill & e_{i+m,j}(z)E^{ij}(z),K1,dL_0\hfill \end{array}`$
defines a representation of $`g\mathrm{}(m|n)^\widehat{}`$ (of level $`1`$) in the space $`F`$.
2. Consider the standard $`g\mathrm{}(m|n)`$-module $`^{m|n}`$ and its contragredient module $`^{m|n}`$. Consider the corresponding $`[t,t^1]g\mathrm{}(m|n)`$-modules $`[t,t^1]^{m|n}`$ and $`[t,t^1]^{m|n}`$, and let $`v(z)=_k(t^kv)z^{k1}`$ for $`v^{m|n}`$ or $`^{m|n}`$. Then the linear maps $`\nu `$ and $`\nu ^{}`$ given by $`(i=1,\mathrm{},m;j=1,\mathrm{},n)`$:
$`\begin{array}{cc}v_i(z)\psi ^i(z),\hfill & v_{j+m}(z)\phi ^j(z)\text{ and}\hfill \\ v_i^{}(z)\psi ^i(z),\hfill & v_{j+m}^{}(z)\phi ^j(z)\hfill \end{array}`$
are equivariant, i.e., they have the following property:
$`\nu (A(z)v(w))`$ $`=`$ $`[\sigma (A(z)),\nu (v(w))],v^{m|n},`$
$`\nu ^{}(A(z)v^{}(w))`$ $`=`$ $`[\sigma (A(z)),\nu ^{}(v^{}(w))],v^{}^{m|n}.`$
###### Proof.
Both statements follow from the corresponding OPE, which are easily derived from Wick’s formula. Below we give the less trivial OPE needed for the proof of (a).
$`E^{ij+}(z)E^k\mathrm{}(w)`$ $``$ $`{\displaystyle \frac{\delta _{jk}a^{i\mathrm{}+}(w)+\delta _i\mathrm{}a^{kj}(w)}{zw}}+{\displaystyle \frac{\delta _i\mathrm{}\delta _{jk}}{(zw)^2}}`$
$`a^{ij\pm }(z)E^{k\mathrm{}\pm }(w)`$ $``$ $`{\displaystyle \frac{\delta _{jk}E^{i\mathrm{}\pm }(w)}{zw}},`$
$`a^{ij\pm }(z)E^k\mathrm{}(w)`$ $``$ $`{\displaystyle \frac{\delta _\mathrm{}iE^{kj}(w)}{zw}},`$
$`a^{ij\pm }(z)a^{k\mathrm{}\pm }(w)`$ $``$ $`{\displaystyle \frac{\delta _{jk}a^{i\mathrm{}\pm }(w)\delta _\mathrm{}ia^{kj\pm }(w)}{zw}}\pm {\displaystyle \frac{1}{(zw)^2}}`$
$`a^{ij\pm }(z)a^k\mathrm{}(w)`$ $``$ $`0.`$
Introduce the *total charge operator*
$$a_0=\sigma (I),\text{ where }I=\underset{i=1}{\overset{m+n}{}}e_{ii}g\mathrm{}(m|n).$$
Its eigenvalues are called *charges* of the corresponding eigenvectors. It is clear from Proposition 3.1 that the total charge can be calculated from the following relations:
(3.7)
$$charge|0=0,charge(\psi _k^{(i)},\phi _k^{(j)})=1,charge(\psi _k^{(i)},\phi _k^{(j)})=1.$$
Consider the charge decomposition of $`F`$, i.e., its decomposition in eigenspaces of $`a_0`$:
(3.8)
$$F=_sF_s.$$
Since $`a_0`$ commutes with $`\sigma (g\mathrm{}(m|n)^\widehat{})`$, we conclude that (3.8) is a decomposition in a direct sum of $`g\mathrm{}(m|n)^\widehat{}`$-modules.
It is clear that $`L_0`$ commutes with $`a_0`$, hence each $`F_s`$ is $`L_0`$-invariant, and since all eigenvalues of $`L_0`$ in $`F`$ lie in $`\frac{1}{2}_+`$, the same holds for eigenvalues of $`L_0`$ in $`F_s`$, $`s`$. Note also that $`L_0`$ commutes with $`\sigma (g\mathrm{}(m|n))`$. It is because all fields $`a^{ij\pm }(z)`$ and $`E^{ij\pm }(z)`$ have conformal weight $`1`$. It follows that each eigenspace of $`L_0`$ in $`F_s`$ is a $`g\mathrm{}(m|n)`$-module. The following proposition describes the lowest energy subspace $`F_s^{low}`$ and the lowest weight vector $`|s`$ in each $`F_s`$.
###### Proposition 3.2.
1. Let $`s_+`$. Then, as a $`g\mathrm{}(m|n)`$-module, $`F_s^{low}`$ is isomorphic to $`\mathrm{\Lambda }^s^{m|n}`$. Furthermore, any highest weight vector of $`g\mathrm{}(m|n)^\widehat{}`$ in $`F_s`$ lies in $`F_s^{low}`$ and is proportional to the vector
$$|s=\psi _{\frac{1}{2}}^{(1)}\mathrm{}\psi _{\frac{1}{2}}^{(s)}|0\text{ with weight }\mathrm{\Lambda }_0+ϵ_1+\mathrm{}+ϵ_s\frac{s}{2}\delta ,$$
provided that $`sm`$, and to the vector
$$|s=\left(\phi _{\frac{1}{2}}^{(1)}\right)^{sm}|m\text{ with weight }\mathrm{\Lambda }_0+ϵ_1+\mathrm{}+ϵ_m+(sm)ϵ_{m+1}\frac{s}{2}\delta $$
provided that $`sm`$.
2. Let $`s_+`$. Then, as a $`g\mathrm{}(m|n)`$-module, $`F_s^{low}`$ is isomorphic to $`\mathrm{\Lambda }^s(^{m|n})^{}`$. Furthermore, any highest weight vector of $`g\mathrm{}(m|n)^\widehat{}`$ in $`F_s`$ lies in $`F_s^{low}`$ and is proportional to the vector
$$|s=\left(\phi _{\frac{1}{2}}^{(n)}\right)^s|0\text{ with weight }\mathrm{\Lambda }_0+sϵ_{m+n}+\frac{s}{2}\delta .$$
###### Proof.
It is clear that, if $`s`$ (resp. $`s`$) $`_+`$, then $`F_s^{low}`$ consists of homogeneous polynomials of degree $`|s|`$ in anticommuting operators $`\psi _{\frac{1}{2}}^{(i)}`$ (resp. $`\psi _{\frac{1}{2}}^{(i)}`$) and commuting operators $`\phi _{\frac{1}{2}}^{(j)}`$ (resp. $`\phi _{\frac{1}{2}}^{(j)}`$), applied to $`|0`$. This proves (a) (resp. (b)), due to Proposition 3.1.
###### Remark 3.1.
The lowest energy in $`F_s`$ is $`\frac{1}{2}|s|`$ and the spectrum of $`L_0`$ in $`F_s`$ is $`\frac{1}{2}|s|+_+`$.
###### Remark 3.2.
Denote by $`\mathrm{\Lambda }_{(s)}`$ the weight of $`|s`$. When restricted to $`s\mathrm{}(m|n)^\widehat{}`$, $`\mathrm{\Lambda }_{(s)}`$ is given by the following formulas:
$`\begin{array}{cc}\mathrm{\Lambda }_s\frac{s}{2}\delta \hfill & \text{ if }0sm,\hfill \\ (1+sm)\mathrm{\Lambda }_m+(sm)\mathrm{\Lambda }_{m+1}\frac{s}{2}\delta \hfill & \text{ if }sm,\hfill \\ (1s)\mathrm{\Lambda }_0s\mathrm{\Lambda }_{m+n1}+\frac{s}{2}\delta \hfill & \text{ if }s0.\hfill \end{array}`$
We identify here $`\mathrm{\Lambda }_{m+1}`$ with $`\mathrm{\Lambda }_0`$ in the case $`n=1`$.
The following theorem is the central result of this section.
###### Theorem 3.1.
Suppose that $`m1`$. Then each $`g\mathrm{}(m|n)^\widehat{}`$-module $`F_s`$, $`s`$, is an irreducible integrable highest weight module of level $`1`$.
###### Remark 3.3.
The $`g\mathrm{}(0|n)^\widehat{}`$-modules $`F_s`$ are not irreducible. For example, one can show that in the case $`(m,n)=(0,2)`$, one has the following decomposition as $`g\mathrm{}(2)^\widehat{}`$-modules (in the standard notation of \[K3\]):
$$chF_s=\underset{j=0}{\overset{\mathrm{}}{}}chL((1+2j+|s|)\mathrm{\Lambda }_0+(2j+|s|)\mathrm{\Lambda }_1)q^{j^2+(|s|+1)j+|s|/2}.$$
E. Frenkel informed one of us that he had found this decomposition too.
The proof of Theorem 3.1 is based on the (super) boson-fermion correspondence, which we shall now recall (cf. \[K4\]).
For each $`i=1,\mathrm{},m`$ there exists a unique invertible odd operator $`e^{ϵ_i}`$ with inverse $`e^{ϵ_i}`$ satisfying the following three properties:
(3.10) $`[e^{ϵ_i},\psi ^j(z)]=0\text{ if }ij,[e^{ϵ_i},\phi ^j(z)]=0\text{ for all }j,`$
(3.11) $`e^{ϵ_i}\psi _k^{(i)}e^{ϵ_i}=\psi _{k1}^{(i)},e^{ϵ_i}\psi _k^{(i)}e^{ϵ_i}=\psi _{k+1}^{(i)},`$
(3.12) $`e^{ϵ_i}|0=\psi _{\frac{1}{2}}^{(i)}|0,e^{ϵ_i}|0=\psi _{\frac{1}{2}}^{(i)}|0.`$
It is easy to see that $`e^{ϵ_i}e^{ϵ_j}=e^{ϵ_j}e^{ϵ_i}`$ if $`ij`$.
We let for short $`(i=1,\mathrm{},m;j=1,\mathrm{},n)`$:
$$ϵ^i(z)=a^{ii+}(z)=\underset{k}{}ϵ_k^{(i)}z^{k1},ϵ^{j+m}(z)=a^{jj}(z)=\underset{k}{}ϵ_k^{(j+m)}z^{k1}.$$
Then we have:
(3.13)
$$[ϵ_k^{(i)},e^{ϵ_j}]=\delta _{ij}\delta _{k0}e^{ϵ_j},i=1,\mathrm{},m+n;j=1,\mathrm{},m.$$
For each $`i=1,\mathrm{},m+n`$ introduce the following fields:
$$\mathrm{\Gamma }_{ϵ_i}^+(z)=e^{_{k=1}^{\mathrm{}}\frac{z^k}{k}ϵ_k^{(i)}},\mathrm{\Gamma }_{ϵ_i}^{}(z)=e^{_{k=1}^{\mathrm{}}\frac{z^k}{k}ϵ_k^{(i)}},$$
and for a linear combination with integer coefficients $`\alpha =_{i=1}^ms_iϵ_i`$ we let $`\mathrm{\Gamma }_\alpha ^\pm (z)=\mathrm{\Pi }_i(\mathrm{\Gamma }_{ϵ_i}^\pm )^{s_i}`$ (recall that all $`ϵ_k^{(i)}`$ commute and all $`ϵ_k^{(i)}`$ commute for $`k1`$, see Proposition 3.1a). The central fact of the classical boson-fermion correspondence is the following formula, see e.g. \[K4\] $`(i=1,\mathrm{},m)`$
(3.14)
$$\psi ^i(z)=e^{ϵ_i}z^{ϵ_0^{(i)}}\mathrm{\Gamma }_{ϵ_i}^+(z)\mathrm{\Gamma }_{ϵ_i}^{}(z),\psi ^i(z)=e^{ϵ_i}z^{ϵ_0^{(i)}}\mathrm{\Gamma }_{ϵ_i}^+(z)\mathrm{\Gamma }_{ϵ_i}^{}(z).$$
The key formulas of the super boson-fermion correspondence are the following \[KL\], \[K4\]($`j=1,\mathrm{},n`$):
(3.15) $`\phi ^j(z)`$ $`=`$ $`z^{ϵ_0^{(i)}}e^{ϵ_i}\mathrm{\Gamma }_{ϵ_i}^+(z)E^{ji}(z)\mathrm{\Gamma }_{ϵ_i}^{}(z)`$
$`\phi ^j(z)`$ $`=`$ $`z^{ϵ_0^{(i)}}e^{ϵ_i}\mathrm{\Gamma }_{ϵ_i}^+(z)E^{ij+}(z)\mathrm{\Gamma }_{ϵ_i}^{}(z),`$
for each $`i=1,\mathrm{},m`$ (we assume here that $`m1`$).
###### Proof of Theorem 3.1.
Since the eigenspaces of $`L_0`$ in $`F_s`$ are finite-dimensional and $`L_0`$ commutes with $`g\mathrm{}(m|n)`$, it follows that $`F_s`$ is a direct sum of finite-dimensional $`g\mathrm{}(m|n)`$-modules, hence $`g\mathrm{}(m|n)`$ acts locally finitely on $`F_s`$. Furthermore, we have:
$$F=F^{fermi}F^{bose},$$
where $`F^{fermi}`$ (resp. $`F^{bose}`$) is the vertex algebra generated by the $`\psi ^i(z),\psi ^i(z)`$ (resp. $`\phi ^j(z),`$$`\phi ^j(z)`$), and the subalgebra $`g\mathrm{}(m)^\widehat{}`$ of $`g\mathrm{}(m|n)^\widehat{}`$ acts on $`F`$ via $`\pi 1`$, where the representation $`\pi `$ of $`g\mathrm{}(m)^\widehat{}`$ on $`F^{fermi}`$ is known to be integrable of level $`1`$ (see \[KP1\]). Thus, the representation of $`g\mathrm{}(m|n)^\widehat{}`$ in each $`F_s`$ is integrable.
The irreducibility of $`F_s`$, provided that $`m1`$, is proved using (3.14) and (3.15) in exactly the same fashion as the proof of Theorem 5.8a from \[K4\].
###### Remark 3.4.
We have got along the way the following vertex operator construction of $`g\mathrm{}(m|n)^\widehat{}`$. For each $`\alpha =_{i=1}^ms_iϵ_i,s_i`$, introduce the usual vertex operator
$$\mathrm{\Gamma }_\alpha =e^\alpha z^{\alpha _0}\mathrm{\Gamma }_\alpha ^+\mathrm{\Gamma }_\alpha ^{}.$$
Then the following map defines an irreducible integrable highest weight module of level $`1`$ in each $`F_s`$:
$`\begin{array}{cc}e_{ii}(z)ϵ_i(z)\hfill & (i=1,\mathrm{},m),K1,\hfill \\ e_{ij}(z)\mathrm{\Gamma }_{ϵ_iϵ_j}\hfill & (i,j=1,\mathrm{},m),\hfill \\ e_{i+m,j+m}(z):\phi ^i(z)\phi ^j(z):\hfill & (i,j=1,\mathrm{},n),\hfill \\ e_{j+m,i}(z)\mathrm{\Gamma }_{ϵ_i}(z)\phi ^j(z)\hfill & (i=1,\mathrm{},m;j=1,\mathrm{},n),\hfill \\ e_{i,j+m}(z)\mathrm{\Gamma }_{ϵ_i}(z)\phi ^j(z)\hfill & (i=1,\mathrm{},m;j=1,\mathrm{},n).\hfill \end{array}`$
Next, we give a standard derivation of a “quasiparticle” character formula for the $`g\mathrm{}(m|n)^\widehat{}`$-modules $`F_s`$, $`s`$.
Given $`a=(a_1,\mathrm{},a_m)`$, $`b=(b_1,\mathrm{},b_m)_+^m`$ and $`c=(c_1,\mathrm{},c_n)`$, $`d=(d_1,\mathrm{},d_n)_+^n`$, denote by $`F(a,b,c,d)`$ the linear span of vectors in $`F`$ obtained from the vacuum vector $`|0`$ by applying all monomials in the $`\psi _k^{(i)}`$, $`\psi _k^{(i)}`$, $`\phi _k^{(i)}`$, $`\phi _k^{(i)}`$ which contain $`a_1`$ factors of the form $`\psi _k^{(1)}`$, $`k\frac{1}{2}+,\mathrm{},a_m`$ factors of the form $`\psi _k^{(m)},b_1`$ factors of the form $`\psi _k^{(1)},\mathrm{},b_m`$ factors of the form $`\psi _k^{(m)}`$, $`c_1`$ factors of the form $`\phi _k^{(1)},\mathrm{},d_n`$ factors of the form $`\phi _k^{(n)}`$. These states lie in $`F_s`$ iff the following condition holds:
(3.17)
$$|a||b|+|c||d|=s,$$
where $`|a|=a_i`$, etc.
It is clear that the state of minimal energy in $`F(a,b,c,d)`$ is (up to a constant factor) the following vector:
$`v(a,b,c,d)`$ $`=`$ $`(\psi _{(a_1\frac{1}{2})}^{(1)}\mathrm{}\psi _{\frac{3}{2}}^{(1)}\psi _{\frac{1}{2}}^{(1)})\mathrm{}(\psi _{(a_m\frac{1}{2})}^{(m)}\mathrm{}\psi _{\frac{1}{2}}^{(m)})`$
$`\times (\psi _{(b_1\frac{1}{2})}^{(1)}\mathrm{}\psi _{\frac{1}{2}}^{(1)})\mathrm{}(\psi _{(b_m\frac{1}{2})}^{(m)}\mathrm{}\psi _{\frac{1}{2}}^{(m)})`$
$`\times (\phi _{\frac{1}{2}}^{(1)})^{c_1}\mathrm{}(\phi _{\frac{1}{2}}^{(n)})^{c_n}(\phi _{\frac{1}{2}}^{(1)})^{d_1}\mathrm{}(\phi _{\frac{1}{2}}^{(n)})^{d_n}|0.`$
All other basis elements from $`F(a,b,c,d)`$ are obtained from $`v(a,b,c,d)`$ by adding to the lower indices of the factors arbitrary non-negative integers. Hence we have (since weight $`|0=\mathrm{\Lambda }_0`$):
(3.18a) $`chF(a,b,c,d)`$ $`=`$ $`e^{weight(v(a,b,c,d))}/\mathrm{\Pi }(q),\text{ where}`$
(3.18b) $`\mathrm{\Pi }(q)`$ $`=`$ $`(q)_{a_1}\mathrm{}(q)_{a_m}(q)_{b_1}\mathrm{}(q)_{b_m}(q)_{c_1}\mathrm{}(q)_{c_n}(q)_{d_1}\mathrm{}(q)_{d_n}.`$
Here and further we use the usual notation and assumptions:
$$(q)_a=(1q)\mathrm{}(1q^a),q=e^\delta \text{ and }|q|<1.$$
Noticing that
(3.19) $`weight(\psi _k^{(i)})`$ $`=`$ $`ϵ_i+k\delta ,weight(\psi _k^{(i)})=ϵ_i+k\delta ,`$
$`weight(\phi _k^{(i)})`$ $`=`$ $`ϵ_{m+i}+k\delta ,weight(\phi _k^{(i)})=ϵ_{m+i}+k\delta ,`$
we obtain from (3.17) and (3.18) the “quasiparticle” character formula for $`F_s`$:
(3.20)
$$chF_s=e^{\mathrm{\Lambda }_0}\underset{\begin{array}{c}a,b_+^{m+n}\\ |a||b|=s\end{array}}{}\frac{e^{_{i=1}^{m+n}(a_ib_i)ϵ_i}q^{\frac{1}{2}_{i=1}^m(a_i^2+b_i^2)+\frac{1}{2}_{i=m+1}^{m+n}(a_i+b_i)}}{\mathrm{\Pi }_{i=1}^{m+n}(q)_{a_i}(q)_{b_i}}.$$
Another formula, which we call a theta function type character formula, is derived as follows. Let
$$chF=\underset{s}{}z^schF_s.$$
Using (3.7) and (3.19), we obtain:
(3.21)
$$chF=e^{\mathrm{\Lambda }_0}\mathrm{\Pi }_{k=1}^{\mathrm{}}\frac{\mathrm{\Pi }_{i=1}^m(1+ze^{ϵ_i}q^{k1/2})(1+z^1e^{ϵ_i}q^{k1/2})}{\mathrm{\Pi }_{j=1}^n(1ze^{ϵ_{m+j}}q^{k1/2})(1z^1e^{ϵ_{m+j}}q^{k1/2})}.$$
In order to compute the coefficient of $`z^s`$, we use the Jacobi triple product identity
(3.22)
$$\mathrm{\Pi }_{k=1}^{\mathrm{}}(1+zq^{k\frac{1}{2}})(1+z^1q^{k\frac{1}{2}})=\frac{1}{\phi (q)}\underset{m}{}z^mq^{\frac{1}{2}m^2},$$
and also the following well-known identity which can be derived from the super boson-fermion correspondence \[K4\]:
$`\mathrm{\Pi }_{k=1}^{\mathrm{}}(1+zq^{k\frac{1}{2}})^1(1+z^1q^{k\frac{1}{2}})^1={\displaystyle \frac{1}{\phi (q)^2}}{\displaystyle \underset{m}{}}(1)^m{\displaystyle \frac{q^{\frac{1}{2}m(m+1)}}{1+zq^{m+\frac{1}{2}}}}`$
$`=`$ $`\phi (q)^2\left({\displaystyle \underset{m,k0}{}}{\displaystyle \underset{m,k<0}{}}\right)((1)^{m+k}z^kq^{\frac{1}{2}m(m+1)+(m+\frac{1}{2})k}).`$
Here and further
$`\phi (q)`$ $`=`$ $`\mathrm{\Pi }_{j=1}^{\mathrm{}}(1q^j).`$
Substituting (3.22) and (3) in (3.21), we get:
$`chF`$ $`=`$ $`{\displaystyle \frac{e^{\mathrm{\Lambda }_0}}{\phi (q)^{m+2n}}}{\displaystyle \underset{k^m}{}}\left({\displaystyle \underset{p_1,a_10}{}}{\displaystyle \underset{p_1,a_1<0}{}}\right)\mathrm{}\left({\displaystyle \underset{p_n,a_n0}{}}{\displaystyle \underset{p_n,a_n<0}{}}\right)`$
$`(1)^{|r|}z^{|k|+|p|}e^{_ik_iϵ_i+_jp_jϵ_{m+j}}q^{\frac{1}{2}_ik_i^2+\frac{1}{2}_j(a_j(a_j+1)+p_j(a_j+1/2))},`$
where $`k=(k_1,\mathrm{},k_m)^m,p=(p_1,\mathrm{},p_n),a=(a_1,\mathrm{},a_n)^n`$, and $`|k|=_ik_i`$.
The coefficient of $`z^s`$ is a rather complicated expression for $`chF_s`$, which, after letting $`r=p+a^n`$, can be written as follows:
(3.25) $`chF_s={\displaystyle \frac{e^{\mathrm{\Lambda }_0+sϵ_1}q^{s^2/2}}{\phi (q)^{m+2n}}}{\displaystyle \underset{k^{m1}}{}}\left({\displaystyle \underset{r_1p_10}{}}{\displaystyle \underset{r_1<p_1<0}{}}\right)\mathrm{}\left({\displaystyle \underset{r_np_n0}{}}{\displaystyle \underset{r_n<p_n<0}{}}\right)`$
$`(1)^{|r+p|}e^{_ik_i(ϵ_iϵ_1)+_jp_j(ϵ_{m+j}ϵ_1)}q^{\frac{1}{2}|k|^2+\frac{1}{2}_ik_i^2+\frac{1}{2}_jr_j(r_j+1)+_{i<j}p_ip_j+|k||p|s(|k|+|p|)},`$
where $`k=(k_2,\mathrm{},k_m)^{m1},p,r^n`$.
We rewrite (3.25) using translation operators $`t_\alpha `$, $`\alpha 𝔥^{}`$, defined by $`(\lambda \widehat{𝔥}^{})`$:
(3.26)
$$t_\alpha (\lambda )=\lambda +(\lambda |\delta )\alpha (\frac{1}{2}(\alpha |\alpha )(\lambda |\delta )+(\lambda |\alpha ))\delta .$$
Let $`M^\mathrm{\#}=_{i=1}^{m1}\alpha _i`$; recall that $`M^\mathrm{\#}`$ acts on $`\widehat{𝔥}^{}`$ via $`\alpha t_\alpha `$ and the image of this action is the translation subgroup of the Weyl group of $`s\mathrm{}(m)^\widehat{}`$. It is straightforward to show that (3.25) can be rewritten as follows:
(3.27) $`chF_s={\displaystyle \frac{1}{\phi (q)^{m+2n}}}\left({\displaystyle \underset{r_1p_10}{}}{\displaystyle \underset{r_1<p_1<0}{}}\right)\mathrm{}\left({\displaystyle \underset{r_np_n0}{}}{\displaystyle \underset{r_n<p_n<0}{}}\right)`$
$`(1)^{|r|+|p|}q^{\frac{1}{2}_jr_j(r_j+1)+_{i<j}p_ip_j|p|+a_s}{\displaystyle \underset{\alpha M^\mathrm{\#}}{}}e^{t_\alpha (\mathrm{\Lambda }_{(s)}_{j=1}^np_j(ϵ_1ϵ_{m+j}))},`$
where $`\mathrm{\Lambda }_{(s)}`$ is the weight of $`|s`$ and $`a_s=s(r_np_n+1)+|p|`$ if $`s0,=0`$ if $`0<sm`$, and $`=(sm)(r_1p_1)`$ if $`sm`$.
In the case $`n=1`$ formula (3.25) can be simplified by making use of the following lemma.
###### Lemma 3.1.
Let $`a,b`$. Then $`(j,k,n)`$:
$`\begin{array}{cc}\text{(a)}\hfill & (_{kja}_{k<j<a})(1)^{j+k}x^jq^{bk}q^{k(k+1)/2}=\frac{x}{1+x}\phi (q)\mathrm{\Pi }_{n=1}^{\mathrm{}}(1+x^1q^{nb1})(1+xq^{n+b}).\hfill \\ \text{(b)}\hfill & \mathrm{\Pi }_{n=1}^{\mathrm{}}(1+x^1q^{nb1})(1+xq^{n+b})=x^bq^{b(b+1)/2}\mathrm{\Pi }_{n=1}^{\mathrm{}}(1+x^1q^{n1})(1+xq^n)\hfill \\ & =x^{b1}q^{b(b+1)/2}\mathrm{\Pi }_{n=1}^{\mathrm{}}(1+x^1q^n)(1+xq^{n1}).\hfill \end{array}`$
###### Proof.
If $`ka`$ (resp. $`k<a`$), we have:
$$\underset{j=a}{\overset{k}{}}(1)^{j+k}x^j(\text{resp. }\underset{j=k+1}{\overset{a1}{}}(1)^{j+k}x^j)=\frac{x^{k+1}(1)^kx^a}{1+x}.$$
Hence the LHS of (a) is equal to
$$\frac{1}{1+x}\underset{k}{}x^{k+1}q^{bk+k(k+1)/2}\frac{x^a}{1+x}\underset{k}{}(1)^kq^{bk+k(k+1)/2}.$$
Noticing that the second summand is zero and applying to the first summand the Jacobi triple product identity (3.22), we obtain (a). In the proof of (b) we assume that $`b>0`$, the case $`b<0`$ being similar:
$$\mathrm{\Pi }_{n=1}^{\mathrm{}}(1+x^1q^{nb1})=\mathrm{\Pi }_{m1b}(1+x^1q^{m1})=\mathrm{\Pi }_{m=1}^{\mathrm{}}(1+x^1q^{m1})\mathrm{\Pi }_{m=1b}^0(1+x^1q^{m1}).$$
The second product on the RHS is equal to
$$\mathrm{\Pi }_{n=1}^b(1+x^1q^n)=\mathrm{\Pi }_{n=1}^bx^1q^n(1+xq^n)=x^bq^{b(b+1)/2}\mathrm{\Pi }_{n=1}^b(1+xq^n).$$
Next, we have:
$$\mathrm{\Pi }_{n=1}^{\mathrm{}}(1+xq^{n+b})=\mathrm{\Pi }_{m=1+b}^{\mathrm{}}(1+xq^m)=\mathrm{\Pi }_{m=1}^{\mathrm{}}(1+xq^m)/\mathrm{\Pi }_{m=1}^b(1+xq^m).$$
These equalities prove (b). ∎
Let now $`n=1`$. Then (3.25) reads:
(3.29) $`chF_s`$ $`=`$ $`{\displaystyle \frac{e^{\mathrm{\Lambda }_0+sϵ_1}q^{s^2/2}}{\phi (q)^{m+2}}}{\displaystyle \underset{k^{m1}}{}}\psi (k)e^{_ik_i(ϵ_iϵ_1)}q^{\frac{1}{2}|k|(|k|2s)+\frac{1}{2}_ik_i^2},`$
where
$`\psi (k)`$ $`=`$ $`\left({\displaystyle \underset{tj0}{}}{\displaystyle \underset{t<j<0}{}}\right)(1)^{j+t}\left(q^{|k|s}e^{ϵ_{m+1}ϵ_1}\right)^jq^{t(t+1)/2}.`$
By Lemma 3.1a for $`x=q^{|k|s}e^{ϵ_{m+1}ϵ_1},b=0`$, we have:
$$\psi (k)=\frac{q^{|k|s}e^{ϵ_{m+1}ϵ_1}\phi (q)}{1+q^{|k|s}e^{ϵ_{m+1}ϵ_1}}\mathrm{\Pi }_{n=1}^{\mathrm{}}\left(1+e^{ϵ_1ϵ_{m+1}}q^{n|k|+s1}\right)\left(1+e^{ϵ_{m+1}ϵ_1}q^{n+|k|s}\right).$$
By Lemma 3.1b for $`x=e^{ϵ_{m+1}ϵ_1},b=|k|s`$, we rewrite this as follows:
$$\psi (k)=\frac{x^bq^{b(b1)/2}\phi (q)}{1+q^bx}\mathrm{\Pi }_{n=1}^{\mathrm{}}(1+x^1q^n)(1+xq^{n1}).$$
Substituting this in (3.29), we obtain:
$`chF_s`$ $`=`$ $`{\displaystyle \frac{e^{\mathrm{\Lambda }_0+sϵ_{m+1}}q^{\frac{s}{2}}}{\phi (q)^{m+1}}}\mathrm{\Pi }_{n=1}^{\mathrm{}}(1+e^{ϵ_1ϵ_{m+1}}q^n)(1+e^{ϵ_{m+1}ϵ_1}q^{n1})`$
$`\times `$ $`{\displaystyle \underset{k^{m1}}{}}{\displaystyle \frac{e^{_{i=2}^mk_i(ϵ_iϵ_{m+1})}}{1+q^{|k|s}e^{ϵ_{m+1}ϵ_1}}}q^{\frac{1}{2}_{i=2}^mk_i(k_i+1)}.`$
This formula agrees with the one obtained in \[KL\] (see also \[K4\]) for $`g\mathrm{}(1|1)^\widehat{}`$.
In the next section we use this formula in the case $`m2`$ in order to derive character formulas for all integrable level $`1`$ $`s\mathrm{}(m|1)^\widehat{}`$-modules in terms of the theta function and the (multivariable) Appell’s functions, and to obtain their high temperature asymptotics.
## 4 Theta function type character formula for integrable level $`1`$ $`s\mathrm{}(m|1)^\widehat{}`$-modules and Appell’s function.
Recall that Appell’s function is defined by the following series (cf. \[A\] and \[P\]):
$$A(x,z,q)=\underset{k}{}\frac{q^{\frac{1}{2}k^2}z^k}{1+xq^k},$$
which converges to a meromorphic function in the domain $`x,z,q`$, $`|q|<1`$. The classical theta function in one variable is a special case of this function:
$$\mathrm{\Theta }(z)\mathrm{\Theta }(z;q)=A(0,z,q).$$
Note that by (3.22) we have a product expansion:
(4.1)
$$\mathrm{\Theta }(zq^{1/2};q)=\phi (q)\mathrm{\Pi }_{k=1}^{\mathrm{}}(1+zq^k)(1+z^1q^{k1}).$$
We shall need also the following multivariable generalization of Appell’s function. Let $`B`$ be an $`N\times N`$ symmetric matrix such that $`ReB`$ is positive definite and let $`\mathrm{}`$ be a linear function of $`^N`$. We define the series
$$A_{B,\mathrm{}}(x;z_1,\mathrm{},z_N;q)=\underset{k^N}{}\frac{q^{\frac{1}{2}k^TBk}z_1^{k_1}\mathrm{}z_N^{k_N}}{1+xq^{\mathrm{}(k)}},$$
which converges to a meromorphic function provided that $`|q|<1`$. Again, letting $`x=0`$, we get the multivariable theta function.
Consider now the $`g\mathrm{}(m|1)^\widehat{}`$-modules $`F_s`$, $`s`$, and assume in this section that $`m2`$. We have (see Remark 3.2):
$`\mathrm{\Lambda }_{(s)}=\{\begin{array}{cc}\mathrm{\Lambda }_s\frac{s}{2}\delta \hfill & \text{if }0sm\hfill \\ (sm)\mathrm{\Lambda }_0+(1+sm)\mathrm{\Lambda }_m\frac{s}{2}\delta \hfill & \text{if }sm\hfill \\ (1s)\mathrm{\Lambda }_0+s\mathrm{\Lambda }_m+\frac{s}{2}\delta \hfill & \text{if }s0\hfill \end{array}.`$
For $`m2`$, we have: $`g\mathrm{}(m|1)^\widehat{}=s\mathrm{}(m|1)^\widehat{}+g\mathrm{}(1)^\widehat{}`$ (sum of ideals), hence:
$$chF_s=chL(\mathrm{\Lambda }_{(s)})\phi (q)^1,$$
where $`L(\mathrm{\Lambda }_{(s)})`$ denotes the irreducible $`s\mathrm{}(m|1)^\widehat{}`$-module with highest weight $`\mathrm{\Lambda }_{(s)}`$. Hence formula (3) gives us the following expression for $`chL(\mathrm{\Lambda }_{(s)})`$ in terms of the theta function $`\mathrm{\Theta }(z;q)`$ (we use (4.1)) and the (multivariable) Appell’s function:
(4.3)
$$chL(\mathrm{\Lambda }_{(s)})=\frac{e^{\mathrm{\Lambda }_0+sϵ_{m+1}}q^{\frac{s}{2}}}{\phi (q)^{m+1}}\mathrm{\Theta }(z_1q^{\frac{1}{2}};q)A_{I,\mathrm{}}(z_1^1q^s;z_2q^{\frac{1}{2}},\mathrm{},z_mq^{\frac{1}{2}};q),$$
where $`z_i=e^{ϵ_iϵ_{m+1}}`$ $`(i=1,\mathrm{},m)`$, $`B=I`$ is the $`(m1)\times (m1)`$ identity matrix and $`\mathrm{}(k)=_ik_i`$, $`k^{m1}`$. (Note that in the simplest case $`m=2`$ we get in this expression the classical Appell’s function.)
Next, we derive yet another character formula for the $`s\mathrm{}(m|1)^\widehat{}`$-module $`L(\mathrm{\Lambda }_0)`$ in the case $`m2`$, in terms of classical theta functions and certain “half” modular forms. We use for this (3.27) for $`n=1`$:
(4.4) $`chL(\mathrm{\Lambda }_0)={\displaystyle \frac{1}{\phi (q)^{m+1}}}\left({\displaystyle \underset{rp0}{}}{\displaystyle \underset{r<p<0}{}}\right){\displaystyle \underset{\alpha M^\mathrm{\#}}{}}(1)^{r+p}e^{t_\alpha (\mathrm{\Lambda }_0p(ϵ_1ϵ_{m+1}))}q^{\frac{1}{2}r(r+1)}.`$
Introduce the following elements of $`M^\mathrm{\#}`$:
$$\beta _k=kϵ_1\underset{i=1}{\overset{k}{}}ϵ_1,k=1,\mathrm{},m,$$
and the element $`\mu =ϵ_{m+1}\frac{1}{m}(ϵ_1+\mathrm{}+ϵ_m)𝔥^{}`$, which is orthogonal to $`M^\mathrm{\#}`$. The even part of $`s\mathrm{}(m|1)^\widehat{}`$ is a sum of ideals $`s\mathrm{}(m)^\widehat{}`$ and $`(\mu )^\widehat{}`$.
Write $`p=jmk`$, where $`j`$, $`1km`$ and note that
(4.5)
$$\mathrm{\Lambda }_0p(ϵ_1ϵ_{m+1})+j\beta _m\beta _k=\dot{\mathrm{\Lambda }}_k+p\mu .$$
where $`\dot{\mathrm{\Lambda }}_k`$ denote fundamental weights of $`s\mathrm{}(m)^\widehat{}`$ and we identify $`\dot{\mathrm{\Lambda }}_m`$ with $`\dot{\mathrm{\Lambda }}_0`$. Adding to $`\alpha `$ the element $`j\beta _m\beta _k`$ in (4.4), and using (4.5), we rewrite (4.4) as follows:
$`chL(\mathrm{\Lambda }_0)`$ $`=`$ $`{\displaystyle \frac{1}{\phi (q)^{m+1}}}{\displaystyle \underset{k=1}{\overset{m}{}}}\left({\displaystyle \underset{\begin{array}{c}j,r\\ r+kjm\\ j>0\end{array}}{}}{\displaystyle \underset{\begin{array}{c}j,rZ\\ r+k<jm\\ j0\end{array}}{}}\right)`$
$`(1)^{r+k+jm}{\displaystyle \underset{\alpha M^\mathrm{\#}}{}}e^{t_\alpha (\dot{\mathrm{\Lambda }}_k)+(jmk)\mu }q^{\frac{r(r+1)}{2}\frac{(jmk)(jmkj)}{2}\frac{k(j1)}{2}}.`$
Denoting by $`\dot{L}(\dot{\mathrm{\Lambda }},a)`$ the irreducible $`s\mathrm{}(m)^\mathrm{^}+(\mu )^\mathrm{^}`$-module with highest weight $`\dot{\mathrm{\Lambda }}+a\mu `$, and recalling that (\[K3\], Proposition 12.13):
(4.6)
$$ch\dot{L}(\dot{\mathrm{\Lambda }}_k,a)=\frac{1}{\phi (q)^m}\underset{\alpha M^\mathrm{\#}}{}e^{t_\alpha (\dot{\mathrm{\Lambda }}_k)+a\mu },$$
we obtain:
(4.7) $`chL(\mathrm{\Lambda }_0)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \underset{\begin{array}{c}p\\ p+k|m\end{array}}{}}b_{k,p}(q)ch\dot{L}(\dot{\mathrm{\Lambda }}_k,p),`$
where
$`b_{k,p}(q)`$ $`=`$ $`{\displaystyle \frac{q^{(\frac{1}{2}\frac{1}{2m})p^2+\frac{1}{2m}k^2+\frac{1}{2}k}}{\phi (q)}}\times {\displaystyle \underset{rp}{}}(1)^{rp}q^{\frac{1}{2}r(r+1)}(\text{resp. }\times {\displaystyle \underset{r<p}{}}(1)^{rp1}q^{\frac{1}{2}r(r+1)}),`$
if $`p0`$ (resp. $`p<0`$).
Thus, the branching functions $`b_{k,p}(q)`$ are “half” modular functions, in a sharp contrast with the case of affine Lie algebras \[K3\]. Recalling that the series $`_{\alpha M^\mathrm{\#}}e^{t_\alpha (\mathrm{\Lambda })}`$ converges to a classical theta function \[K3\], we see that the character of the “basic” $`s\mathrm{}(m|1)^\widehat{}`$-module is a finite linear combination of classical theta functions with coefficients “half” modular functions.
The basic specialization of (3) gives the specialized character formulas for $`s\mathrm{}(m|1)^\widehat{}`$-modules $`L(\mathrm{\Lambda }_{(s)})`$:
(4.9)
$$tr_{L(\mathrm{\Lambda }_{(s)})}q^{L_0}=2q^{\frac{s}{2}}\frac{\phi (q^2)^2}{\phi (q)^{m+2}}\underset{k^{m1}}{}\frac{q^{\frac{1}{2}_ik_i(k_i+1)}}{1+q^{|k|s}}.$$
where, as before, $`|k|=_ik_i`$.In the remainder of this section we discuss asymptotics of (4.9).
Given a positive definite quadratic form $`B(x)`$ on $`^N`$, an affine linear function $`\mathrm{}(x)`$ on $`^N`$ and an element $`\alpha `$ of $`^N`$, consider the following series, where $`q=e^{2\pi i\tau }`$:
$$f_{B,\mathrm{},\alpha }(\tau )=\underset{\gamma ^N+\alpha }{}\frac{q^{\frac{1}{2}B(\gamma )}}{1+q^{\mathrm{}(\gamma )}}.$$
This series converges on the upper half plane $`Im\tau >0`$ to a specialization of the multivariable Appell’s function. From the transformation properties of theta series one gets (see e.g. \[K3\]):
(4.10)
$$f_{B,0,\alpha }(i\beta )=\frac{1}{2}(detB)^{1/2}\beta ^{N/2}+o(\beta )\text{ as }\beta 0(Re\beta >0).$$
In order to get the asymptotics of the functions $`f_{B,\mathrm{},\alpha }(\tau )`$, let $`f_{B,\mathrm{},\alpha }^\pm (\tau )={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\gamma ^N+\alpha \\ \pm \mathrm{}(\gamma )>0\end{array}}{}}q^{\frac{1}{2}B(\gamma )}`$. It is easy to derive from (4.10) by induction on $`n`$ the following asymptotics:
(4.11)
$$f_{B,\mathrm{},\alpha }^\pm (i\beta )=\frac{1}{4}(detB)^{1/2}\beta ^{N/2}+p_\pm (\beta ^{1/2})+o(\beta ),$$
as $`\beta 0,\beta _+`$, where $`p_\pm (x)`$ is a polynomial in $`x`$ of degree strictly less than $`N`$.
The idea of the following lemma is due to A. Polishchuk.
###### Lemma 4.1.
$`|f_{B,\mathrm{},\alpha }(i\beta )f_{B,0,\alpha }(i\beta )|<p(\beta ^{1/2})`$ for $`\beta ,0<\beta <a`$, where $`a`$ is a positive number and $`p(x)`$ is a polynomial in $`x`$ of degree strictly less than $`N`$.
###### Proof.
Let $`g(\beta )=f_{B,\mathrm{},\alpha }(i\beta )f_{B,0,\alpha }(i\beta )`$. We have: $`g(\beta )=g^+(\beta )g^{}(\beta )`$, where
$$g^\pm (\beta )=\underset{\begin{array}{c}\gamma ^N+\alpha \\ \pm \mathrm{}(\gamma )>0\end{array}}{}e^{\pi \beta B(\gamma )}\frac{1e^{2\pi \beta \mathrm{}(\gamma )}}{1+e^{2\pi \beta \mathrm{}(\gamma )}}.$$
Furthermore, we have:
$$0g^\pm (\beta )\underset{\begin{array}{c}\gamma ^N+\alpha \\ \pm \mathrm{}(\gamma )>0\end{array}}{}e^{\pi \beta B(\gamma )}\underset{\begin{array}{c}\gamma ^N+\alpha \\ \pm \mathrm{}(\gamma )>0\end{array}}{}e^{\pi \beta (B(\gamma )\pm 2\mathrm{}(\gamma ))}.$$
The first sum on the right is just $`f_{B,\mathrm{},\alpha }^\pm (i\beta )`$, which has asymptotics (4.11). But the second sum on the right has asymptotics of this form too since it can be written as a product of a power of $`q`$ and a function $`f_{B,\mathrm{}^{},\alpha }^\pm (i\beta )`$ for some other affine linear function $`\mathrm{}^{}`$, by “completing the squares”. The lemma is proved. ∎
We shall write $`f(\tau )g(\tau )`$ if $`lim_{\begin{array}{c}\beta 0\\ \beta _+\end{array}}f(i\beta )/g(i\beta )=1`$. Lemma 4.1 and (4.10) imply:
(4.12)
$$f_{B,\mathrm{},\alpha }(\tau )\frac{1}{2}(detB)^{1/2}\beta ^{N/2}.$$
Since
(4.13)
$$\eta (\tau )\beta ^{1/2}e^{\pi /12\beta },$$
we deduce from (4.12) and (4.9) the following asymptotics along the imaginary axis $`\tau =i\beta ,\beta _+`$:
(4.14)
$$tr_{L(\mathrm{\Lambda }_{(s)})}q^{L_0}\frac{1}{2}\beta ^{\frac{1}{2}}e^{\frac{\pi }{12\beta }(m+1)}.$$
## 5 A Weyl type character formula for integrable level $`1`$ $`g\mathrm{}(m|n)^\widehat{}`$-modules.
In this section we derive a Weyl type character formula (5.13) for principal integrable level $`1`$-modules over $`g\mathrm{}(m|n)^\widehat{}`$ provided that $`mn`$. We use for that formula (3.21) for $`chF`$ and the denominator identity for $`s\mathrm{}(m+1|n)^\widehat{}`$. In order to compare these two formulas, we consider the labelings of simple roots of $`g\mathrm{}(m|n)^\widehat{}`$ and $`s\mathrm{}(m+1|n)^\widehat{}`$ given below:
Putting
(5.1)
$$z=e^{ϵ_1}q^{\frac{1}{2}}y\text{ and }y=e^\alpha _{},$$
we can rewrite formula (3.21) as follows:
(5.2)
$$e^{\mathrm{\Lambda }_0}chF=\mathrm{\Pi }_{k=1}^{\mathrm{}}\frac{\mathrm{\Pi }_{i=0}^{m1}(1e^{\alpha _{}+\alpha _1+\mathrm{}+\alpha _i}q^k)(1e^{\alpha _{}\alpha _1\mathrm{}\alpha _i}q^{k1})}{\mathrm{\Pi }_{j=0}^{n1}(1+e^{\alpha _{}+\alpha _1+\mathrm{}+\alpha _{m+j}}q^k)(1+e^{\alpha _{}\alpha _1\mathrm{}\alpha _{m+j}}q^{k1})}.$$
Denote by $`W^\mathrm{\#}`$ (resp. $`\stackrel{~}{W}^\mathrm{\#}`$) the subgroup of the Weyl group of $`g\mathrm{}(m|n)`$ (resp. $`s\mathrm{}(m+1|n)`$) generated by reflections $`r_\alpha `$ in roots $`\alpha =\alpha _1,\mathrm{},\alpha _{m1}`$ (resp. $`\alpha _{},\alpha _1,\mathrm{},\alpha _{m1}`$) and by $`M^\mathrm{\#}`$ (resp. $`\stackrel{~}{M}^\mathrm{\#}`$) the subgroups of these Weyl groups generated by translations $`t_\alpha `$ in integral linear combinations of these roots. Let $`\stackrel{~}{\rho }`$ denote a Weyl vector for $`s\mathrm{}(m+1|n)^\widehat{}`$ and let $`R`$ denote the denominator for $`g\mathrm{}(m|n)^\widehat{}`$. It is clear by (5.2) that the denominator $`\stackrel{~}{R}`$ of $`s\mathrm{}(m+1|n)^\widehat{}`$ is given by
(5.3)
$$\stackrel{~}{R}=e^{\mathrm{\Lambda }_0}RchF.$$
In order to write down the denominator identity for $`s\mathrm{}(m+1|n)^\widehat{}`$, introduce the roots $`\beta _{ij}=\alpha _i+\alpha _{i+1}+\mathrm{}+\alpha _j`$, $`1ijm+n2`$ (here $`\beta _{ii}=\alpha _i`$), and let $`\beta _i=\beta _{mn+i,m+i1}`$, $`i=1,\mathrm{},n`$. Let
(5.4)
$$\stackrel{~}{\rho }^{}=\stackrel{~}{\rho }+\underset{\begin{array}{c}mn+2imjm+n2\\ jin2\end{array}}{}\beta _{ij}.$$
Then the denominator identity for $`s\mathrm{}(m+1|n)^\widehat{}`$ looks as follows:
(5.5)
$$e^{\mathrm{\Lambda }_0}e^{\stackrel{~}{\rho }^{}}RchF=\underset{w\stackrel{~}{W}^\mathrm{\#}\stackrel{~}{M}^\mathrm{\#}}{}ϵ(w)w\frac{e^{\stackrel{~}{\rho }^{}}}{\mathrm{\Pi }_{j=1}^n(1+e^{\beta _j})}.$$
This identity can be derived from the denominator identity given in \[KW\] by making use of odd reflections as follows.
The denominator identity for $`s\mathrm{}(m+1|n)^\widehat{}`$ in \[KW\] is given for the choice of the set of simple roots with a maximal number of grey nodes:
Let $`\gamma _1,\mathrm{}\gamma _n`$ be the (unique subset) of the set of simple roots without $`\alpha _0`$ such that $`(\gamma _i|\gamma _j)=0`$ for all $`i,j`$, and let $`\stackrel{~}{\rho }^{\prime \prime }`$ be its Weyl vector. Then the identity reads:
(5.6)
$$e^{\stackrel{~}{\rho }^{\prime \prime }}\stackrel{~}{R}=\underset{w\stackrel{~}{W}^\mathrm{\#}\stackrel{~}{M}^\mathrm{\#}}{}ϵ(w)w\frac{e^{\stackrel{~}{\rho }^{\prime \prime }}}{\mathrm{\Pi }_{j=1}^n(1+e^{\gamma _j})}.$$
In order to derive (5.5) from (5.6), we apply a sequence of odd reflections which transforms the initial diagram with two grey nodes to the above final diagram with $`2n`$ grey nodes. In order to explain this sequence, denote by $`\mathrm{\Pi }^{(i,j)}`$ a set of simple roots containing $`\beta _{i,m+j}`$ and by $`\mathrm{\Pi }^{(i,j+1)}`$ the set of simple roots obtained from it by the odd reflection in $`\beta _{i,m+j}`$. Denoting by $`\stackrel{~}{\mathrm{\Pi }}`$ and $`\stackrel{~}{\mathrm{\Pi }}^{}`$ the initial and the final sets of simple roots, we have the following sequence:
$`\stackrel{~}{\mathrm{\Pi }}`$ $`=`$ $`\mathrm{\Pi }^{(m,0)}\mathrm{\Pi }^{(m,1)}\mathrm{}\mathrm{\Pi }^{(m,n2)}\mathrm{\Pi }^{(m,n1)}`$
$`=`$ $`\mathrm{\Pi }^{(m1,0)}\mathrm{\Pi }^{(m1,1)}\mathrm{}\mathrm{\Pi }^{(m1,n2)}`$
$`=`$ $`\mathrm{\Pi }^{(m2,0)}\mathrm{\Pi }^{(m2,1)}\mathrm{}\mathrm{\Pi }^{(m2,n3)}`$
$`=`$ $`\mathrm{\Pi }^{(m3,0)}\mathrm{\Pi }^{(m3,1)}\mathrm{}`$
$`\mathrm{}`$ $`=`$ $`\mathrm{\Pi }^{(mn+2,0)}\mathrm{\Pi }^{(mn+2,1)}=\stackrel{~}{\mathrm{\Pi }}^{}.`$
Using Lemma 1.3, one sees that $`\stackrel{~}{\rho }^{}`$ and $`\stackrel{~}{\rho }`$ are related by formula (5.4) and using Lemma 1.4a, we see that
(5.7)
$$\mathrm{\Lambda }_0^{}=\mathrm{\Lambda }_0.$$
Using the decompositions
$$\stackrel{~}{W}^\mathrm{\#}=W^\mathrm{\#}(_{j=0}^{m1}W^\mathrm{\#}r_{\alpha _{}+\alpha _1+\mathrm{}+\alpha _j})\text{ and }\stackrel{~}{M}^\mathrm{\#}=\alpha _{}+M^\mathrm{\#},$$
we obtain from (5.5):
(5.8) $`e^{\mathrm{\Lambda }_0}RchF`$ $`=`$ $`I_0{\displaystyle \underset{i=0}{\overset{m1}{}}}II_i,`$
where
$`I_0`$ $`=`$ $`e^{\stackrel{~}{\rho }^{}}{\displaystyle \underset{k}{}}{\displaystyle \underset{\alpha M^\mathrm{\#}}{}}{\displaystyle \underset{wW^\mathrm{\#}}{}}ϵ(w)t_{k\alpha _{}}t_\alpha w{\displaystyle \frac{e^{\stackrel{~}{\rho }^{}}}{\mathrm{\Pi }_{j=1}^n(1+e^{\beta _j})}}`$
$`II_i`$ $`=`$ $`e^{\stackrel{~}{\rho }^{}}{\displaystyle \underset{k}{}}{\displaystyle \underset{\alpha M^\mathrm{\#}}{}}{\displaystyle \underset{wW^\mathrm{\#}}{}}ϵ(w)t_{k\alpha _{}}t_\alpha wr_{\alpha _{}+\alpha _1+\mathrm{}+\alpha _i}{\displaystyle \frac{e^{\stackrel{~}{\rho }^{}}}{\mathrm{\Pi }_{j=1}^n(1+e^{\beta _j})}}.`$
In order to compute $`chL(\mathrm{\Lambda }_0)`$ of the $`g\mathrm{}(m|n)^\widehat{}`$-module $`L(\mathrm{\Lambda }_0)`$ we compare the constant terms (i.e., $`y^0`$-terms) in the decomposition of both sides of (5.8) is the series in powers of $`y`$.
We will show that
(5.10) $`\text{constant term of }II_i`$ $`=`$ $`0\text{ for all }i,`$
(5.11) $`\text{constant term of }I_0`$ $`=`$ $`{\displaystyle \underset{wW^\mathrm{\#}M^\mathrm{\#}}{}}ϵ(w)w{\displaystyle \frac{e^{\stackrel{~}{\rho }^{}}}{\mathrm{\Pi }_{j=1}^n(1+e^{\beta _j})}}.`$
Using that $`\stackrel{~}{\rho }^{}`$ when restricted to the Cartan subalgebra of $`g\mathrm{}(m|n)^\widehat{}`$ coincides with $`\mathrm{\Lambda }_0+\rho ^{}`$, where $`\rho ^{}`$ is related to the Weyl vector $`\rho `$ of $`g\mathrm{}(m|n)^\widehat{}`$ by (5.4) with $``$ removed, we obtain from (5.5), (5.10) and (5.11):
(5.12)
$$e^\rho ^{}RchL(\mathrm{\Lambda }_0)=\underset{wW^\mathrm{\#}M^\mathrm{\#}}{}ϵ(w)w\frac{e^{\mathrm{\Lambda }_0+\rho ^{}}}{\mathrm{\Pi }_{j=1}^n(1+e^{\beta _j})}.$$
Applying to this formula odd reflections as above, we obtain an equivalent character formula for the choice of the diagram with $`2n`$ grey nodes as above (where $`\rho `$ is a Weyl vector for this choice of the diagram):
(5.13)
$$e^\rho RchL(\mathrm{\Lambda }_0)=\underset{wW^\mathrm{\#}M^\mathrm{\#}}{}ϵ(w)w\frac{e^{\mathrm{\Lambda }_0+\rho }}{\mathrm{\Pi }_{j=1}^n(1+e^{\alpha _{mn+2j1}})}.$$
The proof of (5.10) and (5.11) is straight-forward, and we explain it in the case $`n=1`$. We have:
$$e^{\stackrel{~}{\rho }}t_{k\alpha _{}}r_{\alpha _{}+\alpha _1+\mathrm{}+\alpha _{m1}}\left(\frac{e^{\stackrel{~}{\rho }}}{1+e^{\alpha _m}}\right)=\frac{e^{mk\alpha _{}}e^{m(\alpha _{}+\alpha _1+\mathrm{}+\alpha _{m1})}q^{k^2mk+mk}}{1+e^{(\alpha _{}+\alpha _1+\mathrm{}+\alpha _m)}q^k}.$$
Hence:
$$e^{\stackrel{~}{\rho }}\underset{k}{}t_{k\alpha _{}}r_{\alpha _{}+\alpha _1+\mathrm{}+\alpha _{m1}}\left(\frac{e^{\stackrel{~}{\rho }}}{1+e^{\alpha _m}}\right)=\left(\underset{k,s0}{}\underset{k,s<0}{}\right)(1)^sy^{m(k+1)+s}(\mathrm{}),$$
where $`(\mathrm{})`$ doesn’t involve $`y`$. But the constant term of the last expression is $`0`$ since$`m(k+1)+s>0`$ if $`k,s0`$ and $`m(k+1)+ss<0`$ if $`k,s<0`$. Thus, constant term of $`II_{m1}=0`$. Furthermore, we have for $`0im2`$:
$$e^{\stackrel{~}{\rho }}t_{k\alpha _{}}r_{\alpha _{}+\alpha _1+\mathrm{}+\alpha _i}\frac{e^{\stackrel{~}{\rho }}}{1+e^{\alpha _m}}=y^{mk+i+1}(\mathrm{}).$$
Since $`mk+i+10`$ if $`0im2`$, we see that constant term of $`II_i`$ is $`0`$ for all $`i`$. Finally:
$$e^{\stackrel{~}{\rho }}t_{k\alpha _{}}\frac{e^{\stackrel{~}{\rho }}}{1+e^{\alpha _m}}=y^{mk}\frac{q^{k^2m+k}}{1+e^{\alpha _m}},$$
hence the constant term of this expression is equal $`\delta _{k,0}(1+e^{\alpha _m})^1`$, which proves (5.11).
Using odd reflections one may derive from (5.12) or (5.13) the Weyl-type character formulas for all other level $`1`$ principal integrable modules. For example, in the case $`\widehat{𝔤}=s\mathrm{}(m|1)^\widehat{}`$, we let $`k=mj+s_+`$, where $`j_+`$ and $`0sm1`$; then
(5.14) $`e^\rho Rch((k+1)\mathrm{\Lambda }_0k\mathrm{\Lambda }_m)`$ $`=`$ $`{\displaystyle \underset{w\widehat{W}}{}}ϵ(w)w{\displaystyle \frac{e^{\mathrm{\Lambda }+\rho }}{1+e^{j\delta \alpha _{ms}\mathrm{}\alpha _m}}},`$
(5.15) $`e^\rho Rch(k\mathrm{\Lambda }_0+(k+1)\mathrm{\Lambda }_m)`$ $`=`$ $`{\displaystyle \underset{w\widehat{W}}{}}ϵ(w)w{\displaystyle \frac{e^{\mathrm{\Lambda }+\rho }}{1+e^{j\delta \alpha _0\mathrm{}\alpha _s}}},`$
and for $`1jm1`$ we have:
(5.16)
$$e^\rho Rch(\mathrm{\Lambda }_j)=\underset{w\widehat{W}}{}ϵ(w)w\frac{e^{\mathrm{\Lambda }_j+\rho }}{1+e^{\alpha _m}}.$$
## 6 Classification of integrable highest weight modules over affine superalgebras.
In this section we consider affine superalgebras of type $`A(m,n)^\widehat{}`$, $`B(m,n)^\widehat{}`$, $`C(n)^\widehat{}`$, $`D(m,n)^\widehat{}`$, $`D(2,1;a)^\widehat{}`$, $`F(4)^\widehat{}`$, $`G(3)^\widehat{}`$. We shall exclude from consideration the well understood case of $`B(0,n)^\widehat{}`$ (see \[K2\] and Section 9.5). In all cases except for $`A(n,n)^\widehat{}`$ these are the affine superalgebras $`\widehat{𝔤}`$ defined by (0.1), (0.2) for $`𝔤=A(m,n)(mn),B(m,n),C(n),D(m,n),D(2,1;a),F(4)`$ and $`G(3)`$ respectively (see \[K1\] for a construction of the simple finite-dimensional Lie superalgebras $`𝔤`$). In the $`A(n,n)^\widehat{}`$ case it is more convenient to take $`𝔤=s\mathrm{}(n+1|n+1)`$ in order not to lose the most interesting modules. The Lie superalgebra $`𝔤`$ carries a unique, up to a constant factor, non-zero invariant bilinear form $`(.|.)`$. This form extends to $`\widehat{𝔤}`$ by formula (0.3)and it is normalized by the values of $`(\alpha _i|\alpha _i)`$, given in Table 1 (see below).
It is convenient to depict Cartan matrices of affine (super)algebras by (generalized) Dynkin diagrams (cf. \[K1\]). We shall assume that the diagonal entries of a Cartan matrix $`\widehat{A}`$ are always $`2`$ or $`0`$ (one can achieve this by rescaling simple coroots). The Dynkin diagram of $`\widehat{A}`$ is a graph whose nodes label the index set $`\widehat{I}=\{0,1,2,\mathrm{}\}`$ and are of the form $`0`$, $``$ or corresponding to cases $`a_{ii}=2`$, $`i\widehat{I}_1`$; $`a_{ii}=0`$ (then $`i\widehat{I}_1`$); and $`a_{ii}=2`$, $`i\widehat{I}_1`$, respectively. These nodes are called white, grey and black respectively, so that $`\widehat{I}_1`$ consists of non-white nodes. We let $`I=\widehat{I}\backslash \{0\}`$, $`I_1=\widehat{I}_1\backslash \{0\}`$. As usual, $`I`$ labels simple roots $`\alpha _1,\alpha _2,\mathrm{}`$ of $`𝔤`$, $`I_1`$ labels odd simple roots of $`𝔤`$, and $`\alpha _0=\delta \theta `$, where $`\theta `$ is the highest root of $`𝔤`$. In the cases $`a_{ij}=a_{ji}=0,ij`$, the $`i`$$`^{\text{th}}`$ and $`j`$$`^{\text{th}}`$ nodes are not connected. In the cases $`a_{ii}=a_{jj}=2`$, $`ij`$, the nodes are, as usual, connected by $`|a_{ij}a_{ji}|`$ edges with an arrow pointing to $`j`$$`^{\text{th}}`$ node if $`|a_{ji}|>1`$. In the remaining cases the nodes are joined as follows:
$$=(\begin{array}{cc}\hfill 0& \hfill a\\ \hfill 1& \hfill 2\end{array}),=\left(\begin{array}{cc}\hfill 0& a\\ \hfill 2& 2\end{array}\right)=\text{},=(\begin{array}{cc}0& a\\ b& 0\end{array}),\text{}=\left(\begin{array}{cc}\hfill 2& \hfill 1\\ \hfill 2& \hfill 2\end{array}\right).$$
In Table 1 below we list the Dynkin diagrams of the symmetrizable Cartan matrices of the affine Lie superalgebras $`\widehat{𝔤}`$ under consideration. The labels against the nodes $`i`$ are $`(\alpha _i|\alpha _i)`$, and the labels against the edges connecting $`i`$ and $`j`$ are $`(\alpha _i|\alpha _j)`$. Recall that $`\alpha _i^{}=2\alpha _i/(\alpha _i|\alpha _i)`$ if $`a_{ii}0`$; we let $`\alpha _i^{}=\alpha _i`$ if $`a_{ii}=0`$. We also give the coefficients of the decomposition of the root $`\delta `$ in terms of simple roots. The nodes are numbered by $`\widehat{I}=\{0,1,\mathrm{}\}`$ in increasing order from left to right, except when it is impossible to do, in which case nodes are numbered by the subscripts of their labels.
###### Remark 6.1.
Recall the definition of the orthosymplectic Lie subalgebra $`osp(M|N)`$ \[K1\]. Let $`V=V_{\overline{1}}V_{\overline{0}}`$ be a superspace, where $`dimV_{\overline{1}}=N,dimV_{\overline{0}}=M`$, and let $`(.|.)`$ be a non-degenerate bilinear form on $`V`$ such that $`(V_{\overline{0}}|V_{\overline{1}})=0`$, the restriction of $`(.|.)`$ to $`V_{\overline{1}}`$ is symmetric and to $`V_{\overline{0}}`$ is skewsymmetric, so that $`N=2n`$ is even; let $`m=[M/2]`$. Then $`(\alpha =0,1)`$:
$`osp(M|N)_{\overline{\alpha }}=\{ag\mathrm{}(M|N)_{\overline{\alpha }}|(a(x)|y)+(1)^{\alpha p(x)}(x|a(y))=0,x,yV\}.`$
For the definition in a matrix form, consider the following $`(M+N)\times (M+N)`$ matrices:
$$C=\left(\begin{array}{cc}C_1\hfill & 0\hfill \\ 0\hfill & C_2\hfill \end{array}\right),F=\left(\begin{array}{cc}I_M\hfill & 0\hfill \\ 0\hfill & I_N\hfill \end{array}\right),$$
where $`C_1`$ (resp. $`C_2`$) is a $`M\times M`$ (resp. $`N\times N`$) symmetric (resp. skewsymmetric) matrix. Then
$$osp(M|N)_{\overline{\alpha }}=\{ag\mathrm{}(M|N)_{\overline{\alpha }}|F^\alpha a^{}C+Ca=0\},\alpha =0,1.$$
Recall that $`B(m,n)=osp(2m+1|2n)`$, $`C(n)=osp(2|2n)`$ and $`D(m,n)=osp(2m|2n)`$. The invariant bilinear form on $`osp(M|N)`$ that is used in Table 1 and throughout the paper is
$$(a|b)=\frac{1}{2}strab.$$
The even parts $`𝔤_{\overline{0}}`$ of the Lie superalgebras $`𝔤`$ are listed in Table 2. In the case of $`D(2,1;a)`$, the subalgebra $`D_2`$ corresponds to $`\alpha _2`$ and $`\alpha _3`$ (see Table 1). We denote by $`𝔤_{\overline{0}}^{}`$ (resp. $`𝔤_{\overline{0}}^{\prime \prime }`$) the first (resp. second) non-zero summand of $`𝔤_{\overline{0}}`$ in the decomposition of Table 2. Note that the invariant bilinear form $`(.|.)`$ (which can be read off from Table 1) is normalized in such a way that it is positive definite on $`𝔤_{\overline{0}}^{}`$ and negative definite on $`𝔤_{\overline{0}}^{\prime \prime }`$ (except that for $`D(2,1;a)`$ we should assume that $`a_>`$), and the maximal square length of a root is $`2`$, except for the cases $`B(1,n)`$ when it is $`1`$ and $`D(2,1;a)`$ when it is max $`(2,2a)`$. If the Killing form on $`𝔤`$ is non-degenerate, then the form $`(.|.)`$ is a positive (resp. negative) multiple of the Killing form, in the cases $`𝔤s\mathrm{}(m|n)`$ with $`m>n`$, $`osp(m|n)`$ with $`m>n+2`$, $`F(4)`$ and $`G(3)`$ (resp. $`s\mathrm{}(m|n)`$ with $`m<n`$ and $`osp(m|n)`$ with $`m<n+2`$).
An irreducible highest weight module over $`\widehat{𝔤}`$ is called *principal* (resp. *subprincipal*) *integrable* module if it is integrable with respect to $`\widehat{𝔤_{\overline{0}}^{}}`$ (resp. $`\widehat{𝔤_{\overline{0}}^{\prime \prime }}`$) and locally finite with respect to $`𝔤`$ (cf. Definition 0.1). As we shall see, the non-trivial principal (resp. subprincipal) highest weight modules have positive (resp. negative) level, except for the cases $`\widehat{𝔤}=A(0,n)^\widehat{}`$ and $`C(n)^\widehat{}`$ (resp. $`A(n,0)^\widehat{}`$). It is easy to see that in these cases the only conditions of integrability are $`k_i_+`$ if $`iI\backslash I_1`$; we shall exclude these cases from further considerations.
Let $`\theta ^{}`$ be the highest root of $`𝔤_{\overline{0}}^{}`$; in cases $`D(2,n)`$ and $`D(2,1;a)`$, which are the only cases when $`𝔤_{\overline{0}}^{}`$ is not simple, we have: $`𝔤_{\overline{0}}^{}=A_1+A_1`$, and the highest roots are $`\theta _+^{}=\alpha _{n+1}`$ and $`\theta _{}^{}=\alpha _{n+2}`$ (where $`n=1`$ for $`D(2,1;a)`$). The root $`\theta ^{}`$ (resp. root $`\theta _\pm ^{}`$) gives rise to a simple root $`\alpha _0^{}=\delta \theta ^{}`$ (resp. simple roots $`\alpha _{0\pm }^{}=\delta \theta _\pm ^{}`$) of $`𝔤_{\overline{0}}^{}`$. The corresponding coroot is $`\alpha _0^{}=2\alpha _0^{}/(\alpha _0^{}|\alpha _0^{})`$.
In all cases except for $`A(m,n)^\widehat{}`$ and $`C(n)^\widehat{}`$ there is a (unique) simple root $`\alpha ^{\prime \prime }`$ of $`𝔤_0^{\prime \prime }`$ (which is a simple Lie algebra), which is not a simple root of $`𝔤`$ \[K1\]. As we have seen in §2, the principal integrability in the case of $`A(m,n)^\widehat{}`$ follows from local nilpotency of the root vector $`f^{}`$ attached to the root $`\alpha _0^{}`$. In all other cases one has to check in addition the local nilpotency of the root vector $`f^{\prime \prime }`$ attached to the root $`\alpha ^{\prime \prime }`$ (in order to ensure the local finiteness with respect to $`𝔤_{\overline{0}}^{\prime \prime }`$; that with respect to $`𝔤_{\overline{0}}^{}`$ follows automatically from the integrability with respect to $`\widehat{𝔤_{\overline{0}}^{}}`$). For that reason, as we have seen in §1, it is important to introduce the following numbers, where $`\rho `$ is a Weyl vector for $`\widehat{𝔤}`$:
$$b^{}=\rho ,\alpha _0^{},b_\pm ^{}=\rho ,\alpha _{0^\pm }^{},b^{\prime \prime }=\rho ,\alpha ^{\prime \prime }.$$
The values of the numbers $`b^{}`$ and $`b_\pm ^{}`$ (resp. $`b^{\prime \prime }`$) are given in Table 3 (resp. 4), the first line for $`D(2,\mathrm{})^\widehat{}`$ in Table 4 being for $`b_+^{}`$ and the second for $`b_{}^{}`$. Table 3 contains also the formula for the *partial levels*
$$k^{}=\mathrm{\Lambda },\alpha _0^{},k_\pm ^{}=\mathrm{\Lambda },\alpha _{0\pm }^{},$$
and the level $`k=\mathrm{\Lambda },K`$ of a weight $`\mathrm{\Lambda }`$ in terms of its labels $`k_i`$ and $`k^{},k_\pm ^{}`$. Table 4 contains also a formula for
$$k^{\prime \prime }=\mathrm{\Lambda },\alpha ^{\prime \prime }$$
and the level $`k`$ in terms of the $`k_i`$ and $`k^{\prime \prime }`$.
###### Theorem 6.1.
For an affine superalgebra from Table 1 (recall that $`A(0,n)^\mathrm{^}`$, $`B(0,n)^\mathrm{^}`$ and $`C(n)^\mathrm{^}`$ are excluded) the labels $`\{k_i\}_{i=\widehat{I}}`$ of the highest weight of a principal integrable irreducible highest weight module $`L(\mathrm{\Lambda })`$ are characterized by the following four series of conditions:
1. $`k_i_+`$ if $`iI\backslash I_1`$,
2. $`k^{}`$ (resp. $`k_\pm ^{}`$ in case $`D(2,\mathrm{})^\widehat{})_+`$, $`k^{\prime \prime }_+`$ (see Tables 3 and 4),
3. if $`k^{}`$ (resp. one of the $`k_\pm ^{}`$) $`b^{}`$ (see Table 4), there are the supplementary conditions:
* $`A(m,n)^\widehat{},m1`$: there exists $`s_+`$, $`sk^{}`$, such that:
$$k_{m+1}=k_{m+2}+\mathrm{}+k_{m+1+s}+s,k_{m+s+2}=\mathrm{}=k_{m+s+1+nk^{}}=0,$$
* $`B(m,n)^\widehat{}`$, $`m2`$, and $`D(m,n)^\widehat{},m3`$: one of the four possibilities hold:
1. there exist $`r,s_+`$, $`r<s`$, such that
$$k^{}=r+s,k_j=0\text{ for }r+1jm+n\text{ and }js,k_s=1,$$
2. there exist $`r,s_+`$, $`rs`$, such that
$$k^{}=r+s,k_j=0\text{ for }r+1jm+n,k_r0,$$
3. there exist $`r_+`$ such that
$$k^{}=n+r,k_j=0\text{ for }r+1jn1,k_n0,k_n+k_{n+1}+1=0,$$
4. there exist $`r_+`$ such that
$$k^{}n+r,k_j=0\text{ for }jr+1,k_r0.$$
* $`B(1,n)^\widehat{}`$: the same as for $`B(m,n)^\widehat{}`$ with $`m>1`$ with the following changes: $`k^{}`$ is replaced by $`\frac{1}{2}k^{}`$ everywhere and $`k_{n+1}`$ is replaced by $`\frac{1}{2}k_{n+1}`$ in (iii),
* $`D(2,n)^\widehat{}`$: the same as for $`D(m,n)^\widehat{}`$ with $`m>2`$ with the following additions:
$$k_+^{}=k_{}^{},k_n+k_{n+2}+1=0\text{ in (iii)},$$
* $`D(2,1;a)^\widehat{}`$: one of the four possibilities holds:
1. $`k_+^{}k_{}^{}=0`$, then all $`k_i=0`$,
2. $`a_{>0}`$, $`a,a^1`$ and one of $`k_\pm ^{}`$ equals $`1`$, then $`k_+^{}=k_{}^{}=1`$
and one has
+ $`k_0=(a+1)^1r1,k_1=r,k_2=r1,k_3=a^1r1`$ for some $`ra`$,
3. $`a^1`$ and $`k_+^{}=1`$, then either (\*) holds or $`k_0=(a+1)^1,k_1=k_2=k_3=0`$,
4. $`a`$ and $`k_{}^{}=1`$, then either (\*) holds or $`k_0=a(a+1)^1,k_1=k_2=k_3=0`$,
* $`F(4)^\widehat{}`$: one of the two possibilities holds:
1. $`k^{}=0`$, then all $`k_i=0`$,
2. $`k^{}=1`$, then $`k_0=\frac{2}{3},k_1=k_2=k_3=k_4=0`$,
* $`G(3)^\widehat{}`$: one of the two possibilities holds:
1. $`k^{}=0`$, then all $`k_i=0`$,
2. $`k^{}=1`$, then $`k_0=\frac{3}{4},k_1=k_2=k_3=0`$,
4. if $`k^{\prime \prime }b^{\prime \prime }`$ (see Table 4), there are the supplementary conditions:
* $`B(m,n)^\widehat{}`$: $`k_j=0`$ for all $`jn+k^{\prime \prime }+1`$,
* $`C(n)^\widehat{}`$: $`k_0=k_1=0`$,
* $`D(m,n)^\widehat{}`$: one of the two possibilities holds:
1. $`k^{\prime \prime }m2`$, then $`k_j=0`$ for all $`jn+k^{\prime \prime }+1`$,
2. $`k^{\prime \prime }=m1`$, then $`k_{m+n1}=k_{m+n}`$,
$`D(2,1;a)^\widehat{}`$: one of the two possibilities holds:
1. $`k^{\prime \prime }=0`$, then $`k_1=k_2=k_3=0`$,
2. $`k^{\prime \prime }=1`$, then $`a`$ and $`k_2+1=|a|(k_3+1)`$,
* $`F(4)^\widehat{}`$: one of the three possibilities holds:
1. $`k^{\prime \prime }=0`$, then $`k_1=k_2=k_3=k_4=0`$,
2. $`k^{\prime \prime }=2`$, then $`k_2=k_4=0`$,
3. $`k^{\prime \prime }=3`$, then $`k_2=2k_4+1`$,
$`G(3)^\widehat{}`$: one of the two possibilities holds:
1. $`k^{\prime \prime }=0`$, then $`k_1=k_2=k_3=0`$,
2. $`k^{\prime \prime }=2`$, then $`k_2=0`$.
###### Proof.
In the case $`\widehat{𝔤}=A(m,n)^\widehat{}`$, the theorem follows from Theorem 2.1. In general, the proof is based on similar arguments. Below we shall give details in the case $`\widehat{𝔤}=B(1,n)^\widehat{}`$; in the rest of the cases arguments are the same.
The even part of $`B(1,n)`$ is $`A_1+C_n`$ and its simple roots are $`\alpha _{n+1}`$ for $`A_1`$ and $`\{\alpha _1,\alpha _2,\mathrm{},\alpha _{n1}`$,$`\alpha ^{\prime \prime }=2(\alpha _n+\alpha _{n+1})\}`$ for $`C_n`$. The simple roots of $`\widehat{A}_1`$ are $`\{\alpha _0^{}=\delta \alpha _{n+1},\alpha _{n+1}\}`$. Due to Lemma 1.5a, the local finiteness (resp. integrability) with respect to $`C_n`$ (resp. $`\widehat{A}_1`$) implies that $`k_1,\mathrm{},k_{n1},k^{\prime \prime }_+`$ (resp. $`k_{n+1},k^{}_+`$). Hence, conditions (1) and (2) are necessary. Furthermore, it follows from Proposition 1.2, that in the cases $`k^{}>b^{}=4n1`$ (resp. $`k^{\prime \prime }>b^{\prime \prime }=\frac{1}{2}`$) the element $`f^{}`$ (resp. $`f^{\prime \prime }`$) is locally nilpotent. It remains to show that in the case of inequality
(6.1)
$$k^{}4n1$$
the element $`f^{}`$ is locally nilpotent iff condition (3) holds, and in the case $`k^{\prime \prime }=0`$, $`f^{\prime \prime }`$ is locally nilpotent iff (4) holds. We shall concentrate on the first claim, the second being easier (cf. also \[K1\]).
Introduce the following isotropic roots:
$$\beta _j=\underset{i=nj}{\overset{n}{}}\alpha _i(j=0,\mathrm{},n),\beta _{n+j}=\beta _n+\underset{i=1}{\overset{j}{}}\alpha _i(j=1,\mathrm{},n1).$$
We have: $`\beta _j^{}=\beta _j`$ for all $`j`$, and
(6.4) $`(\beta _i|\alpha _{n+1})=1\text{ for all }i,(\beta _i|\beta _{i+1})=\{\begin{array}{ccc}1& \text{if}& 0in2,\hfill \\ 2& \text{if}& i=n1.\hfill \end{array}`$
Let $`\mathrm{\Pi }^{(0)}=\mathrm{\Pi }`$, then $`\beta _0\mathrm{\Pi }^{(0)}`$ and we let $`\mathrm{\Pi }^{(1)}=r_{\beta _0}\mathrm{\Pi }^{(0)}`$. Similarly $`\beta _1\mathrm{\Pi }^{(1)}`$ and we let $`\mathrm{\Pi }^{(2)}=r_{\beta _1}\mathrm{\Pi }^{(1)},\mathrm{},\mathrm{\Pi }^{(2n)}=r_{\beta _{2n1}}\mathrm{\Pi }^{(2n1)}`$. We have:
(6.5)
$$\alpha _0^{}\mathrm{\Pi }^{(2n)},\alpha _0^{}=2\alpha _0^{}.$$
Let $`\mathrm{\Lambda }^{(j)}`$ denote the highest weight of $`L(\mathrm{\Lambda })`$ with respect to $`𝔫_+^{(j)}`$. It can be computed by making use of Lemma 1.4. Introduce the following numbers:
$$u_j=\mathrm{\Lambda }^{(j)},\alpha _0^{}(j=0,\mathrm{},2n),t_j=\mathrm{\Lambda }^{(j)},\beta _j(j=0,\mathrm{},2n1).$$
Using (6.4), $`\mathrm{\Lambda },\alpha _0^{}=k^{}`$ and Lemma 1.4, we get the following recurrent formula for the $`u_j`$’s:
(6.6)
$$u_0=k^{},u_{j+1}=u_j2(\text{ resp. }=u_j)\text{ if }t_j0(\text{ resp. }=0).$$
In view of Lemma 1.5, the local nilpotency of $`f^{}`$ follows from $`u_{2n}_+`$. This, clearly, holds if $`k^{}4n`$ (by (6.6)), which again shows that in this case there are no supplementary conditions.
¿From now on we may assume that $`k^{}4n1`$. We may also assume that conditions (1), (2) and (4) hold. We shall derive a recurrent formula for the $`t_i`$. Using that, by Lemma 1.4, $`\mathrm{\Lambda }^{(i+1)}=\mathrm{\Lambda }^{(i)}\beta _i`$ (resp. $`=\mathrm{\Lambda }^{(i)}`$) if $`t_i0`$ (resp. $`=0`$) and that $`\alpha _{ni1}\mathrm{\Pi }^{(i)}`$, we obtain:
(6.7a)
$$t_{i+1}=\{\begin{array}{ccc}t_ik_{ni1}1\hfill & \text{if}\hfill & t_i0\hfill \\ t_ik_{ni1}\hfill & \text{if}\hfill & t_i=0\hfill \end{array}(0in2)$$
(6.7b)
$$t_n=t_{n1}2k_02(\text{resp. }=t_{n1}2k_0)\text{ if }t_{n1}0(\text{resp. }=0).$$
For $`t_j,jn`$, the recurrent formula involves numbers
$$s_i=(\mathrm{\Lambda }^{(n)}|\alpha _i)\text{ for }1in1,s_n=(\mathrm{\Lambda }^{(n)}|\alpha _n+\alpha _{n+1}),$$
which, using the above arguments, can be expressed in terms of the labels of $`\mathrm{\Lambda }`$ as follows:
(6.8a)
$$s_i=\{\begin{array}{ccc}k_i\hfill & \text{if}\hfill & t_{ni1}t_{ni}0\text{ or }t_{ni1}=t_{ni}=0(1in1),\hfill \\ k_i1\hfill & \text{if}\hfill & t_{ni1}0,t_{ni}=0,\hfill \\ k_i+1\hfill & \text{if}\hfill & t_{ni1}=0,t_{ni}0,\hfill \end{array}$$
(6.8b)
$$s_n=k^{\prime \prime }+1(\text{resp. }k^{\prime \prime }),\text{ where }k^{\prime \prime }=k_n\frac{1}{2}k_{n+1},\text{ if }k_n0(\text{resp. }=0).$$
Then we have $`(0in1)`$:
(6.9a)
$$t_{n+i+1}=t_{n+i}+s_{i+1}1(\text{resp. }t_{n+i}+s_{i+1})\text{ if }t_{n+i}0(\text{resp. }=0),$$
where we let
(6.9b)
$$t_{2n}=\frac{1}{2}u_{2n}.$$
Note that $`t_0=k_n0`$ since
$$k^{\prime \prime }=k_n\frac{1}{2}k_{n+1}_+\text{ and }k_{n+1}_+.$$
Since $`k_i_+`$ for $`i0,n`$, formulae (6.7) imply
(6.10)
$$0t_0t_1\mathrm{}t_{n1}.$$
Furthermore, we have
(6.11)
$$t_nt_{n+1}\mathrm{}t_{2n}=\frac{1}{2}u_{2n}.$$
In order to show this, it suffices to prove that $`s_i0`$, $`1in`$. But, due to (6.8a), $`s_i>0`$ can take place for $`1in1`$ only when $`t_{ni1}=0,t_{ni}0,k_i=0`$, which is impossible by (6.7a). Also, $`s_n=k^{\prime \prime }+1(`$resp. $`k^{\prime \prime })`$ if $`k_n0`$ (resp. $`=0`$) cannot be positive since in this case $`k^{\prime \prime }=0`$, which implies that $`k_n=0`$ (see supplementary conditions (4)).
Suppose now that $`f^{}`$ is locally nilpotent. Then $`u_{2n}0`$ (by Lemma 1.5), hence we have from (6.6):
(6.12)
$$t_j=0\text{ for some }0j2n1.$$
Due to (6.12), (6.10) and (6.11), we have the following three possibilities for some $`0i_0n1`$ and $`nj_02n1`$:
$`\begin{array}{cc}(\alpha )& t_0=\mathrm{}=t_{i_0}=0,t_{i_0+1}0,\mathrm{},t_{2n1}0,\hfill \\ (\beta )& t_00,\mathrm{},t_{j_01}0,t_{j_0}=\mathrm{}=t_{2n}=0,\hfill \\ (\gamma )& t_0=\mathrm{}=t_{i_0}=0,t_{i_0+1}0,\mathrm{},t_{j_01}0,t_{j_0}=\mathrm{}=t_{2n}=0.\hfill \end{array}`$
The possibilities (i), (ii), (iii) and (iv) of the supplementary conditions (3) correspond respectively to the following cases:
1. $`(\gamma )`$ when $`i_0+j_0<2n1`$, where we put $`i_0=ns1,j_0=n+r`$,
2. $`(\gamma )`$ when $`i_0+j_02n1`$, where we put $`i_0=nr1,j_0=n+s`$,
3. $`(\beta )`$, where we put $`j_0=n+r`$,
4. $`(\alpha )`$, where we put $`i_0=nr1`$.
We consider in detail only case (ii), the treatment of all other cases being similar. We have:
$$t_0=\mathrm{}=t_{nr1}=0,t_{nr}0,\mathrm{},t_{n+s1}0,t_{n+s}=\mathrm{}=t_{2n1}=0$$
for some integer $`s`$ such that $`0rsn1`$. Hence, by (6.7a) we have:
(6.14)
$$k_r0,k_{r+1}=\mathrm{}=k_n=0,$$
and, by (6.8a), we have:
$$s_i=k_i\text{ if }ir,\mathrm{\hspace{0.17em}\hspace{0.17em}1}in1,s_r=k_r+1,s_n=0.$$
The recurrent formulas (6.7) and (6.9) can be now rewritten respectively as follows:
$`\begin{array}{cccc}& t_{nr}& =& k_r,\hfill \\ & t_{nr+1}& =& t_{nr}k_{r1}1,\hfill \\ & \mathrm{}& & \\ & t_{n1}& =& t_{n2}k_11,\hfill \\ & t_n& =& t_{n1}2k_02,\hfill \\ & t_{n+1}& =& t_nk_11,\hfill \\ & \mathrm{}& & \\ & t_{n+r1}& =& t_{n+r2}k_{r1}1,\hfill \\ & t_{n+r}& =& t_{n+r1}k_r,\hfill \\ & t_{n+r+1}& =& t_{n+r}k_{r+1}1,\hfill \\ & \mathrm{}& & \\ 0=\hfill & t_{n+s}& =& t_{n+s1}k_s1+\delta _{r,s}.\hfill \end{array}`$
Summing up these equalities, we get: $`2{\displaystyle \underset{i=0}{\overset{r}{}}}k_i=(r+s)`$, which, in view of (6.14), implies $`k^{}=2(r+s)`$.
Suppose now that conditions (1)–(4) hold. We have to show that $`u_{2n}0`$. As before, we may assume that $`k^{}4n1`$, hence, due to (3), (6.12) holds. Since (6.10) and (6.11) hold, we again have only the possibilities $`(\alpha )`$, $`(\beta )`$ and $`(\gamma )`$. In cases $`(\beta )`$ and $`(\gamma )`$, $`u_{2n}=2t_{2n}=0`$, hence only case $`(\alpha )`$ remains. This case corresponds to (3)iv when we have: $`t_0=k_n=0`$,$`t_1=t_0k_{n1}=0,\mathrm{},t_{nr1}=t_{nr2}k_{r+1}=0`$. Hence $`v:=\mathrm{\#}\{0j2n1|t_j0\}n+r`$ and $`u_{2n}=u_02vk^{}2(n+r)0`$ (we have used here (6.6) and (3)iv).
###### Theorem 6.2.
For an affine superalgebra $`\widehat{𝔤}`$ from Table 4 the labels $`\{k_i\}_{i\widehat{I}}`$ of the highest weight of a subprincipal integrable irreducible highest weight module are characterized by the following conditions:
1. $`k_i_+`$ if $`i\widehat{I}\backslash I_1`$, $`k^{\prime \prime }_+`$,
2. if $`k^{\prime \prime }b^{\prime \prime }`$ (see Table 4), there are supplementary conditions described by (4) of Theorem 6.1.
###### Proof.
The only simple root of $`\widehat{𝔤_{\overline{0}}^{\prime \prime }}`$ which is not simple for $`\widehat{𝔤}`$ is $`\alpha ^{\prime \prime }`$. Hence the proof of Theorem 6.1 proves Theorem 6.2 as well. ∎
###### Remark 6.2.
It follows from Theorem 6.1 that the level $`k`$ of a principal integrable $`\widehat{𝔤}`$-module $`L(\mathrm{\Lambda })`$ is a non-negative number which is an integer in all cases, except for $`B(1,n)^\widehat{}`$, when it is a half-integer; moreover, if $`k=0`$, then all labels of $`\mathrm{\Lambda }`$ are $`0`$, hence $`L(\mathrm{\Lambda })`$ is $`1`$-dimensional; also, if $`k>0`$, then $`k1`$.
###### Remark 6.3.
It is easy to see that, when restricted to the derived subalgebra $`[\widehat{𝔤},\widehat{𝔤}]`$ of $`\widehat{𝔤}`$ the module $`L(\mathrm{\Lambda })`$ remains irreducible. Two $`\widehat{𝔤}`$-modules are called *essentially equivalent* if they are equivalent as $`[\widehat{𝔤},\widehat{𝔤}]`$-modules. For example, the modules $`L(\mathrm{\Lambda })`$ and $`L(\mathrm{\Lambda }+a\delta )`$ are essentially equivalent for any $`a`$. Theorem 6.1 gives the following complete list of principal integrable modules of level $`1`$ up to essential equivalence:
1. $`A(m,n)^\widehat{},m1`$: $`\mathrm{\Lambda }_s(1sm)`$, $`(a+1)\mathrm{\Lambda }_{m+1}+a\mathrm{\Lambda }_{m+2}(a_+)`$, and $`(a+1)\mathrm{\Lambda }_0+a\mathrm{\Lambda }_{m+n+1}(a_+)`$,
2. $`B(m,1)^\widehat{}`$ and $`D(m,1)^\widehat{}`$: $`\frac{1}{2}\mathrm{\Lambda }_0`$ and $`\frac{3}{2}\mathrm{\Lambda }_0\mathrm{\Lambda }_1`$,
3. $`B(m,n)^\widehat{}`$ and $`D(m,n)^\widehat{},n2`$: $`\frac{1}{2}\mathrm{\Lambda }_0`$ and $`\frac{3}{2}\mathrm{\Lambda }_0+\mathrm{\Lambda }_1`$,
4. $`D(2,1;a)^\widehat{},a^1`$: $`(a+1)^1\mathrm{\Lambda }_0`$ and $`\frac{a+2}{a+1}\mathrm{\Lambda }_0\mathrm{\Lambda }_1+\frac{1a}{a}\mathrm{\Lambda }_3`$ $`a^1`$,
5. $`F(4)^\widehat{}`$: $`\frac{2}{3}\mathrm{\Lambda }_0`$,
6. $`G(3)^\widehat{}`$: $`\frac{3}{4}\mathrm{\Lambda }_0`$.
One can show (cf. Remark 2.5) that in all cases, all weights are conjugate to each other by odd reflections. Thus, for each of the affine superalgebras $`A(m,n)^\widehat{}(m1)`$, $`B(m,n)^\widehat{}`$, $`D(m,n)^\widehat{}`$, $`D(2,1;a)^\widehat{}(a_>)`$, $`F(4)^\widehat{}`$ and $`G(3)^\widehat{}`$ all, up to essential equivalence, principal integrable modules of level $`1`$, can be obtained from one of them by making different choices of the set of positive roots. Note also that in all cases the “basic” module $`L(u\mathrm{\Lambda }_0)`$, where $`u`$ is such that $`u\mathrm{\Lambda }_0`$ has level $`1`$, is a principal integrable module.
###### Remark 6.4.
Using the symmetry of $`A(m,n)^\widehat{}`$ which exchanges the subalgebras $`\widehat{A}_m`$ and $`\widehat{A}_n`$, one gets the classification of the subprincipal integrable modules $`L(\mathrm{\Lambda })`$ for this affine subalgebra:
$$k_i_+\text{ for }i\widehat{I}\backslash \widehat{I}_1,k^{\prime \prime }:=\underset{i=0}{\overset{m+1}{}}k_i_+,$$
and there exists $`s_+`$, $`sk^{\prime \prime }`$, such that
$$k_0+k_1+\mathrm{}+k_s+s=0\text{ and }k_{s+1}=\mathrm{}=k_{s+mk^{\prime \prime }+1}=0.$$
One has: $`k=(k^{\prime \prime }+{\displaystyle \underset{i=n+2}{\overset{m+n+1}{}}}k_i)`$.
###### Remark 6.5.
All principal integrable highest weights of level $`2`$ (up to essential equivalence) for $`B(1,2)^\widehat{}`$ are $`(1+a)\mathrm{\Lambda }_0+a\mathrm{\Lambda }_1`$, where $`a_+`$. Thus, in sharp contrast to the level $`1`$ case, there are infinitely many essentially inequivalent principal integrable highest weight modules of level $`2`$.
###### Remark 6.6.
It follows from Theorem 6.2 and Remark 6.3 that the level $`k`$ of a subprincipal integrable $`\widehat{𝔤}`$-module $`L(\mathrm{\Lambda })`$ is a non-positive number, provided that $`a>1`$ for $`D(2,1;a)`$; moreover, $`dimL(\mathrm{\Lambda })=1`$ if $`k=0`$. Thus, in view of Theorem 6.1, the only $`L(\mathrm{\Lambda })`$ over $`\widehat{𝔤}A(m,0)^\widehat{}`$, $`A(0,n)^\widehat{}`$ or $`C(n)^\widehat{}`$, which are integrable over $`\widehat{𝔤_{\overline{0}}}`$ are $`1`$-dimensional.
###### Remark 6.7.
Using the same arguments, one can show that the non-symmetrizable “twisted” affine superalgebra of type $`Q`$ (which is the universal central extension of the Lie superalgebra $`{\displaystyle \underset{n}{}}(Q(n)_{\overline{0}}t^{2n}+Q(n)_{\overline{1}}t^{2n+1})`$), with the Cartan matrix
$$\left(\begin{array}{ccccc}0& 1& 0& \mathrm{}& 1\\ 1& & & & \\ 0& & & A_n& \\ \mathrm{}& & & & \\ 0& & & & \\ 1& & & & \end{array}\right)$$
has no non-trivial integrable (with respect to its even part) highest weight modules.
###### Remark 6.8.
Consider the $`/2`$-gradation of $`F(4)`$ of type $`(0,0,0,1,0)`$ and that of $`G(3)`$ of type $`(0,0,0,1)`$, cf. Table 1 and \[K3\]. The $`0`$$`^{\text{th}}`$ piece in the first (resp. second) case is isomorphic to $`D(2,1;1/2)A_1`$ (resp. to $`D(2,1;1/3)`$), and its representation on the $`1`$$`^{\text{st}}`$ piece is the module $`^{10}^2`$ (resp. $`^{14}`$), where $`^{10}`$ (resp. $`^{14}`$) is the lowest-dimensional non-trivial module over $`D(2,1;1/2)`$ (resp. $`D(2,1;1/3)`$). This reduces to some extent the construction of the principal integrable level $`1`$ module over $`F(4)^\widehat{}`$ and $`G(3)^\widehat{}`$ to that of $`D(2,1;a)^\widehat{}`$. The free field construction of the principal integrable level $`1`$ modules over $`osp(m,n)^\widehat{}`$ (covering the $`BCD`$ cases) will be given in Section 7.
## 7 Free field realization of level $`1`$ integrable modules over $`osp(M|N)^\widehat{}`$.
Let $`V`$ be the superspace and let $`(.|.)`$ be the bilinear form on $`V`$ considered in Remark 6.1. Recall an equivalent definition of $`osp(M|N)`$ via the Clifford superalgebra:
$$C\mathrm{}V=T(V)/[x,y](x|y)1|x,yV.$$
The Lie superalgebra $`osp(M|N)`$ is identified with the $``$-span of all quadratic elements of $`C\mathrm{}V`$ of the form:
$$:\alpha \beta :\alpha \beta +(1)^{p(\alpha )p(\beta )}\beta \alpha ,\text{ where }\alpha ,\beta V.$$
Such an element is identified with an operator from $`osp(M|N)`$ by the formula:
(7.1)
$$(:\alpha \beta :)v=[:\alpha \beta :,v],vV.$$
Denote by $`\mathrm{\Phi }_V`$ the vertex algebra generated by pairwise local fields $`\gamma (z)`$, where $`\gamma V_{\overline{0}}V_{\overline{1}}`$ and $`\gamma (z)`$ is even (resp. odd) if $`\gamma V_{\overline{0}}`$ (resp. $`V_{\overline{1}}`$), subject to the following OPE:
$$\gamma (z)\gamma ^{}(\omega )\frac{(\gamma |\gamma ^{})}{zw}.$$
This is called the vertex algebra of free superfermions in \[K4\].
###### Remark 7.1.
The vertex algebra $`F`$ considered in § 3 is isomorphic to $`\mathrm{\Phi }_V`$, where $`dimV_{\overline{0}}=2n,dimV_{\overline{1}}=2m`$ and the bilinear form is given by:
$`(\phi ^i|\phi ^j)`$ $`=`$ $`(\phi ^i|\phi ^j)=\delta _{ij}(i,j=1,\mathrm{},n),`$
$`(\psi ^i|\psi ^j)`$ $`=`$ $`(\psi ^i|\psi ^j)=\delta _{ij}(i,j=1,\mathrm{},m),\text{ all other inner products }=0.`$
Furthermore, in the case when $`dimV_{\overline{1}}=2m+1`$ the vertex algebra $`\mathrm{\Phi }_V`$ is isomorphic to $`F\mathrm{\Phi }`$, where $`\mathrm{\Phi }`$ is a vertex algebra generated by one odd field $`\psi (z)`$ with the OPE $`\psi (z)\psi (w)\frac{1}{zw}`$. This corresponds to adding an odd vector $`\psi `$ with $`(\psi |\psi )=1`$ orthogonal to all the above basis vectors.
As in § 3, we construct the Virasoro field $`L(z)_jL_jz^{j2}`$ with respect to which all $`\gamma (z)`$ are primary of conformal weight $`1/2`$. Choose a basis $`\phi ^i,\phi ^i(i=1,\mathrm{},n)`$ of $`V_{\overline{0}}`$, and a basis $`\psi ^i,\psi ^i(i=1,\mathrm{},m)`$ and $`\psi `$ if $`M`$ is odd, with inner products described by Remark 7.1. Then $`L(z)`$ is given by formula (3) if $`M`$ is even. In the case $`M`$ is odd, one should add to the expression (3) the term $`\frac{1}{2}:\psi (z)\psi (z):`$. As in § 3, we shall write $`\gamma (z)={\displaystyle \underset{k\frac{1}{2}+}{}}\gamma _kz^{k1/2},\gamma V_{\overline{0}}V_{\overline{1}}`$.
We shall need also the following well-known fact (see e.g. \[K4\], formula (5.1.5)).
###### Lemma 7.1.
Let $`\psi ^+,\psi ^{}V_{\overline{1}}`$ be such that $`(\psi ^\pm |\psi ^\pm )=0,(\psi ^+|\psi ^{})=1`$. Let $`\alpha (z){\displaystyle \underset{n}{}}\alpha _nz^{n1}`$$`=:\psi ^+(z)\psi ^{}(z):`$. Then one has:
$$:\alpha (z)\alpha (z):=:\psi ^+(z)\psi ^{}(z):+:\psi ^{}(z)\psi ^+(z):.$$
Consequently, the fields $`\psi ^\pm (z)`$ are primary of conformal weight $`1/2`$ with respect to the Virasoro field $`\mathrm{}(z){\displaystyle \underset{n}{}}\mathrm{}_nz^{n2}=\frac{1}{2}:\alpha (z)\alpha (z):`$. In particular, we have
(7.2)
$$[\mathrm{}_0,\psi _n^\pm ]=n\psi _n^\pm .$$
Note that
$$\gamma _j|0=0\text{ for }j>0,\gamma V.$$
Hence $`\mathrm{\Phi }_V`$ is obtained by applying polynomials in the $`\gamma _j`$, $`\gamma V,j>0`$, to the vacuum vector $`|0`$. We have the decomposition
(7.3)
$$\mathrm{\Phi }_V=\mathrm{\Phi }_V^+\mathrm{\Phi }_V^{},$$
where $`\mathrm{\Phi }_V^+`$ (resp. $`\mathrm{\Phi }_V^{}`$) is obtained by applying even (resp. odd) degree polynomials in the $`\gamma _j`$ to $`|0`$.
###### Theorem 7.1.
1. Consider the affine superalgebra $`osp(M|N)^\widehat{}`$ and let $`:\alpha \beta :(z)=`$$`{\displaystyle \underset{k}{}}(t^k:\alpha \beta :)z^{k1}`$ for $`:\alpha \beta :osp(M|N)`$. Then the linear map $`\sigma `$ given by $`(\alpha ,\beta V)`$:
$$:\alpha \beta :(z):\alpha (z)\beta (z):,K1,dL_0$$
defines a principal integrable representation of $`osp(M|N)^\widehat{}`$ of level $`1`$ in the space $`\mathrm{\Phi }_V`$ for which $`\mathrm{\Phi }_V^+`$ and $`\mathrm{\Phi }_V^{}`$ are submodules.
2. The $`osp(M|N)^\widehat{}`$-modules $`\mathrm{\Phi }_V^+`$ and $`\mathrm{\Phi }_V^{}`$ are irreducible highest weight modules isomorphic to $`L(\frac{1}{2}\mathrm{\Lambda }_0)`$ and $`L(\frac{1}{2}\mathrm{\Lambda }_0\frac{1}{2}\alpha _0)`$ respectively, provided that $`(M,N)(1,0)`$ or $`(2,0)`$.
###### Proof.
The proof that $`\sigma `$ is a representation is, as usual, a straightforward use of Wick’s formula. The proof of integrability of $`\sigma `$ is the same as in the proof of Theorem 3.1. This establishes (a).
Note that, as before, $`L_0`$ commutes with $`osp(M|N)`$, and the spectrum of $`L_0`$ on $`\mathrm{\Phi }_V^+`$ (resp. $`\mathrm{\Phi }_V^{}`$) is $`_+`$ (resp. $`\frac{1}{2}+_+`$), the lowest eigenvalue eigenspace being $`S^\pm =|0`$ (resp. $`S^{}=\{\gamma _{\frac{1}{2}}|0|\gamma V\}`$), which is the trivial $`1`$-dimensional (resp. the standard) representation of $`osp(M|N)`$. Provided that $`\mathrm{\Phi }_V^\pm `$ are irreducible $`osp(M|N)^\widehat{}`$-modules, (b) follows.
In order to prove irreducibility of $`\mathrm{\Phi }_V^\pm `$, pick elements $`\psi ^+,\psi ^{}V_{\overline{1}}`$ as in Lemma 7.1 and define the field $`\mathrm{}(z)`$ as in that lemma. Let $`\psi V_{\overline{1}}`$ be an element orthogonal to both $`\psi ^+`$ and $`\psi ^{}`$, and consider the field $`\beta (z)=:\psi ^+(z)\psi (z):{\displaystyle \underset{n}{}}\beta _nz^{n1}`$, so that $`\beta _n={\displaystyle \underset{j\frac{1}{2}+}{}}:\psi _j^+\psi _{nj}:`$. Since $`\mathrm{}_0`$ commutes with $`\psi (z)`$, we have by (7.2):
(7.4)
$$[\mathrm{}_0,\beta _n]=\underset{j\frac{1}{2}+}{}j:\psi _j^+\psi _{nj}:.$$
Let $`U\mathrm{\Phi }_V^\pm `$ be an invariant with respect to $`osp(M|N)^\widehat{}`$ subspace. It follows from Lemma 7.1 and (7.4) that $`vU`$ implies that $`((ad\mathrm{}_0)^s\beta _n)vU`$, $`s_+`$. Hence $`U`$ is invariant with respect to all operators $`\psi _j^+\psi _k`$, where $`\psi ^+,\psi V`$ are such that $`(\psi ^+|\psi ^+)=0=(\psi ^+|\psi )`$ and $`j,k\frac{1}{2}+`$.
Hence, provided that $`M3`$, $`U`$ contains a non-zero purely bosonic element, i.e., an element obtained by applying a polynomial in the $`\gamma _j(\gamma V_{\overline{0}})`$ to $`|0`$. Thus we reduced the problem to the purely bosonic case, i.e., the case when $`M=0`$. In this case the irreducibility was proved in \[L\] using the character formula for modular invariant representations of $`\widehat{C}_n`$ from \[KW1\] and formula (12.13) from \[K3\] (the reference to (13.13) in \[L\] is a misprint).
The remaining cases, when $`M=1`$ or $`2`$ and $`N=2n`$ is even $`2`$ can be reduced again to the purely bosonic case by a direct calculation. We give below details in the $`M=1`$ case, the $`M=2`$ case being similar.
The simple root vectors of $`osp(1,N)^\widehat{}=B(0,n)^\widehat{}`$ are as follows:
$$\begin{array}{ccccc}e_0\hfill & =& (\phi ^1(z)\phi ^1(z))_1\hfill & =& _s\phi _{s+1/2}^{(1)}\phi _{s+1/2}^{(1)},\hfill \\ e_i\hfill & =& (\phi ^i(z)\phi ^{i+1}(z))_0\hfill & =& _s\phi _{s1/2}^{(i)}\phi _{s+1/2}^{(i+1)}(i=1,\mathrm{},n1),\hfill \\ e_n\hfill & =& (\phi (z)\psi (z))_0\hfill & =& _s\phi _{s1/2}^{(n)}\psi _{s+1/2}.\hfill \end{array}$$
Then the simple root vectors of $`sp(N)^\widehat{}=\widehat{C}_n`$ are $`e_0,e_1,\mathrm{}e_{n1}`$ and $`e_n^{}=[e_n,e_n]=`$$`_s\phi _{s1/2}^{(n)}\phi _{s+1/2}^{(n)}`$.
Any vector $`v`$ of $`\mathrm{\Phi }_V`$ can be uniquely written in the form:
$$v=\underset{i_1<\mathrm{}<i_k}{}\psi _{i_1}\mathrm{}\psi _{i_k}u_{i_1,\mathrm{},i_k},$$
where $`u_{i_1,\mathrm{},i_k}`$ are purely bosonic elements (i.e., obtained by applying polynomials in the $`\phi `$’s to $`|0`$). Now, if $`v`$ is a singular vector, i.e., $`e_iv=0`$ for all $`i=0,\mathrm{},n`$, then, in particular, $`e_n^{}v=0`$, and since all $`e_1,\mathrm{},e_{n1},e_n^{}`$ commute with the $`\psi `$’s, we get:
$$\underset{i_1<\mathrm{}<i_k}{}\psi _{i_1}\mathrm{}\psi _{i_k}(e_iu_{i_1,\mathrm{},i_k})=0,\underset{i_1<\mathrm{}<i_k}{}\psi _{i_1}\mathrm{}\psi _{i_k}(e_n^{}u_{i_1,\mathrm{},i_n})=0.$$
It follows that all $`u_{i_1,\mathrm{},i_n}`$ are purely bosonic singular with respect to $`\widehat{C}_n`$ vectors, hence, due to irreducibility of $`\mathrm{\Phi }_V^\pm `$ for $`M=0`$ mentioned above, we obtain that all $`u_{i_1\mathrm{}i_k}`$ are linear combinations of elements $`|0`$ and $`\phi _{1/2}^{(1)}|0`$. Hence
$$v=\underset{i_1<\mathrm{}<i_k}{}a_{i_1,\mathrm{},i_k}\psi _{i_1}\mathrm{}\psi _{i_k}|0+\underset{i_1<\mathrm{}<i_k}{}b_{i_1,\mathrm{},i_k}\psi _{i_1}\mathrm{}\psi _{i_k}\phi _{1/2}^{(1)}|0.$$
Using that $`e_nv=0`$, we obtain:
$`{\displaystyle \underset{r=1}{\overset{k}{}}}{\displaystyle \underset{i_1<\mathrm{}<i_k}{}}(1)^{r1}a_{i_1,\mathrm{},i_k}\psi _{i_1}\mathrm{}\widehat{\psi }_{i_r}\mathrm{}\psi _{i_k}\phi _{i_r}^{(n)}|0`$
$`+{\displaystyle \underset{r=1}{\overset{k}{}}}{\displaystyle \underset{i_1<\mathrm{}<i_k}{}}(1)^{r1}b_{i_1,\mathrm{},i_k}\psi _{i_1}\mathrm{}\widehat{\psi }_{i_r}\mathrm{}\psi _{i_k}\phi _{i_r}^{(n)}\phi _{1/2}^{(1)}|0=0,`$
which implies that $`a_{i_1,\mathrm{},i_k}`$ (resp. $`b_{i_1,\mathrm{},i_k})=0`$ if $`k>0`$. Thus, the only singular vectors in $`\mathrm{\Phi }_V^+`$ (resp. $`\mathrm{\Phi }_V^{}`$) are scalar multiples of $`|0`$ (resp. $`\phi _{1/2}^{(1)}|0`$).
To conclude that the $`B(0,n)^\widehat{}`$-modules $`\mathrm{\Phi }_V^\pm `$ are irreducible, note that $`\mathrm{\Phi }_V`$ carries a unique non-degenerate Hermitian form $`H(.,.)`$ such that the square length of $`|0`$ is $`1`$ and the adjoint operators of $`\phi _k^{(j)}`$ and $`\psi _k`$ are $`\phi _k^{(j)}`$ and $`\psi _k`$, respectively. The absence of non-trivial singular vectors in $`\mathrm{\Phi }_V^+`$ (resp. $`\mathrm{\Phi }_V^{}`$) implies that the $`B(0,n)^\widehat{}`$-submodules $`\mathrm{\Phi }_V^+`$ (resp. $`\mathrm{\Phi }_V^{}`$) generated by $`|0`$ (resp. $`\phi _{1/2}^{(1)}|0`$) is irreducible, hence the restriction of $`H`$ to it is non-degenerate. Hence the orthogonal complement to $`\mathrm{\Phi }_V^+`$ (resp. $`\mathrm{\Phi }_V^{}`$) is a complementary submodule which has no non-zero singular vectors, hence it is zero, and $`\mathrm{\Phi }_V^\pm `$ are irreducible. ∎
###### Remark 7.2.
The irreducibility in the purely fermionic case was established in \[KP1\] by making use of the Weyl-Kac character formula. An argument, using Virasoro operators, was given in \[F\]. The method of using Virasoro operators to prove irreducibility apparently works only in the presence of fermions (cf. Remark 3.3). It is shown in \[L\] that the irreducibility claims of \[FF\], based on the use of Virasoro operators, are false for the constructions of $`A_{2\mathrm{}1}^{(2)}`$ and $`A_2\mathrm{}^{(2)}`$-modules.
Using Theorem 7.1, it is straightforward to write down the characters and supercharacters for the integrable level $`1`$ $`osp(M|N)^\widehat{}`$-modules. We have:
(7.5)
$$ch\mathrm{\Phi }_V^+\pm ch\mathrm{\Phi }_V^{}=e^{\frac{1}{2}\mathrm{\Lambda }_0}\mathrm{\Pi }_{k=1}^{\mathrm{}}\frac{(1\pm q^{k1/2})^{p(M)}\mathrm{\Pi }_{i=1}^m(1\pm e^{ϵ_i}q^{k1/2})(1\pm e^{ϵ_i}q^{k1/2})}{\mathrm{\Pi }_{j=1}^n(1e^{ϵ_{j+m}}q^{k1/2})(1e^{ϵ_{j+m}}q^{k1/2})},$$
where $`p(M)=0`$ (resp. $`1`$) if $`M`$ is even (resp. odd). A similar formula for supercharacters is obtained by reversing signs in the numerator of the right-hand side of (7.5).
Letting all $`ϵ_i`$ and $`\mathrm{\Lambda }_0`$ equal $`0`$ in (7.5), we obtain:
(7.6)
$$tr_{\mathrm{\Phi }_V^+}q^{L_0}\pm tr_{\mathrm{\Phi }_V^{}}q^{L_0}=\mathrm{\Pi }_{k=1}^{\mathrm{}}\frac{(1\pm q^{k\frac{1}{2}})^M}{(1q^{k\frac{1}{2}})^N}.$$
Noticing that
(7.7)
$$\mathrm{\Pi }_{k=1}^{\mathrm{}}(1q^{k\frac{1}{2}})=\frac{\phi (q^{\frac{1}{2}})}{\phi (q)}\text{ and }\mathrm{\Pi }_{k=1}^{\mathrm{}}(1+q^{k\frac{1}{2}})=\frac{\phi (q)^2}{\phi (q^{\frac{1}{2}})\phi (q^2)}$$
and using the asymptotics (4.13) of $`\eta (\tau )`$, we obtain the following asymptotics as $`\tau 0`$:
(7.8)
$$tr_{\mathrm{\Phi }_V^\pm }q^{L_0}\frac{1}{2^{n+1}}e^{\frac{\pi i}{12\tau }(\frac{1}{2}M+N)}.$$
###### Remark 7.3.
The right-hand side of (7.6) multiplied by $`q^{(NM)/48}`$ is a modular function equal to a product of powers of functions $`\eta (\frac{1}{2}\tau )/\eta (\tau )`$ and $`\eta (\tau )^2/\eta (\frac{1}{2}\tau )\eta (2\tau )`$, and the same holds if we replace $`tr`$ by $`str`$. It is well known (and easy to see) that the above two modular functions along with the modular function $`\eta (2\tau )/\eta (\tau )`$ are transitively permuted (with some constant factors) under the action of $`SL(2,)`$. Thus, the normalized by $`q^{(NM)/48}`$ characters and supercharacters of integrable level one $`osp(M|N)^\widehat{}`$-modules are modular functions, but their $``$-span is not $`SL(2,)`$-invariant.
## 8 On classification of modules over the associated vertex algebras
Define numbers $`u`$ and $`h^{}`$ (the dual Coxeter number) by
(8.1)
$$\text{level }(ku\mathrm{\Lambda }_0)=k,\text{level }(\rho )=h^{}.$$
Their values for all affine superalgebras are given in Table 5.
The following proposition is an immediate corollary of Theorems 6.1 and 6.2.
###### Proposition 8.1.
1. If $`\widehat{𝔤}=A(m,n)^\widehat{}`$ with $`m1`$, $`B(m,n)^\widehat{}`$ with $`m2`$, $`D(m,n)^\widehat{}`$, $`F(4)^\widehat{}`$ or $`G(3)^\widehat{}`$, then the $`\widehat{𝔤}`$-module $`L(ku\mathrm{\Lambda }_0)`$ is principal integrable iff $`k_+`$.
The $`B(1,n)^\widehat{}`$-module $`L(ku\mathrm{\Lambda }_0)`$ is principal integrable iff $`k_+`$ $`\{2n\frac{1}{2}+_+\}`$. The $`D(2,1;a)^\widehat{}`$-module $`L(ku\mathrm{\Lambda }_0)`$ is principal integrable iff $`k_+a_+`$.
2. If $`\widehat{𝔤}=B(m,n)^\widehat{}`$, $`C(n)^\widehat{}`$, $`D(m,n)^\widehat{}`$, $`D(2,1;a)^\widehat{}`$, $`F(4)^\widehat{}`$ or $`G(3)^\widehat{}`$, then the $`\widehat{𝔤}`$-module $`L(k_0\mathrm{\Lambda }_0)`$ is subprincipal integrable iff $`k_0_+`$.
Recall that the $`\widehat{𝔤}`$-module $`V_k:=L(ku\mathrm{\Lambda }_0)`$ has a canonical structure of a vertex algebra for any $`k`$ (see e.g. \[K4\]). It is well known that any irreducible $`V_k`$-module is one of the (irreducible) $`\widehat{𝔤}`$-modules $`L(\mathrm{\Lambda })`$ of level $`k`$, and it is an important problem of vertex algebra theory to find out which of these $`L(\mathrm{\Lambda })`$ are actually $`V_k`$-modules. A necessary condition is given by
###### Proposition 8.2.
Suppose that $`k`$ is such that $`L(ku\mathrm{\Lambda }_0)`$ is a principal (resp. subprincipal) integrable $`\widehat{𝔤}`$-module. If a $`\widehat{𝔤}`$-module $`L(\mathrm{\Lambda })`$ of level $`k`$ is a $`V_k`$-module, then it must be a principal (resp. subprincipal) integrable.
###### Proof.
Denote by $`\widehat{𝔤}^0`$ the subalgebra $`\widehat{𝔤}_{\overline{0}}^{}`$ (resp. $`\widehat{𝔤}_{\overline{0}}^{\prime \prime }`$) of $`\widehat{𝔤}`$ (see § 6). This is an affine Lie algebra. Denote by $`V^0`$ the vertex subalgebra $`U(\widehat{𝔤}^0)v_{ku\mathrm{\Lambda }_0}`$ of $`V_k`$. Since, by definition, $`V^0`$ is an integrable $`\widehat{𝔤}^0`$-module, it follows that it is $`\widehat{𝔤}^0`$-irreducible \[K3\], hence $`V^0`$ is a simple affine vertex algebra of non-negative integral level. But one knows \[Z\] that all irreducible modules over such a vertex algebra are integrable $`\widehat{𝔤}^0`$-modules. Using the complete reducibility of $`\widehat{𝔤}^0`$-modules \[K3\], we deduce that any $`V`$-module, viewed as a $`V^0`$-module, is a direct sum of irreducible integrable $`\widehat{𝔤}^0`$-modules, which proves the proposition.
Let $`\widehat{𝔤}_+=[t]_{}𝔤+d`$ and consider a $`1`$-dimensional module $`_k`$ $`(k)`$ over $`\widehat{𝔤}_++K`$ on which $`\widehat{𝔤}_+`$ acts trivially and $`K=k`$. Then $`L(ku\mathrm{\Lambda }_0)`$ is a quotient of the induced $`\widehat{𝔤}`$-module $`\stackrel{~}{V}_k=U(\widehat{𝔤})_{U(\widehat{𝔤}_++K)}_k`$ by a left ideal $`I_k`$ of $`U(\widehat{𝔤})`$ applied to $`11`$. Suppose that $`k`$ is such that $`L(ku\mathrm{\Lambda }_0)`$ is a principal integrable $`\widehat{𝔤}`$-module. As we have seen in the proof of Proposition 8.2, viewed as a $`\widehat{𝔤}_0^{}`$-module, $`L(ku\mathrm{\Lambda }_0)`$ is a direct sum of irreducible integrable highest weight modules. All these modules have the same level $`\mathrm{}`$ (resp. $`\mathrm{}_+,\mathrm{}_{}`$ when $`𝔤_{\overline{0}}^{}`$ has two simple components) given in terms of $`k`$ as follows:
$`\mathrm{}`$ $`=`$ $`\mathrm{}_+=k\text{ if }𝔤B(1,n),\mathrm{}=2k\text{ if }𝔤=B(1,n),`$
$`\mathrm{}_{}`$ $`=`$ $`k\text{ if }𝔤=D(2,n),\mathrm{}_{}=a^1k\text{ if }𝔤=D(2,1;a).`$
In particular, $`I_k`$ contains the element
(8.2)
$$e_\theta ^{}(1)^{\mathrm{}+1}\text{( resp. elements }e_{\theta _+^{}}(1)^{\mathrm{}_+^{}+1}\text{ and }e_\theta _{}^{}(1)^{\mathrm{}_{}^{}+1}).$$
If elements (8.2) generate the left ideal $`I_k`$, it follows that a $`\widehat{𝔤}`$-module $`L(\mathrm{\Lambda })`$ of level $`k`$ is a $`V_k`$-module iff the field $`e_\theta ^{}(z)^{\mathrm{}+1}`$ (resp. fields $`e_{\theta _+^{}}^{\mathrm{}_+^{}+1}(z)`$ and $`e_\theta _{}^{}^{\mathrm{}_{}^{}+1}(z)`$) annihilate $`L(\mathrm{\Lambda })`$. The latter property implies that, viewed as a $`\widehat{𝔤}_{\overline{0}}^{}`$-module, $`L(\mathrm{\Lambda })`$ is a direct sum of irreducible integrable modules and therefore $`L(\mathrm{\Lambda })`$ is a principal integrable $`\widehat{𝔤}`$-module. We thus established a sufficient condition for a $`\widehat{𝔤}`$-module $`L(\mathrm{\Lambda })`$ to be a $`V_k`$-module:
###### Proposition 8.3.
Let $`k`$ be such that $`L(ku\mathrm{\Lambda }_0)`$ is a principal integrable $`\widehat{𝔤}`$-module and suppose that the left ideal $`I_k`$ is generated by (8.2). Let $`L(\mathrm{\Lambda })`$ be a principal integrable $`\widehat{𝔤}`$-module of level $`k`$. Then $`L(\mathrm{\Lambda })`$ is a $`V_k`$-module.
###### Proposition 8.4.
Let $`k`$ be such that $`L(ku\mathrm{\Lambda }_0)`$ is a principal integrable $`\widehat{𝔤}`$-module.
1. Suppose that the highest weight $`ku\mathrm{\Lambda }_0`$ is the only singular weight of the $`\widehat{𝔤}`$-module $`\stackrel{~}{V}_k`$ which is principal integrable. Then elements (8.2) generate the left ideal $`I_k`$.
2. The assumption of (a) holds if
(8.3)
$$k+h^{}0,$$
and for any principal integrable weight $`\mathrm{\Lambda }`$ of level $`k`$ one has:
(8.4)
$$\mathrm{\Lambda }ku\mathrm{\Lambda }_0\widehat{Q}\backslash \delta ,$$
where $`\widehat{Q}=_{i\widehat{I}}\alpha _i`$ is the root lattice of $`\widehat{𝔤}`$.
###### Proof.
Let $`I_k^{}`$ $`(I_k)`$ denote the left ideal of $`U(\widehat{𝔤})`$ generated by elements (8.2). Then the $`\widehat{𝔤}`$-module $`V_k^{}=\stackrel{~}{V}_k/(I_k^{}(11))`$ is principal integrable, hence each of its singular weights $`\mathrm{\Lambda }`$ is integrable. Hence, if the condition of (a) holds, the $`\widehat{𝔤}`$-module $`V_k^{}`$ is irreducible, and therefore $`I_k^{}=I_k`$.
Furthermore, obviously, $`\mathrm{\Lambda }ku\mathrm{\Lambda }_0\widehat{Q}`$, hence (8.4) implies that $`\mathrm{\Lambda }=ku\mathrm{\Lambda }_0+j\delta `$ for some $`j`$. Using the Casimir operator \[K3\], we obtain:
$$(ku\mathrm{\Lambda }_0+\rho |ku\mathrm{\Lambda }_0+\rho )=(ku\mathrm{\Lambda }_0+\rho j\delta |ku\mathrm{\Lambda }_0+\rho j\delta ),$$
which is equivalent to $`j(k+h^{})=0`$. But then (8.3) implies that $`j=0`$, proving (b). ∎
###### Theorem 8.1.
Let $`\widehat{𝔤}`$ be one of the affine superalgebras $`A(m,n)^\widehat{}`$ with $`m1`$, $`B(m,n)^\widehat{}`$ with $`m1`$, $`D(m,n)^\widehat{}`$, $`D(2,1;a)^\widehat{}`$ with $`a^1`$, $`F(4)^\widehat{}`$ or $`G(3)^\widehat{}`$. Then all integrable $`\widehat{𝔤}`$-modules $`L(\mathrm{\Lambda })`$ of level $`1`$ are $`V_1`$-modules (the complete list of these $`\mathrm{\Lambda }`$’s is given by Remark 6.3).
###### Proof.
Note that in the $`A(m,n)^\widehat{}`$ case $`V_1`$ is a subalgebra of the vertex subalgebra $`F_0`$ of $`F`$ (constructed in § 3), while the highest component of the $`F_0`$-module $`F_s`$ restricted to $`V_1`$ is $`L(\mathrm{\Lambda }_{(s)})`$. Since the $`\mathrm{\Lambda }_{(s)}`$ exhaust all integrable highest weights of level $`1`$, by Proposition 8.1, they give a complete list of irreducible $`V_1`$-modules.
In the $`B(m,n)^\widehat{}`$ and $`D(m,n)^\widehat{}`$ cases we note that $`V_1`$ is isomorphic to the vertex algebra $`\mathrm{\Phi }_V^+`$ (see Theorem 7.1), $`\mathrm{\Phi }_V^{}`$ is its irreducible module, and these two modules produce all integrable highest weights of level $`1`$.
The cases $`F(4)^\widehat{}`$ and $`G(3)^\widehat{}`$ are obvious since $`V_1`$ is the only irreducible integrable module of level $`1`$ (see Remark 6.3).
It remains to show that $`L\left(\frac{a+2}{a+1}\mathrm{\Lambda }_0\mathrm{\Lambda }_1+\frac{1a}{a}\mathrm{\Lambda }_3\right)`$ is a $`V_1`$-module in the $`D(2,1;a)^\widehat{}`$ case. But $`\frac{a+2}{a+1}\mathrm{\Lambda }_0\mathrm{\Lambda }_1+\frac{1a}{a}\mathrm{\Lambda }_3=\frac{1}{a+1}\mathrm{\Lambda }_0\left(\frac{1}{2}\alpha _0+\frac{a1}{2a}\alpha _3\right)`$, hence the difference of this weight and $`u\mathrm{\Lambda }_0`$ does not lie in the root lattice; we also have: $`k=1`$ and level $`h^{}=0`$. Hence we may apply Propositions 8.4 and 8.3.
###### Remark 8.1.
The lowest energy $`D(2,1;a)`$-submodule of the module $`L\left(\frac{a+2}{a+1}\mathrm{\Lambda }_0\mathrm{\Lambda }_1+\frac{1a}{a}\mathrm{\Lambda }_3\right)`$ is the module $`\overline{L}(\overline{\mathrm{\Lambda }}_1+(a^11)\overline{\mathrm{\Lambda }}_3)`$. It has dimension $`4a^1+2`$. For $`a=1`$ this is the defining module of $`D(2,1)`$; for $`a=\frac{1}{2}`$ (resp. $`\frac{1}{3}`$) this is the $`10`$\- (resp. $`14`$-) dimensional module mentioned in Remark 6.8. As a $`D(2,1;a)`$-module, the even (resp. odd) part of this module is isomorphic to the irreducible $`s\mathrm{}(2)+s\mathrm{}(2)+s\mathrm{}(2)`$-module
$$^{a^1}^2(\text{resp. }^2^{a^1+1}).$$
###### Remark 8.2.
Let $`V`$ be a vertex algebra with a conformal vector such that $`L_0`$ is diagonizable with finite-dimensional eigenspaces and rational eigenvalues. It is a general belief that if $`V`$ has finitely many irreducible modules, then the character $`tr_Mq^{L_0}`$ of each of these modules $`M`$ becomes a modular function when normalized, i.e., multiplied by a suitable power of $`q`$. The example of the vertex algebra $`V_1`$ for $`B(m,n)^\widehat{}`$ and $`D(m,n)^\widehat{}`$ confirms this conjecture and leads to believe that the same is true for $`D(2,1;a)^\widehat{}`$ (with $`a^1`$), $`F(4)^\widehat{}`$ and $`G(3)^\widehat{}`$.
The vertex algebra $`V`$ is called rational if $`L_0`$ has integral eigenvalues, the number of irreducible $`V`$-modules is finite and any $`V`$-module is completely reducible. It follows from the above discussion that the vertex algebra $`V_1`$ for $`B(m,n)^\widehat{}`$, $`D(m,n)^\widehat{}`$, $`D(2,1;a))`$, $`F(4))^\widehat{}`$ and $`G(3)^\widehat{}`$ is a rational vertex algebra, and that, moreover, the corresponding Zhu algebra \[Z\] is finite-dimensional semisimple (and even $`1`$-dimensional in the $`F(4)^\widehat{}`$ and $`G(3)^\widehat{}`$ cases).
It was proved by Zhu \[Z\] under certain technical assumptions that the $``$-span of normalized characters of irreducible modules over a rational vertex algebra is $`SL(2,)`$-invariant, and it was believed by many that the technical assumptions may be removed. However, the above mentioned rational vertex algebra $`V_1`$ shows that this is not the case.
###### Remark 8.3.
There are only two cases where there exists only a finite number of essentially inequivalent subprincipal integrable $`\widehat{𝔤}`$-modules of a given non-zero level $`k`$:
$$\widehat{𝔤}=F(4)^\widehat{},k=\frac{3}{2}\text{ and }\widehat{𝔤}=G(3)^\widehat{},k=\frac{4}{3}.$$
In both cases the only subprincipal integrable $`\widehat{𝔤}`$-module is $`L(\mathrm{\Lambda }_0)`$. In both cases the associated vertex algebra is rational with a unique irreducible module and the Zhu algebra is $`1`$-dimensional.
## 9 Some remarks and open problems
### 9.1
The calculation of characters of integrable highest weight modules of arbitrary level $`k`$ over affine superalgebras seems to be a very difficult problem. One may expect that the case of the “critical” level $`k=h^{}`$ should be rather different from other levels (as for the affine Lie algebras). However, the construction of level $`1`$ integrable modules over $`osp(m|n)^\widehat{}`$ given in § 7 is the same for all values $`m`$ and $`n`$ though $`1`$ is the critical level iff $`mn=1`$.
Formula (5.13) leads us to believe in the following conjecture:
Consider a principal integrable highest weight module $`L(\mathrm{\Lambda })`$ over an affine superalgebra $`\widehat{𝔤}`$ and suppose that one can choose a set of simple roots $`\widehat{\mathrm{\Pi }}`$ such that it contains a maximal $`\mathrm{\Lambda }+\rho `$-isotropic subset $`S_\mathrm{\Lambda }`$ of roots (i.e., all roots from $`S_\mathrm{\Lambda }`$ are pairwise orthogonal and orthogonal to $`\mathrm{\Lambda }+\rho `$ \[KW\]). Let $`\widehat{W}^\mathrm{\#}`$ be the Weyl group of the integrable part $`\widehat{𝔤}_{\overline{0}}^{}`$ of $`\widehat{𝔤}_{\overline{0}}`$. We conjecture that the following character formula holds:
(9.1)
$$e^\rho RchL(\mathrm{\Lambda })=\underset{w\widehat{W}^\mathrm{\#}}{}ϵ(w)w\frac{e^{\mathrm{\Lambda }+\rho }}{\mathrm{\Pi }_{\beta S_\mathrm{\Lambda }}(1+e^\beta )}.$$
Note that the assumptions of this conjecture exclude the critical level, and include the level $`1`$ integrable modules over exceptional affine superalgebras.
### 9.2
In the papers \[KW1\] and \[KW2\] we proved character formulas for a class of modules $`L(\mathrm{\Lambda })`$ over affine Lie algebras, called admissible modules, which includes integrable modules. These character formulas imply that the normalized specialized characters of admissible modules are modular functions (and we conjecture that this property characterizes admissible modules). Of course, these character formulas break down in the Lie superalgebra case. However, in certain exceptional situations, when the character admits a simple product expansion in the Lie algebra case (see \[KW2\], Theorem 3.2), it seems that a similar product formula holds in the Lie superalgebra case as well.
Concretely, let $`u`$ be a positive integer and let
(9.2)
$$k=h^{}(u^11)$$
(recall that in general the level $`k`$ of an admissible module is $`h^{}(u^11)`$). Let $`y`$ be an automorphism of the root lattice $`\widehat{Q}`$ such that all roots $`\gamma _i=y((u1)\delta _{i0}K+\alpha _i^{})`$ are positive $`(i\widehat{I})`$. The weights of the form $`y.k\mathrm{\Lambda }_0`$, where, as usual, $`y.\lambda =y(\lambda +\rho )\rho `$, are called *admissible*. We conjecture that the following analog of formula (3.3) from \[KW2\] holds:
(9.3) $`chL(y.(k\mathrm{\Lambda }_0))=e^{y.(k\mathrm{\Lambda }_0)}\left({\displaystyle \frac{\phi (q^u)}{\phi (q)}}\right)^{\mathrm{}}\mathrm{\Pi }_{\begin{array}{c}\alpha \overline{\mathrm{\Delta }}_{\overline{0}}\\ n\end{array}}{\displaystyle \frac{1q^{un}e^{y.\alpha }}{1q^ne^\alpha }}/\mathrm{\Pi }_{\begin{array}{c}\alpha \overline{\mathrm{\Delta }}_{\overline{1}}\\ n\end{array}}{\displaystyle \frac{1+q^{un}e^{y.\alpha }}{1+q^ne^\alpha }},`$
where $`\mathrm{}`$ is the rank of $`𝔤`$, $`\overline{\mathrm{\Delta }}_{\overline{0}},\overline{\mathrm{\Delta }}_{\overline{1}}`$ are the sets of even and odd roots of $`𝔤`$, and $`q=e^\delta `$.
This conjecture agrees with formula (7.5) in the case $`\widehat{𝔤}=osp(2|2)^\widehat{}=s\mathrm{}(2|1)^\widehat{}`$. In this case $`k=1/2`$ and $`h^{}=1`$, so that (9.2) holds for $`u=2`$. All the admissible weights of level $`\frac{1}{2}`$ are as follows:
$$\frac{1}{2}\mathrm{\Lambda }_i(i=0,1,2),\frac{1}{2}\mathrm{\Lambda }_0\frac{1}{2}\alpha _0,$$
where the Dynkin diagram is chosen such that $`\alpha _1`$ and $`\alpha _2`$ are odd roots (and $`\alpha _0`$ is even). Character formula (7.5) gives:
(9.6) $`\begin{array}{c}ch(\frac{1}{2}\mathrm{\Lambda }_0)\\ ch(\frac{1}{2}\mathrm{\Lambda }_0\frac{1}{2}\alpha _0)\end{array}=\frac{1}{2}e^{\frac{1}{2}\mathrm{\Lambda }_0}\left({\displaystyle \frac{\mathrm{\Psi }(u^1vq^{\frac{1}{2}};q)}{\mathrm{\Psi }(uvq^{\frac{1}{2}};q)}}\pm {\displaystyle \frac{\mathrm{\Psi }(u^1vq^{\frac{1}{2}};q)}{\mathrm{\Psi }(uvq^{\frac{1}{2}};q)}}\right),`$
where $`u=e^{\frac{1}{2}\alpha _1},v=e^{\frac{1}{2}\alpha _2}`$ and
(9.7) $`\mathrm{\Psi }(z;q)=\mathrm{\Pi }_{k=1}^{\mathrm{}}(1+zq^{k1})(1+z^1q^k),`$
whereas formula (9.3) gives:
(9.8a) $`ch(\frac{1}{2}\mathrm{\Lambda }_0)`$ $`=`$ $`e^{\frac{1}{2}\mathrm{\Lambda }_0}{\displaystyle \frac{\mathrm{\Psi }(u^2q;q^2)\mathrm{\Psi }(v^2q;q^2)\phi (q^2)^2}{\mathrm{\Psi }(uvq;q^2)\phi (q)^2}}`$
(9.8b) $`ch(\frac{1}{2}\mathrm{\Lambda }_0\frac{1}{2}\alpha _0)`$ $`=`$ $`e^{\frac{1}{2}\mathrm{\Lambda }_0\frac{1}{2}\alpha _0}{\displaystyle \frac{\mathrm{\Psi }(u^2;q^2)\mathrm{\Psi }(v^2;q^2)\phi (q^2)^2}{\mathrm{\Psi }(uvq;q^2)\phi (q)^2}}.`$
However, the seemingly different expressions in the right-hand sides of (9.6) and (9.8) actually coincide due to one of the addition theta function formulas (cf. \[M\], formula (6.6) and notation on p. 17):
(9.9)
$$\theta _{00}(\tau ,z_1)\theta _{00}(\tau ,z_2)=\theta _{00}(2\tau ,z_1+z_2)\theta _{00}(2\tau ,z_1z_2)+\theta _{10}(2\tau ,z_1+z_2)\theta _{10}(2\tau ,z_1z_2),$$
if we let $`u=e^{2\pi iz_1},v=e^{2\pi iz_2}`$.
Using \[KW\], formula (6.1), it is immediate to show that the span of supercharacters of the four admissible $`s\mathrm{}(2|1)^\widehat{}`$-modules of level $`1/2`$ is $`SL(2,)`$-invariant. Thus, it is natural to conjecture that this modular invariance property of admissible characters holds for any affine superalgebra $`\widehat{𝔤}`$ and any $`k`$ given by (9.2).
Two other very interesting examples are provided by Remark 8.3: $`\widehat{𝔤}=F(4)^\widehat{}`$ with $`u=2`$, $`y=1`$ and $`\widehat{𝔤}=G(3)^\widehat{}`$ with $`u=3,y=1`$.
### 9.3
The first case not covered in § 5, that when $`m=n1`$, is very interesting. It connects the level $`1`$ modules over $`𝔤\mathrm{}(n1|n)^\widehat{}`$ (or, equivalently the “critical” level $`1`$ modules over $`g\mathrm{}(n|n1)\widehat{}`$) to the denominator identity for $`s\mathrm{}(n|n)^\widehat{}`$, which is unknown. Analyzing this connection, we arrived at the following $`s\mathrm{}(2|2)^\widehat{}`$ denominator identity:
(9.10)
$$e^\rho R=\underset{w\widehat{W}^\mathrm{\#}}{}ϵ(w)w\frac{e^\rho }{(1+e^{\alpha _0})\mathrm{\Pi }_{j=1}^{\mathrm{}}(1+q^je^{\alpha _2})(1+q^{j1}e^{\alpha _2})}$$
where, as before, $`q=e^\delta `$. Here we use the Dynkin diagram with the grey nodes $`\alpha _0`$ and $`\alpha _2`$. If all four nodes are grey we get the same identity with $`\rho `$ replaced by $`0`$; in yet another form (9.10) can be written as follows:
(9.11) $`{\displaystyle \frac{R}{\phi (q)}}\mathrm{\Pi }_{n=1}^{\mathrm{}}(1q^{2n})(1+q^{2n1}e^{\alpha _1+\alpha _3})(1+q^{2n1}e^{\alpha _1\alpha _3})`$
$`={\displaystyle \underset{w\widehat{W}^\mathrm{\#}}{}}ϵ(w)w(\mathrm{\Pi }_{n=1}^{\mathrm{}}(1+q^ne^{\alpha _1})(1+q^ne^{\alpha _3})(1+q^{n1}e^{\alpha _1})(1+q^{n1}e^{\alpha _3})).`$
The latter identity is equivalent to the following identity in $`u=e^{\alpha _1},x=e^{\alpha _2},v=e^{\alpha _3}`$ and $`q`$ (where $`\mathrm{\Psi }`$ is defined by (9.7)):
$`\mathrm{\Psi }(uvq;q^2)\mathrm{\Psi }(ux;q)\mathrm{\Psi }(vx;q)`$
$`=`$ $`\mathrm{\Psi }(uv^1q;q^2)\mathrm{\Psi }(x;q)\mathrm{\Psi }(uvx;q)x\mathrm{\Psi }(uvx^2q;q^2)\mathrm{\Psi }(u;q)\mathrm{\Psi }(v;q).`$
In the notation of \[M\] this identity can be rewritten in terms of theta functions as follows (if we let $`u=e^{2\pi iz_1}`$, $`v=e^{2\pi iz_2}`$, $`x=e^{2\pi iz_3}`$):
(9.12) $`\theta _{00}(2\tau ,z_1+z_2)\theta _{11}(\tau ,z_1+z_3)\theta _{11}(\tau ,z_2+z_3)`$
$`+\theta _{00}(2\tau ,z_1z_2)\theta _{10}(\tau ,z_3)\theta _{11}(\tau ,z_1+z_2+z_3)`$
$`=\theta _{00}(2\tau +z_1+z_2+2z_3)\theta _{10}(\tau ,z_1)\theta _{10}(\tau ,z_2).`$
Identity (9.3) can be derived from (9.7) as follows. Replacing $`z_i`$ by $`z_i+\frac{1}{2}\tau `$ (resp. $`z_i+\frac{1}{2}(1+\tau )`$) in (9.7), we obtain:
(9.13a)
$$\theta _{10}(\tau ,z_1)\theta _{10}(\tau ,z_2)=\theta _{00}(2\tau ,z_1+z_2)\theta _{10}(2\tau ,z_1z_2)+\theta _{10}(2\tau ,z_1+z_2)\theta _{00}(2\tau ,z_1z_2)$$
(9.13b)
$$\theta _{11}(\tau ,z_1)\theta _{11}(\tau ,z_2)=\theta _{00}(2\tau ,z_1+z_2)\theta _{10}(2\tau ,z_1z_2)\theta _{10}(2\tau ,z_1+z_2)\theta _{00}(2\tau ,z_1z_2).$$
Substituting (9.13b) (resp. (9.13a)) in the first (resp. second) summand of the left-hand side of (9.3), we obtain the product of $`\theta _{00}(2\tau ,z_1+z_2+2z_3)`$ and the right-hand side of (9.13a), and, substituting its left-hand side, we obtain the left-hand side of (9.3).
We also have a conjectural formula for an $`s\mathrm{}(3|3)^\widehat{}`$ denominator identity, but it is too cumbersome to be reproduced here. We have no conjectures as how the $`s\mathrm{}(n|n)^\widehat{}`$ denominator identity should look for $`n>3`$.
Using the connection of the $`s\mathrm{}(2|2)^\widehat{}`$ denominator identity to level $`1`$ modules over $`s\mathrm{}(2|1)^\widehat{}`$ we deduce from (9.10) $`(k_+)`$:
(9.14) $`chL(k\mathrm{\Lambda }_0(k+1)\mathrm{\Lambda }_1)={\displaystyle \frac{\phi (q)}{e^\rho R}}{\displaystyle \underset{wr_0}{}}ϵ(w)w{\displaystyle \underset{j_+}{}}t_{j\alpha _0}{\displaystyle \frac{e^{\mathrm{\Lambda }+\rho }}{\mathrm{\Pi }_{n=1}^{\mathrm{}}(1+q^{n1}e^{\alpha _2})(1+q^ne^{\alpha _2})}}.`$
Here the Dynkin diagram is chosen in such a way that $`\alpha _0`$ is even and $`\alpha _1,\alpha _2`$ are odd simple roots.
### 9.4
Let $`k`$ be such that $`V_k=L(uk\mathrm{\Lambda }_0)`$ is a (principal or subprincipal) integrable $`\widehat{𝔤}`$-module of level $`k`$. Is it always true that any integrable $`\widehat{𝔤}`$-module of level $`k`$ can be extended to a module over the vertex algebra $`V_k`$? Of course, this question is closely related to the description of generators of the left ideal $`I_k`$. In the principal integrable case $`I_k`$ contains elements (8.2), and the answer to the above question in this case would be positive if $`I_k`$ were generated by these elements.
Is it true that the normalized (by a power of $`q`$) characters $`trq^{L_0}`$ (where $`L_0`$ is given by the Sugawara construction \[K4\]) of irreducible $`V_k`$-modules are modular functions, provided that there are finitely many of them and $`k+h^{}0`$. Is it true that for $`k`$ of the form (9.2), all $`V_k`$-modules are admissible modules?
### 9.5
A few examples that we have worked out in the paper indicate that the theory of integrable highest weight modules over affine Lie superalgebras is dramatically different from that in the Lie algebra case. The only exception is the case of $`\widehat{𝔤}=B(0,n)^\widehat{}`$. The integrability conditions are (see Table 1 for its Dynkin diagram).
$$k_i_+\text{ for all }i,k_n2_+,$$
hence the level $`k=2k_0+\mathrm{}+2k_{n1}+k_n`$ is a non-negative even integer and the number of integrable highest weight $`B(0,n)^\widehat{}`$-modules is finite for each $`k`$. Moreover, each of these modules extends to an irreducible $`V_k`$-module since $`I_k`$ is generated by $`e_\theta (1)^{k+1}`$ for each $`k`$, and these are all irreducible modules over the vertex algebra $`V_k(k2_+)`$.
Furthermore, the character formula for all integrable $`B(0,n)^\widehat{}`$-modules $`L(\mathrm{\Lambda })`$ is known (see \[K2\]), and it is given by the same expression as that for the twisted affine algebra $`A_{2n}^{(2)}`$ (replacing the black node by a white one). In order to derive the transformation formula of $`B(0,n)^\widehat{}`$ supercharacters from that of $`A_{2n}^{(2)}`$ characters, we need to go from the $`A_{2n}^{(2)}`$ coordinates, which we call $`A`$-coordinates, to the $`B(0,n)^\widehat{}`$ coordinates, which we will call $`B`$-coordinates. This calculation is explained below.
The $`B`$-coordinates $`(\tau ,z_B,u_B)`$ of $`h\widehat{𝔥}`$ are defined by
$$h=2\pi i(\tau \frac{1}{2}\mathrm{\Lambda }_0+z_B+u_B\delta ),\text{ where }z_B𝔥.$$
Let $`\beta =\frac{1}{2}_{j=0}^{n1}(nj)\alpha _j^{}`$. Then $`t_\beta (\mathrm{\Lambda }_i)=\mathrm{\Lambda }_i`$ for $`i=1,\mathrm{},n1`$ and $`t_\beta (\frac{1}{2}\mathrm{\Lambda }_0),\alpha _i^{}=\delta _{0n}`$, hence we may take $`\mathrm{\Lambda }_n=t_\beta (\frac{1}{2}\mathrm{\Lambda }_0)`$ for the $`0`$$`^{\text{th}}`$ fundamental weight of $`A_{2n}^{(2)}`$. Hence the $`A`$-coordinates are expressed via $`B`$-coordinates by
(9.15)
$$\stackrel{~}{h}=t_\beta (h)=2\pi i(\tau \mathrm{\Lambda }_n+z_A+u_A\delta ).$$
Recall that $`SL(2,)`$ acts on functions in $`\tau ,z,u`$ by the formula \[K3, Chapter 13\]:
$$F(\tau ,z,u)|_{\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)}=j(\tau )^nF(\frac{a\tau +b}{c\tau +d},\frac{z}{c\tau +d},u\frac{c(z|z)}{2(c\tau +d)}).$$
Furthermore, defining a new function $`F^{\alpha ,\beta }`$ by
$$F^{\alpha ,\beta }(h)=F(t_\beta (h)+2\pi i\alpha \pi i(\alpha |\beta )\delta ),$$
we have \[KP2\]:
(9.16)
$$F^{\alpha ,\beta }|_{\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)}=(F|_{\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)})^{d\alpha b\beta ,a\beta c\alpha }.$$
We shall use the following connection between supercharacters of $`B(0,n)^\widehat{}`$ and characters of $`A_{2n}^{(2)}`$, which follows from (9.13) and definitions:
(9.17)
$$schL(\mathrm{\Lambda })(h)=chL(\mathrm{\Lambda })(\stackrel{~}{h})^{\beta ,\beta }=(1)^{nk/2}chL(\mathrm{\Lambda })(\stackrel{~}{h})^{\beta ,\beta }.$$
Recall that the normalized $`A_{2n}^{(2)}`$ character $`\stackrel{~}{\chi }_\mathrm{\Lambda }`$ and the normalized $`B(0,n)^\widehat{}`$ supercharacter $`\chi _\mathrm{\Lambda }`$ are defined by:
$$\stackrel{~}{\chi }_\mathrm{\Lambda }=q^{\stackrel{~}{m}_\mathrm{\Lambda }}chL(\mathrm{\Lambda }),\chi _\mathrm{\Lambda }=q^{m_\mathrm{\Lambda }}schL(\mathrm{\Lambda }),$$
where in $`B`$-coordinates:
$$m_\mathrm{\Lambda }=\frac{|\mathrm{\Lambda }+\rho |^2}{2(k+h^{})}\frac{|\rho |^2}{2h^{}}=\frac{(\mathrm{\Lambda }+2\rho |\mathrm{\Lambda })}{2(k+h^{})}\frac{c_k}{24},c_k=\frac{ksdimB(0,n)}{k+h^{}},$$
and $`\stackrel{~}{m}_\mathrm{\Lambda }`$ is defined by a similar formula in $`A`$-coordinates. (Hence $`\chi _\mathrm{\Lambda }(\tau ,0,0)=str_{L(\mathrm{\Lambda })}q^{L_0c_k/24}`$, as it should be.)
Let $`S=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)SL(2,)`$. We shall denote by $`S_A`$ (resp. $`S_B`$) the action of $`S`$ in $`A`$\- (resp. $`B`$-coordinates). We have by (9.17):
(9.18)
$$\chi _\mathrm{\Lambda }|_{S_B}=\stackrel{~}{\chi }_\mathrm{\Lambda }^{\beta ,\beta }=(\stackrel{~}{\chi }_\mathrm{\Lambda }|_{S_A})^{\beta ,\beta },$$
where the last equality holds due to (9.16).
But one has (see \[KP2\],\[K3\], Theorem 13.8a):
$$\stackrel{~}{\chi }_\mathrm{\Lambda }|_{S_A}=\underset{MP_+^kmod\delta }{}S_{\mathrm{\Lambda }M}\stackrel{~}{\chi }_M,$$
where $`(S_{\mathrm{\Lambda }M})`$ is an explicitly known matrix. Hence, continuing the calculation (9.18), we get, using (9.17):
$$\chi _\mathrm{\Lambda }|_{S_B}=\underset{MP_+^kmod\delta }{}S_{\mathrm{\Lambda }M}\stackrel{~}{\chi }_M^{\beta ,\beta }=(1)^{nk/2}\underset{M}{}S_{\mathrm{\Lambda }M}\stackrel{~}{\chi }_M^{\beta ,\beta }.$$
Using again (9.17), we obtain the final transformation formula:
(9.19)
$$\chi _\mathrm{\Lambda }|_{S_B}=(1)^{nk/2}\underset{MP_+^kmod\delta }{}S_{\mathrm{\Lambda }M}\chi _M.$$
It is clear from the above calculation that the $`SL(2,)`$-invariance of normalized $`B(0,n)^\widehat{}`$ characters $`\stackrel{~}{\chi }_\mathrm{\Lambda }`$ does not hold, but the span of $`\{\stackrel{~}{\chi }_\mathrm{\Lambda },\stackrel{~}{\chi }_\mathrm{\Lambda }^{\beta ,0},\stackrel{~}{\chi }_\mathrm{\Lambda }^{0,\beta }\}_{\mathrm{\Lambda }P_+^kmod\delta }`$ is $`SL(2,)`$-invariant.
### 9.6
We use this opportunity to make some corrections to \[KW\].
Due to computer error the following lines disappeared from the paper:
1. Bottom of page 418:
and $`4(n+1)^2`$, respectively (given by Theorem 4.2; see also Examples 5.3 and
2. Bottom of page 421:
(this is independent of the choice of $`B`$), and let $`W^\mathrm{\#}`$ denote the subgroup of $`W`$ generated by reflections $`r_\alpha `$ with respect to all $`\alpha \mathrm{\Delta }_0^\mathrm{\#}`$. Denote by $`(.|.)`$ the even
Also, the diagrams $`B(0,n)`$, $`D(m,n)`$ and $`D(2,1;\alpha )`$ on p. 429 should be as follows:
$`\begin{array}{cc}B(0,n)\hfill & \mathrm{}\text{}\hfill \\ D(m,n),m2\hfill & \begin{array}{ccc}& \text{}\mathrm{}\hfill & _1\hfill \\ \hfill \mathrm{}_2_2\mathrm{}_2& & \\ & \text{}\mathrm{}\hfill & _1\hfill \end{array}\hfill \\ D(2,1;\alpha )\hfill & \begin{array}{ccc}& \text{}\mathrm{}\hfill & _1\hfill \\ \hfill {}_{}{}^{2}& & \\ & \text{}\mathrm{}\hfill & _1\hfill \end{array}\hfill \end{array}`$
Furthermore, the following corrections should be made:
* page 417, line 13$``$: $`\mathrm{}(q)`$
* page 432, line 12$``$: $`M_\mathrm{\Lambda }:=\{\alpha \overline{\mathrm{\Delta }}_0|\alpha S_\mathrm{\Lambda }\}`$,
* page 434, lines 5,7$``$: $`\alpha _2`$ should be replaced by $`\alpha _1`$,
* page 435, line 6$``$: $`(_{j>0}t^j𝔤)`$,
* page 438, line 1$``$: $`:(\alpha _1|\alpha _1)=0`$,
* page 449, line 3$``$: $`(\alpha _2|\alpha _2)=2`$,
* page 450, line 4$``$: $`\widehat{R}=R_m\mathrm{\Pi }_{n=1}^{\mathrm{}}\mathrm{}`$, where $`R_m`$ is the denominator of $`A_m`$, not the one defined by (7.1),
* page 453: Theorem 8.1(a) as stated holds for the subprincipal integrable modules (cf. Theorem 6.2 of the present paper).
It is appropriate to mention here that the specialization (7.2) of Conjecture 7.2 has been proved recently independently by S. Milne (by combinatorial methods) and by D. Zagier (using cusp forms). D. Zagier also proved Conjecture 7.2 in the first unknown case m=2.
In a slightly different form than in \[KW\], Conjecture 7.2 reads:
$`\mathrm{\Pi }_{n=1}^{\mathrm{}}\left(\left({\displaystyle \frac{1q^{2n}}{1q^{2n1}}}\right)^{2s}\mathrm{\Pi }_{\alpha \mathrm{\Delta }}{\displaystyle \frac{1q^{2n}e^\alpha }{1q^{2n1}e^\alpha }}\right)={\displaystyle \underset{\begin{array}{c}n_1,\mathrm{},n_s0\\ k_1\mathrm{}k_s0\end{array}}{}}chL({\displaystyle \underset{i=1}{\overset{s}{}}}k_i\gamma _i,A_m)`$
$`\times q^{\underset{i=1}{\overset{s}{}}k_i(2n_i+1)+(m2i+2)n_i}.`$
Here $`\mathrm{\Delta }`$ is the set of roots of $`A_m,s=\left[\frac{m+1}{2}\right]`$ and $`\{\gamma _1,\mathrm{},\gamma _s\}`$ is the set of positive pairwise orthogonal roots, $`\gamma _1`$ being the highest root.
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# Strong Non-Ultralocality of Ginsparg-Wilson Fermionic Actions
## 1 Introduction
Proper inclusion of chiral dynamics in the framework of lattice-regularized gauge theory has long been one of the main themes in this nonperturbative approach. While there are several useful angles to look at this, the many years of conceptual confusion were primarily caused by the simple fact that lattice theories with naive chiral symmetry can not reproduce the physics of chiral anomaly . At the same time, the departure from naive chiral symmetry to accommodate anomalies typically mutilates other aspects of chiral dynamics.
However, the set of acceptable lattice fermionic actions is rather large, and so it is not implausible that there are actions without naive chiral symmetry, where both the anomalies and the chiral dynamics are accommodated simultaneously. Guided by renormalization group arguments, Ginsparg and Wilson suggested that one should examine lattice Dirac theories in which the chirally nonsymmetric part of the massless propagator is local (coupling decays at least exponentially at long distance) . This is very plausible since the important aspects of chiral dynamics are typically associated with the long distance behavior of the propagator, which is then very mildly affected by the chirally nonsymmetric part of the action. Lattice fermionic actions with this property became known as Ginsparg–Wilson (GW) actions and vectorlike theories defined by them indeed have all the important ingredients of chiral dynamics . More importantly, otherwise acceptable GW actions can be explicitly constructed as demonstrated by the class of operators of Neuberger’s type .
Formally appealing feature of the GW approach is that it can be defined through the symmetry principle . The underlying Ginsparg–Wilson–Lüscher (GWL) symmetry operation is rather unconventional and depends on the action itself. It has been proved that if the action is symmetric, then the corresponding infinitesimal GWL transformation can not be ultralocal . In other words, the transformation must rearrange the system in such a way that a given variable is mixed with infinitely many other variables on an unbounded lattice. This is a universal property of GWL symmetry in the presence of hypercubic symmetries (weak non-ultralocality), and it shows that to enforce features of chiral dynamics without naive chiral symmetry requires a delicate cooperation of many degrees of freedom. Aside from conceptual value, weak non-ultralocality has a well defined physical meaning, because it is associated with non-ultralocality of Noether currents corresponding to GWL symmetry.
While understanding general structural properties of GW actions is of prime importance for both conceptual and practical reasons, the progress has been rather slow. This is mainly due to the fact that the defining condition for GW actions is nonlinear, which causes the related problems to be quite nontrivial. An important basic question here is whether GW actions can be ultralocal. This issue has been settled (with negative answer) for commonly studied class of GW actions some time ago , and the most general result published so far actually follows from weak non-ultralocality. In particular, it was shown that all GW actions for which chirally nonsymmetric part of the propagator is ultralocal (and nonzero), are non-ultralocal in the presence of hypercubic symmetries .
The purpose of this paper is to demonstrate strong non-ultralocality as formulated in Ref. , i.e. non-ultralocality of all doubler-free GW fermionic actions respecting the symmetries of the hypercubic lattice. Our proof will follow the strategy developed in Ref. ,<sup>1</sup><sup>1</sup>1An attempt to prove some new non-ultralocality result using ideas of was also announced in . However, as pointed out in Ref. , the claims there do not appear to be substantiated. which suggests that it is sufficient to consider two-dimensional momentum restrictions for free theories and investigate the consequences of the GW property at the origin of the Brillouin zone. This leads to an algebraic problem reflecting factorization properties of certain class of polynomials . We will adopt the general formalism of Ref. and develop some basic notions of algebraic geometry that will be needed to keep the presentation self-contained.
## 2 Lattice Dirac Operators and Ultralocality
Consider an infinite $`d`$-dimensional hypercubic lattice, where $`d`$ is an even integer. Let $`\psi ,\overline{\psi }^T`$ be the vectors of fermionic variables associated with lattice sites. We will be interested in quadratic expressions (actions) $`\overline{\psi }𝐆\psi `$, defined by the linear operator $`𝐆`$ which acts in the corresponding linear space. The usual gauge–flavor structure of the linear space will be ignored since it will be sufficient for our purposes to consider only operators in unit gauge background. Consequently, the fermionic variables have $`2^{d/2}`$ spinorial components, and arbitrary matrix $`𝐆`$ can be uniquely expanded in the form
$$𝐆_{m,n}=\underset{a=1}{\overset{2^d}{}}𝐆_{m,n}^a\mathrm{\Gamma }^a$$
(1)
Here $`m,n`$ label the lattice points, $`𝐆^a`$ denotes a matrix with position indices, and $`\mathrm{\Gamma }^a`$ is the element of the Clifford basis. Clifford basis is built on gamma–matrices satisfying $`\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu ,\nu }\mathrm{𝕀}`$. We will refer to the operators $`𝐆^a`$ as Clifford components of $`𝐆`$.
It is useful to define the chirally symmetric ($`𝐆_C`$) and chirally nonsymmetric ($`𝐆_N`$) part of arbitrary operator $`𝐆`$, using the $`\gamma _5`$ element of the Clifford algebra, namely
$$𝐆=\frac{1}{2}\gamma _5\{\gamma _5,𝐆\}+\frac{1}{2}\gamma _5[\gamma _5,𝐆]𝐆_N+𝐆_C$$
(2)
This is a unique decomposition such that $`\{𝐆_C,\gamma _5\}=0`$, and $`[𝐆_N,\gamma _5]=0`$.
Important role in our discussion will be played by the subset $`𝔊^{sl}`$ of operators that are local and respect the symmetries of the hypercubic lattice structure. By locality we mean at least exponential decay of interaction at large distances, and the symmetries of hypercubic lattice include translation invariance and invariance under hypercubic rotations and reflections (hypercubic invariance). The elements $`𝐆𝔊^{sl}`$ can be equivalently represented by their Fourier image which is diagonal
$$G(p)\underset{a=1}{\overset{2^d}{}}G^a(p)\mathrm{\Gamma }^aG^a(p)\underset{n}{}g_n^ae^{ipn}g_n^a𝐆_{0,n}^a$$
(3)
Here the functions $`G^a(p)`$ of lattice momenta $`p(p_1,\mathrm{},p_d)`$ are complex–valued, periodic with $`2\pi `$ in all variables, analytic, <sup>2</sup><sup>2</sup>2As is well known from the theory of Fourier series, locality (as defined in Appendix A) implies that the functions $`G^a(p)`$, extended to complex momenta, are analytic in the region $`\{p:|Im(p_\mu )|<\delta ,\mu \}`$ with some positive $`\delta `$. The Brillouin zone of real momenta is embedded in this complex region. We illustrate this in the simplest case of one dimension in Appendix B. and further restricted by hypercubic invariance. The concepts discussed in this paragraph are properly defined in Appendix A.
The set of lattice Dirac operators that will be considered can now be defined. We require locality, symmetries of the hypercubic lattice, and correct classical continuum limit.
###### Definition 1
(Set $`𝔇`$) Let $`D(p)`$ be a local symmetric operator (element of $`𝔊^{sl}`$), and let in the vicinity of $`p=0`$ its Clifford components $`D^a(p)`$ satisfy
$$D^a(p)=\{\begin{array}{cc}ip_\mu +O(p^2),\hfill & \text{if }\mathrm{\Gamma }^a=\gamma _\mu \text{;}\hfill \\ O(p^2),\hfill & \text{if }\mathrm{\Gamma }^a\gamma _\mu ,\mu \hfill \end{array}$$
(4)
Collection $`𝔇𝔊^{sl}`$ of such elements $`D(p)`$ defines the set of lattice Dirac operators.
Note that there are elements both with and without doublers in the set $`𝔇`$.
Finally, our analysis will focus on the property of ultralocality. This notion reflects the fact that the operator doesn’t couple variables beyond some fixed finite distance:
###### Definition 2
(Ultralocality) Let $`𝒞_N`$ denotes the set of all lattice sites contained in the hypercube of side $`2N`$, centered at $`n=0`$, i.e. $`𝒞_N\{n:|n_\mu |N,\mu =1,\mathrm{},d\}`$. Operator $`𝐆`$ is said to be be ultralocal if there is a positive integer $`N`$, so that
$$𝐆_{m,n}^a=\mathrm{\hspace{0.33em}0},m,n:(mn)𝒞_N,a$$
If an ultralocal operator $`𝐆`$ is translation invariant, then Clifford components $`G^a(p)`$ of its Fourier image are functions with finite number of Fourier terms, i.e. $`G^a(p)_{n𝒞_N}g_n^ae^{ipn}`$.
## 3 Formulation of the Problem
Defining property of GW actions is the locality of chirally nonsymmetric part of the propagator. In other words, $`𝐃𝔇`$ is a GW operator if $`(𝐃^1)_N𝔊^{sl}`$. This is very restrictive since the full massless propagator $`𝐃^1`$ is nonlocal. GWL symmetry thus requires that all this nonlocality is contained in the chirally symmetric part $`(𝐃^1)_C`$. Actions with naive chiral symmetry can be viewed as the limiting case in the set of GW actions as $`(𝐃^1)_N`$ completely vanishes.<sup>3</sup><sup>3</sup>3 Strictly speaking, since $`𝐃`$ has zero eigenvalues (for zero momentum states), all the manipulations involving $`𝐃^1`$ should be performed with small mass term added to $`𝐃`$, with massless limit taken at the end. This is always implicitly assumed.
In Fourier space the GW property means that all Clifford components of $`(D^1(p))_N`$ are analytic functions on the whole Brillouin zone. This is automatic for arbitrary action with naive chiral symmetry, and can be also rather easily arranged even for ultralocal $`D(p)`$ with $`(D(p))_N0`$. However, in all such known cases $`D^1(p)`$ has unwanted poles at nonzero momenta, and the theory thus suffers from doubling. This suggests that with only finite number of Fourier components at our disposal, it is perhaps impossible to add to the chirally symmetric operator $`D_C`$ a nonsymmetric part $`D_N`$, that would not interfere with the behavior of the propagator at large distances, and at the same time remove all the unwanted poles originating from $`D_C`$.
We show in the rest of this paper that this is indeed true. More precisely, we will prove the following theorem:
###### Theorem 1
There is no $`D(p)𝔇`$ such that the following three requirements are satisfied simultaneously:
$`D(p)`$ is ultralocal.
$`(D^1(p))_N𝔊^{sl}`$.
$`(D^1(p))_C`$ has no poles except when $`p_\mu =0(mod2\pi ),\mu `$.
## 4 Auxiliary Statements
We first discuss several auxiliary statements that will be needed.
### 4.1 Two-Dimensional Restrictions on the Brillouin Zone
The key ingredient that leads to the successful demonstration of Theorem 1 rests upon the realization that one should concentrate on the consequences of analyticity of $`(D^1)_N`$ at the origin of the Brillouin zone . This is indeed nontrivial because both $`(D^1)_C`$ and $`(D^1)_N`$ are functions of all Clifford components $`D^a`$, while at the same time, the former is singular and the latter is analytic at the origin. We will show that to capture the required analytic structure of the functions involved, it is sufficient to consider two-dimensional restrictions on the Brillouin zone. We thus start with the definition of such restrictions in general.
###### Definition 3
(Restriction $`\mathrm{\Sigma }^\rho `$) Let $`\rho \{1,2,\mathrm{}d\}`$ and let $`\overline{p}`$ denotes the restriction of the momentum variable $`p`$ defined through
$$\overline{p}_\mu =\{\begin{array}{cc}q_\mu ,\hfill & \text{if }\mu =1,\mathrm{},\rho \text{;}\hfill \\ 0,\hfill & \text{if }\mu =\rho +1,\mathrm{},d\hfill \end{array}$$
Map $`\mathrm{\Sigma }^\rho `$ that assigns to arbitrary function $`f(p)`$ of $`d`$ real variables a function $`\overline{f}(q)`$ of $`\rho `$ real variables $`q(q_1,q_2,\mathrm{},q_\rho )`$ through
$$\mathrm{\Sigma }^\rho \left[f(p)\right]\overline{f}(q)f(\overline{p})$$
(5)
will be referred to as restriction $`\mathrm{\Sigma }^\rho `$.
Functions with definite transformation properties under hypercubic group simplify or vanish when restricted through $`\mathrm{\Sigma }^\rho `$. The elements of $`𝔊^{sl}`$ simplify accordingly. In particular, for $`\mathrm{\Sigma }^2`$ we have the following.
###### Lemma 1
Let $`G(p)𝔊^{sl}`$, and let $`\overline{G}(q)`$ be its restriction under $`\mathrm{\Sigma }^2`$ defined through
$$\overline{G}(q)=\underset{a=1}{\overset{2^d}{}}\overline{G}^a(q)\mathrm{\Gamma }^a\overline{G}^a(q)=\mathrm{\Sigma }^2\left[G^a(p)\right]$$
Then $`\overline{G}(q)`$ can be written in the form
$$\overline{G}(q)=X(q)\mathrm{𝕀}+\underset{\mu =1}{\overset{2}{}}Y_\mu (q)\gamma _\mu +Z(q)\gamma _1\gamma _2$$
(6)
where $`X(q)`$,$`Y_\mu (q)`$ and $`Z(q)`$ are analytic functions of two real variables, periodic with $`2\pi `$, and satisfying
$$X(q_1,q_2)=X(q_1,q_2)=X(q_1,q_2)=X(q_2,q_1)$$
(7)
$$Y_1(q_1,q_2)=Y_1(q_1,q_2)=Y_1(q_1,q_2)=Y_2(q_2,q_1)$$
(8)
$$Z(q_1,q_2)=Z(q_1,q_2)=Z(q_1,q_2)=Z(q_2,q_1)$$
(9)
Lemma 1 reflects the naive expectation that $`\overline{G}(q)`$ formally looks as an element of $`𝔊^{sl}`$ in two dimensions. The proof is given in Appendix C.
The distinctive property of ultralocal elements of $`𝔊^{sl}`$ is that its Clifford components are polynomials (or are related to polynomials) in suitably chosen variables. In particular, for the restrictions under $`\mathrm{\Sigma }^2`$ we have the following result proved in Appendix D.
###### Lemma 2
Let $`G(p)`$ be an ultralocal element of $`𝔊^{sl}`$. Then the Clifford components (7-9) of its restriction $`\overline{G}(q)`$ under $`\mathrm{\Sigma }^2`$ can be written in the form
$`X(q_1,q_2)`$ $`=`$ $`P_X(\mathrm{cos}q_1,\mathrm{cos}q_2)`$
$`Y_1(q_1,q_2)`$ $`=`$ $`\mathrm{sin}q_1P_Y(\mathrm{cos}q_1,\mathrm{cos}q_2)=Y_2(q_2,q_1)`$ (10)
$`Z(q_1,q_2)`$ $`=`$ $`\mathrm{sin}q_1\mathrm{sin}q_2P_Z(\mathrm{cos}q_1,\mathrm{cos}q_2)`$
where $`P_X`$ is a symmetric polynomial in its variables, $`P_Z`$ an antisymmetric polynomial, and $`P_Y`$ a polynomial.
### 4.2 Useful Algebraic Results
While the GW property is associated with analytic properties of the Fourier transform of the propagator, the proof of Theorem 1 that we give in the next section is essentially algebraic. Lemma 3 below is a well-known result which will turn out to reflect this connection to algebra for ultralocal operators. Lemma 4 and Corollary 1 are rather interesting new results that we will need.
###### Lemma 3
Let $`P(x,y)`$, $`Q(x,y)`$ be two polynomials over the field of complex numbers such that $`Q(0,0)=0`$, and that the rational function $`R(x,y)P(x,y)/Q(x,y)`$ is analytic in some neighborhood of the origin. Then $`P`$ and $`Q`$ have common polynomial factor $`F`$, such that if we write $`P=F\stackrel{~}{P}`$, $`Q=F\stackrel{~}{Q}`$, then $`\stackrel{~}{Q}(0,0)0`$. The polynomial $`F`$ of minimal degree is unique up to a constant multiplicative factor, and depends only on $`Q`$.
###### Lemma 4
Let $`G_1(x,y)`$, $`G_2(x,y)`$ be polynomials over complex numbers, such that $`G_1(0,0)0`$, $`G_2(0,0)0`$. Let further
$$B(x,y)x(1x)G_1^2(x,y)+y(1y)G_2^2(x,y)$$
If $`F(x,y)P(x,y)=B(x,y)`$ is arbitrary polynomial factorization such that
$$F(0,0)=0P(0,0)0$$
then $`F(x,y)`$ must have at least one additional zero on the set $`𝒵_0\{(x,y):x,y\{0,1\}\}`$.
###### Corollary 1
Let $`G_1(x^2,y^2)`$, $`G_2(x^2,y^2)`$ be polynomials over complex numbers, such that $`G_1(0,0)0`$, $`G_2(0,0)0`$. Let further
$$B(x^2,y^2)x^2(1x^2)G_1^2(x^2,y^2)+y^2(1y^2)G_2^2(x^2,y^2)$$
If $`F(x,y)P(x,y)=B(x^2,y^2)`$ is arbitrary polynomial factorization such that
$$F(0,0)=0P(0,0)0$$
then $`F(x,y)`$ must have at least two additional zeroes on the set $`𝒵_1\{(x,y):x,y\{1,0,1\}\}`$.
We prove these auxiliary statements in the subsections below.<sup>4</sup><sup>4</sup>4 Reader not interested in these proofs can proceed directly to the proof of strong non-ultralocality in Sec. 5. It should be pointed out however, that the proofs of Lemma 4 and Corollary 1 are important for understanding the “heart of the matter”. It turns out that it is quite useful and natural for our purposes to use the established language of algebraic geometry. To make this paper sufficiently self-contained, we will first summarize some basic notions and state some standard results that will be relevant for our purposes.
#### 4.2.1 Preliminaries from Algebraic Geometry
We will be repeatedly concerned here with polynomials over the field of complex numbers $`\mathrm{}`$. Let us thus denote the set of such polynomials in indeterminates $`x_1,x_2,\mathrm{},x_n`$ as $`\mathrm{}[x_1,x_2,\mathrm{},x_n]`$.<sup>5</sup><sup>5</sup>5All of the discussion that follows can be carried over for polynomials (and corresponding algebraic sets) over arbitrary algebraically closed field. However, we will avoid such unnecessary generality here. The degree of arbitrary polynomial $`P`$ will be denoted as $`degP`$. Also, if some expression $`Q(x_1,\mathrm{},x_n,y_1,\mathrm{},y_m)`$ can be viewed as a power series in variables $`x_1,\mathrm{},x_n`$, then $`ord_{x_1,\mathrm{},x_n}(Q)`$ will denote the minimal degree of monomials appearing in such series. This obviously applies also if $`Q`$ is a polynomial. For example, we will work with the power series like $`\pi (t)_{i=0}^{\mathrm{}}p_it^i`$, in which case $`ord_t(\pi (t))`$ is the least $`i`$ such that $`p_i0`$.
###### Definition 4
(Affine Algebraic Curve) Let $`P\mathrm{}[x,y]`$ and $`degP>0`$. The set of points in the complex affine plane $`\mathrm{\Gamma }_P\{(x,y):P(x,y)=0\}`$ defines an affine algebraic curve.
An algebraic curve is thus a geometric figure defined by the locus of zeroes of a polynomial. Arbitrary neighborhood of any point on the curve contains infinitely many other points belonging to the curve, and there are no isolated points. If $`\mathrm{\Gamma }`$ is an algebraic curve and there exists an irreducible polynomial $`P`$ such that $`\mathrm{\Gamma }=\mathrm{\Gamma }_P`$, then $`\mathrm{\Gamma }`$ is called an irreducible algebraic curve. If $`P(x,y)`$ is reducible and $`P(x,y)=_{i=1}^NP_i^{n_i}(x,y)`$ is a factorization into irreducible factors, then obviously $`\mathrm{\Gamma }_P=\mathrm{\Gamma }_{P_1}\mathrm{\Gamma }_{P_2}\mathrm{}\mathrm{\Gamma }_{P_N}`$, and the irreducible curves $`\mathrm{\Gamma }_{P_i}`$ are referred to as the components of $`\mathrm{\Gamma }_P`$.
Introducing the geometrical viewpoint into algebraic problem frequently turns out to be quite useful. Note, for example, that the content of Lemma 4 can be interpreted in the following way. To every polynomial $`B(x,y)`$ constructed as prescribed we can assign the curve $`\mathrm{\Gamma }_B`$, and always $`𝒵_0\mathrm{\Gamma }_B`$. Similarly, to any factor $`F(x,y)`$ such that $`F(0,0)=0`$, we can assign a curve $`\mathrm{\Gamma }_F\mathrm{\Gamma }_B`$, running through the origin. The claim is that every $`\mathrm{\Gamma }_F`$ has to run through one additional point of $`𝒵_0`$. To demonstrate this result, we will need some basic tools from the theory of intersections of algebraic curves which, in turn, is more elegantly discussed in the projective plane rather than in the affine plane.
A. Projective Curves
Two parallel lines in the affine plane do not have any points in common. However, it is intuitively quite acceptable to say that they meet at the “point at infinity”. To handle such “intersections” easily, it is useful to think about algebraic curves in the projective plane.
###### Definition 5
(Projective Plane) Consider the set $`S=\mathrm{}^3\{(0,0,0)\}`$ and define two points $`𝐚,𝐛S`$ to be equivalent iff $`𝐚=\lambda 𝐛`$ for some $`\lambda \mathrm{}`$. The set of equivalence classes of $`S`$ under this equivalence relation is called the projective plane.
A point in the projective plane has thus infinitely many triples of numbers (not simultaneously zero) associated with it. Any one of these fully represents the point and is referred to as a triple of homogeneous coordinates for that point. The coordinates $`\{x,y,z\}`$ on $`\mathrm{}^3`$ are thus used to define the sets of homogeneous coordinates for the points of projective plane. The subset of projective plane consisting of points represented by homogeneous coordinates with $`z0`$ is isomorphic to the affine plane under the map
$$\varphi :(x,y,z)(\frac{x}{z},\frac{y}{z})$$
(11)
Similarly, the subsets of projective plane with $`x0`$ and $`y0`$ are isomorphic to the affine plane under analogously defined maps. However, in our discussion we will use the convention fixed by eq. (11).
To understand why projective plane is useful for dealing with points at infinity, consider an arbitrary line in the affine plane. This is specified by linear equation $`ax+by+c=0`$, where $`a,b`$ are not simultaneously zero. For definiteness, assume that $`b0`$ so that the line can be parametrized through $`x`$ as $`(x,xa/bc/b)`$. In view of map $`\varphi `$, this corresponds in the projective plane to the set of points represented by homogeneous coordinates $`(x,xa/bc/b,1)`$. Concentrating on the part of the line when $`|x|`$ is very large, we can multiply by $`b/x`$ and use the homogeneous coordinates $`(b,ac/x,b/x)`$ instead. Considering now the limiting case $`|x|\mathrm{}`$, it is natural to associate the point at infinity of the affine plane corresponding to the above line with the point in the projective plane with homogeneous coordinates $`(b,a,0)`$. Notice that this point is not part of the isomorphism realized through (11) and there are no infinities associated with this point in the projective plane. Thus, roughly speaking, the projective plane is an affine plane together with its points at infinity.
###### Definition 6
(Projective Algebraic Curve) Let $`P\mathrm{}[x,y,z]`$ be homogeneous of nonzero degree. The set of points in the complex projective plane represented by homogeneous coordinates from the set $`\mathrm{\Gamma }_P\{(x,y,z):P(x,y,z)=0\}`$ defines a projective algebraic curve.
We deal with homogeneous polynomials in projective plane because they simultaneously vanish for all homogeneous coordinates representing the same point of the projective plane. The concepts of irreducibility and component are defined analogously to the affine case, i.e. through irreducibility of the defining homogeneous polynomial.
There is a natural one to one correspondence between affine curves and projective curves that do not have $`z=0`$ (line at infinity) as a component. To see that, it is useful to consider the map $`\mathrm{\Phi }`$ which assigns to arbitrary $`P\mathrm{}[x,y]`$ of degree $`m`$, the homogeneous polynomial $`P^{}\mathrm{}[x,y,z]`$ of the same degree through
$$\mathrm{\Phi }:P(x,y)P^{}(x,y,z)z^mP(\frac{x}{z},\frac{y}{z})$$
(12)
The image $`P^{}`$ of $`P`$ under $`\mathrm{\Phi }`$ is usually referred to as the canonical homogenization of $`P`$. The following obvious theorem holds (see e.g. Ref. )
###### Theorem 2
$`\mathrm{\Phi }`$ is a bijective map (one to one and onto) between $`\mathrm{}[x,y]`$ and the set $`\mathrm{}[x,y,z]`$ of homogeneous polynomials not divisible by $`z`$. The corresponding inverse map is given by $`P^{}(x,y,z)P(x,y)P^{}(x,y,1)`$. Polynomial $`P(x,y)`$ is irreducible if and only if its image $`P^{}(x,y,z)`$ is irreducible.
In this way bijection $`\mathrm{\Phi }`$ also defines the correspondence between affine curve $`\mathrm{\Gamma }_P`$ and projective curve $`\mathrm{\Gamma }_P^{}`$, where the points are associated through map $`\varphi `$ of eq. (11). For any point $`(x,y)`$ of the affine curve $`\mathrm{\Gamma }_P`$ there is a corresponding point on projective curve $`\mathrm{\Gamma }_P^{}`$, represented for example by $`(x,y,1)`$. Conversely, for arbitrary point of the projective curve $`\mathrm{\Gamma }_P^{}`$ with homogeneous coordinates $`(x,y,z),z0`$, there is a corresponding point $`(x/z,y/z)`$ on $`\mathrm{\Gamma }_P`$. The points on $`\mathrm{\Gamma }_P^{}`$ with homogeneous coordinates $`(x,y,0)`$ do not have the corresponding points in the finite affine plane and they are the points at infinity of $`\mathrm{\Gamma }_P`$. In this sense, the projective curve $`\mathrm{\Gamma }_P^{}`$ is frequently referred to as the projective closure of $`\mathrm{\Gamma }_P`$. Conversely, $`\mathrm{\Gamma }_P`$ is an affine representative of $`\mathrm{\Gamma }_P^{}`$.
B. Intersections of Curves
Our discussion of curve intersections will focus on understanding the most basic result here, namely the Bezout’s Theorem. The logic and formalism of our approach will be almost exclusively guided by its utility in the proof of Lemma 4.
We start with the notion of intersection multiplicity, which is a positive integer assigned to any point where two algebraic curves meet. The corresponding definition is motivated by the concept of multiplicity for the root of a polynomial. Let us first consider a simple example in the affine plane. The curve $`\mathrm{\Gamma }_P`$ cut out by the polynomial $`P(x,y)=yF(x)`$ and the curve $`\mathrm{\Gamma }_L`$ with $`L(x,y)=y`$ (line) meet at some point $`(\lambda ,0)`$. In this case one can easily define the intersection multiplicity at that point as the multiplicity of $`\lambda `$ as a root of $`F`$. Knowledge of intersection multiplicity assigned this way clearly adds a useful geometrical detail on how the two curves actually meet. It is instructive to think of this definition of intersection multiplicity also in a different way. Consider the parametrization $`t(t+\lambda ,0)`$ of the line $`\mathrm{\Gamma }_L`$. Since $`t=0`$ corresponds to the point $`(\lambda ,0)`$, such parametrization is usually referred to as parametrization “at the point” $`(\lambda ,0)`$. Upon insertion in $`P(x,y)`$ we get $`P(\lambda +t,0)=F(\lambda +t)`$, and we can conclude that the intersection multiplicity of $`\mathrm{\Gamma }_P`$ and $`\mathrm{\Gamma }_L`$ at the point $`(\lambda ,0)`$ is given by $`ord_t(F(\lambda +t))=ord_t(P(\lambda +t,0))`$. We have just used an elementary fact that if $`F(x)`$ is a polynomial in $`\mathrm{}[x]`$, then the multiplicity of $`\lambda `$ as a root of $`F`$ is given by $`ord_t(F(t+\lambda ))`$.
Our aim is to generalize the above method for arbitrary curves. To do this, it is first necessary to put the concept of parametrization at a point of a curve on a firm footing. The following fundamental statement holds
###### Theorem 3
Let $`\mathrm{\Gamma }_P`$ be arbitrary projective algebraic curve, and let $`(x_0,y_0,z_0)`$ be arbitrary point of this curve. There exists a finite set of triples of functions $`(\mu _j(t),\nu _j(t),\xi _j(t))`$, analytic in the neighborhood of $`t=0`$, such that
(1) $`(\mu _j(0),\nu _j(0),\xi _j(0))=(x_0,y_0,z_0),j`$
(2) $`P(\mu _j(t),\nu _j(t),\xi _j(t))=\mathrm{\hspace{0.17em}0},j`$
(3) For every point $`(x,y,z)(x_0,y_0,z_0)`$ of $`\mathrm{\Gamma }_P`$ in a suitable neighborhood of $`(x_0,y_0,z_0)`$, there is exactly one $`j`$ and a unique value of $`t`$ such that $`(\mu _j(t),\nu _j(t),\xi _j(t))=(x,y,z)`$.
It should be emphasized that the set of triples $`(\mu _j(t),\nu _j(t),\xi _j(t))`$ of Theorem 3 is not unique. In particular, we can perform an analytic change of variables $`t=\alpha _1\tau +\alpha _2\tau ^2+\mathrm{},\alpha _10`$, to obtain an equally valid parametric representation (in $`\tau `$) satisfying conditions $`(13)`$. For fixed $`j`$, the equivalence class of such representations is usually referred to as a branch. The point $`(\mu _j(0),\nu _j(0),\xi _j(0))`$ is said to be the center of the branch, and the branch is said to be centered at that point. Completely analogous considerations hold for affine curves. As a simple example of multiple branches centered at a point, we can consider the curve defined by $`P(x,y)=y^2x^2x^3`$, which has two branches at the origin given by $`t(t,t(1+t)^{1/2})`$ and $`t(t,t(1+t)^{1/2})`$.
###### Definition 7
(Intersection Multiplicity) Let $`𝐚(x_0,y_0,z_0)`$ be the intersection point of projective algebraic curves $`\mathrm{\Gamma }_P`$ and $`\mathrm{\Gamma }_Q`$. Let further $`\mathrm{\Gamma }_P`$ have $`N`$ branches $`\gamma _j(\mu _j(t),\nu _j(t),\xi _j(t))`$ centered at $`𝐚`$. The intersection multiplicity of $`\mathrm{\Gamma }_P`$ and $`\mathrm{\Gamma }_Q`$ at point $`𝐚`$ is defined as
$$i(\mathrm{\Gamma }_P,\mathrm{\Gamma }_Q,𝐚)\underset{j=1}{\overset{N}{}}i(\gamma _j,\mathrm{\Gamma }_Q,𝐚)\underset{j=1}{\overset{N}{}}\underset{t}{ord}\left(Q(\mu _j(t),\nu _j(t),\xi _j(t))\right)$$
where $`i(\gamma _j,\mathrm{\Gamma }_Q,𝐚)`$ is the intersection multiplicity of $`\gamma _j`$ and $`\mathrm{\Gamma }_Q`$ at point $`𝐚`$.
It can be shown that the above definition of $`i(\mathrm{\Gamma }_P,\mathrm{\Gamma }_Q,𝐚)`$ is symmetric in $`\mathrm{\Gamma }_P`$, $`\mathrm{\Gamma }_Q`$, and that it is not dependent on the specific choice of parametrization from the corresponding equivalence classes defining the branches.<sup>6</sup><sup>6</sup>6Intersection multiplicity at a point is frequently introduced differently (through resultants for example) and then the equivalence to Definition 7 is shown.
We are now finally in the position to state the fundamental theorem about intersections of curves in the projective plane .
###### Theorem 4
(Bezout’s Theorem) Let $`\mathrm{\Gamma }_P`$ and $`\mathrm{\Gamma }_Q`$ be two projective algebraic curves without a common component. The total number of their intersections counted with multiplicities is given by $`i(\mathrm{\Gamma }_P,\mathrm{\Gamma }_Q)=degPdegQ`$.
#### 4.2.2 Proof of Lemma 3
Proof. Consider the factorization of $`Q`$ into irreducible factors, and splitting these factors in two groups. In particular, we label the factors that vanish at the origin as $`F_i`$, and the factors that do not vanish at the origin as $`\stackrel{~}{Q}_i`$. Using this, define $`F(x,y)_{i=1}^nF_i(x,y)`$ and $`\stackrel{~}{Q}(x,y)_{i=1}^m\stackrel{~}{Q}_i(x,y)`$. Note that while there is always at least one $`F_i`$, there doesn’t necessarily have to be any $`\stackrel{~}{Q}_i`$. In that case we define $`\stackrel{~}{Q}(x,y)1`$. We thus always have $`Q=F\stackrel{~}{Q}`$, with $`F(0,0)=0`$, and $`\stackrel{~}{Q}(0,0)0`$.
Let us look at an arbitrary but fixed $`F_i`$, defining an irreducible affine complex algebraic curve $`\mathrm{\Gamma }_{F_i}`$, running through the origin. Consider an arbitrary (complex) neighborhood of the origin, in which $`R(x,y)`$ is analytic. Any such neighborhood contains infinitely many points of $`\mathrm{\Gamma }_{F_i}`$. These points represent locations of potential singularities of $`R(x,y)`$, and thus must be canceled by zeroes of $`P(x,y)`$. However, by Bezout’s theorem, any two algebraic curves without common component have only finite number of points in common so this is only possible if $`\mathrm{\Gamma }_{F_i}`$ is the component of $`\mathrm{\Gamma }_P`$. Consequently, $`F_i`$ must be an irreducible factor of $`P`$.
By the above argument, $`F`$ is the desired common factor of $`P`$ and $`Q`$. From the unique factorization theorem it follows that $`F`$ is the only factor (up to a multiplicative constant) of minimal degree. Also, its definition only depends on $`Q`$ as claimed.
#### 4.2.3 Proof of Lemma 4
We will use the formalism described in section 4.2.1 to view Lemma 4 as the problem of algebraic geometry. In particular, we will assign to the polynomials $`G_1(x,y)`$, $`G_2(x,y)`$ and $`B(x,y)`$ their canonical homogenizations defined by the map $`\mathrm{\Phi }`$ of equation (12). In view of Theorem 2 we can then reformulate the problem equivalently in projective terms as follows.
###### Proposition 1
Let $`G_1(x,y,z),G_2(x,y,z)\mathrm{}[x,y,z]`$ be homogeneous polynomials, such that $`G_1(0,0,1)0`$, $`G_2(0,0,1)0`$. Let further
$$B(x,y,z)x(zx)G_1^2(x,y,z)+y(zy)G_2^2(x,y,z)$$
If $`F(x,y,z)P(x,y,z)=B(x,y,z)`$ is arbitrary polynomial factorization such that
$$F(0,0,1)=0P(0,0,1)0$$
then the projective algebraic curve $`\mathrm{\Gamma }_F`$ passing through $`(0,0,1)`$ must pass through at least one additional point from the set $`\stackrel{~}{𝒵}_0\{(x,y,1):x,y\{0,1\}\}`$.
Proof. For convenience, let us denote $`F_1(x,y,z)=x(zx)`$ and $`F_2(x,y,z)=y(zy)`$. We will concentrate on the intersections of $`\mathrm{\Gamma }_F`$ with $`\mathrm{\Gamma }_{F_1}`$. It is sufficient to consider factors $`F`$ that do not have common polynomial factor with $`F_1`$. Indeed, $`x`$ can not be a factor of $`B`$, and hence it can not be a factor of $`F`$. On the other hand, $`zx`$ can be a factor of $`F`$ but if it is, than the claim of the proposition is obviously true.
(I) Since $`F_1`$ and $`F`$ are mutually coprime, it follows from Bezout’s theorem that the total number of their intersections $`i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1})`$ counted with multiplicities is even.
(II) Let $`𝐚`$ be the intersection point of $`\mathrm{\Gamma }_F`$ and $`\mathrm{\Gamma }_{F_1}`$ such that $`𝐚\mathrm{\Gamma }_{F_2}`$. Then the intersection multiplicity $`i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1},𝐚)`$ is even. Indeed, let $`\gamma (\mu (t),\nu (t),\xi (t))`$ be arbitrary branch of $`\mathrm{\Gamma }_F`$ centered at $`𝐚`$. Then $`B(\mu (t),\nu (t),\xi (t))=0`$, or equivalently
$$F_1F_2G_1^2=F_2^2G_2^2\text{for}(x,y,z)(\mu (t),\nu (t),\xi (t))$$
(13)
Because of the squares, we can conclude that upon the substitution, $`ord_t(F_1F_2)`$ is even. At the same time, since $`𝐚\mathrm{\Gamma }_{F_2}`$, we have $`F_2(\mu (0),\nu (0),\xi (0))0`$, and thus $`ord_t(F_2)=0`$. But $`ord_t(F_1F_2)=ord_t(F_1)+ord_t(F_2)`$, and so $`ord_t(F_1)=i(\gamma ,\mathrm{\Gamma }_{F_1},𝐚)`$ is a positive even integer. Since this is true for all branches $`\gamma `$ of $`\mathrm{\Gamma }_F`$ centered at $`𝐚`$, we can conclude that $`i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1},𝐚)`$ is even as claimed.
(III) If $`𝐚=(0,0,1)`$, then $`i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1},𝐚)=1`$. Indeed, $`\mathrm{\Gamma }_{F_1}`$ has a single branch centered at $`𝐚`$, namely the line that can be parametrized as $`\lambda =(0,t,1)`$. Since $`P(0,0,1)0`$ and $`G_2(0,0,1)0`$ we have
$$i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1},𝐚_1)=\underset{t}{ord}\left(F(0,t,1)\right)=\underset{t}{ord}\left(F(0,t,1)P(0,t,1)\right)=\underset{t}{ord}\left(t(1t)G_2^2(0,t,1)\right)=\mathrm{\hspace{0.17em}1}$$
For (I-III) to be satisfied simultaneously, $`\mathrm{\Gamma }_F`$ and $`\mathrm{\Gamma }_{F_1}`$ must necessarily intersect at one additional point from the set of points where $`\mathrm{\Gamma }_{F_1}`$ and $`\mathrm{\Gamma }_{F_2}`$ meet. However, this set is precisely represented by the set of homogeneous coordinates from $`\stackrel{~}{𝒵}_0`$. Consequently, $`\mathrm{\Gamma }_F`$ has to pass through one additional point from $`\stackrel{~}{𝒵}_0`$ as claimed.
#### 4.2.4 Proof of Corollary 1
Proof. Let us consider an arbitrary but fixed polynomial $`B(x^2,y^2)`$ of the prescribed form, and the set of all factorizations $`B(x^2,y^2)=F(x,y)P(x,y)`$ with required properties. It will be sufficient to prove the statement for the case when $`F(x,y)`$ does not contain any nontrivial polynomial factor with nonzero constant term. Indeed, if $`F(x,y)`$ has such a factor $`P_1(x,y)`$, then redefining $`F(x,y)F(x,y)/P_1(x,y)`$, $`P(x,y)P(x,y)P_1(x,y)`$ can only make the set of zeroes for new $`F(x,y)`$ smaller. As a consequence of the unique factorization theorem, the factors $`F`$, $`P`$ are unique up to a trivial rescaling by complex number if this additional condition is imposed.
Assuming the above, consider the transformation $`xx`$, which does not change $`B(x^2,y^2)`$. Since all possible nontrivial irreducible factors $`F_i(x,y)`$ of $`F(x,y)`$ satisfy $`F_i(0,0)=0`$, while all irreducible factors $`P_i(x,y)`$ of $`P(x,y)`$ have $`P_i(0,0)0`$, it follows from the unique factorization theorem that $`F(x,y)F(x,y)`$ and $`P(x,y)P(x,y)`$ under the above reflection. The same is true under $`yy`$ and, as a result, $`F=F(x^2,y^2)`$, $`P=P(x^2,y^2)`$. This implies that we can change the variables $`x^2x`$, $`y^2y`$ in the whole problem. In new variables, we have
$$B(x,y)x(1x)G_1^2(x,y)+y(1y)G_2^2(x,y)=F(x,y)P(x,y)$$
with properties at the origin preserved. However, according to Lemma 4, $`F(x,y)`$ must have an additional zero on the set $`𝒵_0\{(x,y):x,y\{0,1\}\}`$. Consequently, $`F(x^2,y^2)`$ must vanish for at least two additional points from the set $`𝒵_1\{(x,y):x,y\{1,0,1\}\}`$ which completes the proof.
## 5 Proof of Strong Non-Ultralocality (Theorem 1)
We now use the formalism of Sec. 2 and auxiliary statements of Sec. 4 to prove “strong non-ultralocality” of GW actions (Theorem 1).
Proof. We will proceed by contradiction, and thus first assume that there actually exists an element $`D(p)𝔇`$ such that the requirements $`(\alpha \gamma )`$ are satisfied. To this $`D(p)`$ we will assign its restriction $`\overline{D}(q)`$ under $`\mathrm{\Sigma }^2`$ which, according to Lemma 1, has the form $`(q(q_1,q_2))`$
$$\overline{D}(q)=\overline{A}(q)\mathrm{𝕀}+i\underset{\mu =1}{\overset{2}{}}\overline{B}_\mu (q)\gamma _\mu +\overline{C}(q)\gamma _1\gamma _2$$
(14)
with the corresponding propagator being
$$\overline{D}^1=\frac{\overline{A}}{\overline{Q}}\mathrm{𝕀}i\underset{\mu =1}{\overset{2}{}}\frac{\overline{B}_\mu }{\overline{Q}}\gamma _\mu \frac{\overline{C}}{\overline{Q}}\gamma _1\gamma _2\overline{Q}\overline{A}^2+\overline{B}_\mu \overline{B}_\mu +\overline{C}^2$$
(15)
Note that according to conditions $`(\beta )`$ and $`(\gamma )`$, no Clifford component of $`\overline{D}^1`$ has a pole away from the origin of the Brillouin zone. Consequently, the function
$$\left(\frac{\overline{A}}{\overline{Q}}\right)^2+\left(\frac{\overline{B}_\mu }{\overline{Q}}\right)\left(\frac{\overline{B}_\mu }{\overline{Q}}\right)+\left(\frac{\overline{C}}{\overline{Q}}\right)^2=\frac{\mathrm{\hspace{0.33em}1}}{\overline{Q}}$$
(16)
has also no poles away from the origin.
Since $`\overline{D}(q)`$ is ultralocal, we can use Lemma 2, to conclude that its Clifford components have the following structure
$`\overline{A}(q_1,q_2)`$ $`=`$ $`P_{\overline{A}}(\mathrm{cos}q_1,\mathrm{cos}q_2)`$
$`\overline{B}_1(q_1,q_2)`$ $`=`$ $`\mathrm{sin}q_1P_{\overline{B}}(\mathrm{cos}q_1,\mathrm{cos}q_2)=\overline{B}_2(q_2,q_1)`$ (17)
$`\overline{C}(q_1,q_2)`$ $`=`$ $`\mathrm{sin}q_1\mathrm{sin}q_2P_{\overline{C}}(\mathrm{cos}q_1,\mathrm{cos}q_2)`$
where $`P_{\overline{A}}`$ is a symmetric polynomial, $`P_{\overline{C}}`$ an antisymmetric polynomial, and $`P_{\overline{B}}`$ a polynomial without definite symmetry properties. The classical continuum limit (4) requires that $`P_{\overline{A}}(1,1)=0`$, $`P_{\overline{B}}(1,1)=1`$, and also $`P_{\overline{C}}(1,1)=0`$ by virtue of its antisymmetry.
The GW property $`(\beta )`$ demands that the Clifford components of chirally nonsymmetric part of the propagator ($`\overline{A}/\overline{Q}`$ and $`\overline{C}/\overline{Q}`$) are analytic in some complex region containing the (real) Brillouin zone. We will concentrate on the consequences of analyticity in the vicinity of the origin. To this end, it is convenient to introduce a change of variables such as
$$x=\mathrm{sin}\frac{q_1}{2}y=\mathrm{sin}\frac{q_2}{2}$$
(18)
which is invertible in the vicinity of the origin, maps the real region $`[\pi ,\pi ]\times [\pi ,\pi ]`$ onto the square $`[1,1]\times [1,1]`$, and preserves analyticity in new variables except possibly on the boundary. Using (5), we have
$`\overline{A}(q_1,q_2)`$ $`=`$ $`A(x^2,y^2)`$
$`\overline{B}_1(q_1,q_2)`$ $`=`$ $`x\sqrt{1x^2}G(x^2,y^2)=\overline{B}_2(q_2,q_1)`$ (19)
$`\overline{C}(q_1,q_2)`$ $`=`$ $`\sqrt{1x^2}\sqrt{1y^2}xyC(x^2,y^2)`$
where $`A`$ is a symmetric polynomial such that $`A(0,0)=0`$, $`C`$ is an antisymmetric polynomial, and $`G`$ is a polynomial such that $`G(0,0)=2`$. It is also convenient to introduce symmetric polynomials
$$B(x^2,y^2)x^2(1x^2)G^2(x^2,y^2)+y^2(1y^2)G^2(y^2,x^2)$$
(20)
$$Q(x^2,y^2)A^2(x^2,y^2)+B(x^2,y^2)+x^2y^2(1x^2)(1y^2)C^2(x^2,y^2)$$
(21)
corresponding to $`\overline{B}_\mu \overline{B}_\mu `$, and $`\overline{Q}`$ respectively. With this notation, the GW property $`(\beta )`$ implies that the functions
$$R_1(x^2,y^2)\frac{A(x^2,y^2)}{Q(x^2,y^2)}R_2(x,y)\sqrt{1x^2}\sqrt{1y^2}xy\frac{C(x^2,y^2)}{Q(x^2,y^2)}$$
are analytic in the vicinity of the origin.
Let us first consider $`R_1`$ which is a rational function whose denominator vanishes at the origin. Consequently, according to Lemma 3, the polynomials $`A`$ and $`Q`$ have a common polynomial factor $`F(x,y)`$ with zero at the origin, so that the possible singularity is removed. Let us fix $`F`$ to be of minimal degree, and thus unique up to a constant multiplicative factor. Next, consider $`R_2`$. Since $`\sqrt{1x^2}`$ and $`\sqrt{1y^2}`$ are analytic and nonzero near the origin, the rational function $`xyC/Q`$ is also analytic and, moreover, the same denominator $`Q`$ is involved as in $`R_1`$. Consequently, the polynomial $`F`$ divides $`xyC`$. In summary, we thus have that $`F`$ divides $`A`$,$`xyC`$ and $`Q`$. Hence, from the form (21) of $`Q`$ it follows that $`F`$ also divides $`B`$.
According to Corollary 1, an intriguing property of any polynomial $`B(x^2,y^2)`$ of the form (20) is that if it is divisible by polynomial $`F(x,y)`$ vanishing at the origin, then $`F`$ has at least two additional zeroes on the set of points $`𝒵_1=\{(x,y):x,y\{1,0,1\}\}`$. Since $`F`$ divides $`Q`$, it follows that $`Q`$ also has these zeroes on $`𝒵_1`$. Returning back to momentum variables, we thus came to the conclusion that in addition to the origin of the Brillouin zone, the function $`\overline{Q}(q)`$ of equation (15) has at least two additional zeroes on the set $`\stackrel{~}{𝒵}_1=\{(q_1,q_2):q_1,q_2\{\pi ,0,\pi \}\}`$. However, this contradicts the conclusion of (16) that $`1/\overline{Q}`$ has no poles away from the origin of the Brillouin zone. That completes the proof.
We have thus demonstrated that every ultralocal GW action with hypercubic symmetries has a doubler. The argument given above can be refined using the statement of Proposition 3, formulated and proved in Appendix E. This implies that there is always a doubler at the corner $`(\pi ,\pi )`$ of the restricted Brillouin zone.
## 6 Conclusions
As a result of several breakthroughs in the last decade, lattice field theory appears to have finally reached the stage when it can deal with all the important symmetries relevant in particle physics. In particular, it is now conceptually quite clear how to formulate at least vectorlike lattice models simultaneously possessing the lattice counterparts of gauge symmetry (Wilson gauge symmetry), Poincaré symmetry (symmetries of hypercubic lattice) and chiral symmetry (GWL symmetry), so that the continuum limit with desired field-theoretic properties can be taken comfortably. The path of developments leading to this point essentially coincides with the attempts to better understand lattice fermionic actions. Fig. 1 is a graphical representation of the basic knowledge that we now have , including the result on strong non-ultralocality demonstrated here.
The base set $`A`$ in Fig. 1 represents all actions quadratic in fermionic variables, that have “easy symmetries” (gauge and hypercubic symmetries), and proper classical continuum limit. In other words, it is the set of fermionic actions described by lattice Dirac kernels that are covariant under symmetry transformations of hypercubic lattice, gauge covariant, and with classical limit corresponding to continuum Dirac operator. The base set $`A`$ is split into three parts $`A=A^UA^LA^N`$, representing the actions that are ultralocal, local but not ultralocal, and nonlocal respectively. Highlighted is also the subset of actions without doubling of species. Obviously, of prime interest for physics applications is the set $`A^{ND}(A^UA^L)`$.
As a consequence of Nielsen-Ninomiya theorem, the relation of the subset $`A^C`$ of actions with naive chiral symmetry to the above defined sets is as indicated. While it is quite easy to construct nonlocal elements of $`A^C`$ without doublers, there is no intersection of $`A^C`$ with $`A^{ND}`$ on the local part of the diagram. The suggestion of Ginsparg and Wilson was that there might be a larger set of actions contained in $`A`$, respecting the chiral dynamics properly. Recent results in the field confirmed this idea and, more importantly, lead to the conclusion that fermion doubling is not a definite property of local GW actions, i.e. $`A^LA^{GW}A^{ND}\mathrm{}`$ (see the filled area in Fig. 1). However, as a consequence of strong non-ultralocality, doubling is a definite property of ultralocal GW actions and $`A^UA^{GW}A^{ND}=\mathrm{}`$. Unless the set of actions with acceptable chiral dynamics can be further enlarged, non-ultralocality can thus be viewed as a necessary condition to reconcile chiral dynamics with proper anomaly structure in lattice gauge theories respecting the symmetries of the hypercubic lattice.
We would also like to stress that weak non-ultralocality and strong non-ultralocality are two independent statements in the sense that one does not follow from the other. In this respect it should be pointed out that while non-ultralocality of GWL transformations applies regardless of doubling and does not apply for actions with naive chiral symmetry, non-ultralocality of GW actions is only necessary for doubler-free theories and actions with naive chiral symmetry represent no exception for this case.
Acknowledgements: I. H. and C. T. B. are grateful to William Fulton for his interest, help and support regarding Lemma 4. R. M. would like to acknowledge useful input and conversations with Dan Burghelea. I. H. benefited from communications with Pavel Bóna, Martin Niepel, František Marko, Arthur Mattuck and Andrei Rapinchuk, as well as from many pleasant discussions with Ziad Maassarani and Hank Thacker. We also thank Ziad for reading the manuscript.
## Appendix A Local Symmetric Operators
The elements of the set $`𝔊^{sl}`$ satisfy three requirements defined below:
###### Definition 8
(Locality) Operator $`𝐆`$ is said to be local if there are positive real constants $`c`$, $`\delta `$ such that all its Clifford components $`𝐆^a`$ satisfy
$$|𝐆_{m,n}^a|<ce^{\delta |mn|}m,n$$
Here $`|mn|`$ denotes the Euclidean norm of $`mn`$.
###### Definition 9
(Translation Invariance) Operator $`𝐆`$ is said to be translationally invariant if all its Clifford components $`𝐆^a`$ satisfy
$$𝐆_{m,n}^a=𝐆_{0,nm}^ag_{nm}^am,n$$
###### Definition 10
(Hypercubic Invariance) Let $``$ be an element of the hypercubic group in defining representation and $`H`$ the corresponding element of the representation induced on hypercubic group by spinorial representation of $`O(d)`$. Operator $`𝐆`$ is said to have hypercubic invariance if for arbitrary $``$, $`m`$, $`n`$ we have
$$𝐆_{n,m}=H^1𝐆_{n,m}H$$
In the Fourier space, we can directly define:
###### Definition 11
(Set $`𝔊^{sl}`$) Let $`G^a(p),a=1,2,\mathrm{}2^d`$, are the complex valued functions of real variables $`p_\mu `$, and let $`G(p)`$ be the corresponding matrix function constructed as in Eq. (3). We say that $`G(p)`$ belongs to the set $`𝔊^{sl}`$ if:
Every $`G^a(p)`$ is an analytic function with period $`2\pi `$ in all $`p_\mu `$.
For arbitrary hypercubic transformation $``$ it is true identically that
$$G(p)=\underset{a=1}{\overset{2^d}{}}G^a(p)\mathrm{\Gamma }^a=\underset{a=1}{\overset{2^d}{}}G^a(p)H^1\mathrm{\Gamma }^aH$$
(22)
Arbitrary hypercubic transformation $``$ can be decomposed into products of reflections of single axis ($`_\mu `$) and exchanges of two different axis ($`𝒳_{\mu \nu }`$). Transformation properties of the elements of the Clifford basis are determined by the fact that $`\gamma _\mu `$ transforms as a vector. In particular
$$R_\nu ^1\gamma _\mu R_\nu =\{\begin{array}{cc}\gamma _\mu ,\hfill & \text{if }\mu =\nu \text{;}\hfill \\ \gamma _\mu ,\hfill & \text{if }\mu \nu \hfill \end{array}$$
and
$$X_{\rho \sigma }^1\gamma _\mu X_{\rho \sigma }=\{\begin{array}{cc}\gamma _\sigma ,\hfill & \text{if }\mu =\rho \text{;}\hfill \\ \gamma _\rho ,\hfill & \text{if }\mu =\sigma \text{;}\hfill \\ \gamma _\mu ,\hfill & \text{otherwise}\hfill \end{array}$$
where $`R_\mu ,X_{\mu \nu }`$ are the spinorial representations of $`_\mu ,𝒳_{\mu \nu }`$. The elements of the Clifford basis naturally split into groups with definite transformation properties and the hypercubic symmetry thus translates into definite algebraic requirements on functions $`G^a(p)`$.
## Appendix B Locality and Analyticity
###### Proposition 2
Let $`G(p)`$ be a complex function of single variable periodic with $`2\pi `$, and forming a Fourier pair with the sequence of complex numbers $`\{g_n;n=0,\pm 1,\pm 2\mathrm{}\}`$, i.e.
$$G(p)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}g_ne^{inp}g_0+g_1e^{ip}+g_1e^{ip}+g_2e^{i\mathrm{\hspace{0.17em}2}p}+g_2e^{i\mathrm{\hspace{0.17em}2}p}+\mathrm{}$$
and
$$g_n=\frac{1}{2\pi }_\pi ^\pi 𝑑pe^{inp}G(p)$$
Then the following two statements are equivalent:
(a) There exist $`c>0`$,$`\delta >0`$ such that $`g_n<ce^{\delta |n|},n`$.
(b) There exists $`\mathrm{\Delta }>0`$, such that $`G(p)`$ is analytic in the complex strip $`Im(p)<\mathrm{\Delta }`$.
Proof. We first show that (a) implies (b). Consider the change of variables $`z=e^{ip}`$. Then $`G=g(e^{ip})`$ where the function $`g`$ is defined by
$$g(z)=g_0+g_1z+g_1z^1+g_2z^2+g_2z^2+\mathrm{}$$
We split this up as $`g(z)=h(z)+k(z)`$ where
$$h(z)=g_0+g_1z+g_2z^2+\mathrm{}k(z)=g_1z^1+g_2z^2+\mathrm{}$$
Condition (a) simply says that $`|g_n|<c\beta ^{|n|}`$ where $`\beta >1`$, implying that $`h(z)`$ has a radius of convergence strictly greater than $`1`$. Similarly, if we further change the variable $`z=1/t`$ in $`k(z)`$, so that $`l(t)=k(1/t)`$, we can infer that $`l(t)`$ too has a radius of convergence strictly greater than $`1`$. The above conclusions prove that $`g(z)`$ is analytic in an annulus containing the unit circle of the complex plane. Consequently, $`G(p)=g(e^{ip})`$ is a composition of two analytic maps, implying analyticity in the complex strip as desired in (b).
Next we show that (b) implies (a). Consider arbitrary $`0<\rho <\mathrm{\Delta }`$. We will integrate in the complex plane along the rectangular contour $`𝒞`$ with the following line segments: $`𝒞_1`$ from point $`(\pi ,0)`$ to point $`(\pi ,0)`$, $`𝒞_2`$ from $`(\pi ,0)`$ to $`(\pi ,\rho )`$, $`𝒞_3`$ from $`(\pi ,\rho )`$ to $`(\pi ,\rho )`$, and $`𝒞_4`$ from $`(\pi ,\rho )`$ to $`(\pi ,0)`$.
Let us assume that $`n0`$. According to Cauchy theorem, we have
$$_𝒞𝑑pe^{inp}G(p)=\mathrm{\hspace{0.33em}0}=_{𝒞_1}+_{𝒞_2}+_{𝒞_3}+_{𝒞_4}$$
Since $`G(p)`$ is periodic with $`2\pi `$, the contributions from the integrals along $`𝒞_2`$ and $`𝒞_4`$ will cancel each other and we thus obtain
$`{\displaystyle _{𝒞_1}}𝑑pe^{inp}G(p)`$ $`=`$ $`{\displaystyle _\pi ^\pi }𝑑te^{int}G(t)={\displaystyle _{𝒞_3}}𝑑pe^{inp}G(p)=`$
$`=`$ $`{\displaystyle _\pi ^\pi }𝑑te^{in(ti\rho )}G(ti\rho )=e^{n\rho }{\displaystyle _\pi ^\pi }𝑑te^{int}G(ti\rho )`$
Finally, $`G(p)`$ is bounded on the path of the last integral due to analyticity, and we thus have
$$g_n=\frac{1}{2\pi }_\pi ^\pi 𝑑te^{int}G(t)e^{n\rho }\underset{t[\pi ,\pi ]}{\mathrm{max}}G(ti\rho )$$
Similarly if $`n<0`$, then we will use the rectangular integration contour in the upper half of the complex plane, yielding an analogous bound. Together, this then implies (a) as claimed.
## Appendix C Proof of Lemma 1
Proof. Let us denote the set of indices $`u^\rho \{1,2\mathrm{}\rho \}`$ for arbitrary positive integer $`\rho `$. Clifford basis can be subdivided into non–intersecting subsets $`\mathrm{\Gamma }=_j\mathrm{\Gamma }_{(j)}`$, where $`\mathrm{\Gamma }_{(j)}`$, $`j=0,1,\mathrm{},d`$, contains the elements that can be written as the product of $`j`$ gamma-matrices. In particular, $`\mathrm{\Gamma }_{(0)}=\{\mathrm{𝕀}\}`$, $`\mathrm{\Gamma }_{(1)}=\{\gamma _\mu ,\mu u^d\}`$, and so on. With the appropriate convention on ordering of gamma–matrices in the definition of $`\mathrm{\Gamma }^a`$, we can then rewrite the Clifford decomposition of $`G(p)`$ in the form
$$G(p)=\underset{j=0}{\overset{d}{}}\underset{\genfrac{}{}{0pt}{}{\mu _1,\mu _2\mathrm{}\mu _j}{\mu _1<\mu _2\mathrm{}<\mu _j}}{}F_{\mu _1,\mu _2\mathrm{}\mu _j}(p)\gamma _{\mu _1}\gamma _{\mu _2}\mathrm{}\gamma _{\mu _j},\mu _iu^d$$
(23)
We will concentrate on contributions to $`\overline{G}(q)`$ originating from subsets $`\mathrm{\Gamma }_{(j)}`$, where $`j2`$. Consider a single term in decomposition (23) from this group, specified by the set of indices $`v\{\mu _1,\mu _2\mathrm{}\mu _j\}u^2`$. Then there exists element $`\mu v`$, such that $`\mu u^2`$. Under reflection $`_\mu `$ through the corresponding axis, we have $`R_\mu ^1\gamma _{\mu _1}\gamma _{\mu _2}\mathrm{}\gamma _{\mu _j}R_\mu =\gamma _{\mu _1}\gamma _{\mu _2}\mathrm{}\gamma _{\mu _j}.`$ Hypercubic symmetry of $`G(p)`$ then requires that $`F_{\mu _1,\mu _2\mathrm{}\mu _j}(_\mu p)=F_{\mu _1,\mu _2\mathrm{}\mu _j}(p).`$ However, since $`\mu u^2`$, the restricted variable $`\overline{p}`$ under $`\mathrm{\Sigma }^2`$ satisfies $`_\mu \overline{p}=\overline{p}`$, and hence
$$\overline{F}_{\mu _1,\mu _2\mathrm{}\mu _j}(q)F_{\mu _1,\mu _2\mathrm{}\mu _j}(\overline{p})=F_{\mu _1,\mu _2\mathrm{}\mu _j}(\overline{p})=\mathrm{\hspace{0.33em}0}$$
Consequently, the only Clifford element contributing from this group is $`\gamma _1\gamma _2`$ (when $`v=u^2`$), and the form (6) follows.
Analyticity, periodicity and relations (7-9) for functions $`X(q),Y_\mu (q),Z(q)`$ follow from corresponding properties of unrestricted operator.
## Appendix D Proof of Lemma 2
Proof. Since $`G(p)`$ is ultralocal, the Clifford components of $`\overline{G}(q)`$ have finite number of Fourier terms. Then there exists a non-negative integer $`N`$, such that when grouping together the Fourier terms related by reflection properties in (7-9) of Lemma 1, the Fourier expansions can be written in the form
$`X(q_1,q_2)`$ $`=`$ $`{\displaystyle \underset{n_1=0}{\overset{N}{}}}{\displaystyle \underset{n_2=0}{\overset{N}{}}}x_{n_1n_2}\mathrm{cos}n_1q_1\mathrm{cos}n_2q_2`$
$`Y_1(q_1,q_2)`$ $`=`$ $`{\displaystyle \underset{n_1=1}{\overset{N}{}}}{\displaystyle \underset{n_2=0}{\overset{N}{}}}y_{n_1n_2}\mathrm{sin}n_1q_1\mathrm{cos}n_2q_2=Y_2(q_2,q_1)`$
$`Z(q_1,q_2)`$ $`=`$ $`{\displaystyle \underset{n_1=1}{\overset{N}{}}}{\displaystyle \underset{n_2=1}{\overset{N}{}}}z_{n_1n_2}\mathrm{sin}n_1q_1\mathrm{sin}n_2q_2`$
Furthermore, the exchange properties in (7) and (9) imply that $`x_{n_1,n_2}=x_{n_2,n_1}`$, while $`z_{n_1,n_2}=z_{n_2,n_1}`$. Using the formulas for trigonometric functions of multiple arguments, namely
$`\mathrm{sin}nx=\mathrm{sin}x[\mathrm{\hspace{0.17em}2}^{n1}\mathrm{cos}^{n1}x`$ $``$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{n2}{1}}\right)2^{n3}\mathrm{cos}^{n3}x+\left({\displaystyle \genfrac{}{}{0pt}{}{n3}{2}}\right)2^{n5}\mathrm{cos}^{n5}x`$
$``$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{n4}{3}}\right)2^{n7}\mathrm{cos}^{n7}x+\mathrm{}]`$
and
$`\mathrm{cos}nx=\mathrm{\hspace{0.33em}2}^{n1}\mathrm{cos}^nx`$ $``$ $`{\displaystyle \frac{n}{1}}\mathrm{\hspace{0.17em}2}^{n3}\mathrm{cos}^{n2}x+{\displaystyle \frac{n}{2}}\left({\displaystyle \genfrac{}{}{0pt}{}{n3}{1}}\right)2^{n5}\mathrm{cos}^{n4}x`$
$``$ $`{\displaystyle \frac{n}{3}}\left({\displaystyle \genfrac{}{}{0pt}{}{n4}{2}}\right)2^{n7}\mathrm{cos}^{n6}x+\mathrm{}`$
the forms (10) directly follow, together with the exchange symmetry properties.
## Appendix E Refined Position of the Doubler
###### Proposition 3
Let $`G(x,y,z)\mathrm{}[x,y,z]`$ be homogeneous polynomial, such that $`G(0,0,1)0`$. Let further
$$B(x,y,z)x(zx)G^2(x,y,z)+y(zy)G^2(y,x,z)$$
If $`F(x,y,z)P(x,y,z)=B(x,y,z)`$ is arbitrary polynomial factorization such that
$$F(0,0,1)=0P(0,0,1)0$$
then the projective algebraic curve $`\mathrm{\Gamma }_F`$ passing through $`(0,0,1)`$ must also pass through $`(1,1,1)`$.
Proof. We are dealing with the special case of Proposition 1, given by $`G_1(x,y)=G_2(y,x)=G(x,y)`$, and all the arguments of the the proof in section 4.2.3 are valid here as well. Using the notation defined there and setting $`𝐚_1=(0,0,1)\stackrel{~}{𝒵}_0`$ we have in particular
$$\underset{𝐚\stackrel{~}{𝒵}_0}{}i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1},𝐚)=\mathrm{\hspace{0.33em}0}(mod2)i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1},𝐚_1)=1$$
(24)
We will show below that for $`𝐚_2=(1,0,1),𝐚_3=(0,1,1)`$ it is now true in addition that
$$i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1},𝐚_2)+i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1},𝐚_3)=\mathrm{\hspace{0.17em}0}(mod2)$$
(25)
However, denoting $`𝐚_4=(1,1,1)`$ we have $`\stackrel{~}{𝒵}_0=\{𝐚_1,𝐚_2,𝐚_3,𝐚_4\}`$ and so (24), (25) imply that $`i(\mathrm{\Gamma }_F,\mathrm{\Gamma }_{F_1},𝐚_4)=1(mod2)`$. That proves the claim of the proposition.
To show (25), it is sufficient to concentrate on irreducible $`\mathrm{\Gamma }_F`$. This is because such $`\mathrm{\Gamma }_F`$ is unique and all the reducible ones will thus contain it. The uniqueness of irreducible $`\mathrm{\Gamma }_F`$ follows from the fact that if there were at least two, then $`ord_{x,y}(B(x,y,z))`$ would also have to be at least two, which is not the case. This also means that the corresponding $`F(x,y,z)`$ is symmetric in $`x,y`$. Indeed, if it were not, then the polynomial $`F(y,x,z)`$ would define another component of $`\mathrm{\Gamma }_B`$ running through the point $`(0,0,1)`$ thus contradicting the uniqueness.
Assuming the above, consider arbitrary branch $`\gamma =(\mu (t),\nu (t),1)`$ of $`\mathrm{\Gamma }_F`$, centered at $`𝐚_2`$. Due to symmetry in $`x,y`$ there is a corresponding branch $`\gamma ^{}=(\nu (t),\mu (t),1)`$ centered at $`𝐚_3`$. Using the fact that $`ord_t(\mu )=ord_t(1\nu )=0`$ and $`ord_t(\nu )>0,ord_t(1\mu )>0`$ we have
$$i(\gamma ,\mathrm{\Gamma }_{F_1},𝐚_2)=\underset{t}{ord}(\mu (1\mu ))=\underset{t}{ord}(1\mu )i(\gamma ^{},\mathrm{\Gamma }_{F_1},𝐚_3)=\underset{t}{ord}(\nu (1\nu ))=\underset{t}{ord}(\nu )$$
while at the same time
$$\underset{t}{ord}(1\mu )+\underset{t}{ord}(\nu )=\underset{t}{ord}(\mu (1\mu )\nu (1\nu ))=\underset{t}{ord}(F_1F_2)=\mathrm{\hspace{0.17em}0}(mod2)$$
The last two equalities relate to substituting $`(x,y,z)(\mu (t),\nu (t),1)`$ in $`F_1,F_2`$, and using eq. (13) which is valid for arbitrary branch $`\gamma `$ of $`\mathrm{\Gamma }_F`$. It follows that $`i(\gamma ,\mathrm{\Gamma }_{F_1},𝐚_2)+i(\gamma ^{},\mathrm{\Gamma }_{F_1},𝐚_3)`$ is even for arbitrary pair $`(\gamma ,\gamma ^{})`$, implying (25).
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# The Rotation Curve of the Large Magellanic Cloud and the Implications for Microlensing
## 1 Introduction
If the Galactic dark halo were composed entirely of massive compact halo objects (MACHOs), the instantaneous probablity of microlensing (the “optical depth”) towards the Large Magellanic Cloud (LMC) would have been $`\tau 5\times 10^7`$ (Paczýnski 1986). The most recent observational result for MACHOs with masses up to about 1 $`M_{}`$ is $`\tau =1.2_{0.3}^{+0.4}\times 10^7`$, which comes from an analysis of 5.7 years of LMC microlensing survey data (Alcock et al. 2000). Thus an all-MACHO Galactic dark halo is ruled out. However, the observed optical depth is significantly higher than the estimate for known Galactic and LMC stellar populations, $`\tau \stackrel{<}{}4\times 10^8`$ (Alcock et al. 2000). The excess LMC microlensing signal may be telling us the MACHO fraction of the Galactic dark halo or revealing gaps in our understanding of the essential structure of the LMC and Galaxy.
Critical discussions of the microlensing optical depth from known stellar populations have developed along two lines. One debate has been about the possible existence of a “new” stellar population which would account for the LMC microlensing result. Some example suggestions include an intervening dwarf galaxy or tidal stream (Zhao 1998; Zaritsky & Lin 1997, Zaritsky et al. 1999), and a very warped Galactic disk (Evans et al. 1998). These hypotheses have been tested, and in some cases ruled out (Alcock et al. 1997; Beaulieu & Sackett 1998; Bennett 1998; Gould 1998, 1999; Gyuk, Flynn, & Evans 1999; Ibata, Germaint, & Beaulieu 1999; Johnston 1998). The second debate has been about the importance of LMC star-star “self-lensing” (Sahu 1994; Wu 1994). In particular, although the LMC is well represented by an exponential disk model which often serves as the basis for the self-lensing optical depth calculation, it may exhibit some important detailed structures. Our paper is motivated by this LMC self-lensing problem.
The line-of-sight velocity dispersion of LMC stars yields a strong constraint on star-star self-lensing from the virialized LMC disk: $`\tau \stackrel{<}{}1\times 10^8`$ (Gould 1995). Gould’s elegant limit is probably uncertain by no more than a factor of $``$2. However, several LMC models have been devised which would increase the self-lensing optical depth over the Gould limit. One example of a detailed structure that might increase the self-lensing optical depth is a highly-inclined and flared LMC disk (Zhao 1999). It has also been suggested that the oldest LMC disk stars have a very large characteristic scale height, or that the LMC harbors an as yet unseen but massive stellar spheroid<sup>1</sup><sup>1</sup>1The distinction between a massive stellar spheroid or very thick disk and an LMC dark halo is a matter of the characteristic mass to light ratios (i.e. M/L $``$ 2-4 for the former, and a very large M/L for the latter), and the density profiles. (e.g. Aubourg et al. 1999; Salati et al. 1999; Evans & Kerins 2000). Finally, it is possible that a nonvirialized stellar component (i.e., a shroud of tidal debris; Weinberg 2000) acts to increase the self-lensing optical depth (Zhao 1998).
Kinematic studies play a critical role in testing these various proposed LMC models. For example, the flare of the LMC disk may be inferred from its velocity dispersion at different radii, and the existence of a non-equilibrium stellar component lying near the LMC might be proven with a large kinematic survey (Zhao 1999; Graff et al. 1999; see also Zaritsky et al. 1999; Ibata et al. 1999). In addition, one could search directly for an LMC stellar population with a velocity dispersion of $``$50 km s<sup>-1</sup>, the prediction for a spheroid or very thick disk. These latter components must be identified before their importance to the self-lensing optical depth can be assessed. The LMC rotation curve provides a framework for detailed kinematic studies such as these.
The decomposition of the LMC rotation curve into disk and dark halo components has several implications for microlensing. First, as emphasized by Gyuk et al. (1999), the total mass of the LMC is an important constraint on self-lensing optical depth calculations. A high-mass model will typically predict a high self-lensing optical depth, with an additional dependence on whether the mass lies mostly in the disk or in a halo. Unfortunately, recent analyses of the LMC rotation curve have lead to mass estimates that range over a factor of $``$4, which can be attributed primarily to different assumptions about a dark halo (Kim et al. 1998, Kunkel et al. 1997, Schommer et al. 1992). Second, a massive LMC dark halo might significantly affect the LMC disk kinematics (e.g. Bahcall 1984), with possible consequence for constraints on self-lensing or determining the flare of the disk. The influence of an LMC dark halo on the LMC disk kinematics has not been investigated in detail. We note that the interplay between the disk and dark halo is also of general interest for studies of galaxy formation and evolution. Finally, if the LMC has a dark halo with a physical makeup similar to that of the Galactic dark halo (i.e. composed partly of MACHOs), the MACHO fraction of the Galactic dark halo implied by the observed optical depth would be lowered (Alcock et al. 2000).
Motivated by the above considerations, we present a new analysis of the LMC rotation curve. In this work we are particularly concerned with the flare of the LMC disk. However, other issues pertinent to microlensing are discussed as our calculations and analyses permit. We refer extensively to the H I rotation curve recently presented by Kim et al. (1998) and the impressive radial velocity dataset for LMC carbon stars summarized by Kunkel, Demers, Irwin & Albert (1997; hereafter KDIA). A significant subset of these latter data are public (Kunkel, Irwin & Demers 1997), and archived electronically at the CDS<sup>2</sup><sup>2</sup>2Centre de Donnes astronomique de Strasbourg; located at http://cdsweb.u-strasbg.fr.. We also refer extensively to the analysis of the LMC space motion presented by Kroupa & Bastian (1997).
Our paper is organized as follows. In §2, we present the impetus for this work, a comparison of constant-thickness exponential disk models (with no dark halos) to the LMC disk velocity dispersions reported by KDIA. In §3, we reanalyze the LMC rotation curve and velocity dispersion curve using archived carbon-star radial-velocity data. In §4, we present a multi-mass component kinematic model for the LMC. In §5, we compare our carbon-star rotating disk solution to our model. We discuss the microlensing implications of our analyses in §6 and conclude in §7.
## 2 A Flared LMC Disk?
### 2.1 Theoretical Models for Disk Galaxies
Disks in galaxies often have exponential surface brightness profiles from which we infer exponential mass surface density profiles of the form
$$\mathrm{\Sigma }(R)=\mathrm{\Sigma }_0e^{R/\mathrm{\Lambda }}$$
(1)
by assuming constant mass-to-light ratios ($`M/L`$). In equation (1), $`\mathrm{\Sigma }_0`$ is the central mass surface density, $`\mathrm{\Lambda }`$ is the radial scale length, and $`R`$ is a cylindrical radial coordinate. For an infinitely thin disk described by equation (1), Freeman (1970) calculated that the rotation curve reaches a maximum circular velocity at $`R2.2\mathrm{\Lambda }`$
$$V_{max}=0.623\left(\frac{GM_{disk}}{\mathrm{\Lambda }}\right)^{1/2}$$
(2)
where $`M_{disk}`$ is the total mass of the disk, which is related to the central surface density by
$$M_{disk}=2\pi \mathrm{\Sigma }_0\mathrm{\Lambda }^2$$
(3)
For the case of a disk with a modest finite thickness, $`V_{max}`$ will be decreased by $`5\%`$ (van der Kruit & Searle 1982). A useful model with a finite thickness is the exponential disk (also known as a double-exponential disk), which has a density distribution
$$\rho (R,z)=\rho _0e^{|z|/h}e^{R/\mathrm{\Lambda }}$$
(4)
In equation (4), the variable $`h`$ is the vertical, or what we will sometimes refer to as the “z” scale height. If we choose the spatial density normalization $`\rho _0=\mathrm{\Sigma }_0/2h`$, the exponential disk projects to the surface density of the infinitely thin Freeman disk, i.e. equation (1). The z-velocity dispersion, $`\sigma _z`$, for the exponential disk is given by Wainscoat, Freeman and Hyland (1989),
$$\sigma _z(R,z)=2h\left[\pi G\rho _0e^{R/\mathrm{\Lambda }}\left(1\frac{1}{2}e^{|z|/h}\right)\right]^{1/2}$$
(5)
which assumes that the disk is virialized. Equation (5) is strictly valid for a radially infinite disk. The radial dependence (i.e., the $`e^{R/\mathrm{\Lambda }}`$ term) in equation (5) scales the local density normalization, as prescribed by others (van der Kruit & Searle 1982; van der Kruit & Freeman 1984).
A model extensively discussed by van der Kruit & Searle (1981, 1981b, 1982) is the isothermal disk. The vertical distribution of stars in an isothermal disk is described by the sech<sup>2</sup> function, as originally derived by Spitzer (1942; see also §4.2 of this paper). The spatial density for this disk model is,
$$\rho (R,z)=\rho _0\mathrm{sech}^2(z/h)e^{R/\mathrm{\Lambda }}$$
(6)
which projects to the surface density of the infinitely thin Freeman disk if $`\rho _0=\mathrm{\Sigma }_0/2h`$. Assuming virialization, the z-velocity dispersion of the isothermal disk is proportional to the z-scale height and independent of z:
$$\sigma _z(R)=h\left[2\pi G\rho _0e^{R/\mathrm{\Lambda }}\right]^{1/2}$$
(7)
where the radial dependence scales the local density normalization. Equation (7) is strictly valid for a radially infinite disk. We note that the “z-scale height” in the isothermal disk is not strictly a scale height (and not the same as $`h`$ in the double-exponential model), but it is a similar parameter. Further comparisons of the exponential and isothermal disk models are found in Wainscoat, Freeman & Hyland (1989), and van der Kruit (1988).
The systematic error in the velocity dispersion predicted by the radially infinite isothermal disk model has been shown to be only $``$10% in the range $`1<R/\mathrm{\Lambda }<4`$ using numerical studies of truncated isothermal disk models (van der Kruit & Searle 1982). We will use the isothermal disk model in our analyses, because these numerical calculations provide us with an estimate of the model accuracy, and will later allow us to approximately account for the finite radial extent of the LMC disk (see §4.3 of this paper). Observations of edge-on spirals indicate that both the exponential and isothermal models are viable representations of real disks (de Grijs, Peletier, & van der Kruit 1997). The specific choice of model does not affect the main results of this work.
It is reasonable to compare the model z-velocity dispersion to the observed velocity dispersion at different projected radii in the LMC for the case of uniform anisotropy ($`\sigma _Z/\sigma _R`$ constant at different radii; van der Kruit & Freeman 1984). Although we do not know if this is true for the LMC, the assumption of uniform anisotropy is supported by kinematic studies of the old Galactic disk (Lewis & Freeman 1989). Gould (1995) notes that if $`\sigma _Z/\sigma _R<`$ 1 as found elsewhere, the line of sight velocity dispersion measured for an inclined disk overestimates $`\sigma _Z`$. Therefore, our model z-velocity dispersion is likely an upper-limit on the true z-velocity dispersion.
Combining equations (2), (3) and (7) yields
$$\sigma _z^2(R)1.288\left(\frac{h}{\mathrm{\Lambda }}\right)V_{max}^2e^{R/\mathrm{\Lambda }}$$
(8)
which is most accurate for disks in the range $`1<R/\mathrm{\Lambda }<4`$. The z-scale height has been observed to be constant over a large range of $`R`$ in numerous disk galaxies (van der Kruit & Freeman 1984; de Grijs, Peletier, & van der Kruit 1997), implying that $`\sigma `$ decreases with increasing $`R`$ in these galaxies. Indeed, this trend of decreasing velocity dispersion has been observed in over a dozen spiral/disk galaxies (Bottema 1993), including our own.
In summary, if the disk velocity dispersions are the same at different radii, the disk is likely flared (e.g. Zhao 1999). However, this inference assumes that (1) the disk is virialized, (2) the radial scale length is constant, (3) the disk velocity dispersions exhibit uniform anisotropy, and (4) the galaxy has no dark halo (or that the dark halo has a negligible affect on the disk velocity dispersions). Finally, one must be careful to compare the observed velocity dispersions to the simple model described above over a restricted range of radii, because real galactic disks are finite in radial extent. We compare to the LMC in the next section.
### 2.2 Comparison of Disk Models to the LMC
In this section, we will adopt a plausible value for the product $`hV_{max}^2`$ in equation (8), and compare that model to LMC disk velocity dispersion data. Following KDIA and many others, we adopt a distance to the LMC of 50.1 kpc, and an inclination of $`i`$ = 33 degrees. For comparison, wide-field UV polarimetric image data have yielded a precise but model-dependent estimate of the inclination of the LMC disk: $`i=36_5^{+2}`$ degrees (Cole, Wood, & Nordsieck 1999). Westerlund (1997) summarizes a number of other inclination measurements.
We adopt $`\mathrm{\Lambda }=1.6`$ kpc for the radial scale length of the LMC disk. The surface brightness profile of the LMC disk is well-fit by an exponential with a scale length $`\mathrm{\Lambda }1.6`$ kpc (de Vaucouleurs 1957, Bothun & Thompson 1988). Moreover, the surface density profile of intermediate-age long period variables (LPVs) in the LMC is also fit by an exponential with a scale length of $`\mathrm{\Lambda }1.6`$ kpc (Hughes, Wood & Reid 1991). This result directly associates our kinematic dataset with an LMC population that follows the $`\mathrm{\Lambda }`$ = 1.6 kpc exponential profile, because some of the intermediate-age LPVs are also carbon stars<sup>3</sup><sup>3</sup>3The KDIA carbon stars have a mean velocity dispersion similar to that of the Hughes et al. (1991) intermediate-age LPV sample. Therefore, as KDIA also note, these two samples of stars probably have similar ages and conform to the dynamics of the same inclined disk.. Finally, the RR Lyrae stars in the LMC also appear to lie in a $`\mathrm{\Lambda }1.6`$ kpc exponential disk (Alcock et al. 2000b). Since the RR Lyraes are presumably much older than the carbon stars, the carbon stars were likely born into the same disk. We note that no stellar population in the LMC has yet been shown to be inconsistent with a $`\mathrm{\Lambda }=1.6`$ kpc exponential profile. We assume that the carbon stars observed by KDIA properly represent this LMC disk.
The KDIA dataset consists of radial velocity measurements for 759 carbon stars spanning a true LMC radius of 2 to 10 kpc. KDIA summarize the results of 11 different zonal solutions for a rotating disk. In Figure 1, we plot the KDIA velocity dispersions as a function of radial distance, which we take as the central value for each zonal solution given in KDIA’s Table 1. We show data only for $`R`$ = 3 to 6 kpc. This is the range of radii where equation (8) will most accurately predict the velocity dispersions. The dashed line in Figure 1 is the mean of the plotted, observed velocity dispersions, $`\sigma =`$ 13.7 km s<sup>-1</sup>. Assuming a constant value of $`\sigma `$ equal to the mean, we find $`\stackrel{~}{\chi }^2`$ = 0.4, consistent with the data. The solid line plotted in Figure 1 shows the prediction of equation (8) assuming values of $`V_{max}=70`$ km s<sup>-1</sup> and $`h`$ = 0.35 kpc at $`R=2\times \mathrm{\Lambda }`$. This model predicts $`\sigma _z13.7`$ km s<sup>-1</sup> at $`R=2\times \mathrm{\Lambda }3.2`$ kpc. The fit for this model distribution is $`\stackrel{~}{\chi }^2`$ = 8.7, which can be rejected with high significance. Models with larger scale heights intersecting the data at larger $`R`$ give similarly poor fits<sup>4</sup><sup>4</sup>4These values of $`\stackrel{~}{\chi }^2`$ are not strictly apppropriate because the zonal solutions given by KDIA are not all independent. However, subsets of the $`\sigma (R)`$ measurements which are independent also give poor fits to the distribution given by equation (7).. One possible interpretation of the observed constancy of $`\sigma `$ with $`R`$ is that the z-scale height increases as $`e^{+R/2\mathrm{\Lambda }}`$, i.e. the LMC disk is flared.
## 3 The LMC Rotation Curve Revisited
An alternate interpretation of Figure 1 might be a constant-thickness disk under the dynamical influence of a dark halo (Bahcall 1984). Indeed, Kim et al. (1998) compared their H I rotation curve with the KDIA carbon star rotation curve and, on this basis, argued that an LMC dark halo is dynamically significant at a radius of $``$4 kpc. However, KDIA favored a different interpretation; they attributed the rising outer portion of their carbon star rotation curve to tidal effects, and not a dark halo<sup>5</sup><sup>5</sup>5In any case, the effect of the LMC dark halo (if it exists) would be small at the radius we chose to normalize the model $`\sigma (R)`$ curve in Figure 1, supporting this aspect of our analysis in §2.. Before concluding that the LMC disk is flared, it is worth testing disk plus dark halo kinematic models in a self-consistent manner. With this model comparison in mind, let us consider the KDIA analysis in greater detail.
We have several specific concerns with the KDIA analysis. First, and perhaps most importantly, KDIA excluded some carbon stars from their solutions for their association with a polar ring. This procedure might have artificially lowered the disk velocity dispersions. By retaining these stars in our solutions, we guarantee that our results will be properly comparable to our model (which will not include a polar ring). KDIA also corrected for the LMC transverse motion by forcing the position angle of the kinematic line of nodes to match the photometric line of nodes in each zonal solution. This is a questionable procedure. Moreover, it is inconsistent with the H I analysis by Kim et al. (1998). We will correct the carbon-star velocities for the transverse motion of the LMC in a manner consistent with the H I rotation curve analysis by Kim et al. (1998). In fact, we will make our carbon star rotation curve analysis similar in other ways to the H I rotation curve analysis (i.e., by adopting the same kinematic center), in order to lend maximum credibility to detailed comparisons of the two rotation curves.
### 3.1 Rotation Curve Solution
#### 3.1.1 The Carbon Star Radial Velocity Data
We assemble positions (right ascension and declination; equinox 1950) and galactocentric radial velocities for Magellanic Cloud carbon stars from Kunkel, Irwin & Demers (1997). We discard all Small Magellanic Cloud (SMC) carbon stars, “inter-cloud” carbon stars (their Table 15), and carbon stars located near the center of the LMC (their Table 17; see also §3.3). We also discard one rogue carbon star (C0433$``$6607) for its highly discrepant velocity. The resulting homogeneous dataset contains 422 carbon stars, representing a significant subset of the data analyzed in KDIA. Ongoing observational campaigns to obtain carbon star radial velocities will likely increase the sample size in the coming years (Suntzeff 1998). The archival dataset we have assembled is sufficient for the purposes of this study.
#### 3.1.2 Transverse Motion Radial Velocity Correction
Feast, Thackeray & Wesselink (1961) have discussed the apparent rotation induced by the transverse motion of the LMC. Each of our 422 carbon star velocities are corrected for this projected radial velocity gradient as follows. We adopt the LMC space motion calculated by Kroupa & Bastian (1997), which is derived from an average of the LMC proper motion they measure with Hipparcos data and the LMC proper motion measured by Jones, Klemola, & Lin (1994). The Hipparcos and Jones et al. (1994) measurements of the LMC proper motion agree within their respective error bars. We note that Kim et al. (1998) make their transverse motion correction using the proper motion measurement by Jones et al. (1994).
It is useful to establish a Galactic coordinate system for the vector analysis that follows. Following Kroupa & Bastian (1997), the galactocentric coordinates used here are such that the Galactic Center and the Sun are located at 0 and $`8.5`$ kpc along the $`x`$-axis, respectively. The positive $`z`$-axis points toward the north Galactic pole, and the positive $`y`$-axis points in the direction of Galactic rotation at the position of the Sun. In these coordinates, the position vectors of the Sun and LMC are
$$R_{}=(8.5,0.0,0.0)$$
(9)
$$R_{LMC}=(1.0,40.5,26.6)$$
(10)
in units of kpc. The velocity vector of the LMC is
$$V_{LMC}=(+41\pm 44,200\pm 31,+169\pm 37)$$
(11)
in units of km s<sup>-1</sup> (Kroupa & Bastian 1997). The galactocentric radial velocity of the LMC is the vector dot product $`V_{LMC}{}_{}{}^{.}\widehat{R}_{LMC}^{}=(41,200,169){}_{}{}^{.}(1.0,40.5,26.6)/48.5=74`$ km s<sup>-1</sup>, where $`\widehat{R}_{LMC}`$ is a normalized unit vector. The vector between the Sun and LMC is $`X=R_{LMC}R_{}=(7.5,40.5,26.6)`$, and the radial velocity component of the LMC seen from the local standard of rest (accounting for the Sun’s peculiar velocity and Galactic rotation) at the position of the Sun is $`V_{LMC}{}_{}{}^{.}\widehat{X}_{LMC}^{}=80`$ km s<sup>-1</sup>. The vector $`X_{LMC}`$ points toward the “center” of the LMC, which is assumed to be at $`\mathrm{}=280.46`$, $`b=32.89`$ (Kroupa & Bastian 1997). The unit vector connecting the Sun to any position in the sky in galactocentric coordinates is
$$\widehat{X}(b,\mathrm{})=(\mathrm{cos}b\mathrm{cos}\mathrm{},\mathrm{cos}b\mathrm{sin}\mathrm{},\mathrm{sin}b)$$
(12)
The transverse motion velocity correction that we seek is the difference between the projected radial velocity toward the LMC center and the projected radial velocity toward each carbon star. If we define the vector $`\widehat{S}=\widehat{X}_{LMC}\widehat{X}(b,\mathrm{})`$, the apparent rotation induced by the LMC space motion is $`\delta V(\mathrm{},b)=\widehat{S}{}_{}{}^{.}V_{LMC}^{}`$. We make our correction by subtracting $`\delta V(\mathrm{},b)`$ from the observed carbon star velocities. Small differences between the Kroupa & Bastian (1997) model used here for the transverse motion radial velocity correction and assumptions (i.e., for the LMC center and distance) in subsequent sections have a negligible effect on our results.
#### 3.1.3 Zonal Solutions
First, we convert each carbon star’s right ascension and declination into spherical coordinates in units of radians. We designate these $`\alpha `$ and $`\delta `$, respectively. We adopt the center of the LMC at right ascension $`05^h17^m.6`$ and declination $`69^002^{^{}}.0`$ (equinox 2000) following Kim et al. (1998). We designate the LMC center with $`\alpha _0`$ and $`\delta _0`$ (also in units of radians and precessed to equinox 1950). The distribution of carbon stars is shown in Figure 2 (where North is up and East is to the right). The SMC is located well outside of the boundaries of the figure in a South-Westerly direction (toward the lower left). We remind that $`\alpha `$ and $`\delta `$ are spherical coordinates plotted on a rectilinear scale. Therefore, the distribution of carbon stars shows a $``$cos($`\delta `$) distortion.
Our solution proceeds as follows. We calculate the position angle ($`\mathrm{\Theta }`$) and radial distance from the center of the LMC ($`S`$ in units of radians) for each carbon star using the formulae
$$S=\mathrm{cos}^1\left[\mathrm{sin}\delta _0\mathrm{sin}\delta +\mathrm{cos}\delta _0\mathrm{cos}\delta \mathrm{cos}(\alpha \alpha _0)\right]$$
(13)
$$\mathrm{\Theta }^{}=\mathrm{sin}^1\left[\frac{\mathrm{cos}\delta _0\mathrm{sin}(\alpha \alpha _0)}{\mathrm{sin}S}\right]$$
(14)
$$\mathrm{\Theta }=2\pi \mathrm{\Theta }_{NE}^{},\mathrm{\Theta }=\pi +\mathrm{\Theta }_{SE}^{},\mathrm{\Theta }=\mathrm{\Theta }_{NW}^{},\mathrm{\Theta }=\pi +\mathrm{\Theta }_{SW}^{}$$
(15)
Equation (15) assigns the correct sign and phase to the position angle depending on which quadrant each star resides in Figure 2. We have defined $`\mathrm{\Theta }`$ to increase West of North. We deproject each radial distance on the sky to true LMC radius ($`R`$ in units of kpc) using,
$$R=50.1\left|\mathrm{tan}^1\left(S\mathrm{cos}(\mathrm{\Theta }+\mathrm{\Theta }_0)\left[1+\mathrm{tan}^2(\mathrm{\Theta }+\mathrm{\Theta }_0)\mathrm{sec}^2i\right]^{1/2}\right)\right|$$
(16)
We take the absolute value of the inverse tangent in equation (16) to guarantee a positive value for $`R`$. We accommodate a twisting disk model by introducing $`\mathrm{\Theta }_0`$. The parameter $`\mathrm{\Theta }_0`$ allows for a changing kinematic line of nodes in different radial zone solutions. We note that $`\mathrm{\Theta }_0`$ is sensitive to the adopted LMC transverse velocity (Meatheringham et al. 1989). We define our different zones by making cuts in true radius. Each zonal solution is derived by fitting the function $`V(\mathrm{\Theta })`$ to the distribution of galactocentric radial velocities,
$$V(\mathrm{\Theta })=\pm \left(\frac{V_c\mathrm{sin}i}{\left[1+\mathrm{tan}^2(\mathrm{\Theta }+\mathrm{\Theta }_0)\mathrm{sec}^2i\right]^{1/2}}+V_{sys}\right)$$
(17)
We must choose the correct sign for $`V(\mathrm{\Theta })`$ at the appropriate position angle, which is the meaning of the ‘$`\pm `$’ in equation (17). We choose ‘$`+`$’ for $`(\mathrm{\Theta }+\mathrm{\Theta }_0)`$ = $`\pi /2`$ to $`\pi /2`$, and ‘$``$’ for $`(\mathrm{\Theta }+\mathrm{\Theta }_0)`$ = $`\pi /2`$ to $`3\pi /2`$. The fitted model could be generalized to account for a warped and decentered disk by letting $`i`$ and ($`\alpha _0`$, $`\delta _0`$) vary in each zonal solution. However, we fix $`i`$, $`\alpha _0`$, and $`\delta _0`$ to the values given above in all solutions.
We derive solutions in four radial zones which span $`R`$ = 2.5 - 5, 5 - 6, 6 - 7, and 7 - 13 kpc. These were chosen to provide an adequate number of stars and phase coverage in each zone. The inner-most and outer-most boundaries were chosen to be inclusive of the entire dataset when $`\mathrm{\Theta }_0`$ = 0. We begin by solving for $`V_{sys}`$, $`\mathrm{\Theta }_0`$, and $`V_c`$ that yield a minimum dispersion about $`V(\mathrm{\Theta })`$ in each of our zones. We average the values of $`V_{sys}`$, excluding the outermost zone (which was sometimes poorly constrained). In a second round of fitting, we fix $`V_{sys}`$ to the three-zone average value, and search only for values of $`\mathrm{\Theta }_0`$ and $`V_c`$ that minimize the velocity dispersion. This procedure was easily reproduced, and facilitated our bootstrap and Monte Carlo error analyses (described below).
The mean radii in each final zonal solution are $`<R>`$ = 4.0, 5.6, 6.5, and 7.8 kpc. The minimum velocity dispersions occurred at $`\mathrm{\Theta }_0`$ = $``$37, $``$22, $``$8, and 3 degrees . We find circular velocities $`V_c`$ = 72, 68, 65, and 65 km s<sup>-1</sup> and velocity dispersions $`\sigma `$ = 17.7, 14.9, 19.3 and 20.6 km s<sup>-1</sup>, respectively (in order of increasing radius). The four zonal solutions are shown in Figure 3. We indicate the carbon stars with small open circles. Each carbon star is plotted twice at $`\mathrm{\Theta }`$ and ($`\mathrm{\Theta }`$ \+ $`2\pi `$). The best fit curve is plotted as a solid line. The mean radius, circular velocity, and velocity dispersion are also labeled in each panel.
### 3.2 Rotation Curve Error Analysis
Our first estimate of the errors is made with a bootstrap resampling analysis (Barrow, Sonada, & Bhavsar 1984). We create 250 artificial datasets by randomly removing 1-10% of the carbon stars from our full dataset, and then find new rotating disk solutions for each of these. In a bootstrap reasampling analysis, the variances in the mean values of interest averaged over all of the artificial datasets are taken to represent the internal variances of the full dataset. We find errors in $`V_c`$ of 1.3, 1.7, 2.3, and 2.4 km s<sup>-1</sup> for the four radial zones (in order of increasing radii). We find an error in $`\sigma `$ of 0.3, 0.2, 0.4, and 0.5 km s<sup>-1</sup>, and errors in $`\mathrm{\Theta }_0`$ of 1.1, 1.0, 2.4, and 1.5 degrees, respectively.
Our second calculation is designed to estimate the errors associated with our choice of zonal boundaries and the movement of stars between zones due to the changing $`\mathrm{\Theta }_0`$ parameters. For each of the four zones, we find new solutions for 64 different variations in the zonal boundaries. We varied each boundary by no more than $`\pm `$ 0.25 kpc, which yielded roughly comparable numbers of stars and mean radii in each zonal solution (except for the outer boundary of the outer-most zone which was varied by 1.5 kpc in steps of 0.25 kpc). The variances of $`V_c`$, $`\sigma `$, and $`\mathrm{\Theta }_0`$ in the zonal boundary solutions are similar to those found with the bootstrap analysis. Moreover, we find no systematic trends with boundary choice, except for the outermost zone. In the outermost zone, $`V_c`$ decreases systematically (from $``$60 to 30 km s<sup>-1</sup>) as the inner boundary is varied from $``$6.7 to 7.2 kpc.
Finally, we perform a Monte Carlo calculation to propagate the error on the LMC space motion to the rotation curve. We create 500 artificial datasets, where the LMC space motions are drawn randomly from a sample that reproduces the best fit values and errors in equation (11). We assume Gaussian distributions for the space motion errors. This calculation yields errors in $`V_c`$ of 7.1, 8.7, 10.4 and 15.8 km s<sup>-1</sup>, errors in $`\sigma `$ of 0.1, 0.2, 0.3 and 0.4 km s<sup>-1</sup>, and errors in $`\mathrm{\Theta }_0`$ of 6.6, 4.6, 7.6 and 8.4 degrees, in the four zones (in order of increasing radius). Similar values were found with as few as $``$250 artificial datasets, indicating that our calculation is robust. For most of the parameters in our rotating disk solution, the error from the LMC space motion dominates. Our final solutions and adopted errors are summarized in Table 1.
We note that there would be no stars in common between the different solutions if all zones had the same $`\mathrm{\Theta }_0`$, but in our final solutions, $``$5% of the stars are located in two zones (see Table 1). However, since we have accounted for the movement of small numbers of stars across the zonal boundaries in our error analysis, we may consider each zonal solution as independent in subsequent analyses.
### 3.3 Velocity Dispersion in the LMC Bar
We previously discarded 22 radial velocity measurements for carbon stars very near the center of the LMC (Kunkel, Irwin & Demers 1997; their Table 17) because this sample of stars did not span an adequate range of $`\mathrm{\Theta }`$ to yield an interesting zonal rotation curve solution. However, the velocity dispersion of these carbon stars is of considerable interest. This sample of stars at $`<R>=0.5`$ kpc has a velocity dispersion of $`\sigma =22.1`$ km s<sup>-1</sup>. Guided by our error analyses above, we adopt an error of 1 km s<sup>-1</sup>. The average velocity dispersion of carbon stars in our five zones is $`\sigma `$ = 18.9 km s<sup>-1</sup>. We defer further discussion and interpretation of the LMC rotation curve and disk velocity dispersions until §5, after we present our theoretical kinematic model.
## 4 A Multi-Mass Component Model for the LMC
### 4.1 Rotation Curve
We will consider two representations of the contribution of the luminous mass in the LMC disk to the rotation curve. The first is derived empirically by Kim et al. (1998) from the $`R`$-band surface brightness data of de Vaucouleurs (1958). For two assumptions of the disk mass-to-light ratio, we present this stellar rotation curve in Table 2. The second representation of the disk rotation curve we will consider is the finite-thickness, truncated disk model of van der Kruit and Searle (1982). We adopt their model No. 2 (truncated at $`R`$ = 4$`\mathrm{\Lambda }`$) from their Appendix A.1. Table 3 summarizes this disk model rotation curve, which is given in units of the radial scale length ($`\mathrm{\Lambda }`$) and the equivalent maximum circular velocity for an infinitely thin exponential disk ($`V_{max}`$; see §2.1 of this paper). This model will also be used to calculate a correction to the scale heights inferred from the disk velocity dispersions at small radii in the LMC (see §4.3 of this paper).
When comparing either of these two disk rotation curves to our carbon star data points, we will add in quadrature the small additional contribution of gas in the LMC disk to the rotational velocity. Kim et al. (1998) calculated the gas rotation curve using the single-dish H I data of Luks & Rolfs (1992), and assuming a 30% contribution from He I. The gas rotation curve is summarized in Table 4.
Finally, we will account for a pseudo-isothermal dark halo in the traditional manner, assuming a density profile,
$$\rho _{halo}(R)=\frac{\rho _0}{1+\left(R^2/a^2\right)}$$
(18)
where the circular velocity is a function of the integrated total mass
$$V_{halo}^2(R)=4\pi G\rho _0a^2\left[1\frac{a}{R}\mathrm{tan}^1\left(\frac{R}{a}\right)\right]$$
(19)
The psuedo-isothermal dark halo is parameterized by a central density, $`\rho _0`$, and core radius, $`a`$. We add the circular velocities of the disk ($`V_{disk}`$) and psuedo-isothermal halo ($`V_{halo}`$) in quadrature in order to calculate the rotation curve of our model: $`V_c(R)`$.
### 4.2 Disk Velocity Dispersions
The effect of a spherical dark halo on the velocity dispersion of an embedded disk has been extensively discussed by Bahcall (1984), Bahcall & Casertano (1984) and Bottema (1993). The vertical motions of stars in a disk and halo system are governed by Poisson’s equation
$$\frac{^2\varphi }{z^2}=4\pi G\left(\rho _{disk}+\rho _{halo}^{eff}\right)$$
(20)
where
$$\rho _{halo}^{eff}=\rho _{halo}\frac{1}{4\pi GR}\frac{}{R}V_c^2(R)$$
(21)
and the first moment of the Boltzmann equation
$$<\sigma _z^2>\frac{\rho (z)}{z}=\frac{\varphi }{z}\rho (z)$$
(22)
The effective halo density, $`\rho _{halo}^{eff}`$, includes the contribution from the radial deriviative of the total circular velocity, $`V_c`$. If the rotation curve is fairly flat, we may discard the radial deriviative term in equation (22) and equate $`\rho _{halo}^{eff}=\rho _{halo}`$. This formulation of the problem (Bahcall 1984) neglects the $`<\sigma _z\sigma _R>`$ cross terms in Poisson’s equation. With the following redefinitions (Bahcall 1984)
$$z_0=\left(\frac{<\sigma _z^2>}{2\pi G\rho _{disk}(z=0)}\right)^{1/2},x=\frac{z}{z_0},y(x)=\frac{\rho _{disk}(x)}{\rho _{disk}(x=0)}$$
(23)
and
$$ϵ=\frac{\rho _{halo}^{eff}(z=0)}{\rho _{disk}(z=0)}$$
(24)
we rewrite the equations governing the vertical motions of the disk stars as
$$y\frac{d^2y}{dx^2}\left(\frac{dy}{dx}\right)^2=2y^22ϵy^2$$
(25)
with the boundary conditions
$$y(0)=1,\left(\frac{dy}{dx}\right)_{x=0}=0$$
(26)
The solution for $`ϵ`$ = 0 was first published by Spitzer (1942):
$$y=\mathrm{sech}^2x$$
(27)
The solution to equation (25) for $`ϵ>>1`$ is
$$y=e^{ϵx^2}$$
(28)
The differential equation (25) must be solved numerically for intermediate values of $`ϵ`$.
The $`ϵ`$ parameter define in equation (24) is the ratio of effective halo density to the disk density in the plane of the disk. It is particularly useful parameter because it relates ratio of the velocity dispersion in a disk plus halo system, $`\sigma _{disk+halo}`$, to the velocity dispersion in a disk with no halo, $`\sigma _{disk}`$. Following Bottema (1993), and suppressing the subscript $`z`$ for clarity,
$$\frac{\sigma _{disk+halo}}{\sigma _{disk}}=\alpha ^1$$
(29)
where $`\alpha `$ is the surface density
$$\alpha =\underset{0}{\overset{\mathrm{}}{}}y(x)𝑑x$$
(30)
In the limit of $`ϵ>>1`$,
$$\alpha ^1=2\left(\frac{ϵ}{\pi }\right)^{1/2}$$
(31)
We have solved for $`\alpha ^1`$ as a function of $`ϵ`$ numerically<sup>6</sup><sup>6</sup>6Bottema (1993) also made this calculation. We note that Bottema’s Fig. 15 is incorrect. However, the relevant text in Bottema (1993) is correct.; these values are listed in Table 5. In the limit of large $`ϵ`$, these values confirm equation (31). The error in $`\alpha `$ is $``$30% at $`ϵ1`$ and less than 20% at $`ϵ>2`$, in the sense that equation (31) overestimates the numerically calculated value of $`\alpha `$.
Finally, we note that the density profile of the dark halo may be derived from the observed velocity dispersion in the disk at large $`R`$, if the disk scale height is constant. Equating $`\rho _{halo}^{eff}=\rho _{halo}`$ and $`\sigma _{obs}(R)=\sigma _{disk+halo}(R)`$, and combining equations (6), (7), (24), (29), and (31) gives
$$\sigma _{obs}^2(R)=8Gh^2\rho _{halo}(R,z=0)8Gh^2\frac{\rho _0a^2}{R^2}$$
(32)
where the approximate density profile of the dark halo (at large $`R`$) follows from equation (18). Thus, if a galactic disk is embedded in a dark halo, one expects $`\sigma _{obs}(R)R^1`$ at large $`R`$ with the slope proportional to $`a\rho _0^{1/2}`$.
### 4.3 The Numerical z-Force Correction
The velocity dispersion in the LMC bar will yield an accurate estimate of the scale height only if we make a numerical correction to account for the finite extent of the LMC. We note that applying a pure disk model at small radii in the LMC is reasonable because the LMC has no bulge, and the surface brightness profile (even in the bar) is purely exponential (Bothun & Thompson 1988). We appeal to the numerical integrations by van der Kruit & Searle (1982) to estimate a “correction factor” for the disk velocity dispersions given by equation (7). In Table A1 of van der Kruit & Searle (1982), model No. 2, the ratio of the actual z-force ($`K_z`$) to the model z-force is 0.55, 0.82, 0.95, 1.07, 1.06, 0.64, and 0.77 at $`R/\mathrm{\Lambda }`$ = 0, 1, 2, 3, 4, 5 and 6, respectively. We have simply averaged over the tabulated values of $`z/z_0`$ (various other weighted-average schemes give similar values). The correction to the model velocity dispersons is the square root of the z-force ratios (van der Kruit & Searle 1982), and is $``$25% near the LMC center. The scale height inferred from the velocity dispersion in the LMC bar would be underestimated by a factor of $``$2 without the correction.
### 4.4 Summary of Theoretical Calculations
Here we summarize the theoretical calculations presented thus far. In §2.1, we introduced the basic formulae that describe an exponential disk with no dark halo, and a prediction for the velocity dispersion perpendicular to the plane of the disk. The latter calculation was made by assuming a radially infinite disk. For a disk truncated at $`R/\mathrm{\Lambda }`$ = 4 or 5, the error in this predicted velocity dispersion is less than 10% at radii of a few scale lengths (van der Kruit & Searle 1982). In §4.1, we described the rotation curve of an exponential disk embedded in a psuedo-isothermal dark halo. We presented a stellar rotation curve derived from LMC surface brightness data, and a gas rotation curve, both from Kim et al. (1998). We presented a model rotation curve for a finite-thickness disk truncated at $`R/\mathrm{\Lambda }`$ = 4, from van der Kruit & Searle (1982). In §4.2, we calculated the velocity dispersion perpendicular to the plane of a disk that is embedded in a psuedo-isothermal dark halo. At large radii, the dynamical influence of the dark halo can dominate. In this limit, we showed that the disk velocity dispersion can be predicted by an analytic formula. We estimated the error of this approximation using numerical integrations. In §4.3, we give a numerical correction to the velocity dispersion perpendicular to the plane of the disk (the prediction from §2.1) that accounts for truncation of the disk at $`R/\mathrm{\Lambda }`$ = 4, which is adopted from van der Kruit and Searle (1982). Although the LMC disk may not be truncated at precisely four radial scale lengths, by using the same model throughout this work, our analysis is self-consistent.
## 5 Analysis of the LMC Kinematic Structure
### 5.1 Decomposition of the Rotation Curve
In Figure 4, we present the LMC rotation curve and a “maximal” disk decomposition. Our four zonal solutions are indicated with bold dots and error bars. Data points on the H I rotation curve from Kim et al. (1998) are shown with bold cross symbols. For comparison, we plot the carbon star zonal solutions of KDIA with open circles. Also shown in Figure 4 are the gas and stellar disk rotation curves. We assume M/L = 2.2 for the stellar rotation curve, which is similar to that found for the Galactic disk (Bahcall, Flynn & Gould 1992). Disk mass-to-light ratios of $``$2 are favored by dynamical stability arguments (Bottema 1993). Finally, we plot the sum of the gas and stellar disk rotation curves as a bold solid line.
Our four zonal solutions and the solutions of KDIA agree within our respective errors. The model rotation curve and our data points show quite good agreement. In fact, although the overall scaling of the model is a free parameter (the M/L ratio of the stellar disk), the variation with radius is reproduced remarkably well. The disagreement between the model and H I data points would appear less severe if error bars similar to those on the carbon star data points were appropriate. If the lower envelope of the carbon star error bars were more representative of the true LMC rotation curve, then one could adopt a slightly smaller M/L ratio and the agreement with the H I data points would improve. However, since the H I data points appear to deviate from a smooth curve, most notably the dip at $`R`$ 3 kpc, it is possible that all of the H I data points lying below the model curve represent a significant failure of the model. We speculate that a bar in the LMC may be responsible, although we have not attempted any modeling of this effect.
In Figure 5, we show the same carbon star and H I data points as in Figure 4. We now replace the stellar rotation curve derived from surface brightness data with the finite-thickness, truncated model disk rotation curve. The gas rotation curve, and the sum of the gas and model disk rotation curves are also plotted. The model curve assumes a scale length of 1.6 kpc and an infinitely thin disk equivalent maximum circular velocity<sup>7</sup><sup>7</sup>7We clarify that the maximum circular velocity in our rotation curve is 72 km s<sup>-1</sup>. A finite-thickness disk has a maximum circular velocity approximately 5% smaller than an equivalent infinitely thin disk (see Table 3). In this decomposition, we adopt the finite-thickness disk model, but account for the contribution from gas (approximately 5% near maximum), and quote here the maximum circular velocity for an equivalent infinitely thin disk. $`V_{max}`$ = 71 km s<sup>-1</sup>. The model disk rotation curve is clearly consistent with the data, and fits the run of H I and carbon star data points even better than the (semi-empirical) stellar rotation curve shown in Figure 4. The small differences between the model disk rotation curve, the stellar disk rotation curve, and the H I data points as illustrated in Figures 4 and 5 are not critical to our conclusions.
We estimate the mass of the LMC disk using $`V_{max}`$ = 71 km s<sup>-1</sup> and the formulae in §2.1, which yields $`M_{disk}=4.8\pm 1.0`$ $`\times `$ $`10^9M_{}`$. The error is estimated by adopting the uncertainty of the maximum observed circular velocity in our carbon star solutions ($`\pm `$7 km sec<sup>-1</sup>; see Table 1) for the uncertainty of $`V_{max}`$. This maximal disk model has a surface density normalization of $`\mathrm{\Sigma }_0=298\pm 59`$ $`M_{}`$ pc<sup>-2</sup>. Kim et al. (1998) estimate the total mass of gas in the LMC to be 0.5 $`\times `$ $`10^9M_{}`$. Thus, the total mass of the LMC is $`5.3\pm 1.0`$ $`\times `$ $`10^9M_{}`$.
In Figure 6, we present a “minimal” disk decomposition of the LMC rotation curve. We plot the same carbon star and H I data points as in Figures 4 and 5. We plot the stellar rotation curve assuming M/L = 1, the gas rotation curve, and the sum of the stellar and gas rotation curves. We also plot the contribution of a pseudo-isothermal dark halo, and the sum of the disk and halo rotation curves. In this decomposition, we calculate $`M_{disk}=1.1\pm 1.0`$ $`\times `$ $`10^9M_{}`$ and $`\mathrm{\Sigma }_0=68`$ $`M_{}`$ pc<sup>-2</sup>. The shallow and declining run of carbon star data points favors a small core radius for the LMC dark halo. The dark halo shown has $`a`$ = 1 kpc and $`\rho _0`$ = 0.10 $`M_{}`$ pc<sup>-3</sup>. By assuming M/L = 1 in the disk, this decomposition has a maximal halo, and yields a total LMC mass of $``$$`\times `$ $`10^9M_{}`$ (corresponding to an upper limit on the LMC global mass-to-light ratio of $``$4).
### 5.2 Decomposition of the Disk Velocity Dispersions
In Figure 7, we plot our LMC disk velocity dispersions as a function of true radius with bold dots and error bars. We plot the results of KDIA with open circles (see also Fig. 1). The prediction of equation (8) for our maximal disk model is shown as a dotted line and labeled. We plot this model prediction again, but now corrected for the truncation of the LMC disk (the “$`K_Z`$ force correction”) as a bold solid line. We have assumed a constant scale height of $`h=0.5`$ kpc, which normalizes the latter curve to intersect the carbon star data point in the LMC bar. A constant-thickness disk model cannot be reconciled with the data by any choice of $`h`$.
For our minimal disk decomposition, we must account for the effect of the LMC dark halo on the disk velocity dispersions. By combining equations (6), (18), and (24), and equating $`\rho _{halo}=\rho _{halo}^{eff}`$, we may estimate the $`ϵ`$ parameter,
$$ϵ(R)=\frac{\rho _{0,halo}\left(1+\frac{R^2}{a^2}\right)^1}{\rho _{0,disk}e^{R/\mathrm{\Lambda }}}$$
(33)
Assuming a constant scale height of $`h=0.4`$ kpc, the spatial density normalization of the minimal disk is $`\rho _{0,disk}`$ = 0.068 $`M_{}`$ pc<sup>-3</sup>. Adopting the parameters of the maximal LMC dark halo from above, we find $`ϵ`$ 1.1, 1.0, 1.2, 1.9 and 3.4 at integer steps of $`R/\mathrm{\Lambda }`$ = 1 to 5. For these values of $`ϵ`$, the error associated with equation (31) is too large, and thus we appeal directly to our numerical integrations in Table 5.
For this decomposition, we begin with the velocity dispersion according to equation (8), and then make the $`K_Z`$ force correction. We calculate $`ϵ(R)`$ using equation (33), and use spline interpolations of the numerical data in Table 5 to calculate $`\alpha `$. The model disk velocity dispersions are corrected according to equation (29) by a factor of $`\alpha ^1`$ (equating $`\sigma _{obs}=\sigma _{disk+halo}`$). These various curves are plotted in Figure 8. Our choice of $`h`$ = 0.4 kpc normalizes the dark halo-corrected curve to intersect the data point in the LMC bar. A constant-thickness disk in the presence of a maximal dark halo cannot be reconciled with the observed velocity dispersions by any choice of $`h`$. For comparison, a fit to the velocity dispersion data with the minimal disk and an arbitrary pseudo-isothermal dark halo implies a maximum circular velocity in the rotation curve of $``$500 km s<sup>-1</sup>, which is about 7 times higher than observed.
#### 5.2.1 Scale Heights of the Carbon Stars
Our maximal disk model yields $`h`$ = 0.25, 0.93, and 1.62 kpc at $`R`$ = 0.5, 4.0, and 5.6 kpc, respectively. For comparison, the minimal disk yields $`h`$ = 0.55, 1.62, and 2.61 kpc at the same radii. At larger radii ($`R>`$ 6 kpc), where the disk velocity dispersions begin to rise, the implied scale heights for the maximal disk model are $`h`$ = 5.1 and 28.4 kpc (at $`R`$ = 6.5 and 8.2 kpc, respectively). The scale height inferred from the last data point is much larger than the tidal radius of the LMC (Weinberg 2000), and obviously wrong. We will return to the interpretation of these high velocity dispersions at large radii in a moment. A linear regression on the three data points with $`R<6`$ kpc yields:
$$\sigma (R)=1.39(\pm 0.10)\times R+22.89(\pm 0.42)$$
(34)
where $`R`$ is in units of kpc, and $`\sigma (R)`$ is in units of km s<sup>-1</sup>. We adopt this regression to represent the flare of the LMC disk. The change of scale height is represented by a function of the form
$$h(R)=\kappa e^{+R/\beta \mathrm{\Lambda }}$$
(35)
where $`h`$ and $`R`$ are in units of kpc. We find $`\beta =1.4`$, and $`\kappa =0.14`$ kpc. Equation (35) predicts the scale height to within 0.1 kpc of the values estimated above.
Assuming a constant mass-to-light ratio, our flared disk model must project to a surface density that varies as $`e^{R/\mathrm{\Lambda }}`$ in order to reproduce the observed surface brightness profile of the LMC (de Vaucouleurs 1957, Bothun & Thompson 1988). Therefore, the spatial density of the flared disk, $`\rho _f(R)`$, must change with an effective radial scale length that compensates for the variation of scale height with $`R`$. This effective radial scale length is easily calculated by considering
$$\mathrm{\Sigma }(R)=\mathrm{\Sigma }_0e^{R/\mathrm{\Lambda }}=2\rho _f(R)h(R)=2\rho _{0f}e^{R/\gamma \mathrm{\Lambda }}\kappa e^{+R/\beta \mathrm{\Lambda }}$$
(36)
where we require that $`(1/\gamma 1/\beta )=1`$. Substituting $`\beta `$ = 1.4 yields $`\gamma `$ = 0.583 (note that $`\gamma `$ and $`\beta `$ are dimensionless scale factors). It also follows that $`\rho _{0f}`$ = $`\mathrm{\Sigma }_0/(2\kappa )`$ = 1.064 $`M_{}`$ pc<sup>-3</sup>. In summary, our flared disk model has the following spatial density,
$$\rho _f(z,R)=\rho _{0f}e^{R/\gamma \mathrm{\Lambda }}\mathrm{sech}^2\left(\frac{z}{h(R)}\right)$$
(37)
where $`h(R)`$ is given in equation (35).
#### 5.2.2 Dynamical Influence of the Galactic Dark Halo
The dynamical influence of the Galactic dark halo is calculated following the discussion in §4.2 of this paper (see also Aubourg et al. 1999). We designate the mean density of the Galactic dark halo at the distance of the LMC as $`\overline{\rho }`$. If we substitute $`\overline{\rho }`$ for $`\rho _{eff}`$ in equation (20), then $`ϵ`$ from equation (24) yields $`\alpha `$ from Table 5, and thus the correction to the disk velocity dispersions through equation (29). We estimate $`\overline{\rho }`$ following Griest (1991; see his Eqn. ). We assume that the Galactic dark halo has a core radius of 3 kpc, and a solar neighborhood spatial density normalization of 0.0079 $`M_{}`$ pc<sup>-3</sup>. We adopt the distance from the Sun to the Galactic center of 8.5 kpc, and the distance from the Sun to the LMC of 50.1 kpc, which yields $`\overline{\rho }`$ = 0.00025 $`M_{}`$ pc<sup>-3</sup>. We estimate $`ϵ`$ = 0.001, 0.007, 0.040, 0.224 and 1.246 at integer steps of $`R/\mathrm{\Lambda }`$ = 1 to 5.
In Figure 9, we plot the same data points as in Figures 7 and 8. We indicate the regression of equation (34) as a solid line, and this regression corrected for the effect of the Galactic dark halo as a bold solid line. Although the fit is not perfect, the agreement with the last data point is quite good. Thus the influence of the Galactic dark halo on our flared disk model is of the correct magnitude to account for the outer-disk velocity dispersions. In addition, the radius where the Galactic dark halo begins to have a significant dynamical effect (i.e., where the velocity dispersions begin to rise in the disk) is also reasonably reproduced. A more detailed modeling of the outer-disk velocity dispersions is beyond the scope of this paper.
We find that $`ϵ`$ is quite small at all radii of interest for the constant-thickness maximal disk model (accounting for the Galactic dark halo but no LMC dark halo), and for the constant-thickness minimal disk model (accounting for both the Galactic and LMC dark halos). In summary, constant-thickness exponential disk models cannot be reconciled with the observed run of LMC disk velocity dispersions.
As discussed above, we adopted the mean density of the Galactic dark halo at a distance of 50 kpc of $`\overline{\rho }`$ = 0.00025 $`M_{}`$ pc<sup>-3</sup>, or $`\mathrm{log}\rho _{50}3.6`$ (Griest 1991). It is worth considering the range of allowed Galactic dark halo profiles. Assuming flat rotation curves at large radii, many authors have made Galaxy models with pseudo-isothermal density profiles (e.g. Caldwell & Ostriker 1981; Bahcall, Schmidt, & Soneira 1982). Typical core radii range from 2 to 8 kpc. The average value of the density at 50 kpc predicted by these models is $`\mathrm{log}\rho _{50}=3.6`$, with a standard deviation of 0.2 dex. As a final check, we compare the N-body cold dark matter simulation of the Galactic dark halo presented by Dubinski (1994). The initial density profile for this dark halo model was similar to an ellipsoidal Hernquist potential, which is like the universal dark halo profile found in dark matter-dominated galaxies (Kravstov et al. 1998). Dubinski’s (1994) initial density profile was then compressed by the Galactic potential (e.g. Blumenthal et al. 1986) and in its final form predicts $`\mathrm{log}\rho _{50}3.6`$. The density profile of the Galactic dark halo assumed for MACHOs (Alcock et al. 2000) is also pseudo-isothermal; thus our analysis of the LMC disk kinematics supports this dark halo profile over the range of interest ($`8`$ to 50 kpc).
#### 5.2.3 Tidal Debris?
It has been suggested that the LMC is embedded in a “shroud” of tidal debris (Weinberg 2000). Such debris is unlikely to contribute significantly to the LMC self-lensing optical depth because too little mass is involved, and it is likely to be located close to the disk (Weinberg 2000; see §6.2). Nonetheless, it is worth considering the possibility that our sample of true disk carbon stars is contaminated by tidal debris, which might reconcile a constant-thickness disk model with the disk kinematic data. In this case, tidal debris would also affect our interpretation of the outer-disk velocity dispersions. The notion of tidal debris surrounding the LMC is somewhat similar to the suggestion by KDIA that the LMC harbors a polar ring, also of tidal origin. (We note that KDIA analysed a larger sample of carbon stars, in which they claim the kinematic signature of a polar ring is evident.) We have thus far presented an interpretation of the radial velocity data for 422 carbon stars in the LMC that does not require a polar ring, or tidal debris. The nature of tidal debris (in a polar ring or otherwise) is investigated as follows.
If we assume that a constant-thickness maximal disk model represents the true LMC disk, the “excess” velocity dispersion in each zonal solution may be calculated by subtracting the model contribution from the observed dispersions in quadrature, which yields $`\sigma _{ex}`$ = 15.0, 13.7, 18.8 and 20.5 km s<sup>-1</sup> at $`R`$ = 4.0, 5.6, 6.4 and 8.2 kpc, respectively. Next, if we assume that the contaminating debris has a uniform velocity dispersion of $``$50 km s<sup>-1</sup>, the ratio of contaminating stars to total stars in each zonal solution would be of order $``$10-20%. (This ratio is inversely proportional to the velocity dispersion assumed for the tidal debris.) In each of our zones, we estimate the number of $``$50 km s<sup>-1</sup> tidal debris stars would be 12, 17, 10, and 8 (in order of increasing radius), if the LMC disk were of constant thickness.
We calculate that our four disk zones subtend relative sky areas of 0.40/0.16/0.16/1.00, in order of increasing zone radius. A regression of the relative numbers of tidal debris stars with the relative zone areas shows an anti-correlation, significant at the $``$1$`\sigma `$ level. Assuming that the true disk and tidal debris carbon stars are similarly affected by incompleteness in the data, this is not what we naively expect from the tidal debris scenario described by Weinberg (2000). The two outermost zones are illustrative. They have area ratios that differ by a factor of $``$6, yet the estimated numbers of contaminating tidal debris stars are 10 and 8. Therefore, the contaminating debris must have a non-uniform radial distribution of surface density, velocity dispersion, or both. It is also possible, but even more contrived, that a spatially-dependent incompleteness in the carbon star dataset has conspired with a non-uniform spatial distribution of tidal debris to yield the observed velocity dispersions. In light of this analysis, we prefer an interpretation of the LMC rotation curve and disk velocity dispersions that invokes a negligible contamination from tidal debris (i.e., our maximal flared disk model).
The unambiguous detection of nonvirialized LMC stars and the accurate characterization of their spatial distribution are clearly desirable. Observations of such a population/structure would be relevant to our interpretation of the LMC rotation curve and disk kinematics, with possible additional implications for LMC microlensing (e.g. Zhao 1999; Graff et al. 1999; Weinberg 2000).
## 6 Microlensing Implications
### 6.1 Self-Lensing of the Flared LMC Disk
Here we calculate the LMC self-lensing optical depth of our maximal flared disk model by directly integrating the spatial density of stars given in equation (37). We begin with an integral of the form
$$\tau (R_s,\varphi _s,z_s)=\frac{4\pi G}{c^2}\underset{0}{\overset{D_s}{}}\rho (y)\left(1\frac{y}{D_s}\right)y𝑑y\frac{4\pi G}{c^2}\underset{0}{\overset{D_s}{}}\rho (y)y𝑑y$$
(38)
which yields the optical depth to any source star in the LMC at position $`(R_s,\varphi _s,z_s)`$, a cylindrical coordinate system in the plane of the disk. The position angle $`\varphi `$ is defined to be zero at the near side if the inclined disk. The integration variable is the line of sight distance between the source and lens. We use the subscripts “$`s`$” and “$`l`$” to denote the coordinates of the sources and lenses, respectively. The approximation in equation (38) is that $`y`$ is always much smaller than the distance to the source ($`D_s`$ 50.1 kpc), which is reasonable for LMC self-lensing. Note that $`y`$ is simply related to $`D_l`$, the distance from the observer to the lens, and $`D_s`$, the distance from the observer to the source
$$y=D_lD_s=\frac{z_lz_s}{\mathrm{cos}i}$$
(39)
where $`i`$ is the inclination angle of the disk (measured from the plane of the sky). For completeness, we provide the following geometric identity relating the source and lens positions to $`y`$:
$$R_l^2=R_s^2+y^2\mathrm{sin}^2i+2R_sy\mathrm{cos}\varphi _s$$
(40)
It is necessary to integrate over the line of sight a second time, in order to calculate a density-weighted average of $`\tau (R_s,\varphi _s,z_s)`$ over the distribution of source stars
$$\overline{\tau }(R,\varphi )=\frac{\underset{0}{\overset{\mathrm{}}{}}\rho (R_s,z_s)\tau (R_s,\varphi _s,z_s)𝑑D_s}{\underset{0}{\overset{\mathrm{}}{}}\rho (R_s,z_s)𝑑D_s}$$
(41)
which may be properly compared to the observed optical depth. The integration over $`D_s`$ is made in practice by transforming to a variable in our LMC cylindrical coordinate system (i.e., $`dD_s=dz_s/\mathrm{cos}i`$). We truncate the LMC disk at a conservative radius of 10 kpc. Our results are insensitive to the choice of truncation radius at the few percent level. We calculate $`\overline{\tau }(R,\varphi )`$ for multiple lines of sight, designated by ($`R,\varphi `$) in the plane of the disk. Our calculation is the case of “pure” self-lensing (Gould 1995), which yields an upper limit. The effects of dust obscuration and nonlensing mass in the plane of the disk will lower the real value of the optical depth for a fixed total mass of the LMC.
We have calculated the self-lensing optical depths in each of the 30 survey fields in the 5.7-year LMC microlensing analysis by Alcock et al. (2000; see also Gyuk et al. 1999). In the central-most fields, we find $`\overline{\tau }1.7\times 10^8`$, while in the fields lying at the largest true radii, we find $`\overline{\tau }1.2\times 10^8`$. The field-averaged optical depth for LMC disk self-lensing is: $`\overline{\tau }_{30}=1.4\times 10^8`$.
The LMC inclination dominates the uncertainty in $`\overline{\tau }_{30}`$. First, we note that the inclination is probably not a serious concern for the rotation curve and disk kinematics because the shape of the $`V(\mathrm{\Theta })`$ function (Eqn. ) is fairly insensitive to $`i`$ over the range of plausible values. However, Gould (1995) and Gyuk et al. (1999) show that, to first order, $`\tau `$ is proportional to the square of the mass-weighted velocity dispersion and an inclination factor of sec$`{}_{}{}^{2}i`$. Therefore, we simply rescale our calculation of $`\overline{\tau }_{30}`$ for different inclinations,
$$\overline{\tau }_{30}<1.0\times 10^8\mathrm{sec}^2i$$
(42)
which should be quite accurate for the LMC inclinations typically found. As noted in §2.2, Cole et al. (1999) find $`i=36_5^{+2}`$. Thus their $``$2$`\sigma `$ upper limit, $`i<38`$, yields $`\overline{\tau }_{30}<1.6\times 10^8`$. For comparison, Bothun & Thompson (1988) find $`i=45`$, which corresponds to $`\overline{\tau }_{30}<2.0\times 10^8`$. The contribution to the uncertainty in $`\overline{\tau }_{30}`$ from the velocity dispersion follows from Gould’s (1995) formula and our equation (34); it is of order $`(0.42/22.89)^2`$, or a negligible 0.03%. Finally, the $``$10% uncertainty in $`V_C`$ yields a $``$20% uncertainty on the total LMC mass, but this corresponds to a small uncertainty in the factor of $`1.0\times 10^8`$ in equation (42).
The self-lensing optical depth on the near and far sides of the minor axis at the distances spanned by the 30 fields in Alcock et al. (2000) varies by $``$5% due to the inclination of the disk (e.g. Gould 1994). Thus the flaring of the LMC disk has only a minor effect on the spatial distribution of the self-lensing optical depth, which is a potential diagnostic of the lens population (e.g. Alcock et al. 2000). The estimate of a $``$5% minor axis asymmetry in the optical depth would increase for larger adopted values of the inclination.
### 6.2 Discussion of LMC Self-Lensing
Gyuk et al. (1999) recently reviewed different self-lensing calculations in the literature. To first order, self-lensing from the virialized LMC disk is proportional to the following combinations of derivable quantities (Gyuk et al. 1999),
$$\tau \left(\frac{M_{disk}h}{\mathrm{\Lambda }^2}\right)\mathrm{sec}^2i\sigma ^2\mathrm{sec}^2i$$
(43)
Our model accounts for “second order” effects, such as the finite radial extent and flare of the LMC disk. We have shown that uncertainties in the LMC proper motion contribute a $``$1% uncertainty to $`\sigma `$ and a $``$20% uncertainty to $`M_{disk}`$ (via $`V_C`$), assuming a truncated and flared maximal disk model. Assuming this model, our numerical calculation of the self-lensing optical depth (Eqn. ) is appropriate.
The maximal disk model is consistent with LMC dynamics. However, if a model of the LMC with a dark halo is preferred, the halo will probably contribute to LMC self-lensing in addition to the disk. In this case, our equation (42) is not appropriate (see Gyuk et al. 1999). We have adopted a minimal LMC disk, and provided an example maximal dark halo which is consistent with the rotation curve. Measurements of the LMC disk mass to light ratio would better constrain the minimal LMC disk model (our model is extreme). For the more complicated case of an LMC with a dark halo, constraints on self-lensing are weaker. This is due partly to the uncertainty that the LMC proper motion contributes to the rotation curve, which constrains the total LMC mass.
The LMC self-lensing optical depth may depend systematically on the age of the stars whose kinematics are studied. Historically, it has been difficult to prove that old LMC stellar populations have different velocity dispersions. For example, the CH stars in the LMC were originally thought to trace the elusive LMC halo population, by analogy with CH stars in the Galaxy (Hartwick & Cowley 1988, Cowley & Hartwick 1991). However, the CH stars in the LMC are now believed to represent a much younger population (Suntzeff et al. 1993). We note that our average carbon star velocity dispersion (18.9 km s<sup>-1</sup>) is similar to that found for the CH stars (20 km s<sup>-1</sup>).
According to Aubourg et al. (1999) and Salati et al. (1999), the LMC could harbor an old population with a large velocity dispersion. These authors invoke Wielen’s (1977) age-velocity dispersion calibration from the solar neighborhood. We note that Weilen’s calibration was derived for stars younger than $``$3 Gyr, and extrapolation beyond this age may not be appropriate. Moreover, Weinberg (2000) predicts that the velocity dispersions in the LMC disk, under the influence of the Galactic tidal field, will remain constant (or decrease) over time.
Extant observational studies have not yielded a clear picture of the variation of velocity dispersion with age in the LMC. For example, the old globular clusters in the LMC, with ages of $``$12 Gyrs (Olsen et al. 1998), exhibit a velocity dispersion of $``$21-24 km s<sup>-1</sup> (Schommer et al. 1992; see also Freeman, Illingworth, & Oemler 1982). Thus, the LMC carbon stars and ancient clusters support the hypothesis of a constant disk velocity dispersion (for stellar populations older than $``$3 Gyr), in agreement with Weinberg’s (2000) prediction. However, the velocity dispersion of old LPVs in the LMC, with ages of $``$8 Gyr (Hughes, Wood, & Reid 1991; Olszweski, Sunzteff, & Mateo 1996), is $``$35 km s<sup>-1</sup> (Hughes et al. 1991). Further observational studies are needed, particularly given the provocative old LPV result.
Last, we suggest that a shroud of tidal debris or a polar ring, if they exist, would probably make a negligible contribution to the LMC self-lensing optical depth because the total amount of mass involved is small. Weinberg (2000) reached a similar conclusion.
## 7 Conclusion
The rotation of the disk of the LMC has been derived from the radial velocities of 422 carbon stars (Kunkel, Irwin, & Demers 1997). We have propogated the uncertainty in the LMC space motion to the LMC rotation curve with a Monte Carlo calculation. The associated disk velocity dispersions found in each rotating disk zonal solution are less sensitive to systematic uncertainties from the LMC space motion than the circular velocities obtained. We note that our carbon star rotation curve is derived in a manner consistent with the H I rotation curve analysis by Kim et al. (1998); thus our respective results are properly comparable.
We have fit our carbon star rotation curve and the H I rotation curve with a truncated, maximal LMC disk model, yielding a total LMC mass of $`5.3\pm 1.0`$ $`\times `$ $`10^9M_{}`$. We also conclude that the disk of the LMC is flared. Extrapolating the flare inferred at small radii (where tidal perturbations are small) to the outer-disk, and accounting for the influence of the Galactic dark halo, we are able to approximately reproduce the observed disk kinematics. This model favors an isothermal density profile for the Galactic dark halo out to a distance of 50 kpc. Our truncated and flared maximal disk model yields a limit on the spatially-averaged LMC self-lensing optical depth of $`\overline{\tau }_{30}<1.0\times 10^8\mathrm{sec}^2i`$. For plausible values of the LMC inclination, this low self-lensing rate compared to the measured microlensing rate allows for the existence of a dark lens population in the Galactic halo (Alcock et al. 2000).
Finally, we caution that we have not included in our analysis (1) the dynamical effects of an LMC bar, (2) large-scale non-circular motions, (3) non-uniform anisotropy of the disk velocity dispersions, or (4) arbitrary spatial distributions of tidal debris (i.e., a polar ring or stars out of virial equilibrium). Despite these shortcomings, our truncated and flared maximal disk model successfully accounts for the general dynamical characteristics of the LMC, lending inferences from the model high weight.
## 8 Acknowledgments
D.R.A. acknowledges Howard E. Bond for support of work performed at the Space Telescope Science Institute, and Kem H. Cook for support of work performed at the Lawrence Livermore National Laboratory. NASA Research Grant NAG5-6821 “UV, Visible, and Gravitational Astrophysics Research and Analysis” and DOE Contract W7405-ENG-48 are recognized. C.A.N. acknowledges support from National Physical Science Consortium Graduate Student Fellowship. We thank Ken Freeman, David Bennett, Tim Axelrod, Howard Bond, Kailash Sahu, Geza Gyuk, Greg Bothun, and the anonymous referee for their helpful discussions and comments.
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# Formation of a tidal dwarf galaxy in the interacting system Arp 245 (NGC 2992/93)
## 1 Introduction
Research activity in the field of interacting galaxies has increased quite dramatically over the last thirty years (see the recent very comprehensive review by Struck 1999). Galactic collisions trigger a number of phenomena, such as inward–transportation of gas from distances of up to kiloparsecs to the nucleus which is thought to be an efficient means to fuel a central starburst or nuclear activity. The inverse process is the ejection of material into the intergalactic medium by tidal forces. The prominent tidal tails and bridges that emanate from interacting galaxies have proved to be important tools to study the interaction, constraining the orbital parameters (Toomre & Toomre 1972) of the collision. Recently, attempts have even been made to use the formation of tidal tails as a diagnostic of the mass distribution of halos within the framework of Cold Dark Matter cosmologies (Dubinski et al. 1996, Springel & White 1999). Much less attention has been paid to what goes on within these tidal features (see for instance Schombert et al. 1990, Wallin 1990, Hibbard et al. 2000).
Detailed H I maps of a number of interacting systems (e.g., Hibbard & van Gorkom 1996, Kaufman et al. 1997, and references therein) have shown that a large fraction of the gaseous component of colliding galaxies can be expelled into the galactic halos or even into the intergalactic medium as a result of tidal forces. In some systems, up to 90% of the atomic hydrogen is observed outside the optical disk (like in Arp 105, Duc et al. 1997). Even if part of this gas falls back towards its progenitors (Hibbard & Mihos 1995), a significant amount of gas will be lost for the merger remnant for time scales of at least 1–10 Gyr. The stellar/gaseous tidal debris might be dispersed in the intergalactic/intracluster medium where the stellar component then contributes to the diffuse background light observed in clusters (Gregg & West 1998) or recondense within the halo of the merger and form a new generation of galaxies: the so–called tidal dwarf galaxies (TDGs; see the review by Duc & Mirabel 1999). TDGs are typically found at the tip of tidal tails at distances between 20 and 100 kpc from the merging objects, of which at least one should be a gas–rich galaxy. They are gas–rich objects that can be as massive as the Magellanic Clouds, form stars at a rate which might be as high as in blue compact dwarf galaxies (BCDGs) and seem dynamically independent from their parent galaxies. Although the observational evidence for the existence of recycled galaxies has been well established by now, their formation process is as yet not well understood. In particular it is not known when and under which conditions TDGs form during the interaction. This is one of the topics of the current paper.
Dating events is a general problem in extragalactic astronomy that may however be more easily achieved for interacting systems via numerical simulations. The comparison of morphological and/or kinematical features with predictions based on numerical simulations provides strong indications on the age of the collision and therefore sets a constraint on the formation history of TDGs. Interacting galaxies observed just after the first perigalacticon are particularly attractive. Whereas they have already developed long tails and bridges, their disks are still clearly separated and hence their intrinsic properties (i.e., orientation, sense and amplitude of their rotation, which are key input parameters to the numerical models) can be well defined. In contrast, galaxies on course for their first interaction do not show strong perturbations before reaching perigalacticon and hence do not provide a simple handle on any time scale. On the other hand, evolved mergers have lost all memory of the initial properties.
Arp 245, an interacting system consisting of two spiral galaxies, NGC 2992 and NGC 2993, appears to be an interesting testcase as it can be fairly easily modelled. Moreover, it is a relatively nearby system, at an adopted distance of 31 Mpc<sup>1</sup><sup>1</sup>1We use in this paper H$`{}_{0}{}^{}=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup>= 100 h km s<sup>-1</sup> Mpc<sup>-1</sup>. At the distance of Arp 245, 1$`\mathrm{}`$ corresponds to 9 kpc. Its prominent tidal tails host a tidal dwarf galaxy candidate which, because of its proximity, can be studied in detail. Finally Arp 245 has already been well studied. NGC 2992, in particular, has been the object of numerous articles focusing mostly on its active, Seyfert 1.9 nucleus. Radio continuum maps have revealed a striking pair of loops near the nucleus with a “Figure–8” shape (Ulvestad & Wilson 1984) which has later on been detected at other wavelengths as well (Wehrle & Morris 1988, Chapman et al. 2000). Glass (1997) has monitored the AGN in the near–infrared and has reported an outburst. Durret & Bergeron (1987, 1988), Colbert et al. (1996) and Allen et al. (1999) have studied the ionization cone extending from the AGN and the surrounding extended emission line regions. Evidence for outflow has been claimed by Colbert et al. (1996) and Márquez et al. (1998). NGC 2992 has been observed at many wavelengths from the X–ray regime, where it is a strong emitter (Marshall et al. 1981, Gilli et al. 2000), to the centimetre wavelengths (Condon et al. 1982). Single dish H I (Mirabel & Wilson 1984) and CO (Sanders & Mirabel 1985) data are also available. NGC 2993 is mentioned in several catalogs of starburst galaxies but has not yet been the subject of any individual study. Surprisingly, the fact that both galaxies are partaking in a spectacular interaction has largely been ignored, apart from the photometric work by Schombert et al. (1990).
We decided to subject Arp 245 to a comprehensive multiwavelength study. In this first paper, we detail the observations, present in particular the first complete H I map of the system, and assess the status of the interaction with the help of numerical N–body/hydrodynamical simulations. We then focus on the tidal features, emphasizing the properties of the tidal dwarf galaxy candidate. Based on these results, and of similar such systems from the literature, we attempt to better understand the formation process of TDGs. In a second paper (Paper II, in prep.), we will concentrate on the internal properties of NGC 2992, and study in particular its Active Galactic Nucleus (AGN) and the ionization filaments.
## 2 Observations
Table 1 and 2 summarize our multi–wavelength observations<sup>2</sup><sup>2</sup>2Partly carried out at the European Southern Observatory, La Silla, Chile (ESO No 54.A–0606, 56.A–0757, 64.N–0163 and 64.N–0361) of Arp 245 and lists technical details. In the next sections we will discuss these observations in turn.
### 2.1 Optical and near–infrared broad–band imaging
Optical broad band BVR images of Arp 245 have been collected in February 1995 with the 3.5–m NTT at la Silla observatory. The red arm of EMMI has been used. The weather conditions were photometric and the seeing varied between $`0\stackrel{}{\mathrm{.}}9`$ and $`1\stackrel{}{\mathrm{.}}2`$. The $`9\stackrel{}{\mathrm{.}}2\times 8\stackrel{}{\mathrm{.}}7`$ field of view of the R–band image covered the entire interacting system including the tidal features whereas the B<sup>3</sup><sup>3</sup>3The B–like filter, Bb, optimized for the red arm of EMMI, has actually been used. and V band images were somewhat offset and missed a small part of the southern tail. However, B and V images of this part of the system had been previously taken with the PUMA camera on the CFHT<sup>4</sup><sup>4</sup>4The Canada–France–Hawaii Telescope is operated by the National Research Council of Canada, the Centre National de la Recherche Scientifique de France and the University of Hawaii. in February 1992. The weather conditions at CFHT were poor and the seeing was 1$`\mathrm{}`$. Finally BVR images of a field situated to the South-West of NGC 2993 were obtained in January 2000 with the ESO 3.6m telescope. The weather conditions were photometric and the seeing was 1$`\mathrm{}`$. Images from the three telescopes were eventually combined. Landolt fields of photometric standard stars (Landolt 1992) were observed for flux calibration.
Near–infrared JHK images of NGC 2992 were obtained in February 1996 with the IRAC2B camera installed on the MPG/ESO 2.2m. The field of view of each individual image was $`2\stackrel{}{\mathrm{.}}1\times 2\stackrel{}{\mathrm{.}}1`$. Sky images were taken of adjacent fields offset by 2$`\mathrm{}`$. The entire field of NGC 2992, including the tidal tail, was covered with a mosaic technique. The highest effective integration time was reached towards the tail which was observed for a total of 15 min in J, 12 min in H and 13 min in K. Several photometric standard stars from the list of UKIRT faint IR standards (Casali & Hawarden 1992) were observed under photometric conditions. The seeing varied between $`1\stackrel{}{\mathrm{.}}3`$ and $`1\stackrel{}{\mathrm{.}}5`$.
The reduction of the optical/NIR data has been performed within IRAF with standard tasks from the ccdred package complemented by self–written scripts which perform a semi–automatic processing of the NIR data set. An astrometric correction achieving a precision of $`0\stackrel{}{\mathrm{.}}3`$ was performed to each optical/NIR image using the positions of several tens of stars from the USNO A1.0 astrometric catalog (Monet 1996) queried via the ESO skycat browser. The frames were corrected for distortions during this process. They were registered, PSF–matched and combined to produce color maps. An optical “true color” image of the system was constructed from the combination of the NTT/CFHT/3.6m B, V, and R images. It is shown in Figure 1. The monochromatic V–band image is displayed in Figure 2.
Photometry was carried out on the registered images with the digiphot package in IRAF. Polygonal apertures were chosen to enclose respectively the disks of NGC 2992 and NGC 2993, their outer regions only, the bridge between both systems, the tidal tails, and the tidal dwarf galaxy. Because of the complex geometry of the system, the sky level was measured manually and averaged at numerous positions surrounding the objects. This procedure was chosen so as to ensure that sky contamination by stars or galaxies would be minimal. The photometric accuracies take into account variations in the sky background as well as photon noise. The magnitudes, corrected for Galactic extinction are listed in Tables 3, 4, and 5. Circular aperture photometry of both NGC 2992 and NGC 2993 was performed as well, for comparison purposes. In the optical, our V and R magnitudes differ by less than 0.05 mag with data in the literature (Prugniel & Heraudeau 1998) whereas in the B–band, the agreement is not as good. The B magnitudes were systematically too faint by 0.1 mag. A scatter in the color terms in the blue is expected since the EMMI Bb filter which we have used slightly differs from a Bessel B filter. As the system does not show large variations in color, we have applied to all B–band magnitudes a correction corresponding to an offset of 0.1 mag. This post–calibration has been later on validated using the B–band photometric measurements obtained in January 2000. In the near–infrared, comparisons are more difficult due to the intrinsic variability of the nuclear flux. Glass (1997) monitored the AGN and measured excursions of the central K–band flux of up to 0.6 mag due to outbursts. Colors appear to be less affected though and our J–H and H–K colors are within 0.1 mag of those of Alonso-Herrero et al. (1998). On average our optical/NIR photometry might be affected by systematic errors of up to 0.1 mag. Note that the instrumental errors quoted in Tables 4 and 5 do not take this into account.
### 2.2 H<sub>α</sub> imaging
An image with a narrow–band filter centered on the redshifted H<sub>α</sub> line has been obtained with EMMI during the NTT run. The filter which has a width (FWHM) of 66Å includes light from H<sub>α</sub> and \[NII\] . The image covers all of NGC 2992, including its tidal plume and bridge, and the disk of NGC 2993 with part of its tidal tail. Only the southern most region has been missed. The exposure time was 15 min. The continuum was obtained by using a similar narrow–band filter but offset in wavelength. A halftone H<sub>α</sub> image is displayed in Figure 3 and intensity contours of the line emission towards NGC 2993 and towards the northern tidal tail are shown in Figure 4 and Figure 17, respectively. A contour map towards NGC 2992 will be presented in Paper II. Our H<sub>α</sub> image shows a number of structures not seen in the maps previously published (see e.g. Colbert et al. 1996, Allen et al. 1999, for some recent ones), such as long ionization filaments outside the disk of NGC 2992 and H ii regions at the tip of the northern tidal plume. Flux calibration was achieved by observing a spectroscopic standard star from the list of Hamuy et al. (1992) observed with same setup. Our integrated flux towards NGC 2992 agrees to within 20% with the flux measured by Colbert et al. (1996).
### 2.3 Optical spectroscopy
Longslit spectroscopic observations were carried out in February 1995 with the NTT and EMMI. The grism used has an intrinsic resolution of 9Å. Coupled with the red–arm CCD, it covers the wavelength range 3800–8400 Å. Two spectra were obtained along an axis roughly parallel to the morphological axis of NGC 2992. The position angles were $`18^{}`$ and $`20.2^{}`$. The $`1\stackrel{}{\mathrm{.}}5`$ wide slits encompassed the western regions of NGC 2992 and two of the brightest clumps of the tidal dwarf galaxy at the tip of the plume. The total integration time was 45 min divided in three exposures. The spectra of several spectroscopic standard stars from the list of Hamuy et al. (1992) were measured through a wide slit. Weather conditions were photometric. A spectroscopic follow–up was performed in January 2000 with the EFOSC2 instrument installed on the ESO 3.6m telescope. The MOS mode offered by punched masks was used to obtain the optical spectra of several extended objects near NGC 2992 as well as some new condensations in its tidal dwarf. The grism used had a wider wavelength coverage but a slightly lower spectral resolution. Five exposures were obtained and the total integration time was 100 min.
Data reduction, wavelength and flux calibrations were performed with the twodspec package within IRAF following standard procedures. Line fluxes and errors were measured with the IRAF splot task. Data were corrected for extinction using the formula:
$$\frac{I(\lambda )}{I(\text{H}\text{β})}=\frac{F(\lambda )}{F(\text{H}\text{β})}10^{cf(\lambda )}$$
where $`F(\lambda )`$ is the observed line flux, $`f(\lambda )`$, the reddening function taken from Torres-Peimbert et al. (1989) and $`c`$ the logarithmic reddening correction at H<sub>β</sub> obtained from a constant $`\text{H}\text{α}/\text{H}\text{β}`$ Balmer decrement of 2.85. The lines were not corrected for underlying stellar absorption. Table 6 displays the values of $`c`$, the absolute flux and equivalent width of the H<sub>β</sub> line and the observed/corrected fluxes relative to H<sub>β</sub> of the principal lines. The accuracy of the absolute flux calibration has been checked using our H<sub>α</sub> narrow–band image. The integrated H<sub>α</sub> fluxes of several emission line regions extracted from the H<sub>α</sub> image and from the slit spectra agree to within 5% to 15%.
### 2.4 VLA H I observations
We obtained observations with the Very Large Array (VLA)<sup>5</sup><sup>5</sup>5The VLA is operated by the National Radio Astronomy Observatory, a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. of Arp 245 in the 21–cm line of neutral hydrogen (H I) in September 1997, using the CS–array<sup>6</sup><sup>6</sup>6This configuration results in a resolution equivalent to that of C–array, but has some antennae placed at locations usually employed for D–array, with the aim to provide sufficient short spacing information so that separate D–array observations are not necessary.. We used the calibrator 1331+305 for absolute flux calibration and 0902–142 as secondary calibrator. We determined a flux density for 0902–142 of $`2.89\pm 0.01`$ Jy. The data were taken in correlator mode 4 with two passbands (or IFs) measuring both right and left hand polarization. Each band contained 32 channels, and we used an 8–channel overlap between the bands. The first and last three channels of each band were discarded because of the decreasing gain toward the edges of the bandpass. Full details of the observations are presented in Table 2.
Data reduction and calibration were performed using the Astronomical Image Processing System (aips) package, following standard procedures. The uv–data from each of the two IFs were calibrated separately. The visibilities were inspected and bad data points due to interference were removed. Because of the presence of strong solar interference at the lowest spatial frequencies, we excluded baselines shorter than 1 k$`\lambda `$ (or about 210 m) in the calibration. After the standard amplitude and phase calibration we went through one cycle of self–calibration (phase only), improving the final quality of the images. As reference source we used the bright nuclear emission of NGC 2992 which has a flux density of $`190\pm 2`$ mJy.
For each IF we generated “dirty” data cubes of line plus continuum emission. After running several tests we decided that the best maps were those made with the task imagr using the default values for the robust weighting scheme and including the entire observed uv–range. The first 15 and the last 15 channels of IF1 and IF2 respectively, were found to be free of line emission. From each corresponding IF the average of those line–free channels was subtracted to construct data cubes with line emission only. An average of all line free channels from both IFs was used to construct a map of the continuum emission. This map, after cleaning, is shown in Figure 6. Besides the nuclear emission due to NGC 2992 we see emission corresponding to the nucleus of NGC 2993. This emission is resolved and measures $`12.2^{\prime \prime }\times 9.4^{\prime \prime }`$ at a position angle of $`58^{}`$. Its peak flux corresponds to 49 mJy.
After continuum subtraction the maps were inspected and cleaned, and the two IFs were appropriately merged. In order to isolate the H I emission in the maps from the noise, the cube, which has originally a resolution of $`19.4^{\prime \prime }\times 14.4^{\prime \prime }`$ was convolved to a circular beam of $`35^{\prime \prime }`$. We clipped this cube at a level of twice the rms noise. We inspected all features in this cube, retaining only those regions which show emission in at least five consecutive channels. We then went back to the original, high resolution cube and applied conditional blanking, retaining emission in the high resolution cube from only those regions which were preserved in the processed $`35^{\prime \prime }`$ cube. As a fringe benefit this process removed most if not all traces of solar interference. Besides emission, H I absorption is seen towards the strong nuclear source in NGC 2992. We ensured that this was preserved in the conditional blanking process. In addition to a blanked cube at the highest resolution ($`19.4^{\prime \prime }\times 14.4^{\prime \prime }`$) we produced lower resolution cubes as well, such as at $`25^{\prime \prime }\times 25^{\prime \prime }`$ and $`30^{\prime \prime }\times 30^{\prime \prime }`$. As a last step we obtained the moments of the cubes. The conversion factors from mJy beam<sup>-1</sup> to Kelvin brightness temperatures are listed in Table 2.
Figure 7 displays a mosaic of the H I channel maps at a resolution of 25$`\mathrm{}`$, showing every second channel, superimposed on an image obtained from the DSS of the same area. The final H I cube consists of 32 channels with velocities ranging from 2655 km s<sup>-1</sup> to 2006 km s<sup>-1</sup> with a channel spacing of 21 km s<sup>-1</sup>. Emission is found from 2592 km s<sup>-1</sup> to 2048 km s<sup>-1</sup>. Starting at 2267 km s<sup>-1</sup> and extending over some 200 km s<sup>-1</sup> to 2477 km s<sup>-1</sup> we see strong absorption against the nuclear continuum source of NGC 2992.
Maps of the integrated H I distribution at a resolution of $`25^{\prime \prime }`$ overlayed on a V–band and on an H<sub>α</sub> image of the system are shown in Figure 2 and in Figure 3, respectively. In addition to the two main galaxies, H I emission is clearly associated with the tidal dwarf galaxy to the north of NGC 2992, at the tip of the tidal arm. It is further traced along the northern tidal arm and in the bridge between NGC 2992 and NGC 2993. Surprisingly, there is a huge ring–like H I structure extending to the southeast of NGC 2993. In addition to emission from the NGC 2992/93 system, a strong signal is picked up from what turns out to be a companion object, known as FCG 0938 and located to the southwest. Total H I fluxes of the galaxies are 18.6 Jy km $`\mathrm{s}^1`$ for the interacting system, corresponding to an H I mass of $`4.2\times 10^9\text{M}\text{}`$. For the companion we find 3.1 Jy km $`\mathrm{s}^1`$ which corresponds to an H I mass of $`7.0\times 10^8\text{M}\text{}`$.
The VLA integrated H I spectrum of NGC 2992 matches well, both qualitatively and quantitatively, with the single dish H I spectrum obtained by Mirabel & Wilson (1984) at Arecibo. After beam corrections, the fluxes differ by less than 10 %.
### 2.5 CO observations
CO(1–0) and CO(2–1) spectra towards the northern tidal tail of Arp 245 were obtained with the IRAM 30 m antenna in June 1999. These observations are presented in Braine et al. (2000). The CO emission over the entire system, including NGC 2992, NGC 2993, and the tidal features, has subsequently been mapped in November 1999 with the SEST 15–m antenna at la Silla. Details regarding those observations will be given in Paper II.
## 3 Results
### 3.1 Status of the interaction
#### 3.1.1 The galaxies involved
Table 3 summarizes the global properties of all members of the system involved in the interaction: the two spirals of moderate luminosity, NGC 2992 and NGC 2993, and the dwarf galaxy FGC 0938.
The two galaxies NGC 2992 and NGC 2993 have radial velocities which differ by less than 100 km s<sup>-1</sup> and suffer a strong tidal interaction. NGC 2992 is a disturbed Sa spiral seen nearly edge–on at an inclination of $`70^{}`$. On top of having an active Seyfert 1.9 nucleus and a perturbed morphology due to the interaction, NGC 2992 shows a number of peculiarities. A prominent dust lane crosses most of the galaxy from the North–East to the South–West (see Fig. 1). It could be the remnant of a past merger. The H I distribution is rather complex. On top of the expected rotating disk, fast–moving clouds are observed. Finally, the galaxy exhibits several ionization filaments that are located well beyond the ionization cone associated with the AGN. They can be seen in Figure 3. Their kinematics is highly suggestive of outflows. Paper II will be devoted to the study of these phenomena which are likely more directly related to the nuclear activity than to the interaction.
NGC 2993, a face–on Sab, is situated at about 3′(27 kpc) to the South–East. On a true color image (see Fig. 1), NGC 2993 appears much bluer than its companion. Its blue color index ($`\text{B}\text{V}`$= 0.3 mag) and its rather high far–infrared to blue luminosity ratio ($`L_{\text{FIR}}`$/$`L_\text{B}`$= 1.7) is suggestive of an active starburst. The H<sub>2</sub> content is however rather low (Sofue et al. 1993) and the far–infrared to H<sub>2</sub> mass — a measure of the star formation efficiency — is much higher than in classical starburst galaxies (Sanders & Mirabel 1985). Active star–forming regions, as traced by H<sub>α</sub>, are present all over the galaxy (see Fig. 4), but are more concentrated in the clumpy inner regions where the H I peaks as well. Some, red unresolved stellar clumps can be seen in the outskirts of the galaxy, especially to the North. They might be old globular clusters.
Our VLA map (Fig. 2) disclosed a third partner at roughly the same velocity, lying to the south–west which is listed in the Catalog of Flat Galaxies (Karachentsev et al. 1993) as FGC 0938. On our CCD image (Fig. 5), it appears as an almost undisturbed edge–on disk with an ellipticity of 0.83, consistent with an inclination of 80$`\mathrm{°}`$. With $`M_\text{B}=15.5`$, FGC 0938 has the absolute blue magnitude of a dwarf galaxy. Its projected blue central surface brightness is 22.5 $`\mathrm{mag}\mathrm{arcsec}^2`$. Its rotating H I disk looks regular apart from perhaps a small extension to the West. The peak rotational velocity is of order 60 km s<sup>-1</sup>, corroborating its classification as a dwarf galaxy. Despite their different environments, FGC 0938 has many properties in common with the superthin galaxy UGC 7321 recently studied by Matthews et al. (1999). It would be remarkable if a galaxy involved in a tidal interaction would manage to keep a flat disk. This is a strong indication that FGC 0938 is presently heading for its first encounter, falling towards the NGC 2992/93 system.
#### 3.1.2 The tidal features
Arp 245 features three major tidal structures resulting from the interaction between NGC 2992 and NGC 2993: two long tidal tails escaping from both spiral galaxies and a bridge between them. Their properties are summarized in Table 4. The northern tail associated with NGC 2992 is about 16 kpc long, a rather modest extent compared to the 100–kpc long antennae in the prototypical interacting systems NGC 4038/39 and NGC 7252. This might be due to the smaller size and earlier morphological type of the progenitor and to the fact that the interaction is witnessed in its early phase when the tails did not have time to fully develop. The northern tail has roughly the same distribution and extent in both its stellar (optical) and gaseous (H I) components. The optical tail looks like a plume, spreading at its tip where the H I peaks at a column density of $`2\times 10^{21}\text{cm}\text{-2}`$. This sub–structure, detected in all optical and near–infrared bands, suggests the formation of a tidal dwarf galaxy. It hosts several star–forming regions, clearly visible in the H<sub>α</sub> map (see Fig. 17) at the location of the H I clump. The integrated colors of the tail as a whole ($`\text{B}\text{V}`$= 0.57, $`\text{V}\text{K}`$= 0.42) are similar to the colors of the outer regions of NGC 2992 ($`\text{B}\text{V}`$= 0.62, $`\text{V}\text{K}`$= 0.43). Some slightly bluer colors by 0.1 mag are measured within the H ii regions. The stellar population of the tail is hence dominated by old stars pulled out from the parent galaxy with some weak contribution from young stars born in situ in the tidal object. The TDG candidate will be studied in more detail in Sect. 3.3.
A diffuse stellar bridge, with a maximum V surface brightness of 24 $`\mathrm{mag}\mathrm{arcsec}^2`$, connects the two spiral galaxies. The orientation of its H I counterpart seems to be slightly twisted. No stellar clusters no H ii regions or H I substructures are found there. Instances of star–forming bridges do exist (e.g., NGC 6845; Rodrigues et al. 1999) but are in general quite rare among interacting systems.
The tail emanating from NGC 2993 to the East is highly curved in the plane of the sky. This is particularly visible in its H I component which shapes like a ring. The projected total extent of the NGC 2993 tail is 30 kpc. The optical tail appears sharp edged at its base and more diffuse further out. The southern and western parts of the H I ring–like structure have no optical counterparts up to a limit of $`\mu _\text{V}`$=25 $`\mathrm{mag}\mathrm{arcsec}^2`$. No obvious H ii regions are found along the tail, at least over the area covered by our H<sub>α</sub> image. Further out, the absence of bright optical condensations in the broad band images strongly suggests that star formation, if any, is very weak there, contrary to the NGC 2992 tail. The main difference between the two tidal features is their respective H I column density which turns out to be much lower in the NGC 2993 tail. The few H I condensations scattered over the ring have all peak column densities lower than $`3\times 10^{20}\text{cm}\text{-2}`$. Surprisingly, the gobal color of this quiescient tidal feature ($`\text{B}\text{V}=0.39`$) is bluer than that of the northern star–forming tail and similar to that of the outskirts of its parent galaxy, NGC 2993, a starburst spiral. The blue color of the tidal tail might therefore reflect the past star–forming activity in the parent galaxy.
Our VLA map revealed one more feature with, perhaps, a tidal origin and lying east of the NGC 2993 ring, a detached H I cloud with a mass of $`6.5\times 10^7\text{M}\text{}`$. This object, visible in several VLA channels, and therefore likely real, has no optical counterpart.
In total the H I in tidal features contributes to about 45% of the total H I mass measured in Arp 245. The H I mass of the northern tail is about 60% of that of NGC 2992<sup>7</sup><sup>7</sup>7The mass of the HI clouds seen in absorption towards the nucleus of NGC 2992 has not been included whereas the H I mass of the southern ring is 80% of that of NGC 2993. Such large values are typical of gas–rich interacting systems where the contribution of the tails may reach 90% (such as in Arp 105; Duc et al. 1997). On the other hand the tidal features contain a much weaker proportion of stars. The tail associated with NGC 2992 and NGC 2993 contains respectively 17% and 13% of their stellar luminous population. Such values have been determined from the V–band luminosities, assuming that the M/L ratio is the same for all objects.
The 2D kinematics of the southern tail derived from the H I datacube is quite smooth (see the channel map in Fig. 7 and the velocity map in Fig. 11). The ring shows a coherent velocity field which is matched by our simulations (see below), as does the northern tail just south of the tidal dwarf. We will show in Sect. 3.3.6 that the kinematics of the TDG itself is probably decoupled from that of its host tail.
#### 3.1.3 Other unrelated objects
Weedman (1971) reported the presence of a blue stellar object, Weedman 2 or \[BOB94\] 0943-1403 (Bowen et al. 1994), near the tip of the tidal plume of NGC 2992. Based on low resolution spectroscopic data, Burbidge et al. (1972) claimed that it could be a QSO somehow “associated” with NGC 2992. We obtained in January 2000 with EFOSC2 at the ESO 3.6m telescope a spectrum of this object. It turns out to be a star. During the same run, we measured the redshifts of two uncatalogued background galaxies that are situated North of NGC 2992 at resp. $`\alpha =09^\mathrm{h}45^\mathrm{m}45.55^\mathrm{s},\delta =14^{}16^{}23.7^{\prime \prime }`$ and $`\alpha =09^\mathrm{h}45^\mathrm{m}51.14^\mathrm{s},\delta =14^{}15^{}43.8^{\prime \prime }`$ (J2000). Both lie at a redshift of $`z=0.111`$.
### 3.2 A numerical model of the interaction
In this paper, we are mainly interested in a first, approximate model of the NGC 2992/93 system, capturing its essential features while not necessarily providing a perfect fit to all its details. In a future paper, we plan to systematically refine this model and improve on its ability to fit the finer features of the system. The model developed here is particularly useful to date the various phenomena observed in Arp 245, to reconstruct its history and to predict its ultimate fate.
#### 3.2.1 Parameters of the collision
The first encounter models based on restricted N–body simulations date from almost thirty years ago when Toomre & Toomre (1972) were able to approximately match the morphology of four well–studied interacting systems. Since then, numerical methods have become much more sophisticated but the number of interacting systems with tidal tails that have been modeled in a self–consistent way with numerical simulations is still small (see e.g., Salo & Laurikainen 1993, Thomasson & Donner 1993, Mihos et al. 1993, Hibbard & Mihos 1995). The main reason for that is the huge parameter space spanned by the numerous free parameters of any such model (for instance, those defining the orbital plane), precluding simple methods to find a matching solution for an observed system. Moreover the parameters describing the structure of the colliding galaxies might not be well constrained by the observations. In that respect, Arp 245 appears as a particularly attractive system to model. The morphology and dynamics of each of the interacting partners are indeed fairly well known. The large parameter space of possible collision scenarios can thus be reduced drastically by making reasonable assumptions and “educated guesses” about the colliding galaxies.
A sketch of the adopted geometry for our simulations is shown in Figure 8. Coordinates where the orbital plane coincides with the xy–plane have been used while the orientations of the spin vectors of the two disk galaxies have been specified in terms of ordinary spherical coordinates $`(\theta ,\varphi )`$.
The nearly face–on orientation of the disk of NGC 2993, within $`20^{}`$, and the orientation of its tidal arm emanating in the north–west suggest that its spin vector points into the plane of the sky. We have therefore identified the plane of the sky with the plane of the disk of NGC 2993, i.e., the observer is located in the opposite direction of the spin vector of NGC 2993. The disk of NGC 2992 is seen about 70 inclined to the plane of the sky with the upper east side being closer to the observer (Chapman et al. 2000). Hence we have assumed that the spin vectors of the disks of NGC 2992 and NGC 2993 enclose exactly an angle of 70 with each other. The prominence of both tidal tails suggests a largely prograde encounter of the two galaxies, with the vector of the orbital angular momentum presumably lying somewhere within the cone spanned by the two spin vectors, or at least relatively near to it. To further simplify things, we have here assumed that the orbital angular momentum lies in the plane spanned by the two spin vectors.
The velocity field shows that the center of NGC 2992 and its northern tail is blue–shifted with respect to NGC 2993. Also note that the well–developed tails of the two galaxies indicate that the system is seen after its first encounter (as will be confirmed below). In this phase, NGC 2992 and its tail can be expected to move towards the negative y–axes. The tidal response of the stellar disk of NGC 2993 seems relatively weak compared to that of NGC 2992 which is in the process of forming a dwarf galaxy out of its tidal debris. This suggests that NGC 2992 is colliding almost fully prograde. Finally, we have restricted ourselves to parabolic, zero–energy encounters. The angular momentum of the orbit has been expressed in terms of the minimum separation $`b`$ of a corresponding Keplerian collision.
The galaxies themselves were modeled with a massive and extended dark halo with an adiabatically modified NFW–profile, an embedded exponential stellar disk, a stellar bulge, and a gas distribution in the disk (see details on the numerical representations of these models in Springel & White 1999). Note that in many observed disk galaxies the H I gas is substantially more extended than the exponentially distributed stars. To be able to follow the gas at large radii, we have split the available gas into two components of equal mass, one distributed just like the stars and one forming a more extended component with a constant surface mass density.
#### 3.2.2 Numerical simulations
Using fully self–consistent simulations carried out with the Tree–SPH code GADGET, we have coarsely explored the parameter space remaining under the above assumptions until a best–matching solution was found. More specifically, the numerical model presented below is specified by the following parameters. The virial velocities of the two galaxies were set to $`V_\mathrm{c}=120\mathrm{km}\mathrm{s}^1`$, i.e. each of them has a total mass of $`4.0\times 10^{11}h^1\mathrm{M}_{}`$. We have assumed that a fraction $`m_\mathrm{d}=0.05`$ of the total mass resides in the disk, 30% of it in the form of gas, the rest as collisionless stars. A fraction $`m_\mathrm{b}=0.016`$ of the total mass has been put into a central stellar bulge. Adopting a spin parameter<sup>8</sup><sup>8</sup>8The spin parameter $`\lambda `$ is defined as $`\lambda =\frac{L|E|^{1/2}}{GM^{5/2}}`$, where $`L`$ is the angular momentum, $`E`$ the total energy, and $`M`$ the mass of the galaxy. of $`\lambda =0.06`$, the resulting exponential scale length of the disks is $`R_\mathrm{d}=2.7h^1\mathrm{kpc}`$. Half of the gas has been distributed in the disk like the stars, and the other half in a disk of radius $`8R_\mathrm{d}`$ with constant surface mass density. For the orbit of the collision, we adopted the parameters $`\theta _1=50^{}`$, $`\varphi _1=90^{}`$, $`\theta _2=20^{}`$, $`\varphi _2=270^{}`$, and $`b=4h^1\mathrm{kpc}`$ (see Fig. 8). The numerical simulation shown here used 60000 particles for the dark matter, 40000 for the stellar disk, 10000 for the stellar bulge, and 40000 for the gas. Cooling of gas, star formation and supernova feedback were modeled as in Springel (2000). The cooling function used to describe radiative dissipation effectively cuts off at $`10^4`$ K, when the gas becomes neutral. Some of the gas can reach somewhat lower temperatures by adiabatic expansion. Star formation is assumed to convert gas into collisionless stellar material at a rate given by a simple Schmidt-type law, and supernovae feedback is modeled by an effective turbulent pressure term. A full description of the thermodynamic properties of this model is given in Springel (2000).
The simulation begins when the galaxies’ dark halos just start to touch, i.e., at a separation of $`240h^1\mathrm{kpc}`$. It then takes about $`1h^1\mathrm{Gyr}`$ to reach the perigalacticon for the first time, when strong tidal forces eject stars and gas out of the disks. This material forms the pronounced bridges and tails in Arp 245 when the galaxies separate again. Eventually, they fall back together for a second encounter, which takes place about $`700h^1\mathrm{Myr}`$ after the first collision. One hundred $`\mathrm{Myr}`$ later the galaxies completely coalesce and form a single merger remnant. This time evolution of the system is shown in Figure 9.
From Figure 9, it is clear that we observe the system NGC 2992/93 at a time between the first and second encounter. In this phase of the collision, the tidal features are very well defined, whereas they are going to be much more diffuse at later stages of the collision, which is a general result from simulations of other systems. Our numerical simulation provides the best match for the morphology of the system at a time of around $`100h^1\mathrm{Myr}`$ after the first encounter, i.e. $`1.1h^1\mathrm{Gyr}`$ after the start of the simulation<sup>9</sup><sup>9</sup>9Near perigalacticon the timestep in our adaptive timestep code is of order $`0.05h^1\mathrm{Myr}`$.. Figure 10 shows the face–on projection of the numerical model at this time. After matching the morphology it becomes of course interesting to study the velocity field of the gas distribution at this moment of the interaction. As Figure 11 shows, the overall pattern of gas flow is well reproduced by our model. Also, the striking ring of gas around NGC 2993 has nicely formed in our model. Most of this gas is simply stemming from the long tail which is pulled from the extended HI disk. The inclination of the disk relative to the orbital plane helps to curl up the tail in projection to an almost closed ring. There is actually some shocked gas from the region of the bridge that helps closing the ring.
#### 3.2.3 Limitations of the model
The simulation shown in this paper corresponds to the model which so far provided the best match to the morphology of the encounter. While we are confident that the orbital geometry is reasonably well determined in this model, the internal structure of the galaxies is less well defined and may be subject to revision when we refine the model using more detailed comparisons between the simulation and the observational data. In particular, we have here assumed that the two colliding galaxies have the same mass and the same internal structure, an assumption that is unlikely to hold in reality.
Besides, the dwarf galaxy FGC 0938 has not been taken into account. Having preserved a flat disk, this galaxy is most probably plunging into the system and has not yet faced its first encounter with NGC 2992/93. Therefore, this intruder has not yet affected the interaction between the two spirals much. Later on, this galaxy might however slightly perturb the merging history of the system.
Perhaps the main shortcoming of the present model is that it does not form the tidal dwarf galaxy seen in the northern tail of NGC 2992 although there is at least some gaseous overdensity at about the right place. It remains to be seen whether the tidal dwarf can be produced by a simple encounter model like the one examined here or whether additional physics or more sophisticated collision scenarios have to be invoked.
### 3.3 The tidal dwarf galaxy candidate, A245N
In this section we will study in detail the properties of the tidal object observed at the tip of the northern tail of Arp 245. We will henceforth refer to it as A245N.
#### 3.3.1 Morphology and structural parameters
Identifying a tidal dwarf in its host tail is obviously a difficult task. Such a problem which raises the basic question of the definition of a tidal dwarf galaxy will be addressed in Sect. 4.1.1. In a first approach, we have isolated the TDG candidate based on a morphological criterium. We have considered as belonging to A245N all stellar and gaseous material in the tail located inside the isophote $`\mu _\text{B}`$=24.5 $`\mathrm{mag}\mathrm{arcsec}^2`$, north of $`\delta =14\mathrm{°}18\mathrm{}`$ (see Fig. 2). This region contains the bulk of the atomic and ionized gas and appears to be “detached” in optical images (see Fig. 1). Table 5 lists the main properties of A245N. Figure 12 presents images of A245N at different wavelengths from the optical to the near–infrared. The tidal object has been detected in all BVRJHK bands. With an absolute blue magnitude of $`M_\text{B}=17.2`$, A245N belongs to the bright end of the dwarf galaxy population. It is actually as luminous and as extended as the LMC. Its surface brightness profile, shown in Figure 13, has been computed in the B and R bands following Papaderos et al. (1996a, b). It is exponential up to a radial distance of 25$`\mathrm{}`$ (3.7 kpc) and drops beyond. The extrapolated central surface brightness in the B band is 22.4 $`\mathrm{mag}\mathrm{arcsec}^2`$ and the exponential scale length of the disk is 19″ or 2.8 kpc. Put on the archetypal absolute magnitude vs. surface brightness and scale length diagrams, A245N occupies the locus of low surface brightness dwarf irregular galaxies (see Fig. 9 and Fig. 10 in Patterson & Thuan 1996).
#### 3.3.2 Stellar populations
Tidal tails contain two basic types of stellar populations. The first category includes stars older than the age of the interaction. They were originally formed in the parent galaxy from which they have been pulled out. Numerical simulations indicate that the stars now found at the tip of the tidal tails, i.e., in the TDGs, initially belonged to the outskirts of the parent disk. The second category is made of stars younger than the interaction formed in situ in the tails after the collapse of tidally expelled H I material. For nearby systems, deep color–magnitude diagrams would be the ideal tool to disentangle the first and second generation stars. At the distance of Arp 245, the galaxy can unfortunately not be resolved into stars. One has to rely on techniques based on the comparison of integrated broad–band photometric data with predictions of evolutionary synthesis models.
As a first step towards determining the age of the populations of A245N, we have compared its optical color with that of its parent galaxy. The B–R profile of both objects are displayed together in Figure 15. The external region of NGC 2992, beyound $`r=30\mathrm{}`$ (4.5 kpc) appears to have the same color as the tidal object except in its inner 10″. Their spectral energy distributions remain similar over a larger wavelength range, as shown in Figure 16 which presents the SEDs of A245N and of the outskirts of NGC 2992 ($`r>30\mathrm{}`$). Each color index differs by less than 0.1 mag, within the photometric errors. Note that the data have not been corrected for internal extinction as this is highly uncertain (see Sect. 3.3.3). Hence, it is clear that the stellar population of A245N is currently dominated by stars pulled out from the disk of the parent galaxy.
When analyzed in detail however, the stellar population of A245N does not appear to be completely uniform. First, the object hosts several stellar clusters showing a range of colors (see Fig. 12). Their faintness and the large photometric errors prevent us from deriving from this color spread differences in ages or metallicities for the individual clusters<sup>10</sup><sup>10</sup>10Obvious bright Galactic stars have been subtracted from the images. However, some faint foreground stars could have been mistaken for stellar clusters. The brightest blue clump in the B–band image, for which we could obtain a deep optical spectrum, does not exhibit the strong emission lines expected for an H ii region. However the detection of faint absorption lines (in particular, Ca H&K $`\lambda \lambda `$ 3924,3968) at the right velocity confirms that it is a genuine star cluster in the TDG.. Moreover, our spectroscopic data indicate that OB stars are present in A245N and therefore that the tidal object is currently forming stars. Figure 17a displays the H<sub>α</sub> map towards A245N. Several individual clumps are concentrated in the northern region where the H I emission peaks. Two of the five brightest have counterparts in the broad–band images (see Fig. 12) and could be associated with older star clusters. The current star formation episode weakly affects the global photometry of the galaxy. The $`\text{V}\text{K}`$ color map shown in Figure 12 (where the H<sub>α</sub> contours have been superimposed) indicates that star–forming regions are bluer by about 0.1 mag with respect to quiescient regions where no ionized gas is detected. This variation appears to be very small compared to predictions of photometric models. For a pure instantaneous starburst, the $`\text{V}\text{K}`$ index might change in the first 100 Myr by 1 to 3 mag, depending on the metallicity (Leitherer & Heckman 1995). Taking into account the old stellar component, one would get however a smaller color evolution.
We have estimated quantitatively the relative contribution of young to “old” stars using an evolutionary population synthesis model developed at the Göttingen observatory (Krüger et al. 1991). This code has been updated to include the bimodal star formation history of tidal dwarfs (see details in Weilbacher et al. 2000). We have only considered solar metallicities and a Scalo IMF. We have first reproduced the spectral energy distribution (SED) of the underlying population of A245N, assuming that it is similar to the SED of the external regions of its parent galaxy, NGC 2992. The best fit, shown in Figure 16, was obtained by a template for an Sb galaxy with a characteristic age of 5 Gyr. Note that this result, which is valid for the outer disk of the spiral, is consistent with the global morphological type of the galaxy being classified as an Sa. The effect of extinction is indicated in the figure. We have then simulated an additional burst, varying its shape, strength and duration and tried to fit the SED of the TDG candidate. Taking into account the photometric errors and the dispersion between the observed SED and the best model, we found that stars younger than 100 Myr, i.e., stars formed in situ in the tail, cannot contribute for more than 2% to the overall stellar mass of the tidal object that we estimated from our model to be $`3\times 10^9\text{M}\text{}`$. With such a low burst strength, A245N differs from most of the tidal dwarf galaxies studied so far. In particular, the southern TDG in the Arp 105 system appears much bluer than its parent galaxy; young stars could amount in this object to about 25% (Fritze-v. Alvensleben & Duc 1998) of the stellar mass. The very blue tidal dwarfs in the NGC 5291 system also seem to be completely dominated by stars formed in situ in an H I tail (Duc & Mirabel 1998). A possible explanation for the difference between those systems and Arp 245 might be that the latter object suffered its first encounter only 100 Myr ago, which might not have been a long enough time span to go through multiple bursts of star formation in order for the young population to dominate the old one.
#### 3.3.3 Current star formation and extinction
We have obtained longslit spectra of some of the H<sub>α</sub> condensations of A245N. Our two slit positions were roughly aligned along the knots labeled TDG 1 and TDG 4 (see Fig. 17a). A third spectrum, of TDG 3, was obtained during our MOS run at the ESO 3.6m. The spectra of TDG 3 and TDG 4 are shown in Figure 17b and the spectrophotometric data corresponding to the knots with the highest signal to noise are listed in Table 6. The detected emission lines have relative fluxes typical of H ii regions ionized by young OB stars. There is no doubt that the H<sub>α</sub> map towards the tidal dwarf traces regions of star–formation. Despite the strong contribution from stars older than 100 Myr noted before, no absorption lines are visible in the spectra. The low surface brightness of this stellar component might account for this result.
A Balmer decrement H<sub>α</sub>/H<sub>β</sub> of 5.7 $`\pm `$ 1 has been derived from the spectrum of TDG 4. It implies an intrinsic extinction of about $`A_\text{B}=2.6\pm 0.6`$ mag, or E(B–V)=0.6, a value much higher than the typical E(B–V)=0.1 of irregular galaxies (Hunter 1982) and about twice the value towards 30 Doradus in the LMC (Greve et al. 1991). The discrepancy could suggest a very localized obscuration in the H ii regions and hence a clumpy dust distribution. Slightly higher absorptions have been derived in TDG 1 and TDG 3, albeit with a much larger uncertainty. Besides, H ii regions in spiral disks also show comparable values of E(B–V) (Dufour et al. 1980).
The equivalent width of the H<sub>β</sub> line in TDG 4, about $`12\pm 2`$ Å, is consistent with a starburst age of 5–7 Myr, according to the models by Cerviño & Mas-Hesse (1994) for which an instantaneous burst with no underlying old component and a solar metallicity have been assumed. Given the age of the interaction, these newly formed stars were born in situ in the tidal tail.
We have estimated the star formation rate of the TDG from the total line emission extracted from the narrow–band H<sub>α</sub> image. The fluxes have been decontaminated for \[NII\] emission using the mean \[NII\]/H<sub>α</sub> line ratio measured along our long–slit spectra. The calibration of Kennicutt (1998) has been used to derive the SFR from the H<sub>α</sub> luminosity corrected for galactic extinction. It assumes a Salpeter IMF (0.1–100 M) and solar metallicity. We obtain an estimate of 0.03 M/yr. Taking into account the intrinsic extinction derived in the optical from the Balmer decrement, one would get a SFR of 0.13 M/yr. The star formation rate per unit area (in kpc<sup>-2</sup>), $`\mathrm{log}(SFR/A)`$, is -3.1 (-2.5 after correction for extinction). The normalized SFR derived by Hunter (1997) for a sample of Im galaxies covers a large range of 4 dex, between -5 and -1, with a relative peak at -3.5. The SFR we find for A245N is somewhat higher than this and is comparable to that of the LMC/SMC (Kennicutt & Hodge 1986). But despite these similarities, A245N appears much redder than star–forming Irrs and more typical of a spiral. The simple fact that the global photometry of the tidal dwarf does not seem to be much affected by the episode of in situ star formation implies that the latter cannot have lasted long with a constant SFR of $``$ 0.1 M/yr. Given the short time scale available — less than 100 Myr, a constraint set by the simulations — it is likely that the starburst started only late, perhaps less than 10 Myr ago, as suggested by the equivalent width of the H<sub>β</sub> line.
#### 3.3.4 Gas content
A245N is a gas–rich object. Its total H I reservoir of $`9\times 10^8\text{M}\text{}`$ could sustain star–formation episodes at rates of 0.1 M/yr for several Gyr. Its H I mass to blue luminosity of 0.7 is typical of dwarf irregular galaxies (according to Roberts & Haynes 1994, the median value for Sm and Im UGC galaxies is 0.66 ). A245N is the first tidal dwarf where molecular gas has been unambiguously found. Extended CO(1–0) emission has been detected with the IRAM 30 meter antenna over a region of 40$`\mathrm{}`$ centered on the H I peak (Braine et al. 2000). The H<sub>2</sub> mass derived from the combined CO line flux is $`1.4\times 10^8\text{M}\text{}`$. Braine et al. (2000) argue that the molecular gas has formed in situ in the H I clouds, directly fueling the star–formation regions. The very similar width of the CO and H I lines supports this idea.
The star formation efficiency, i.e the ratio of the star formation rate to the available molecular gas, $`L_{\text{H}\alpha }`$/$`M_{\text{H}\text{2}}`$ is 0.006, a value much closer to that of early type spirals (about 0.01 $`L_{\text{H}\alpha }`$/$`M_{\text{H}\text{2}}`$) than to that of irregular galaxies ($``$ 0.05 $`L_{\text{H}\alpha }`$/$`M_{\text{H}\text{2}}`$, Rownd & Young 1999). Smith et al. (1999) detected molecular gas in the extended tail of the interacting galaxy Arp 215 (NGC 2782) where they derived a star formation efficiency three times lower than in A245N. Whereas the molecular gas in Arp 215 was found at the base of the tail and has probably been stripped from the disk of the parent galaxy, the CO detection in the tail of NGC 2992 is more directly related to A245N. We have carried out a complete CO mapping of Arp 245 with the 15m–SEST antenna at la Silla in November 1999. The only CO cloud detected outside the parent disks is associated with A245N at the location of the IRAM source, strongly suggesting that the CO found here is related to the increase in density due to self gravity in the tidal object and subsequent gravitational collapse of the tidal H I material.
#### 3.3.5 Metallicity
We have estimated the oxygen abundance of A245N using the spectrophotometric data of the H ii regions TDG 3 and TDG 4 listed in Table 6. A direct determination of the abundance is not possible as the temperature sensitive \[OIII\]<sub>λ4363</sub> line is not detected. Semi–empirical calibrations were applied instead and several diagnostic diagrams tested. We first used the approximate calibrations of Edmunds & Pagel (1984) based on the $`R23=(\text{[OII]}\text{λ3727}+\text{[OIII]})/\text{H}\text{β}`$ line ratio and applied it to TDG 3. The relation between $`R23`$ and O/H is ambiguous as a single value of $`R23`$ corresponds to two values of the oxygen abundance. However, according to McGaugh (1994), the \[NII\]/\[OII\] line ratio provides a useful diagnostic for choosing between the lower and upper branch of the $`R23`$–O/H relation. log(\[NII\]/\[OII\]) $`>`$ -1 favors the highest metallicities. For TDG 3, we measured log$`(\text{[NII]}/\text{[OII]})=0.4`$, hence we selected the upper branch, and derived from $`R23`$ an oxygen abundance 12+log(O/H) of 8.6 $`\pm `$ 0.2. For TDG 4, the \[OII\] line is outside the instrumental wavelength range. Our estimate for this H ii region is therefore based on the less reliable $`\text{[OIII]}/\text{H}\text{β}`$ line ratio as also calibrated by Edmunds & Pagel (1984). We derived 12+log(O/H) $`=`$ 8.7 $`\pm `$ 0.2. However, these calibrations are not unique and depend very much on the ionization parameter $`U`$ (e.g., McGaugh 1991) and on the temperature of the thermal ionization source $`T`$. We therefore decided to run a photoionization code trying to fit all our data on TDG 3 in a consistent way. For this we used cloudy (Ferland 1996). Our constraints were the \[OII\]/H<sub>β</sub>, \[OIII\]/H<sub>β</sub> and \[NII\]<sub>λ6584</sub>/H<sub>α</sub> observed line ratios. We let the ionization parameter and temperature vary; we assumed a hydrogen density of 100 cm<sup>-2</sup> and selected a metallicity of about one half of solar (12+log(O/H)=8.6). With log(U)=-3.4 and T=56200 K, the code accurately reproduced the oxygen line intensities. The \[NII\]<sub>λ6584</sub>/H<sub>α</sub> line ratio could only be obtained by increasing the relative nitrogen abundance log(N/O) to -1.1. This abundance is quite high compared to classical dIrrs and BCDGs (Kobulnicky & Skillman 1996) but consistent with that measured in the outer regions of spiral disks (Ferguson et al. 1998). Models with solar metallicity failed to reproduce the \[OII\]/H<sub>β</sub> line ratio.
The spectrophotometric study of a sample of tidal dwarf galaxies by Duc & Mirabel (1999) indicates that TDGs have an average oxygen abundance of Z/3 which is independent of their absolute blue magnitude (see Fig. 18). Made of pre–enriched material, TDGs are more metal rich than isolated dwarf galaxies of the same luminosity. This property may be used to identify recycled dwarf galaxies and to investigate the origin of their building material in the disk of their progenitors, as discussed in Sect. 4.2.
#### 3.3.6 Kinematics
At first sight, the mean H I velocity field along the northern tidal tail (see Fig. 11a) appears to roughly follow the model field governed by streaming tidal motions (Fig. 11b). Velocities range between 2270 km s<sup>-1</sup> at the base of the tail to 2140 km s<sup>-1</sup> at its tip. A detailed analysis of the H I datacube shows however a possible evidence for a sub–structure associated with A245N which is visible in position–velocity (PV) diagrams. This is best seen after rotating the HI datacube ($`\alpha ,\delta ,\mathrm{V}_{\mathrm{Hel}}`$) clockwise by 15 degrees so that the tidal tail points upwards. Figure 19 presents three PV–diagrams taken along a direction perpendicular to the tail axis. One cut is taken through the optical center of A245N, the other two are offset along the tail by $`45^{\prime \prime }`$. The cut through the tidal object shows a velocity gradient of $`65\text{km s}\text{-1}`$, peak to peak extending over almost $`50^{\prime \prime }`$ spatially. This might suggest that part of the tidal H I is kinematically decoupled from its host tail showing some evidence for solid body rotation with its kinematical major axis roughly perpendicular to the tail. This component contains the bulk of the H I in the tail (more than 60%) and corresponds spatially to the optical tidal object. Due the scarcity of H ii regions in the tail, we could not determine, based on our longslit data, any velocity profile in the ionized gas component.
Based on our low resolution H I data only we estimate a lower limit for its dynamical – virial – mass (i.e. uncorrected for inclination) of $`9\times 10^8\text{M}\text{}`$, identical within the errors to the H I mass and three times smaller than the estimated stellar mass. This could indicate that a large fraction of the old stellar population present in A245N does not belong to the kinematically decoupled region of the galaxy. If reliable, our low value of the dynamical mass would be consistent with a very low dark matter content, as expected for TDGs (Barnes & Hernquist 1992).
The H I velocity dispersion towards the TDG candidate, 25–30 km s<sup>-1</sup>, is much higher than in a quiescent disk. The turbulence is hence dynamically important.
#### 3.3.7 Conclusions on the properties of A245N
Our detailed study of A245N indicates that the tidal object has apparently properties ranging between those of dwarf irregular galaxies (structural parameters, gas content, star formation rate) and those of spiral disks (metallicity, extinction, star formation efficiency, stellar population). A straightforward explanation would be the intrinsic hybrid nature of the galaxy which is made of disk material but results from the collapse of gaseous clouds that have masses and characteristics of dwarf galaxies.
## 4 Discussion
Tidal dwarf galaxies candidates have now been found in a variety of interacting galaxies, i.e., in disk–disk systems (e.g., NGC 4038/39, Mirabel et al. 1992), disk–spheroid systems (e.g., Arp 105, Duc & Mirabel 1994), gas–rich spheroid–spheroid systems (e.g., NGC 5291, Duc & Mirabel 1998) and advanced mergers (e.g., NGC 7252, Hibbard et al. 1994). These interacting systems have a variety of environments: field (e.g., NGC 4038/39), compact groups (e.g., the Hickson groups, Hunsberger et al. 1996) and clusters (e.g., Arp 105, NGC 5291). On the other hand, not all interacting systems with gaseous tails form dwarf galaxies, as recently shown by Hibbard & Yun (1999). Together with the overall properties of the parent galaxies and history of the collision, the environment — i.e., density of companion galaxies, density of the intergalactic medium and strength of the associated ram–pressure — might play a role in the formation and more significantly in the evolution and survival time of tidal dwarf galaxies. The exact contribution of each process is not yet known. In that respect, Arp 245 is a simple system situated in a relatively poor environment. It may hence give us some clues about the required minimum conditions for the creation of such objects. Incorporating our data on Arp 245N with data from the literature on similar such systems, we will in the following review some of the outstanding questions regarding TDGs, starting with a proper definition.
### 4.1 Identifying Tidal Dwarf Galaxies in tidal features
#### 4.1.1 A working definition of a Tidal Dwarf Galaxy
Luminous, star–forming knots are commonly found along the tidal features of interacting systems. Will all of them become independent from their parent galaxies and hence become true galaxies ? Given their particular environments, one may doubt whether all of them will manage to survive for more than one Gyr. Therefore, in order to distinguish between such a short–lived object and a true “Tidal Dwarf Galaxy” we need a stricter definition. We propose as a working definition for a “Tidal Dwarf Galaxy” an object which is a self–gravitating entity, formed out of the debris of a gravitational interaction.
This restrictive definition ensures that TDGs are not simply agglomerated debris of collisions but that they are active objects that have their own dynamics and a potential well that is strong enough to sustain themselves against internal or external disruption for at least 1 Gyr. Such kinematically decoupled tidal objects and associated rotating gas clouds have been found in the interacting systems Arp 105 (Duc et al. 1997) and NGC 5291 (Duc et al. 1997).
#### 4.1.2 Possible mis-identification of TDGs
Under the definition given above, a TDG should be considered a separate entity based on dynamical rather than morphological criteria. Tidal objects may not always look as “detached” or appear as contrasted entities in optical images, nor should they always have very distinct colors if they have not yet formed a substantial proportion of luminous young stars. The true optical morphology of the galaxy might hence be hidden by an old unbound stellar component pulled out from the disk of the parent galaxy. Computing the integrated properties of a TDG, one should in principle only consider bound tidal material corresponding to the kinematically detached part of the tail. In practice, this is difficult since it would require two–dimensional high–resolution velocity data. The integrated properties of a tidal object might hence well be contaminated.
Moreover, projection effects should be taken into account. Tidal tails are generally curved (see Fig. 10) and seen edge-on, they exhibit at their apparent tip superimposed material from the near and back side. The resulting projected structure might be as luminous as a dwarf galaxy and hence be mistaken as a single object. Even the streaming motions of an expending tail could mimic at its projected bend the dynamical signature of a rotating TDG (Duc et al. 1997). We could see there a range of velocities due to the contribution of several velocity vectors at different angles. It is however expected that, if this is the case, the resulting velocity field should be asymmetric.
Besides, some TDGs could, in principle, be mistaken with dwarf galaxies preexisting the collision, and the tail linking them to a parent galaxy would then be a bridge.
#### 4.1.3 Is A245N a tidal dwarf galaxy ?
Several observational facts are inconsistent with the hypothesis that A245N is a preexisting dwarf galaxy. First of all, its high metallicity is a clear sign of a recycled origin. Its colors are remarkably similar to those of the parent’s outside disk. Finally our numerical simulations show that no third body is required to reproduce the morphology of the interacting system.
The only massive tidal object in Arp 245 is found at the tip of an almost edge–on tail, a case for which projection effects might be important. However the H I, H<sub>α</sub>, and CO lines peak at the same location in A245N and have similar velocities. The projection hypothesis would require that all these components are rather uniformly distributed along the tails which is not seen in the face–on eastern tail, at least for the H I gas. Besides, as discussed in Sect. 4.3.2, ionized gas is only observed above a critical H I column density. The H<sub>α</sub>/H I coincidence is a strong indication that both phases of the gas are physically linked. It is hence more likely that the various gaseous components towards A245N form, at our spatial and velocity resolution, a single entity. The surface brightness profile of A245N, which is very well fit by an exponential profile (see Fig. 13), might suggest the presence of a stellar disk. However, the profile expected for a pure tidal tail without bound objects in it is not yet well known. In our low resolution simulations, the SBP at the tip of the numerical tail seems to diverge from an exponential profile (see Fig. 14). Note however that the SBP measured towards A245N corresponds to that of the old stellar component pulled out from the disk of its parent galaxy. Given the young age of the interaction and the non–dissipative nature of stars, it would be surprising that the stellar population is already relaxed and bound with the gaseous entity.
The kinematical independence of A245N is yet difficult to assess. H I position–velocity diagrams show a fairly localized symmetric velocity gradient reminiscent of solid–body rotation with an axis perpendicular to the tail (see Sect. 3.3.6). But this is at the limits of the spatial and velocity resolution of our VLA observations. Besides, position–velocity diagrams of the simulated system computed at the same locations as the H I ones appear strikingly similar (see Fig. 19 and Fig. 20). Therefore, either a kinematically decoupled tidal object has already formed in our simulations, or the H I PV–diagrams mostly reflect streaming tidal motions. Clearly higher resolution numerical simulations and H I observations will shed more light on this. On the other hand, the high velocity dispersion measured in the H I component of A245N indicates that turbulence is almost as important dynamically as the putative rotation.
Most probably, A245N is a tidal dwarf galaxy observed in the early phases of its formation. Its gas might just be becoming self–gravitating, overcoming turbulence and streaming tidal motions. This would be consistent with the young age of the interaction.
### 4.2 Origin of the TDG building material
TDGs consist of material that has been pulled out from parent galaxies. But where exactly do the tidally expelled stars and gas clouds come from? Numerical simulations by Hibbard & Mihos (1995) show that particles found at the end of the simulated tidal tails where TDGs are usually observed once belonged to the outermost regions of the parent disk. This is confirmed by our numerical simulations of Arp 245 (see Fig. 10). The gas particles towards the tip of the TDG were originally in the outer gas disk, at a radius of about eight times the exponential scale length of the stellar disk. The simulations show however that the inner parts of the tail come from regions at smaller radii in the disk.
The prediction of the simulations can be observationally checked by noting that the metallicity of TDGs reflects that of the region where their building material has come from and been pre–enriched. Spiral galaxies show strong metallicity gradients, from above solar in the core to one tenth of solar in the outer regions beyond the optical disk (Ferguson et al. 1998). An oxygen abundance of one third solar, the median metallicity of the ionized gas in TDGs (see Fig. 18), roughly corresponds to a radius of $`R_{25}`$. The slightly higher metallicity of TDG A245N suggests an even smaller galactic radius.
Hence, it is improbable that most TDGs are made up of material from much further out, unless a strong local enrichment has occured or unless the parent galaxy had an unusual metallicity distribution prior to the collision. Enrichment of the ionized gas by a burst of star formation within a time scale of 100 Myr or less – the maximum age of the TDG A245N as determined from numerical simulations – is unlikely. Studies of classical metal poor dIIrs and BCDGs have indeed shown that enrichment of the ionized gas does not seem to be efficient in low mass galaxies even over longer periods. Parent galaxies with perturbed metallicity gradients, e.g. higher than normal, might be found in the class of Seyfert galaxies to which NGC 2992 belongs. Their active nucleus could pollute the interstellar medium via large–scale outflows. Whereas the absence of a steep abundance gradient has been reported by Evans & Dopita (1987) in the prototype Seyfert 2 galaxy NGC 1068, this property was not found in other Seyfert galaxies (Schmitt et al. 1994, Storchi-Bergmann et al. 1996). Statistical data are so far missing to link any anomalous metallicity gradient distributions with the nuclear activity and related ejection phenomena.
Clearly, precise measurements of the metallicity distribution in the parent galaxies <sup>11</sup><sup>11</sup>11This study is extremely difficult to do in NGC 2992 because of the absence of H ii regions in its disk. Metallicity measurements are very uncertain in the ELRs of active galaxies. and more detailed numerical simulations of interacting systems would be required to further investigate the origin of the building material of TDGs.
### 4.3 Required conditions for the formation of a TDG
#### 4.3.1 Location of TDGs in the tail
Numerical simulations by Barnes & Hernquist (1992), Elmegreen et al. (1993) suggest that TDGs should form from gravitational instabilities that grow in the debris of the collision. These objects appear as bound condensations distributed all along the tidal features. Some tails of interacting systems indeed host numerous faint blue knots showing signs of star formation (Schombert et al. 1990, Weilbacher et al. 2000). One of the two most spectacular objects of that kind are NGC 4676, the Mice, (Hibbard & van Gorkom 1996, Sotnikova & Reshetnikov 1998) and IRAS 19254–7245, the Superantennae (Mirabel et al. 1991, Mihos & Bothun 1998). However, systems containing substructures that comply with our definition of a TDG generally show the star formation being concentrated in a single object located at the tip of the stellar tail, such as TDG A105N (Duc et al. 1997). And this is also, of course, the case for A245N, which is located at the end of the tidal plume stemming from NGC 2992.
So, why is there only one single massive star forming clump as soon as a TDG has been formed and why is it located at the tip of the tidal tail ?
Excluding the projection effects that could lead to a mis–identification of a TDG at the tip of the tail, part of the explanation might have to do with the timescale for tidal material to fall back to the merger. This rate scales as $`t^{\frac{5}{3}}`$ (Hibbard & Mihos 1995) and is highest shortly after perigalacticon. Therefore, bound objects that might have formed at the base of the tail would have already fallen back towards their progenitor, leaving as the only viable region for the formation of a TDG the tips of tails. Modeling the prototype merger NGC 7252 with N–body simulations, Hibbard & Mihos (1995) have found that about half of the tidal material situated at the base of the already formed tail (130 Myr after periapse) falls back within 130 Myr. Arp 245 is observed about 100 Myr after perigalacticon and only about 50 Myr after the tails have developed (see Fig.9). Over a time scale of 50 Myr, about one fourth of the initial tidal material might have fallen back already. A more precise study of the temporal evolution of the return of both gaseous and stellar material, based on our own numerical simulations of Arp 245, will be studied elsewhere.
Another hypothesis could be that the initial clumps, proto–TDGs, might have merged. Detailed high resolution simulations could test whether the formation of TDGs proceeds in a hierarchical way. If true, short–lived tidal features would not offer a suitable environment for the formation and survival of TDGs. In particular, tails that have a life expectancy greater than 1 Gyr are more likely to generate TDGs than bridges which only last for about $`10^8`$ yr (Struck 1999). Another obstacle faced by tidal bridges, illustrated for instance by Arp 245, is their intrinsic low surface density, a critical parameter as we will show in the next section.
#### 4.3.2 Conditions for star formation
The fact that all observed tidal dwarfs, i.e., kinematically independent objects in tidal features, are star–forming objects is partly a selection effect: TDG candidates have generally been selected for their brightness and blue color. On the other hand no bound tidal gas cloud is yet known which is not associated with a stellar counterpart. So the question raised in the previous section — where do TDGs form? — could be addressed by answering another one: what conditions are necessary for star formation (SF) to commence in tidal debris?
H I maps of interacting systems give some clues: all TDGs are found in H I tails and are roughly adjacent to local peaks in the H I column density map. This fact is not specific to TDGs, of course, and actually applies to SF regions in all types of galaxies, from spirals to dwarf irregulars, and might simply reflect the fact that H I is the raw material for any SF event. However, not every H I clump hosts a SF region. The onset of star formation apparently depends on an H I column density threshold, as first noted by Davies et al. (1976) for the LMC. Further studies of dwarf irregulars (Gallagher & Hunter 1984, Skillman 1987, Taylor et al. 1994) and low–surface brightness galaxies (van der Hulst et al. 1993, van Zee et al. 1997) suggest that this threshold is about 0.5 – 1 $`\times 10^{21}\text{cm}\text{-2}`$.
Neither the universality of this threshold, nor the underlying physics is known. It would be surprising if it would be the same in spiral disks, in dwarf galaxies, and in the more chaotic environment of interacting galaxies. And a word of caution, when making comparisons, one should aim for comparable linear resolution as the peak H I column density, because of beam dilution, decreases with decreasing resolution.
In the case of Arp 245 both tidal tails show a clumpy H I structure. Star formation occurs only in the TDG above densities of 2 $`\times 10^{21}\text{cm}\text{-2}`$. All condensations in the southern tail have their peak below $`5\times 10^{20}\text{cm}\text{-2}`$ (see Fig. 3) and no SF is associated with this tail. In turn, towards the spiral NGC 2993 ionized gas is observed above our sensitivity limit of $`0.4\times 10^{16}\text{erg s}\text{-1}\text{ cm}\text{-2}\text{ arcsec}\text{-2}`$ at those locations where the H I density exceeds $`10^{21}\text{cm}\text{-2}`$. The threshold hence appears to be two times higher in the tidal dwarf than in the spiral disk. The connection between H<sub>α</sub> and H I in the other spiral, NGC 2992, is more difficult to assess. Some ionization filaments are observed at large galactic radii where the H I has column densities well below $`5\times 10^{20}\text{cm}\text{-2}`$. However this apparent decoupling between the two phases of the gas is mainly due to the very nature of the ionizing source, probably hard UV radiation from the central AGN or from a shock induced process (Allen et al. 1999, Paper II).
As mentioned, the physical origin for the empirical H I threshold is not well understood (see review by Skillman 1996). One possible explanation is a gravitational instability, as elaborated by Kennicutt (1989) in the case of spiral disks. Because of the limited resolution and the obviously strong tidal forces, it is not possible to test the validity of this explanation for TDGs.
Alternatively, the threshold in $`N_{\text{HI}}`$ might correspond to a minimum required column density of dust to shield the gas from the UV radiation field (Federman et al. 1979, Skillman 1987). Given sufficient shielding, a molecular cloud can form from the H I allowing star formation to proceed. This critical density should be metallicity dependent since the dust–to–gas ratio varies with the element abundance (Franco & Cox 1986). It is then expected that recycled tidal dwarfs should have an $`N_{\text{HI}}`$ threshold lower than that of isolated metal–poor dwarf galaxies and similar to that of spiral disks. Table 7 lists the H I column density data available in the literature for several star–forming TDGs. Column 1 shows the peak $`N_{\text{HI}}`$ — or a range of values in case several TDGs are present — and column 2 the linear resolution. The majority of the TDGs so far studied have a column density of 2 – 3 $`\times 10^{20}\text{cm}\text{-2}`$, indeed three times lower than in dIrrs (Skillman 1987, Taylor et al. 1994). However, the metal rich, CO rich, and presumably dusty TDG belonging to Arp 245 has its H ii clustered in the central core of the H I cloud where the $`N_{\text{HI}}`$ exceeds $`2\times 10^{21}\text{cm}\text{-2}`$. Rather than stating that Arp 245N is the exception, one should perhaps rather question the low star formation thresholds in the other TDGs, noting that beam dilution for the other TDGs, which were observed at much lower resolution, could partly account for this large difference.
### 4.4 Models and numerical simulations of TDGs
Several theoretical models for the formation of TDGs have been proposed. Elmegreen et al. (1993) suggest that massive H I clouds form first in the outer disk of interacting galaxies. According to the Jeans criterion, the local high velocity dispersion induced by the collision allows the formation of clouds with masses much higher than in a quieter environment, up to $`10^8\text{M}\text{}`$. These clouds are then tidally expelled into the IGM where they collapse. Wallin (1990) put forward a geometrical torsion of the tail leading to a local density enhancement. Barnes & Hernquist (1992) propose a local amplification of pre–existing stochastic clumps in the tidal tails that start accreting low $`\sigma `$ material — in particular gas — and become gravitationally bound. The collected material will then collapse until rotation takes over.
Numerical simulations that might help selecting between these different theories of the formation of TDGs should also take into account the observed properties of these objects, and explain:
(a) knots, perhaps star forming, but not necessarily leading to TDGs, along the tidal tails; (b) star–formation whenever the H I clumps reach some critical column density, the origin of which could either be a gravitational instability within the tidally ejected gas, or shielding from the UV–radiation field once clumps get dense enough; (c) that in most systems eventually only one massive object per tail forms, located at the end of it; (d) the possible merging between individual preexisting clumps; (e) the time scale for the formation of a TDG — this can be quite short: in Arp 245, a TDG developed in less than 100 Myr; (f) the about $`1/3`$ of Solar metallicity in the TDGs; (g) the formation of TDGs in pure H I tails which is the case when the interaction involves gas–rich early type galaxies from which stars are much less efficiently pulled out than the gas (e.g., in NGC 5291, Duc & Mirabel 1998).
## 5 Conclusions
We have presented a multi–wavelength study of the interacting system Arp 245 (NGC 2992/93) and a preliminary numerical simulation of the collision. From our observations and models, we obtain the following results:
a) The system is observed at an early stage of the interaction, $`100\mathrm{Myr}`$ after first perigalacticon and 700 Myr before the final merger. At the current time two long tidal tails and a bridge have developed. N–body/SPH simulations reproduce fairly well both the optical and H I features, in particular the ring–like structure observed in the gaseous component. The VLA map reveals a third member to the interacting system: an H I rich dwarf galaxy, seen almost edge–on and looking unperturbed. This object is hence probably falling into the group.
b) H I in emission is found all along the two optical tidal tails where it accounts for respectively 60% and 80% of the H I of their progenitors, NGC 2992 and NGC 2993. The H I shows a peak at the end of the NGC 2992 tail where a gas reservoir of $`10^9\text{M}\text{}`$ supplies a star–forming, tidal object, A245N. H I is seen in absorption towards the active nucleus of NGC 2992. Star formation, as traced by H<sub>α</sub>, occurs in the main body of NGC 2993. In the tidal features, it is restricted to the tip of the northern tail at the location of the TDG A245N, where the H I column density is at least 1.5–2 $`\times 10^{21}\text{cm}\text{-2}`$.
c) The global properties of A245N range between those of dwarf irregular galaxies as far as its structural parameters, gas content and star formation rate are concerned, and those of spiral disks as for its metallicity, extinction, star formation efficiency and stellar population. In particular, although the object is actively forming stars, and has a large atomic and molecular gas reservoir, the bulk of its stellar population is still dominated by the old component pulled out from the parent galaxy. Whereas A245N appears to be a self–gravitating entity, our data lack spatial resolution to probe its proper dynamics. Most probably, A245N is a tidal dwarf galaxy that is still in the process of formation.
We are grateful to the support astronomers and night assistants from the 2p2, NTT and 3p6 teams which have helped us during our different observing runs at la Silla. Special thanks to Pierre Leisy for his precious support, to Vassilis Charmandaris who did the CO observations at SEST and to Polychronis Papaderos who has computed the surface brightness profiles of our objects. This work has greatly benefited from discussions with Jonathan Braine, Uta Fritze–v.Alvensleben, Simon White and from the very useful comments of the referee, Curtis J. Struck. P.–A. D. acknowledges support from the network Formation and Evolution of Galaxies set up by he European Commission under contract ERB FMRX–CT96086 of its TMR program. I.F. M. acknowledges support from CONICET/Argentina. This research has made use of the Lyon–Meudon Extragalactic Database (LEDA) supplied by the LEDA team at the CRAL–Observatoire de Lyon (France), as well as of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.
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# On the Observables Describing a Quantum Reference Frame.
## 1 Introduction
As Mach remarked at the end of the nineteenth century , from the physical point of view a frame of reference is defined by a material object of the same nature as the objects that form the system under investigation and the measuring instruments. Such an idea doesn’t conflict with classical mechanics: for example a rigid body can define a spatial origin and an orientation. The situation becomes more complicated in quantum mechanics: Heisenberg’s uncertanty relations forbid the exact determination of the position and the velocity of a frame. As noted by some authors , such an analysis not only contributes to remove a classical concept from quantum mechanics, but also reveals some surprising physical consequences, such as the “paradox of the quantum frames”. That is, if we have three frames of reference $`F_1,F_2`$ and $`F_3`$, the observables which describe the relation between $`F_1`$ and $`F_2`$ may be not compatible with the observables which describe the relation between $`F_2`$ and $`F_3`$, even if the systems don’t interact.
We follow an operational approach: the mathematical structures involved should have a direct physical meaning. From this point of view, a frame is determined by the procedures which transform an initial frame into the chosen one. The set of the transformations allowed by a relativistic theory is represented by the proper orthochronous Poincaré group $`𝒫`$. Each element of $`𝒫`$ can be identified by means of ten indipendent parameters, indicating the coordinates of the new origin in the Minkowsky space-time, the three components of the velocity and three angles of orientation. From a physical point of view these ten variables can be determined by their measurement performed on the physical object defining the frame: They have to be considered observables. Unfortunately spectral measures, the mathematical structures traditionally associated to the physical concept of observables in quantum mechanics, cannot describe neither simultaneous measurements of position and velocity, nor measurements of time (Pauli’s theorem). According to Gleason’s theorem, the natural generalization of spectral measures, compatible with the “Copenhagen interpretation”, is given by the so-called positive operator value measures (POVMs) .
An observable is often characterized by its transformation properties under a particular symmetry group. We can define a system of covariance as a POVM endowed with its covariance properties under some symmetry group . If a POVM reduces to a spectral measure, the structure so defined is called asystem of imprimitivity .
Following some hints which can be found in ref , in section 2 we illustrate a general construction procedure for covariant observables, which allows one to assign the statistical distribution of the outcomes to the state vector of the system on which the measurement is performed. The cornerstones of our procedure are three theorems. The first one (covariant dilatation) asserts that any system of covariance can be derived from a suitable system of imprimitivity to which it is linked by means of a suitable “intertwinig operator”. The second theorem is Mackey’s imprimitivity theorem, which allows us to find the most general form of a system of imprimitivity. The third one is the “intertwinig operator theorem”, which we derived in order to find the most general intertwining operator connecting the systems of imprimitivity to the unknown system of covariance. Its application is possible when the unitary representations of the symmetry group acting on the involved Hilbert spaces are decomposed into irreducible unitary representations.
We stress that our results are very general and allow one to describe all the possible measurements of a given observable, defined by its spectrum and its trasformation properties under a relevant symmetry group, performed on a physical system which is identified by its covariance properties under the same group. We do not introduce any model, but use general symmetry properties of the measurement, to get all the POVMs describing a choosen observable.
In section 3, the developed procedure is used for a new derivation of the most general POVM on Minkowsky space-time which is covariant with respect to the Poincaré group, found in ref by means of a different method.
Coming back to quantum frames in section 4, their description can be given by a system of covariance on the proper orthochronous Poincaré group, which is in this case both the parameter space and the symmetry group. Harmonic analysis on $`SL(2C)`$ and on the group of translations of $`𝐑^4`$ allows the decomposition of the most general unitary representation of $`𝒫`$ into irreducible unitary representations. In this way, the intertwining operator theorem can be applied and the most general probability distribution can be found.
A comparison with the Poincaré covariant POVM on Minkowsky space-time indicates the existence of some constraints. In particular the so-called baricentric measures cannot be obtained: In other words one cannot require that the coordinates of the origin coincide with the coordinates of the centre of mass of the physical system defining the frame. Moreover through an analysis of our results, one realizes that, for a complete description of a quantum reference frame, kinematical variables are not sufficient and internal degree of freedom of the system have to be involved.
Finally, the formalism we adopted also allows an alternative derivation of the paradox of the quantum frames. Indeed, in section 5, we derive the form of the POVMs describing the relative observables between a generical quantum object and a quantum reference frame. The observables describing the relations between non-interacting quantum frames are just a particular case.
We hope that an analysis of the variances of the probability distributions we have found will give a new class of indeterminacy relations.
## 2 Mathematical tools
In the traditional framework of quantum mechanics, the states of a system are described by rays in a complex separable Hilbert space, or by normalized positive trace class operators, while observables are described by self-adjoint operators or, equivalentely, by spectral measures. It is well known that the last ones cannot describe neither joint measurements of incompatible observables, nor measurements of time, indeed Pauli’s Theorem forbids the description by means of a self-adjoint operator of an observable canonically conjugate to an Hamiltonian with a semibounded spectrum. Nevertheless, as noted by some authors , the description of some measuring instruments requires a different mathematical structure, which can be recognized as a generalization of spectral measures: the so-called positive operator value measures (POVM). Gleason’s theorem assures they are the most general mathematical structures describing observables compatible with the probabilistic interpretation of quantum mechanics.
The analysis of the proof of Pauli’s theorem, shows that it is based on the covariance of time measurements with respect to time translations. This is not accidental, but shows the importance of simmetry in our discussion. Indeed, the requirement of precise covariance properties of the quantum measurement under some simmetry group $`𝒢`$ leads to the following definition of covariance systems .
Notation From now on:
1. $`S`$ will indicate a topological space which is locally compact and has a countable base of its topology. $`S`$ is called the “space of the possible results of the measurement”.
2. $``$ will indicate a complex separable Hilbert space. $``$ is called “the space of states (of the quantum system)”.
3. $`𝒢`$ will indicate a locally compact topological group which has a countable base of its topology. $`𝒢`$ is called the “symmetry group of the theory”.
###### Definition 1
Let $``$ be a space of states and $`S`$ be a space of possible results. A POVM on $`S`$ is a class $`\tau :=\{\tau (I)\}_{I(S)}`$, where $``$ is the $`\sigma `$\- algebra of Borel sets of $`S`$ and each $`\tau (I):`$ is a positive bounded operator such that:
1. $`\tau (I)\tau (\mathrm{})=0I`$
2. $`\tau (I_i)=\tau (I_i)`$
3. $`\tau (S)=1`$
where $`\{I_i\}`$ is a countable collection of disjoint elements of $``$ and the convergence is in the weak topology.
Given a space of states $``$, a space of possible results $`S`$ and a POVM $`\{\tau (I)\}_{I(S)}`$, for any pure state of the system determined by a normalized vector $`\varphi `$, the probability that the outcome of the measurement of the observable described by $`\{\tau (I)\}_{I(S)}`$ belongs to the Borel set $`I`$ is:
$$P(\varphi ,I)=\varphi ,\tau (I)\varphi .$$
(1)
More generally, for any mixed state of the system determined by a normalized positive trace class operator $`\rho `$ on $``$, the probability above is given by:
$$P(\rho ,I)=\mathrm{Tr}[\rho \tau (I)].$$
(2)
Note: If $`\tau (I_1I_2)=\tau (I_1)\tau (I_2)`$ for all $`I_1,I_2`$ then $`\tau `$ is a spectral measure.
The physical requirement that two observers, related by a transformation of the simmetry group $`𝒢`$ and performing the same experiment, get the same statistical distribution of the outcomes of the measurement, leads to a natural covariance condition and, eventually, to the following definition:
###### Definition 2
Let $``$ be a space of states, $`S`$ be a space of possible results and $`𝒢`$ a symmetry group of the theory. Suppose that $`𝒢`$ acts on $`S`$ by means of a representation $`\mathrm{\Lambda }:g\mathrm{\Lambda }(g),g𝒢`$, where $`\mathrm{\Lambda }(g):SS`$ are Borel mappings. Suppose $`𝒢`$ acts on $``$ by means of a strongly continuous unitary representation $`U:gU(g)`$. Let $`\{\tau (I)\}_{I(S)}`$ a POVM on $`S`$, with the property:
$$U(g)\tau (I)U(g)^{}=\tau (\mathrm{\Lambda }(g)I)$$
(3)
for any $`I(S),g𝒢`$. In this case the structure $`(,S,𝒢,\mathrm{\Lambda },U,\tau )`$ so defined is called system of covariance. If, furthermore, $`\tau `$ is a spectral measure, it is called system of imprimitivity .
While spectral measures represent a “property” of the system on which the measurement is performed, generic POVMs can’t describe “definite observables”, but we have to prefer them because they are able to describe simultaneous measurements of incompatible observables and measurements of time.
While there is a unified treatment of imprimitivity systems, mainly due to G. W. Mackey, we cannot say the same for covariance systems. Anyway we can overcame this difficulty by means of the following theorem, which may be recognized as a covariant version of Naimark’s dilatation theorem .
###### Proposition 1
Let $`(,S,𝒢,\mathrm{\Lambda },U,\tau )`$ a system of covariance. Then there is an imprimitivity system $`(^{},S,𝒢,\mathrm{\Lambda },V,E)`$, where $`^{}`$ is an Hilbert space, $`V`$ a strongly continuous unitary representation of the symmetry group $`𝒢`$ acting on $`^{}`$, $`E`$ is a spectral measure on the Borel $`\sigma `$-algebra $``$ of $`S`$, and there is an “intertwining operator” $`A:^{}`$, with the property $`AU(g)=V(g)A,`$ for any $`g𝒢`$, so that the following relation connects the spectral measure $`E`$ to the POVM $`\tau `$:
$$\tau (I)=A^+E(I)A.$$
(4)
Moreover, $`\tau `$ is normalized (i.e. $`\tau (S)=1`$) if and only if $`A^+A=1`$, namely if $`A`$ is isometric.
Finally Mackey’s imprimitivity theorem allows one to find the most general form of a system of imprimitivity.
###### Proposition 2
Let $`(^{},S,𝒢,\mathrm{\Lambda },V,E)`$ be a transitive system of imprimitivity. Let $`qS`$ be a generical element of $`S`$, for any $`xS`$ let $`g_x𝒢`$ be an element of $`𝒢`$ with the property $`x=\mathrm{\Lambda }(g_x)q`$. Let $`H_q`$ be “the little group”, namely the closed subgroup of $`𝒢`$ defined by
$$gH_q\mathrm{\Lambda }(g)q=q.$$
(5)
Then one can represent $`^{}`$ as direct integral of Hilbert spaces on $`S`$
$$^{}=_S^{}^{}(x)𝑑\mu (x),$$
(6)
where $`d\mu (x)`$ is a measure on $`(S)`$ having the same null sets as the spectral measure $`E`$. The vectors $`\varphi ^{}`$ can be represented by “wave functions” $`\psi (x)`$ and the projectors $`E(I)`$ as diagonal operators:
$$(E(I)\psi )(x)=f_I(x)\psi (x),$$
(7)
where $`f_I(x)`$ is the characteristic function of the Borel set $`I`$. Moreover the unitary representation $`V`$ of $`𝒢`$ takes the form of an “induced representation”:
$$[V(g)\psi ](x)=\left[\frac{d\mu (x^{})}{d\mu (x)}\right]^{\frac{1}{2}}R(g_x^1gg_x^{})\psi (x^{}),x^{}=\mathrm{\Lambda }(g^1)x$$
(8)
where $`g_x^1gg_x^{}H_q`$ and $`R`$ is an unitary representation of $`H_q`$.
The introduction of a system of imprimitivity is very advantageous: in this way the probability that the result of the measurement, performed on the state $`\varphi `$, belongs to the Borel set $`IS`$ takes the following simple form:
$$P(\varphi ,I)=\varphi ,\tau (I)\varphi =A\varphi ,E(I)A\varphi =_Sf_I(x)\psi (x)^2𝑑\mu (x),\psi =A\varphi .$$
(9)
In other words the usual concept of probability density, which can be found in the traditional formulation of quantum mechanics, can be re-established even if spectral measures are replaced by generic POVMs.
The last step of our construction procedure is the description of the most general intertwinig operator joining the imprimitivity system found by means of Mackey’s theorem to the unknown covariance system. The following theorem allows one to know when such an operator exists and what its general form is. It is based on a generalization of an argument given in ref and on Schur’s lemma.
###### Proposition 3
Let $`𝒢`$ be a locally compact topological group with a countable base of open sets and of type I. Let $`\widehat{𝒢}`$ be its dual space, namely the space of equivalence classes of its irriducible representations. Let $`U`$ and $`V`$ be two unitary representations of its, defined by their central decompositions:
$$U(g)=_{\widehat{𝒢}}^{}(U_\lambda (g)1_\lambda )𝑑\mu (\lambda ),V(g)=_{\widehat{𝒢}}^{}(U_\lambda (g)1_\lambda ^{})𝑑\mu ^{}(\lambda ),$$
(10)
acting respectively on Hilbert spaces
$$=_{\widehat{𝒢}}^{}_\lambda 𝒦_\lambda 𝑑\mu (\lambda ),^{}=_{\widehat{𝒢}}^{}_\lambda 𝒦_{}^{}{}_{\lambda }{}^{}𝑑\mu ^{}(\lambda ),$$
(11)
where $`U_\lambda `$ are irriducible representation and $`1_\lambda `$ and $`1_\lambda ^{}`$ are the unity operators acting on the Hilbert spaces $`𝒦_\lambda `$ or $`𝒦_{}^{}{}_{\lambda }{}^{}`$.
An isometric intertwinig operator $`A:^{}`$
$$AU(g)=V(g)A,g𝒢,A^+A=1$$
(12)
exists if and only if $`\mu `$ is absolutely continuous with respect to $`\mu ^{}`$ and
$$\mathrm{dim}(𝒦_{}^{}{}_{\lambda }{}^{})\mathrm{dim}(𝒦_\lambda )$$
(13)
almost everywhere whith respect to $`\mu `$. In this case it will assume the following form:
$$[A\varphi ]_\lambda =\left(\frac{d\mu }{d\mu ^{}}\right)^{\frac{1}{2}}(1_\lambda A_\lambda )\varphi _\lambda ,A_\lambda ^+A_\lambda =1,$$
(14)
where $`1_\lambda `$ is the unity operator in $`_\lambda `$ and $`A_\lambda :𝒦_\lambda 𝒦_{}^{}{}_{\lambda }{}^{}`$ is an isometry defined almost everywhere with respect to $`\mu `$.
The theorem reduces all our efforts, once we have the imprimitivity system, to the decomposition of V into I.U.R.s of the simmetry group $`𝒢`$.
The condition for the applicability of the theorem are not too restrictive, since most of the groups of physical interest have the required properties, namely they are locally compact with a countable base of open sets and of type I. However we shall see that the absolute continuity of the measure $`\mu `$ on $`\widehat{𝒢}`$ with respect to $`\mu ^{}`$ leads to interesting physical consequences, namely to a series of constraints on the realizability of some measurements on particular physical systems.
## 3 Localization of events in space-time
The first step necessary for the description of a realistic quantum reference frame is the definition of its origin. From an operational point of view this is the description of the way in which a microscopical object can localize a point of the Minkowsky space-time manifold, namely an instant indicating the beginning of the time scale and a point in space with respect to which position measurements are referred. In other words, how a quantum system can point at a particular event, in a relativistic covariant way.
This kind of measurement can be described by a POVM on the Minkowsky space-time $``$ covariant with respect to the universal covering of the proper orthochronous Poincaré group, which will be indicated by $`𝒫`$. The problem has already been studied in with a different method. We are going to rederive those results by means of the above developed construction procedure.
The first step is the construction of the most general imprimitivity system on $``$ covariant with respect to $`𝒫`$, whose action $`\stackrel{~}{\mathrm{\Lambda }}`$ on $``$ is given by
$$\stackrel{~}{\mathrm{\Lambda }}(y,a)(x)=y+\mathrm{\Lambda }(a)x,(y,a)𝒫,y𝒯_4,aSL(2C),$$
(15)
where $`\mathrm{\Lambda }:a\mathrm{\Lambda }(a)`$ is the representation of $`SL(2C)`$ acting on $``$ by means of the Lorentz matrices.
The system of imprimitivity is transitive. If we choose as representative point of the only orbit in $``$ under the action of $`𝒫`$ the origin $`O=(0,0,0,0)`$, we can recognize the little group in the Lorentz group, or more precisely in $`SL(2C)`$, its universal covering. According to the imprimitivity theorem the unitary representation $`V`$ has the form of an induced representation
$$[V(y,a)\psi ](x)=D(a)\psi (x^{}),x,x^{},$$
(16)
where $`\psi `$ takes its values in a Hilbert space $`\stackrel{~}{}`$, $`D(a)`$ is a unitary representation (not necessarily irreducible) of $`SL(2C)`$ and
$$x^{}=\mathrm{\Lambda }(a^1)(xy).$$
(17)
The projection-valued measure $`E`$ on the homogeneous space $``$ allows one to represent the vectors $`\psi `$ belonging to the Hilbert space $`^{}`$ as square-integrable vector-value function defined on $``$. The Lebesgue measure $`d^4x`$ on $``$, canonically associated to Minkowsky coordinates, is invariant under the action of $`𝒫`$ and the norm of $`\psi `$ assumes the simple form:
$$\psi ^2=_{}\psi (x)^2d^4x,$$
(18)
while the spectral measure $`E`$ on the Borel $`\sigma `$algebra B of $``$ assumes the diagonal form:
$$[E(I)\psi ](x)=f_I(x)\psi (x),I.$$
(19)
The second step is the decomposition of $`V`$ into I.U.R.s of the Poincaré group $`𝒫`$. We perform a Fourier transform on $``$ and pass from the coordinate representation to the momentum representation:
$$\stackrel{~}{\psi }(k)=(2\pi )^2_{}\mathrm{exp}(ikx)\psi (x)d^4x,kx=x^\alpha k_\alpha ,$$
(20)
$$\psi ^2=\stackrel{~}{\psi }(k)^2d^4k.$$
(21)
$`V`$ takes the following form:
$$[V(y,a)\stackrel{~}{\psi }](k)=exp(iky)D(a)\stackrel{~}{\psi }(k^{}),k^{}=\mathrm{\Lambda }(a^1)k.$$
(22)
The physical states of $``$ contain only non-negative energy representations, it follows that if $`A`$ is an intertwinig operator between $`U`$ and $`V`$, then $`A^{\prime \prime }^{}`$, where $`^{\prime \prime }`$ is the invariant subspace of $`^{}`$ which contains the vectors with non-negative energy, namely the wave functions $`\stackrel{~}{\psi }(k)`$ with support in the future cone $`V_+`$. We may disregard the values taken by $`\stackrel{~}{\psi }(k)`$ on the boundary of the cone, which has vanishing Lebesgue measure. In what follows we consider the subrepresentation $`V^{\prime \prime }`$ of $`V`$ acting on $`^{\prime \prime }`$.
Now we introduce for any $`k`$ in the open future cone an element $`a_kSL(2C)`$ defined by
$$k=\mathrm{\Lambda }(a_k)(M,0,0,0),k_{}^{0}{}_{}{}^{2}𝐤^2=M^2,$$
(23)
and the new wave function $`\psi ^{}`$, defined by:
$$\stackrel{~}{\psi }(k)=D(a_k)\psi ^{}(k).$$
(24)
The representation $`V^{\prime \prime }`$ takes the following form:
$$[V^{\prime \prime }(y,a)\psi ^{}](k)=\mathrm{exp}(iky)D(a_k^1aa_k^{})\psi ^{}(k^{}),$$
(25)
where $`a_k^1aa_k^{}=uSU(2)`$. We can now consider the decomposition of $`D`$ into I.U.R.s of $`SL(2C)`$, whose matrix elements we indicate with $`D_{jmj^{}m^{}}^{\rho n}(a)`$. They are identified by two parameters: $`\chi =(\rho ,n)`$. Two different I.U.R.s identified by $`(\rho ,n)`$ and $`=(\rho ^{},n^{})`$ are equivalent if and only if either $`(\rho ,n)=(\rho ^{},n^{})`$, or $`(\rho ,n)=(\rho ^{},n^{})`$. There are two series of I.U.R.s: the principal series with $`\rho `$ real and $`n`$ integer, and the supplementary series with $`\rho `$ imaginary and $`n=0`$ . Moreover one should not forget the trivial one-dimensional representation. The restriction of these representations to the subgroup $`SU(2)`$ is given by
$$D_{jmj^{}m^{}}^\chi (u)=\delta _{jj^{}}R_{mm^{}}^j(u),$$
(26)
where $`R_{mm^{}}^j(u)`$ stands for the matrix elements of the I.U.R. of $`SU(2)`$, labelled by the integer or half-integer index $`j`$, with
$$j=|\frac{n}{2}|,|\frac{n}{2}|+1,\mathrm{}m=j,j+1,\mathrm{},j1,j$$
(27)
Every unitary representation of $`SL(2C)`$ can be decomposed uniquely into primary representations, which are direct sums of I.U.R.s, as $`SL(2C)`$ is a type I group. We consider the direct integral decomposition of the Hilbert space $`\stackrel{~}{}`$ into irreducible spaces labelled by the variable $`\chi =(\rho ,n)`$, and introduce an index $`\alpha `$, which distinguishes the spaces where equivalent I.U.R.s operate:
$$\stackrel{~}{}=_{\widehat{SL(2C)}}^{}\underset{\alpha }{}\stackrel{~}{}_\alpha ^\chi d\omega (\chi ),$$
(28)
$$\psi ^2=_{\widehat{SL(2C)}\times V_+}\underset{\alpha }{}\psi _\alpha (k,\chi )^2d\omega (\chi )d^4k,$$
(29)
where $`\omega `$ is a generic measure on $`\widehat{SL(2C)}`$.
For fixed values of $`\alpha ,M,\chi `$ the Poincaré group $`𝒫`$ acts in the way described by Wigner :
$$[V^{\prime \prime }(y,a)\psi ^{}]_{\alpha ,jm}(k,\chi )=\mathrm{exp}(iky)\underset{m^{}}{}R_{mm^{}}^j(a_k^1aa_k^{})\psi _{\alpha ,jm^{}}^{}(k^{},\chi ),$$
(30)
as
$$a_k^1aa_k^{}=uSU(2)$$
(31)
Every I.U.R. of $`𝒫`$ with positive mass, identified by the variables (M,j), appears in the direct integral decomposition of $`V^{\prime \prime }`$ with a given multiplicity (defined almost everywhere on the positive real axis, i.e. on the $`M`$-axis). The multiplicity of a particular representation (M,j) is strictly positive if the subset of $`\widehat{SL(2C)}`$, whose elements are the I.U.R.s $`\chi =(\rho ,n)`$ of $`SL(2C)`$ with $`n2j`$, has non-vanishing measure $`\omega `$. Then one can always assume that the multiplicity is as large as one needs, allowing the index $`\alpha `$ to take a sufficient number of different values.
As we have seen in the previous section, the intertwinig operator theorem can be applied once $`U`$, the unitary representation of $`𝒫`$ acting on the Hilbert space $``$, is decomposed into direct integral of spaces where I.U.R.s of $`𝒫`$ operate:
$$[U(y,a)\varphi ]_{\alpha jm}(k)=\mathrm{exp}(iky)\underset{m^{}}{}R_{mm^{}}^j(u)\varphi _{\alpha jm^{}}(k^{}),$$
(32)
with
$$\varphi ^2=\underset{\alpha ,j}{}\varphi _\alpha (k,j)^2d\mu (k).$$
(33)
The discrete index $`\alpha `$ distinguishes the spaces where equivalent I.U.R.s operate. Note that the range of the sum on the indices $`\alpha `$ and $`j`$ may depend on $`M`$. The measure $`d\mu (k)`$ gives some informations about the mass spectrum of the system on which the measurement is performed. According to our third theorem, an isometric intertwinig operator $`A`$ between $`U`$ and $`V^{\prime \prime }`$ exists only if $`\mu (k)`$ (and therefore the corresponding measure on the range of $`M`$) is absolutely continuous with respect to the Lebesgue measure $`d^4k`$. This is possible if and only if the physical system on which the measurement is performed has a continuous mass spectrum, so we have to disregard the vacuum state and the one-particle states, whose mass-spectrum has a vanishing Lebesgue measure. Moreover, if a value $`j`$ appears in the decomposition, the measure $`\omega `$ of the set $`I(\widehat{SL(2C)})`$ with $`I=(\{\chi =(\rho ,n)\widehat{SL(2C)},n2j\})`$ has to be strictly positive. Eventually, if these condition are satisfied the most general intertwinig operator takes the following form:
$$\psi _{\alpha jm}^{}(k,\chi )=\underset{\alpha ^{}}{}A_{\alpha \alpha ^{}}^j(M,\chi )\varphi _{\alpha ^{}jm}(k),$$
(34)
assuming that $`\frac{d\mu (k)}{d^4k}=1`$ when $`M`$ belongs to the mass spectrum, with
$$\underset{\alpha }{}\overline{A_{\alpha \alpha ^{}}^j(M,\chi )}A_{\alpha \alpha ^{\prime \prime }}^j(M,\chi )d\omega (\chi )=\delta _{\alpha ^{}\alpha ^{\prime \prime }}.$$
(35)
Finally the most general density of probability on the Minkowsky space-time, describing the measurement of the coordinates of an event individuated by a quantum state described by a vector $`\varphi `$ takes the following form:
$$\rho (x)=\underset{\alpha }{}_{\widehat{SL(2C)}}\psi _\alpha (x,\chi )^2𝑑\omega (\chi ),$$
(36)
where
$$\psi _{\alpha pl}(x,\chi )=(2\pi )^2\mathrm{exp}(ikx)\underset{\alpha ^{},jm}{}D_{pljm}^\chi (a_k)A_{\alpha \alpha ^{}}^j(M,\chi )\varphi _{\alpha ^{}jm}(k)d^4k.$$
(37)
This is the main result of ref .
## 4 Quantum frames of reference
From an operational point of view a reference frame $`F`$ can be defined by the operations which allow one to connect it to an initially fixed frame $`F_0`$. In a relativistic theory the set of the allowed transformation is represented by the Poincaré group $`𝒫`$. Every element of $`𝒫`$ individuates the translation in the Minkowsky space-time and the Lorentz transformation which make the origin and the axes of the two frames coincide. We can also recognize in the four-vector individuating the translation the coordinates of the new origin with respect to the old one, while in the columns of the Lorentz matrix one finds the components of the new four orthogonal axes, relative to the old orthogonal basis. The ten indipendent parameters individuating the Poincaré transformation can also be recognized as relative observables of the two frames, namely relative position, time, velocity and spatial orientation. Their description can be given by a POVM on $`𝒫`$, covariant with respect to $`𝒫`$ itself. In other words the Poincaré group is in this case both the symmetry group and the measure space, endowed with the invariant Haar measure $`\nu `$. The description is simplified by the assumption of the classical nature of the frame $`F_0`$, in other words it will be considered an abstract mathematical tetrad with well-defined position, velocity and orientation. In this case the Poincaré group acts just on the system $`F`$ by means of left translation and the covariance condition assumes the following form:
$$U(g)\tau (I)U(g)^1=\tau (gI),g𝒫,I𝒫.$$
(38)
As shown in the previous sections, the starting point is the application of Mackey’s imprimitivity theorem. In this case it is particulary simple: there is only one orbit and, if we choose for example as representative point the identity $`e𝒫`$, the little group is reduced to the identity and the induced representation is simply the left regular representation or the direct sum of several representations equivalent to the left regular one and distinguished by the index $`\alpha `$:
$$[V(g^{})\psi ]_\alpha (g)=\psi _\alpha (g^1g),$$
(39)
$$[V(y,a)\psi ]_\alpha (x,b)=\psi _\alpha (\mathrm{\Lambda }(a^1)(xy),a^1b).$$
(40)
It can be decomposed into i.u.r.s of $`𝒫`$ by means of the harmonic analysis on $`𝒫`$, defined by
$$\stackrel{~}{\psi }(\gamma )=\psi (g)D^\gamma (g)𝑑\nu (g),$$
(41)
where $`\gamma `$ stands for $`(M,j)`$. The inversion formula is given by:
$$\psi (g)=_{\widehat{𝒫}}\mathrm{Tr}[\stackrel{~}{\psi }(\gamma )D^\gamma (g^1)]𝑑\widehat{\nu }(\gamma ),$$
(42)
where $`d\widehat{\nu }(\gamma )`$ is the Plancherel measure on $`\widehat{𝒫}`$. On the new “wave function ”, defined on $`\widehat{𝒫}`$, the space of equivalence classes of i.u.r.s of $`𝒫`$, the group action assumes the following form:
$$[V(g)\stackrel{~}{\psi }]_\alpha (\gamma )=D^\gamma (g)\stackrel{~}{\psi }_\alpha (\gamma ).$$
(43)
This procedure can be repeated whenever the action of the symmetry group $`𝒢`$ on the measure space $`S`$ is free and transitive, namely if for all $`x,yS`$ exists one and only one $`g𝒢`$ so that $`y=\mathrm{\Lambda }(g)x`$. In this way, for example, we can construct time measurements covariant with respect to time translations and position measurements covariant with respect to space displacements.
In our case we have just to combine the usual Fourier transform on $`R^4`$ and the harmonis analysis on $`SL(2C)`$, which give:
$$\stackrel{~}{\psi }_{\alpha ,jmj^{}m^{}}(k,\rho ,n)=(2\pi )^2\mathrm{exp}(ikx)D_{jmj^{}m^{}}^{\rho n}(a)\psi _\alpha (x,a)𝑑\mu (a)d^4x,$$
(44)
and
$$\psi ^2=_{\widehat{𝒯}\times \widehat{SL(2C)}}\underset{\alpha }{}\mathrm{Tr}[\psi _\alpha ^+(k,\chi )\psi _\alpha (k,\chi )]d^4kd\widehat{\mu }(\chi ),$$
(45)
where $`d\widehat{\mu }(\chi )`$ is the Plancherel measure on $`\widehat{SL(2C)}`$, which is concentrated on the principal series only.
As in the preceding section we consider only $`^{\prime \prime }`$, the invariant subspace of $`^{}`$ which contains only the vectors with non-negative energy, and the subrepresentation $`V^{\prime \prime }`$ acting on $`^{\prime \prime }`$. For every $`k`$ belonging to the open future cone $`V_+`$ we introduce an element $`a_kSL(2C)`$ defined by equation (23) and the new wave function $`\psi ^{}`$, given by:
$$\stackrel{~}{\psi }_{\alpha ,plj^{}m^{}}(k,\chi )=\underset{qs}{}D_{plqs}^\chi (a_k)\psi _{\alpha ,qsj^{}m^{}}^{}(k,\chi ).$$
(46)
In this way the action of $`𝒫`$ on $`^{\prime \prime }`$ for fixed values of $`\alpha ,M,\chi ,j^{},m^{}`$ assumes the form introduced by Wigner :
$$[V^{\prime \prime }(y,a)\psi ^{}]_{\alpha ,jmj^{}m^{}}(k,\chi )=\mathrm{exp}(iky)\underset{m^{}}{}R_{mn}^j(a_k^1aa_k^{})\psi _{\alpha ,jnj^{}m^{}}^{}(k^{},\chi ),$$
(47)
because of equation (31) and (26) Once we have introduced the primary decomposition of $`U`$ shown in equation (32), the most general intertwinig operator can be found if and only if the measure $`d\mu (k)`$, defining the mass spectrum of the physical system on which the measurement is performed, is absolutely continuous with respect to $`d^4k`$. In this case, redefinig if necessary the normalization of the wave function $`\varphi `$ so that $`\frac{d\mu (k)}{d^4k}=1`$ when $`M`$ belongs to the mass spectrum, we can write:
$$\psi _{\alpha jmqs}^{}(k,\chi )=\underset{\alpha ^{}}{}A_{qs,\alpha \alpha ^{}}^j(M,\chi )\varphi _{\alpha ^{}jm}(k),$$
(48)
with
$$\underset{\alpha }{}\underset{q=|\frac{n}{2}|}{\overset{\mathrm{}}{}}\underset{s=q}{\overset{q}{}}\overline{A^j(M,\chi )_{\alpha \alpha ^{}qs}}A^j(M,\chi )_{\alpha \alpha ^{\prime \prime }qs}d\widehat{\mu }(\chi )=\delta _{\alpha ^{}\alpha ^{\prime \prime }}.$$
(49)
Finally the density of probability assumes the following form:
$$\rho (x,b)=\underset{\alpha }{}|\psi _\alpha (x,b)|^2,$$
(50)
where
$$\psi _\alpha (x,b)=(2\pi )^2\mathrm{exp}(ikx)\mathrm{Tr}[D^\chi (b^1a_k)\psi _\alpha ^{}(k,\chi )]d^4k𝑑\widehat{\mu }(\chi )$$
$$=(2\pi )^6d^4k\mathrm{exp}(ikx)_0^+\mathrm{}𝑑\rho \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}(n^2+\rho ^2)$$
$$\times \underset{j,q=|\frac{n}{2}|}{\overset{\mathrm{}}{}}\underset{m=j}{\overset{j}{}}\underset{s=q}{\overset{q}{}}D_{qsjm}^{(\rho ,n)}(b^1a_k)\psi _{\alpha jmqs}^{}(k,\chi ).$$
(51)
It is quite interesting to compare the density of probability on the Minkowsky space-time (36), which was found indipendently in the previous section, and the density of probability which can be found from (50) by integration on $`SL(2C)`$:
$$\rho (x)=_{SL(2C)}\rho (x,a)𝑑\mu (a),$$
(52)
which after some calculations assumes the following form:
$$\rho (x)=\underset{\alpha }{}\mathrm{Tr}[\psi _\alpha ^{}(x,\chi )^+\psi _\alpha ^{}(x,\chi )]𝑑\widehat{\mu }(\chi ),$$
(53)
with
$$\widehat{\psi }_\alpha (x,\chi )=(2\pi )^2\mathrm{exp}(ikx)\stackrel{~}{\psi }_\alpha (k,\chi )d^4k,$$
(54)
where $`d\widehat{\mu }(\chi )`$ is the Plancherel measure on the principal series of the i.u.r.s of $`SL(2C)`$.
As equations (36) and (37) show, the most general density of probability on the Minkowsky space-time admits a generic measure $`d\omega (\chi )`$ on the space of i.u.r.s of $`SL(2C)`$. From these considerations one can guess there are some constraints on the realizability of some measurement, whose properties can be found through an analysis of the physical meaning of the parametres $`\chi =(\rho ,n)`$ in this context. For example, as shown in ref , one can recognize as “baricentric” a measurement of events such that the measure $`\omega `$ on $`(\widehat{SL(2C)})`$ appearing in equation (36) is concentrated on the trivial representation $`D(a)=1`$. Our results show that such a requirement can’t be compatible with the measurement of the further parametres fully describing a reference frame. In other words the origin of the quantum reference frame can never be localized on the world line of the centre of mass of the microscopical system defining it. Moreover, as the same author suggested in a previous paper , for a complete description of a quantum reference frame kinematical variables are not sufficient, but internal degree of freedom must be involved. We can see for example that neither the invariant mass of the system nor its spatial distribution, namely the centre of mass position, can be arbitrarily fixed and disregarded, but have a foundamental role in the whole description.
## 5 The observables relative to a quantum reference frame
The quantum picture is complete if every classical element is disregarded and every kinematical variable of a quantum system $`F_j`$ is referred to a quantum reference frame $`F_i`$, namely a microscopical system with continuous mass spectrum. This can be simply obtained in two steps if the quantum systems $`F_i`$ and $`F_j`$ don’t mutually interact.
First of all let’s introduce as a preliminary tool a classical frame $`F_0`$, with respect to which the parametres of the Poincaré tranformation connecting it to the quantum frame $`F_i`$ and a cinematical variable of the system $`F_j`$ are referred. The first ones are described by a POVM $`\tau _i`$ on the universal covering of the Poincaré group $`𝒫`$, acting on the Hilbert space $`_i`$, while the second ones are described by a POVM $`\tau _j`$ on a measure space $`S`$ acting on the Hilbert space $`_j`$. If there is no interaction the POVMs $`\tau _i`$ and $`\tau _j`$ and the unitary representation of $`𝒫`$ can be extended to the whole Hilbert space $`=_j_i`$ by the relations:
$$U(g)=U_j(g)U_i(g),g𝒫,$$
(55)
$$\widehat{\tau }_i(I)=I\tau _i(I)\widehat{\tau }_j(J)=\tau _j(J)I,$$
(56)
for all Borel subsets $`I𝒫`$ and $`JS`$. One can easily see that $`\widehat{\tau }_i`$ and $`\widehat{\tau }_j`$ are endowed with the right covariance properties with respect to Poincaré transformations:
$$U(g)\widehat{\tau }_i(I)U(g^1)=\widehat{\tau }_i(gI)U(g)\widehat{\tau }_j(J)U(g^1)=\widehat{\tau }_j(\mathrm{\Lambda }(g)J)g𝒫.$$
(57)
Moreover the operators in their ranges are mutually commuting:
$$[\widehat{\tau }_i(I),\widehat{\tau }_j(J)]=0,I𝒫,JS.$$
(58)
If these condition are satisfied the convolution $`\tau _{ij}`$ of the two POVMs $`\widehat{\tau }_i`$ and $`\widehat{\tau }_j`$ can be defined by the relation:
$$\tau _{ij}(J)=f_J(\mathrm{\Lambda }(g^1)x)𝑑\widehat{\tau }_i(g)𝑑\widehat{\tau }_j(x),g𝒫,xS,JS.$$
(59)
It is suitable for the description of the relative observables of the system $`F_j`$ with respect to the quantum frame $`F_i`$. Indeed $`\tau _{ij}`$ is endowed with the properties of a POVM acting on the Hilbert space $``$, namely positivity, $`\sigma `$additivity and normalization. Moreover, as we expected, it is invariant under the action of the Poincaré group $`𝒫`$:
$$U(\stackrel{~}{g})\tau _{ij}(J)U(\stackrel{~}{g}^1)=f_J(\mathrm{\Lambda }(g^1\stackrel{~}{g})\mathrm{\Lambda }(\stackrel{~}{g}^1)x)𝑑\widehat{\tau }_i(g)𝑑\widehat{\tau }_j(x)=\tau _{ij}(J),$$
(60)
in fact a Poincaré transformation will act on both $`F_i`$ and $`F_j`$, changing the “absolute” cinematical variables of the two systems but leaving invariant the relative ones.
The mathematical description of the relations connecting two not-interacting quantum frames $`F_i`$ and $`F_j`$ can be obtained as a special case of this formalism. If in the previous discussion the quantum system $`F_j`$ has a continuous mass spectum, while the measure space $`S`$ coincides with the Poincaré group again, $`\tau _{ij}`$ describes the measurement of the ten parametres of the transformation connecting the two frames. It assumes the following form:
$$\tau _{ij}(I)=f_I(g^1g^{})𝑑\widehat{\tau }_i(g)𝑑\widehat{\tau }_j(g^{}),J𝒫.$$
(61)
The intuition can be helped by a calculation of the density of probability describing the statistics of the measurement, defined by
$$\varphi _i\varphi _j,\tau _{ij}(I)\varphi _i\varphi _j=_𝒫f_I(g)\rho _{ij}(g)𝑑\mu (g),$$
(62)
which assumes the following form:
$$\rho _{ij}(g)=_𝒫\rho _i(\varphi _i,g^{})\rho _j(\varphi _j,g^{}g)𝑑\mu (g^{}),$$
(63)
where $`\rho _i(\varphi _i,g^{})`$ and $`\rho _j(\varphi _j,g^{}g)`$ are the densities of probability describing a measurement of the “absolute” parameters $`g^{}`$ and $`g^{}g`$, namely relative to a classical reference frame $`F_0`$, which were calculated in the previous section. In other words if the Poincaré transformations identified by the elements $`g^{}`$ and $`g^{}g`$ connect the classical frame to the quantum frames $`F_i`$ and $`F_j`$ respectively, then the transformation from $`F_i`$ to $`F_j`$ will be individuated by the element $`g𝒫`$, whatever $`g^{}`$ may be.
The introduction of a third quantum frame $`F_k`$ in the description leads to some surprising consequences, which are commonly called “the paradox of quantum frames”. While the relative observables of $`F_i`$ and $`F_j`$, or of $`F_i`$ and $`F_k`$, can be described respectevely by $`\tau _{ij}=\stackrel{~}{\tau _i}\tau _j`$ or by $`\tau _{ik}=\stackrel{~}{\tau _i}\tau _k`$, the relative observables of $`F_j`$ and $`F_k`$ can’t be obtained by the convolution $`\stackrel{~}{\tau _{ik}}\tau _{ij}`$ as it doesn’t own the necessary properties. The operators in its range could be positive if and only if the POVMs $`\tau _{ik}`$ and $`\tau _{ij}`$ commute, but it can’t be required and it is not generally true. One can easily see that the commutativity of the POVM $`\{\tau _i(I)\}_{I(𝒫)}`$ is a sufficient condition for the commutativity of $`\tau _{ik}`$ and $`\tau _{ij}`$:
$$[\widehat{\tau }_i(I),\widehat{\tau }_i(I^{})]=0[\tau _{ik}(I),\tau _{ij}(I^{})]=0,I,I^{}𝒫,$$
(64)
however the first condition cannot be required. Note that the commutativity of the projectors in the range of the spectral measure $`\{E_i(I)\}_{I(𝒫)}`$ does not involve the commutativity of the POVM $`\tau _i(I)=A^+E_i(I)A`$, unless the range of the intertwinig operator $`^{\prime \prime }=A`$ is an invariant subspace under the action of the projectors $`E_i(I)`$. One can easily see that in this case the positive operators in the range of the POVM $`\tau _i`$ would be projectors too, but this is forbidden by Pauli’s theorem and by the non compatibility of the observables describing a reference frame. In other words sequential measurement of the relative parametres of two quantum frames $`F_j`$ and $`F_k`$ with respect to a third quantum frame $`F_i`$ are not compatible, even if they don’t mutually interact .
## 6 Aknowledgments
I am grateful to M. Toller for his precious suggestions. I also wish to thank V. Moretti.
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# SMALL 𝑋 BEHAVIOUR OF PARTON DISTRIBUTIONS
## Abstract
We investigate the $`Q^2`$ evolution of parton distributions at small $`x`$ values, recently obtained in the case of soft initial conditions. The results are in excellent agreement with deep inelastic scattering experimental data from HERA.
The measurements of the deep-inelastic scattering structure function $`F_2`$ in HERA have permitted the access to a very interesting kinematical range for testing the theoretical ideas on the behavior of quarks and gluons carrying a very low fraction of momentum of the proton, the so-called small $`x`$ region. In this limit one expects that non-perturbative effects may give essential contributions. However, the resonable agreement between HERA data and the NLO approximation of perturbative QCD that has been observed for $`Q^2>1`$GeV<sup>2</sup> (see the recent review in ) indicates that perturbative QCD could describe the evolution of structure functions up to very low $`Q^2`$ values, traditionally explained by soft processes.
Here we illustrate the results obtained recently in . These results are the extension to the NLO QCD approximation of previous LO studies . The main ingredients are:
1. Both, the gluon and quark singlet densities are presented in terms of two components ($`{}_{}{}^{}+_{}^{}`$ and $`{}_{}{}^{}_{}^{}`$) which are obtained from the analytical $`Q^2`$ dependent expressions of the corresponding ($`{}_{}{}^{}+_{}^{}`$ and $`{}_{}{}^{}_{}^{}`$) parton distributions moments.
2. The $`{}_{}{}^{}_{}^{}`$ component is constant at small $`x`$, whereas the $`{}_{}{}^{}+_{}^{}`$ component grows at $`Q^2Q_0^2`$ as
$$\mathrm{exp}\left(2\sqrt{\left[a_+\mathrm{ln}\left(\frac{a_s(Q_0^2)}{a_s(Q^2)}\right)\left(b_++a_+\frac{\beta _1}{\beta _0}\right)\left(a_s(Q_0^2)a_s(Q^2)\right)\right]\mathrm{ln}\left(\frac{1}{x}\right)}\right),$$
where the LO term $`a_+=12/\beta _0`$ and the NLO one $`b_+=412f/(27\beta _0)`$. Here the coupling constant $`a_s=\alpha _s/(4\pi )`$, $`\beta _0`$ and $`\beta _1`$ are the first two coefficients of QCD $`\beta `$-function and $`f`$ is the number of active flavors.
We have analyzed $`F_2`$ HERA data at small $`x`$ from the H1 coll.. The initial scale of the parton distributions was fixed into the fits to $`Q_0^2`$ = 1 $`GeV^2`$, although later it was released to study the sensitivity of the fit to the variation of this parameter. The analyzed data region was restricted to $`x<0.01`$ to remain within the kinematical range where our results are accurate.
Fig. 1 shows $`F_2`$ calculated from the fit with Q<sup>2</sup> $`>`$ 1 GeV<sup>2</sup>. Only the lower $`Q^2`$ bins are shown. One can observe that the NLO result (dot-dashed line) lies closer to the data than the LO curve (dashed line). The lack of agreement between data and lines observed at the lowest $`x`$ and $`Q^2`$ bins suggests that the initial flat behavior should occur at $`Q^2`$ lower than 1 GeV<sup>2</sup>. In order to study this point we have done the analysis considering $`Q_0^2`$ as a free parameter. Comparing the results of the fits (see ) one can notice the better agreement with the experiment of the NLO curve at fitted $`Q_0^2=0.55GeV^2`$ (solid curve) which is more apparent at the lowest kinematical bins.
A.K. was supported by Alexander von Humboldt fellowship and by DIS2000 Orgcommittee. G.P. was supported in part by Xunta de Galicia (XUGA-20602B98) and CICYT (AEN96-1673).
References
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# Nonlinear optical response and spin-charge separation in one-dimensional Mott insulators
## Abstract
We theoretically study the nonlinear optical response and photoexcited states of the Mott insulators. The nonlinear optical susceptibility $`\chi ^{(3)}`$ is calculated by using the exact diagonalization technique on small clusters. From the systematic study of the dependence of $`\chi ^{(3)}`$ on dimensionality, we find that the spin-charge separation plays a crucial role in enhancing $`\chi ^{(3)}`$ in the one-dimensional (1D) Mott insulators. Based on this result, we propose a holon-doublon model, which describes the nonlinear response in the 1D Mott insulators. These findings show that the spin-charge separation will become a key concept of optoelectronic devices.
The charge gap in Mott insulators is a consequence of strong electron correlation. This is completely different from the band insulators, where the charge gap is basically originated from band effects. The nature of charge excitation across the gap is thus essentially different between the two types of insulators. The optical response is used to investigate the charge excitation across the gap. In the response, the linear susceptibility with respect to the applied electric filed, $`\chi ^{(1)}`$, provides information on the dipole-allowed states with odd parity. In addition to $`\chi ^{(1)}`$, the nonlinear susceptibilities detect the odd states together with the dipole-forbidden states with even parity.
Very recently, anomalously enhanced third-order nonlinear optical susceptibility $`\chi ^{(3)}`$ has been reported for one-dimensional (1D) Mott insulators of Cu oxides and Ni halides, as compared with those for the band insulators. In addition, the 1D Cu oxide, Sr<sub>2</sub>CuO<sub>3</sub>, exhibits ultrafast nonlinear optical response ($``$1 ps) at room temperature. These facts strongly suggest a great potential of the 1D Mott insulators as optoelectronic materials with high performance. Furthermore, it is shown that $`\chi ^{(3)}`$ in 2D Mott insulators is smaller than that in the 1D Mott insulators.
Motivated by these striking experiments, we theoretically examine photoexcited states and nonlinear optical response in the Mott insulators, and clarify underlying physics of optical nonlinearity of the 1D Mott insulators. We use the half-filled Hubbard model to describe the Mott insulators. The $`\chi ^{(3)}`$ is calculated by using the numerically exact diagonalization method on finite-size clusters, and compared with that for higher dimensional systems with ladder and 2D geometry. We find that $`\chi ^{(3)}`$ increases with decreasing the dimensionality. It is shown that in the 1D system dipole-allowed (odd) and -forbidden (even) states are almost degenerate in energy, having very large dipole coupling between the two states, while in the 2D system the dipole coupling is rather small in spite of the closeness of the odd and even states. We demonstrate that this peculiar feature in the 1D system is due to the spin-charge separation inherent in the 1D Mott insulators. We propose an effective model that can describe the optical nonlinearity, a holon-doublon model, where holon and doublon represent the charge degree of freedom for photoinduced unoccupied and doubly occupied sites, respectively. The model reproduces very well the characteristic behaviors of the experimental $`\chi ^{(3)}`$. We find that the spin-charge separation becomes a key concept of future optoelectronic devices.
The electric field E of the incident light induces the dielectric polarization $`𝐏_{\mathrm{ind}}`$ in a material, which is described by a power series of nonlinear optical susceptibility $`\chi ^{(n)}`$: $`𝐏_{\mathrm{ind}}=ϵ_0\left(\chi ^{(1)}𝐄+\chi ^{(2)}𝐄^2+\chi ^{(3)}𝐄^3+\mathrm{}\right)`$. The linear susceptibility $`\chi ^{(1)}`$ is given by
$`\chi _{jk}^{(1)}(\omega ;\omega )`$ (1)
$`={\displaystyle \frac{1}{ϵ_0L}}{\displaystyle \frac{e^2}{\mathrm{}}}{\displaystyle \underset{a}{}}\left({\displaystyle \frac{r_{0a}^jr_{a0}^k}{\mathrm{\Omega }_ai\mathrm{\Gamma }_a\omega }}+{\displaystyle \frac{r_{0a}^kr_{a0}^j}{\mathrm{\Omega }_a+i\mathrm{\Gamma }_a+\omega }}\right),`$ (2)
where $`L`$ is the number of sites, $`ϵ_0`$ is the dielectric constant, $`j`$ and $`k`$ are the polarization directions, $`er_{0a}`$ is the dipole moment between the ground state $`|0`$ and excited state $`|a`$ with odd parity, $`\mathrm{\Omega }_a`$ is the energy difference between $`|0`$ and $`|a`$, and $`\mathrm{\Gamma }_a`$ is the damping factor. Due to symmetry restrictions, $`\chi ^{(2)}`$ vanishes in centrosymmetric materials to which Cu oxides and Ni halides belong. The lowest observable nonlinearity is, therefore, $`\chi ^{(3)}`$, which is expressed as
$`\chi _{jklm}^{(3)}(\omega _\sigma ;\omega _1,\omega _2,\omega _3)={\displaystyle \frac{1}{ϵ_0L}}{\displaystyle \frac{e^4}{3!\mathrm{}^3}}𝒫{\displaystyle \underset{a,b,c}{}}`$ (3)
$`{\displaystyle \frac{r_{0a}^jr_{ab}^kr_{bc}^lr_{c0}^m}{(\mathrm{\Omega }_ai\mathrm{\Gamma }_a\omega _\sigma )(\mathrm{\Omega }_bi\mathrm{\Gamma }_b\omega _2\omega _3)(\mathrm{\Omega }_ci\mathrm{\Gamma }_c\omega _3)}},`$ (4)
where $`\omega _\sigma =\omega _1+\omega _2+\omega _3`$, $`b`$ and $`c`$ denote even and odd states, respectively, and $`𝒫`$ represents the sum of permutation on ($`j`$, $`\omega _1`$), ($`k`$, $`\omega _2`$), ($`l`$, $`\omega _3`$), and ($`m`$, $`\omega _\sigma `$). In the present study, by setting all of the polarization directions the same along chains, we examine $`\chi ^{(3)}(\omega ;\omega ,\omega ,\omega )`$ and $`\chi ^{(3)}(\omega ;0,0,\omega )`$, the imaginary parts of which give two-photon absorption (TPA) and electroabsorption spectra, respectively. Hereafter, we set $`e`$, $`\mathrm{}`$, and $`ϵ_0`$ to be unity.
The insulating Cu oxides and Ni halides are known to be charge-transfer (CT)-type Mott insulators, where both 3$`d`$ and 2$`p`$ orbitals exist in the transition metal and ligand ions, respectively, and participate in the electronic states. The values of the gap are predominantly determined by the energy position of the $`p`$ orbitals. However, it is well established that the electronic states of the CT-type insulators can be described by the Hubbard model with single band by mapping a bound state called the Zhang-Rice singlet state onto the lower Hubbard band. Actually, the single-band Hubbard model explains very well the spectral line shape of angle-resolved photoemission and electron-energy loss spectroscopies in the insulating Cu-oxides.
The single-band Hubbard Hamiltonian is defined as
$`H`$ $`=`$ $`t{\displaystyle \underset{i,j,\sigma }{}}(c_{i,\sigma }^{}c_{j,\sigma }+\mathrm{H}.\mathrm{c}.)+U{\displaystyle \underset{i}{}}n_{i,}n_{i,}`$ (5)
$`+`$ $`V{\displaystyle \underset{i,j}{}}n_in_j,`$ (6)
where $`c_{i,\sigma }^{}`$ is the creation operator of an electron with spin $`\sigma `$ at site $`i`$, $`n_i`$=$`n_{i,}`$+$`n_{i,}`$, $`i,j`$ runs over pairs on nearest neighbor sites, $`t`$ is the hopping integral, $`U`$ is the on-site Coulomb interaction, and $`V`$ is the Coulomb interaction between nearest-neighbor sites.
The realistic values of the parameters for cuprates are chosen: $`U/t`$=10, and $`V/t`$=1.5 for a 1D system and $`V/t`$=1 for ladder and 2D systems. Additionally, in the 2D system, a second-nearest-neighbor hopping $`t^{}`$ is considered to simulate a parent compound of high-$`T_c`$ cuprates. The damping factors are assumed to be the same for all excited states of the systems with $`\mathrm{\Gamma }`$=$`0.4t`$.
We numerically examine linear absorption spectra, that are defined by an imaginary part of $`\chi ^{(1)}`$, and TPA spectra for a 12-site chain with open boundary condition. The Lanczos technique is used for the calculation of the ground state. The spectra are calculated by using the correction vector technique. We also calculate them for two-leg ladder, and 2D systems with 12 sites in order to clarify the dependence of $`\chi ^{(3)}`$ on dimensionality. The clusters used in the calculation are shown in Fig. 1(c).
The calculated results are shown in Figs. 1(a) and 1(b). The solid, dashed, and dotted lines denote the results for the 1D, ladder and 2D systems, respectively. The upper panel (a) represents the linear absorption spectra, which detect odd states. The optical gap is estimated to $``$2.4 (=6$`t`$) eV assuming $`t`$$``$0.4 eV, comparable with the experimental optical gap. The linear absorption has a large magnitude for the 1D system with strong enhancement in the low-energy region ($`\omega 7t`$), and decreases with increasing dimensionality from the 1D to 2D systems. These results are consistent with the experimental data of optical conductivity for 1D, ladder, and 2D cuprates.
The lower panel (b) shows the TPA spectra, which detect even states with the same parity as the ground state. We find that the spectral shapes of TPA bear a close resemblance to those of $`\chi ^{(1)}`$ when the energy in TPA is doubled. This means that the energy positions of the even states lie in the same energy region as the odd states shown in Fig. 1(a) \[see the denominator of Eq. (2)\]. These characteristics are independent of the dimensionality, and agree with the results obtained by pump-probe measurements. We also find that the magnitude of $`\chi ^{(3)}`$ is the largest in the 1D system, and decreases with increasing the dimensionality. This indicates that one dimensionality is favorable to the enhancement of $`\chi ^{(3)}`$, which is consistent with the recent experimental data on Sr<sub>2</sub>CuO<sub>3</sub> and Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>.
In order to clarify the TPA spectra in detail, we next examine photoexcited states and dipole moments between these states. The dipole moments are given by $`_i`$odd$`|x_in_i|`$even$``$, where $`|\mathrm{odd}`$ and $`|\mathrm{even}`$ denote eigenstates with odd and even parities for the Hamiltonian of Eq. (3), respectively, $`x_i`$ is the $`x`$-coordinate of the position at the $`i`$ site, and $`n_i`$ is the number operator of the electron.
Figure 2 shows the results of the energy distribution of dipole moments for the 1D (left panel) and 2D (right panel) systems, which are calculated exactly by using clusters shown in the insets. The upper panel \[(a) and (d)\] represents the dipole moments between the ground state and odd states. The middle panel \[(b) and (e)\] shows dipole moments between the first odd state labeled as $`O_1`$ in (a) and (d) and even states. In the 1D case, the $`O_1`$ state strongly couples to an even state with nearly the same energy as the $`O_1`$ state. The magnitude of the dipole moment is about 3. On the other hand, in the 2D case the dipole moment is smaller than that in the 1D system and the magnitude is less than 0.5. In addition, other dipole moments are also very small in the 2D system. The lower panel \[(c) and (f)\] shows dipole moments between another odd state $`O_2`$ and even states. We find again that an even state with very close energy to the $`O_2`$ state yields large dipole moment in the 1D case, whereas the dipole moment in the 2D case is small although there are even states whose energy is very close to the $`O_2`$ state. These results indicate a quantitative difference of the dipole moments between the 1D and 2D systems. The difference plays an important role in the difference of the magnitude of TPA between the 1D and 2D systems.
In general, the dipole moment becomes large when two states which lead to the dipole moment are constructed by similar wave functions. Therefore it is meaningful to examine the nature of the odd and even states in the 1D and 2D systems. We define a quantity $`𝒮`$, which measures the similarity of the wave functions, as follows: $`𝒮_i|\alpha _{o,i}\alpha _{e,i}|`$, where $`\alpha _{o,i}`$ ($`\alpha _{e,i}`$) is the coefficient of a basis which constructs the odd (even) eigenstates of the Hamiltonian of Eq. (3), and $`i`$ is the index of the bases. The quantity $`𝒮`$ becomes unity when two eigenstates are the same except for the phase of each basis.
In the 1D system, $`𝒮`$ between the $`O_1`$ ($`O_2`$) state and an adjacent even state is 0.96 (0.81), which is close to unity. On the other hand, in the 2D system, $`𝒮`$ is 0.66 and 0.64 for the $`O_1`$ and the $`O_2`$ states, respectively. These results show that in the 1D system the two wave functions with close energies are quite similar to each other in contrast to the 2D ones.
Next, we consider the origin of the different behavior of the photoexcited states between the 1D and 2D systems. In the 1D system, a photoinduced carrier separates into two collective modes carrying the spin and charge degrees of freedom (spin-charge separation). In addition, the wave function itself is factorized as the products of spin and charge wave functions when $`U`$ is very large. Therefore, it is reasonable in the 1D system to take into account only the charge degree of freedom of electrons for the optical response. We propose below an effective model to describe the nonlinear optical response in the 1D system. In contrast to the 1D system, the motion of the carriers in the 2D system is known to be strongly influenced by spin background as long as the excitation energy is low, implying that the wave function cannot be simply factorized unlike the 1D system. It is, therefore, natural to suppose that the spin degree of freedom in the 2D system plays a crucial role in the broad spectral weight of $`\chi ^{(1)}`$ and $`\chi ^{(3)}`$ (Fig. 1) as well as the broad distribution of the photoexcited states (Fig. 2).
Based on the spin-charge separation picture, we introduce a holon-doublon model for the 1D Mott insulators. The Hamiltonian is given by
$`H`$ $`=`$ $`t{\displaystyle \underset{i}{}}(h_i^{}h_{i+1}d_i^{}d_{i+1}+\mathrm{H}.\mathrm{c}.)V{\displaystyle \underset{i,j}{}}h_i^{}h_id_j^{}d_j`$ (7)
$`+`$ $`{\displaystyle \frac{U}{2}}{\displaystyle \underset{i}{}}\left(h_i^{}h_i+d_i^{}d_i\right),`$ (8)
with the constraint of no double occupation at site $`i`$, $`(h_i^{}h_i+d_i^{}d_i)1`$. The holon $`h`$ (doublon $`d`$) represents the charge degree of freedom of unoccupied (doubly occupied) sites produced by the photoexcitation. We note that only two particles, i.e., one holon and one doublon, are contained in the present system. Although the model includes the attractive Coulomb interaction, we note that the model is different from a model for standard exciton of 1D semiconductors, since the holon and the doublon cannot occupy the same site due to the hard core constraint unlike the semiconductors. As will be shown below, the constraint plays an important role in the degeneracy of the even and odd photoexcited states.
It is easily shown in the holon-doublon model that the even and odd states are degenerate and thus the overlap integral $`𝒮`$ between them is unity: Let us introduce a set of the bases with odd and even parities written as $`|\mathrm{basis},o`$=($`|hd00\mathrm{}|00\mathrm{}dh`$)$`/\sqrt{2}`$ and $`|\mathrm{basis},e`$=($`|hd00\mathrm{}+|00\mathrm{}dh`$)$`/\sqrt{2}`$, where $`h`$, $`d`$ and 0 denote holon, doublon, and vacant sites, respectively. Since there is no exchange between holon and doublon due to the hard core constraint, the first terms of the bases never couple to the second terms. This means that the matrix elements of the Hamiltonian of Eq. (4) are independent of the sign of the second terms. Therefore, the Hamiltonian matrix with odd parity is equal to that with even parity. This gives rise to the degeneracy between odd and even eigenstates and thus $`𝒮`$=1. This model naturally explains the characteristics in the photoexcited states in the 1D Mott insulators discussed above. We note that the even and odd states are always degenerate regardless of the value of $`V`$ (Ref. 18) or the range of Coulomb interactions.
Figure 3 exhibits the linear and nonlinear absorption spectra in the holon-doublon model for an 80-site chain. The linear absorption (a) and TPA (b) spectra reproduce well those for the 1D Hubbard model shown in Fig. 1. The obtained results are also very similar to those of the experimental data in Sr<sub>2</sub>CuO<sub>3</sub> (Refs. 2 and 3) in the following points: (i) The spectral weights of linear absorption and TPA spectra are concentrated on the same energy region when the energy in TPA is doubled. (ii) The electroabsorption spectra (c) show an oscillating structure at the spectral edge of the linear absorption ($`6t`$). The agreement suggests that this model is a proper one in describing the nonlinear optical response in the 1D Mott insulators.
Finally we discuss the difference of nonlinear optical response between 1D band insulators and Mott insulators. The photoexcitation process in the band insulators creates an electron-hole pair bound by their attractive Coulomb interaction (Mott-Wannier exciton). The pair gives a series of bound states with odd and even parities below the gap. In fact, in the 1D band insulators such as silicon polymers and Pt halides, it has been observed that the odd and even states are concentrated on different energy regions, and their difference is rather large ($``$1 eV). However, such an energy difference between odd and even states is generally unfavorable for obtaining the large dipole moments because dipole moments are inversely proportional to their energy difference. In this respect, the 1D Mott insulators are advantageous to large dipole moments since they have the nearly degenerate odd and even states as shown in Figs. 1 and 3, which leads to the large $`\chi ^{(3)}`$.
In summary, we have clarified the nonlinear optical response and photoexcited states of Mott insulators. We found that $`\chi ^{(3)}`$ is enhanced with decreasing the dimensionality, which is consistent with the recent experiments. In the 1D system, the odd and even photoexcited states are nearly degenerate in energy, and constructed by similar wave functions. We found that these properties, which become important for the large nonlinearity, are originated by the spin-charge separation. We also proposed the effective model in 1D Mott insulators. These findings imply that the concept of the spin-charge separation is not of purely academic interest but will show up as underlying physics of optoelectronic devices in the near future.
The authors thank H. Okamoto, H. Kishida, M. Kuwata-Gonokami, and Y. Tokura for valuable discussions. This work was supported by a Grant-in-Aid for Scientific Research on Priority Areas from the Ministry of Education, Science, Sports and Culture of Japan, CREST and NEDO. The parts of the numerical calculation were performed in the supercomputing facilities of ISSP, University of Tokyo, and IMR, Tohoku University.
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# Effective action for QED4 through 𝜁-function regularization
## I Introduction
In QED, the effective one-loop Lagrangian describes the effective nonlinear interaction of the electromagnetic fields due to a single fermion loop. In two dimensions, its general form has been obtained both through proper time and $`\zeta `$ function regularizations . In four dimensions, on the other hand, only particular field configurations have been studied.
The $`3+1`$ dimensional problem of constant electromagnetic fields was first studied by Euler and Heisenberg and independently by Weisskopf . These authors obtained an integral expression for the one-loop effective Lagrangian in the framework of the electron-hole theory. Later on, Schwinger rederived this integral representation in a field-theoretical scenario, by making use of proper time techniques . In all these references, explicit results were derived in some limits, the most famous being the weak-field one. This and other particular field configurations were subsequently studied through the proper-time regularization by a number of authors (see, for example, ).
More recently, the interest in the subject was renewed, and the Euclidean effective action for constant electromagnetic background configurations was studied through $`\zeta `$ function techniques : In reference analytic expressions were found for the case of purely magnetic fields in any number of dimensions. In this same reference, the case of equal electric and magnetic fields in four Euclidean dimensions was also studied. A step towards more general field configurations was given in , where the authors obtained the effective Lagrangian as a power series in $`\frac{B}{E}`$.
It is the aim of this paper to obtain, through $`\zeta `$ function methods, an explicit non-perturbative expression for the full one-loop effective action of Quantum Electrodynamics in four dimensions in the case of constant, but otherwise arbitrary, electromagnetic fields. To this end, we will work in Euclidean space-time, and define the determinant of the relevant Dirac operator $`\overline{)}D`$ through the derivative of the $`\zeta `$ function of $`\overline{)}D^{}\overline{)}D`$.
The organization of the paper is as follows:
After summarizing some well-known generalities in section II, we devote section III to analytically extending the relevant $`\zeta `$ function to the region $`\mathrm{}s>2`$. (The main point here is the analytic extension of a Barnes $`\zeta `$-function). Its value at $`s=0`$ is also given in this section.
In section IV, a complete analytical expression for the effective action in terms of special functions is given, and the renormalization issue is discussed.
Section V contains a comparison between $`\zeta `$ and proper-time regularizations.
The Appendices A and B contain the derivation of some particular limits for the relevant zeta and for the effective action, thus allowing for the comparison with previous work on less general field configurations.
## II Generalities
We study the effective action for massive Dirac particles in the presence of uniform, but otherwise arbitrary, electromagnetic background fields. We work in four-dimensional Euclidean space. Then, the effective action in the one-loop approximation is given by
$$S_{eff}\left[A_\mu \right]=S_{cl}\left[A_\mu \right]\mathrm{log}Det\left(\overline{)}D\left[A_\mu \right]\right),$$
(1)
where $`S_{cl}\left[A_\mu \right]`$ is the classical Euclidean action and $`\overline{)}D\left[A_\mu \right]=\gamma _\mu \left(_\mu ieA_\mu \right)+im`$ is the Euclidean Dirac operator, $`m`$ being the fermion mass.
Note that, even though $`\overline{)}D`$ is not self adjoint, it is normal; so, the functional determinant appearing in the one-loop correction to the action can be defined through $`\zeta `$ function regularization , which leads to
$$S_{eff}\left[A_\mu \right]=S_{cl}\left[A_\mu \right]+S^{(1)}\left[A_\mu \right]=S_{cl}\left[A_\mu \right]+\frac{1}{2}\frac{}{s}\zeta (s;\overline{)}D^{}\overline{)}D)_{s=0}.$$
(2)
In order to evaluate the one-loop correction $`S^{(1)}`$ in the previous expression, it is necessary to obtain the spectrum of the operator $`\overline{)}D^{}\overline{)}D`$, which is well known in the case of uniform fields . In this particular situation, one can always choose a reference frame such that $`F_{03}=F_{30}=E`$ and $`F_{12}=F_{21}=B`$, while the remaining components of the field tensor vanish. When doing so, the required zeta function turns out to be
$`\zeta (s;\overline{)}D^{}\overline{)}D)=\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}[2{\displaystyle \underset{n_a=1}{\overset{\mathrm{}}{}}}(2n_aa+c)^s+2{\displaystyle \underset{n_b=1}{\overset{\mathrm{}}{}}}(2n_bb+c)^s+`$
$$4\underset{n_a=1}{\overset{\mathrm{}}{}}\underset{n_b=1}{\overset{\mathrm{}}{}}(2n_aa+2n_bb+c)^s+c^s].$$
(3)
Here, $`\mathrm{\Omega }`$ is the volume of the four-dimensional Euclidean space, $`a=\frac{e|E|}{\mu ^2}`$, $`b=\frac{e|B|}{\mu ^2}`$, $`c=\frac{m^2}{\mu ^2}`$, and $`\mu `$ is a parameter with mass dimension, introduced to render the $`\zeta `$ function dimensionless. Note that the series in equation (3) are all convergent for $`\mathrm{}s>2`$, where they define an analytic function of $`s`$.
## III Analytic extension of the $`\zeta `$ function
In this section, we will perform the analytic extension of the relevant $`\zeta `$ function to a region containing $`s=0`$. In particular, we will show it to be finite at $`s=0`$ and give its value at this point.
The first two terms in equation (3) can be rewritten in terms of Hurwitz’ zeta functions, which are well known to be meromorphic functions with a unique simple pole at $`s=1`$. On the other hand, the third term is a zeta function of the Barnes’ type (see also and references therein). In order to analytically extend this term, we write it in integral form. After doing so, we get
$`\zeta (s;\overline{)}D^{}\overline{)}D)=\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}\{{\displaystyle \frac{2}{\left(2a\right)^s}}\zeta (s,{\displaystyle \frac{c}{2a}}+1)+{\displaystyle \frac{2}{\left(2b\right)^s}}\zeta (s,{\displaystyle \frac{c}{2b}}+1)+`$
$$\frac{1}{\mathrm{\Gamma }(s)}_0^{\mathrm{}}dtt^{s1}\frac{4e^{2at}e^{2bt}e^{ct}}{\left(1e^{2at}\right)\left(1e^{2bt}\right)}+c^s\}=\mathrm{A}(s)+\mathrm{B}(s)+\mathrm{C}(s)+\mathrm{D}(s),$$
(4)
where $`\zeta (s,v)`$ is Hurwitz’ zeta function. This expression (invariant under $`ab`$) is, in principle, well defined for $`\mathrm{}s>2`$. Since the analytic structure of $`\mathrm{A}(s)`$ and $`\mathrm{B}(s)`$ is well known, we will concentrate on the Barnes term $`\mathrm{C}(s)`$, which will be extended to $`\mathrm{}s>2`$.
To this end, we will use the expansion
$$\frac{1}{e^{at}e^{at}}=\frac{1}{2at}+at\underset{k=1}{\overset{\mathrm{}}{}}(1)^k\frac{1}{(at)^2+(k\pi )^2},$$
(5)
thus obtaining
$`\mathrm{C}(s)=2\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}{\displaystyle \frac{1}{\mathrm{\Gamma }(s)}}\{{\displaystyle \frac{1}{2a}}{\displaystyle _0^{\mathrm{}}}dtt^{s2}{\displaystyle \frac{e^{(a+b+c)t}}{e^{bt}e^{bt}}}+`$
$`a{\displaystyle _0^{\mathrm{}}}dtt^s{\displaystyle \frac{e^{(a+b+c)t}}{e^{bt}e^{bt}}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(1)^k{\displaystyle \frac{1}{(at)^2+(k\pi )^2}}\}+ab=`$
$$\mathrm{C}_1(s)+\mathrm{C}_2(s).$$
(6)
The first term, $`\mathrm{C}_1(s)`$, can be easily seen to be
$$\mathrm{C}_1(s)=2\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\frac{1}{2a}\frac{1}{(s1)\left(2b\right)^{s1}}\zeta (s1,\frac{a+2b+c}{2b})+ab.$$
(7)
As all the terms we have analytically extended up to this point, $`\mathrm{C}_2(s)`$ in equation (6 )involves an integral which diverges at $`s=0`$. In order to isolate this singularity, we will rewrite this term as
$`\mathrm{C}_2(s)=2\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}{\displaystyle \frac{1}{\mathrm{\Gamma }(s)}}a{\displaystyle _0^{\mathrm{}}}dtt^s{\displaystyle \frac{e^{(a+b+c)t}}{\left(e^{bt}e^{bt}\right)}}\{{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(1)^k[{\displaystyle \frac{1}{(at)^2+(k\pi )^2}}{\displaystyle \frac{1}{(k\pi )^2}}]+`$
$$\underset{k=1}{\overset{\mathrm{}}{}}(1)^k\frac{1}{(k\pi )^2}\}+ab=\mathrm{CF}_2(s)+\mathrm{CD}_2(s).$$
(8)
The integral appearing in $`\mathrm{CD}_2(s)`$ is divergent at $`s=0`$ but, after performing the sum, it can be trivially extended to give
$$\mathrm{CD}_2(s)=\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\frac{a}{6}\frac{s}{\left(2b\right)^{s+1}}\zeta (s+1,1+\frac{a+c}{2b})+ab.$$
(9)
Now, $`\mathrm{CF}_2(s)`$, can be rewritten as
$$\mathrm{CF}_2(s)=2\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\frac{1}{\mathrm{\Gamma }(s)}a^3\underset{k=1}{\overset{\mathrm{}}{}}\frac{(1)^k}{(k\pi )^2}_0^{\mathrm{}}𝑑tt^{s+2}\frac{e^{(a+2b+c)t}}{\left(1e^{2bt}\right)}\frac{1}{(at)^2+(k\pi )^2}+ab.$$
(10)
As is easily seen, this integral converges for $`\mathrm{}s>2`$. We have thus obtained an analytic extension for the $`\zeta `$ of the operator as a meromorphic function with only simple poles. Such extension is valid for $`\mathrm{}s>2`$.
Now, the factor $`\frac{1}{(at)^2+(k\pi )^2}`$ can be written as an integral. In fact
$`{\displaystyle \frac{1}{(at)^2+(k\pi )^2}}={\displaystyle \frac{1}{2ik\pi }}\left[{\displaystyle \frac{1}{at+ik\pi }}{\displaystyle \frac{1}{atik\pi }}\right]={\displaystyle \frac{1}{k\pi }}{\displaystyle _0^{\mathrm{}}}𝑑ue^{atu}\mathrm{sin}(k\pi u).`$
When replaced in equation (10), this gives
$`\mathrm{CF}_2(s)=2\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}{\displaystyle \frac{1}{\mathrm{\Gamma }(s)}}a^3{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{(k\pi )^3}}{\displaystyle _0^{\mathrm{}}}𝑑tt^{s+2}{\displaystyle \frac{e^{(a+2b+c)t}}{\left(1e^{2bt}\right)}}{\displaystyle _0^{\mathrm{}}}𝑑ue^{atu}\mathrm{sin}(k\pi u)+ab`$
or, after interchanging the integrals
$`\mathrm{CF}_2^2(s)=2\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}{\displaystyle \frac{a^3}{\mathrm{\Gamma }(s)}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{(k\pi )^3}}{\displaystyle _0^{\mathrm{}}}𝑑u\mathrm{sin}(k\pi u){\displaystyle \frac{\mathrm{\Gamma }(s+3)}{(2b)^{s+3}}}\zeta (s+3,{\displaystyle \frac{a+2b+c+au}{2b}})+ab.`$
When the $`\zeta `$ function is written in terms of its series development (which is valid for $`\mathrm{}s>2`$) one has (after interchanging this series and the integral)
$`\mathrm{CF}_2(s)=2\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}{\displaystyle \frac{a^3}{\mathrm{\Gamma }(s)}}{\displaystyle \frac{\mathrm{\Gamma }(s+3)}{(2b)^{s+3}}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{(k\pi )^3}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}𝑑u\mathrm{sin}(k\pi u)\left(l+{\displaystyle \frac{a+c+au}{2b}}\right)^{(s+3)}+ab.`$
Finally, after performing the remaining integral and making use of the functional relations between incomplete gamma functions , one gets
$`\mathrm{CF}_2(s)=i\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}{\displaystyle \frac{\mathrm{\Gamma }(s+3)}{\mathrm{\Gamma }(s)}}a^s{\displaystyle \frac{1}{s+2}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{(k\pi )^{1s}}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}[i^{s+2}e^{i\frac{k\pi }{a}(2bl+a+c)}\mathrm{\Gamma }(s1,i{\displaystyle \frac{k\pi }{a}}(2bl+a+c))`$
$$(i)^{s+2}e^{i\frac{k\pi }{a}(2bl+a+c)}\mathrm{\Gamma }(s1,i\frac{k\pi }{a}(2bl+a+c))]+ab.$$
(11)
The replacement of equations (7), (9) and (11) into equation (4) completes the analytic extension of the relevant $`\zeta `$ function. Its value at $`s=0`$ can be easily computed, which gives:
$$\zeta (0;\overline{)}D^{}\overline{)}D)=\frac{\mu ^4\mathrm{\Omega }}{4\pi ^2}\left\{\frac{1}{2}c^2+\frac{a^2+b^2}{3}\right\}.$$
(12)
The agreement with the known results for null and equal fields is shown in Appendix A.
## IV The effective action and its renormalization
This section contains the main result in this paper, i.e., the one-loop correction to the Euclidean effective action. According to equation (2), to obtain such result, one must perform the derivatives at $`s=0`$ of the various terms in equation (4).
We start from $`\mathrm{A}(s)`$, which contributes with
$$\frac{1}{2}\frac{}{s}\mathrm{A}(s)_{s=0}=\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\{\mathrm{log}(2a)(\frac{1}{2}+\frac{c}{2a})+\mathrm{log}\mathrm{\Gamma }(\frac{c}{2a}+1)\frac{1}{2}\mathrm{log}(2\pi )\}.$$
(13)
In a completely analogous way, one has
$$\frac{1}{2}\frac{}{s}\mathrm{B}(s)_{s=0}=\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\{\mathrm{log}(2b)(\frac{1}{2}+\frac{c}{2b})+\mathrm{log}\mathrm{\Gamma }(\frac{c}{2b}+1)\frac{1}{2}\mathrm{log}(2\pi )\}.$$
(14)
It is also through a direct calculation that one gets
$$\frac{1}{2}\frac{}{s}\mathrm{C}_1(s)_{s=0}=\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\frac{1}{2a}\{2b(1+\mathrm{log}(2b))\zeta (1,1+\frac{a+c}{2b})2b\frac{}{s}_{s=0}\zeta (s1,1+\frac{a+c}{2b})\}+ab.$$
(15)
$$\frac{1}{2}\frac{}{s}\mathrm{CD}_2(s)_{s=0}=\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\frac{a}{24b}\{\mathrm{log}(2b)+\mathrm{\Psi }(1+\frac{a+c}{2b})\}+ab.$$
(16)
As regards $`\mathrm{CF}_2(s)`$, due to the presence of $`\mathrm{\Gamma }(s)`$ in the denominator, the required derivative reduces to the product $`\mathrm{\Gamma }(s)\mathrm{CF}_2(s)`$ at $`s=0`$, i.e.,
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{s}}\mathrm{CF}_2(s)_{s=0}={\displaystyle \frac{i}{2}}\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}{\displaystyle }_{k=1}^{\mathrm{}}{\displaystyle \frac{(1)^k}{k\pi }}{\displaystyle }_{l=1}^{\mathrm{}}[e^{i\frac{k\pi }{a}(2bl+a+c)}\mathrm{\Gamma }(1,{\displaystyle \frac{ik\pi }{a}}(2bl+a+c))`$
$$e^{i\frac{k\pi }{a}(2bl+a+c)}\mathrm{\Gamma }(1,\frac{ik\pi }{a}(2bl+a+c))]+ab.$$
(17)
Summarizing, the Euclidean effective action is given by the sum of the partial contributions in equations (13) to (17), plus
$$\frac{1}{2}\frac{}{s}\mathrm{D}(s)_{s=0}=\mu ^4\mathrm{\Omega }\frac{ab}{8\pi ^2}\mathrm{log}(c).$$
(18)
Notice that even though the result is finite, it depends on the arbitrary parameter $`\mu `$. However, this effective action still admits a finite renormalization. We will perform it by adopting the criterium (used, for instance, in reference ), that a very massive field does not fluctuate. Thus, we will subtract the one loop correction to the effective action in the limit $`m\mathrm{}`$, $`\mu \mathrm{}`$, with constant $`c`$. From equation (B6) in Appendix B, the effective action in this limit can be seen to be
$$\mu ^4\mathrm{\Omega }\frac{1}{4\pi ^2}\left\{\left[\frac{3}{8}\frac{1}{4}\mathrm{log}(c)\right]c^2\frac{1}{6}(b^2+a^2)\mathrm{log}(c)\right\}.$$
(19)
After doing this subtraction, all dependence on the parameter $`\mu `$ disappears, and the Euclidean effective action is given by
$`S_{eff}^{Ren}\left[A_\mu \right]={\displaystyle \frac{\mathrm{\Omega }\mu ^4}{2e^2}}(a^2+b^2)+`$
$`\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}\{{\displaystyle \frac{1}{8}}\mathrm{log}\left({\displaystyle \frac{4ab}{c^2}}\right){\displaystyle \frac{1}{24}}{\displaystyle \frac{(a^2+b^2)}{ab}}\mathrm{log}\left({\displaystyle \frac{4ab}{c^2}}\right)+{\displaystyle \frac{c}{4a}}\mathrm{log}\left({\displaystyle \frac{a}{b}}\right){\displaystyle \frac{c^2}{16ab}}\mathrm{log}\left({\displaystyle \frac{4ab}{c^2}}\right)+`$
$`\mathrm{log}\left({\displaystyle \frac{\mathrm{\Gamma }(\frac{c}{2a}+1)}{\sqrt{2\pi }}}\right){\displaystyle \frac{b}{a}}\zeta (1,1+{\displaystyle \frac{a+c}{2b}}){\displaystyle \frac{b}{a}}{\displaystyle \frac{}{s}}_{s=0}\zeta (s1,1+{\displaystyle \frac{a+c}{2b}})`$
$`{\displaystyle \frac{i}{2}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{k\pi }}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left[e^{i\frac{k\pi }{a}(2bl+a+c)}\mathrm{\Gamma }(1,{\displaystyle \frac{ik\pi }{a}}(2bl+a+c))e^{i\frac{k\pi }{a}(2bl+a+c)}\mathrm{\Gamma }(1,{\displaystyle \frac{ik\pi }{a}}(2bl+a+c))\right]+`$
$$\frac{a}{24b}\mathrm{\Psi }(1+\frac{a+c}{2b})\frac{3}{16}\frac{c^2}{ab}\}+ab.$$
(20)
The renormalization performed amounts to subtracting the zero field effective action (thus redefining the cosmological constant), and renormalizing the classical action. As a result, one gets the following running charge relationship
$$\frac{1}{e^2}=\frac{1}{e_0^2}+\frac{1}{12\pi ^2}\mathrm{log}\frac{\mu ^2}{m^2}.$$
(21)
Equivalently, for the fine structure constant one has
$$\alpha =\frac{\alpha _0}{1+\frac{\alpha _0}{3\pi }\mathrm{log}\frac{\mu ^2}{m^2}}.$$
(22)
Note that this expression reduces, in the perturbative limit, to the well known result (see, for example )
$$\alpha =\alpha _0(1\frac{\alpha _0}{3\pi }\mathrm{log}\frac{\mu ^2}{m^2}).$$
(23)
## V Comparison with the proper time result
In Appendix B we show that, in the weak field limit, our result for the $`\zeta `$ regularized effective action coincides, once renormalized, with the Euclidean version of the well known Schwinger’s proper time one.
In this section, we will show that this is also the case for arbitrary field strenghts. In fact, Schwinger’s integral expresion for the one loop correction to the effective action is given, after subtracting the divergent terms, by
$$S_{PT}^{(1)}=\mu ^4\mathrm{\Omega }\{\frac{ab}{8\pi ^2}_0^{\mathrm{}}dtt^{s1}e^{ct}\mathrm{coth}(bt)\mathrm{coth}(at)\frac{1}{8\pi ^2}_0^{\mathrm{}}dtt^{s3}e^{ct}\frac{a^2+b^2}{24\pi ^2}_0^{\mathrm{}}dtt^{s1}e^{ct}\}_{s=0}.$$
(24)
Now, performing the integrals in the last two terms and comparing with equation (4) (with the Hurwitz’s zetas written in integral form), the previous expression can be rewritten as
$$S_{PT}^{(1)}=\frac{1}{2}\{\mathrm{\Gamma }(s)\zeta (s;\overline{)}D^{}\overline{)}D)\frac{\mu ^4\mathrm{\Omega }}{4\pi ^2}(c^{2s}\mathrm{\Gamma }(s2)+\frac{a^2+b^2}{3}c^s\mathrm{\Gamma }(s))\}_{s=0}.$$
(25)
After developing around $`s=0`$, it is easy to see that
$$S_{PT}^{(1)}=S_\zeta ^{(1)}\frac{\mu ^4\mathrm{\Omega }}{4\pi ^2}\left[\frac{3}{8}c^2\left(\frac{c^2}{4}+\frac{a^2+b^2}{6}\right)\mathrm{log}c\right],$$
(26)
where $`S_\zeta ^{(1)}`$ is the $`\zeta `$-regularized one loop correction to the effective action, as defined in equation (2), and the remaining terms are precisely the ones we have subtracted through renormalization. So, the exact agreement between both renormalized effective actions is apparent.
###### Acknowledgements.
We thank Horacio Falomir, Klaus Kirsten and Roberto Soldati for carefully reading the manuscript, and for many useful suggestions. This work was partially supported by UNLP, under Grant No 11/X230, ANPCyT, under Grant PICT00039, and CONICET, under Grant PIP0459.
## A The limits of null and equal fields
In this section, we will show the agreement of our general $`\zeta `$ function with the results obtained by other authors for some particular cases, i.e., the case of a null electric or magnetic field and that of equal electric and magnetic fields .
We will start with the $`B0`$ limit. It is easy to see that $`lim_{b0}\mathrm{A}(s)=0`$. As regards $`lim_{b0}\mathrm{B}(s)`$, it can be studied by making use of the asymptotic expansion for Hurwitz’ $`\zeta `$ function (see, for example, )
$$\zeta (s,v)=\frac{1}{\mathrm{\Gamma }(s)}\left\{v^{1s}\mathrm{\Gamma }(s1)+\frac{1}{2}v^s\mathrm{\Gamma }(s)+\underset{n=1}{\overset{N}{}}B_{2n}\frac{\mathrm{\Gamma }(s+2n1)}{(2n)!}v^{1s2n}\right\}+O(v^{2Ns1}),$$
(A1)
which gives
$$\underset{b0}{lim}\mathrm{B}(s)=\underset{b0}{lim}\left\{\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\frac{2}{(2b)^s}\frac{\mathrm{\Gamma }(s1)}{\mathrm{\Gamma }(s)}\left(\frac{c}{2b}+1\right)^{1s}\right\}=\frac{\mu ^4\mathrm{\Omega }}{4\pi ^2}\frac{a}{s1}c^{1s}.$$
(A2)
The only contribution to $`\mathrm{C}(s)`$ in this limit comes from $`\mathrm{C}_1(s)`$, which gives
$$\underset{b0}{lim}\mathrm{C}(s)=\frac{\mu ^4\mathrm{\Omega }}{4\pi ^2}\frac{(2a)^{2s}}{s1}\left\{\zeta (s1,\frac{c}{2a})\left(\frac{c}{2a}\right)^{1s}\right\}.$$
(A3)
Finally, $`\mathrm{D}(s)`$ vanishes for $`b=0`$. Then, replacing all these partial results into equation (4), one obtains
$$\zeta (s,\overline{)}D^{}\overline{)}D)_{b=0}=\frac{\mu ^4\mathrm{\Omega }}{4\pi ^2}\frac{(2)^{1s}}{s1}a^{2s}\{2\zeta (s1,\frac{c}{2a})\left(\frac{c}{2a}\right)^{1s}\}$$
(A4)
which is in complete agreement with previous results .
Of course, the $`E0`$ limit, gives an analogous expression, which can be obtained by changing $`ab`$ in equation (A4).
We will now study the equal fields limit. In this situation, taking $`a=b`$ in the different terms appearing in the $`\zeta `$ function (4), we have
$`\zeta (s;\overline{)}D^{}\overline{)}D)_{a=b}=\mu ^4\mathrm{\Omega }{\displaystyle \frac{a^2}{4\pi ^2}}\{{\displaystyle \frac{4}{\left(2a\right)^s}}\zeta (s,{\displaystyle \frac{c}{2a}}+1)+c^s+{\displaystyle \frac{2^{2s}a^s}{s1}}\zeta (s1,{\displaystyle \frac{3}{2}}+{\displaystyle \frac{c}{2a}})`$
$`{\displaystyle \frac{1}{6}}(2a)^s\zeta (s+1,{\displaystyle \frac{3}{2}}+{\displaystyle \frac{c}{2a}})i\mathrm{\hspace{0.17em}2}a^s(s+1)s{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^{k+1}}{(k\pi )^{1s}}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}[i^{s+2}e^{ik\pi (2l+1+\frac{c}{a})}\mathrm{\Gamma }(s1,ik\pi (2l+1+{\displaystyle \frac{c}{a}}))`$
$$(i)^{s+2}e^{ik\pi (2l+1+\frac{c}{a})}\mathrm{\Gamma }(s1,ik\pi (2l+1+\frac{c}{a}))]\}.$$
(A5)
In order to compare this expression with the result in , we use the functional relations between incomplete gamma functions, thus getting
$`\zeta (s;\overline{)}D^{}\overline{)}D)_{a=b}=\mu ^4\mathrm{\Omega }{\displaystyle \frac{a^2}{4\pi ^2}}\{{\displaystyle \frac{4}{\left(2a\right)^s}}\zeta (s,{\displaystyle \frac{c}{2a}}+1)+c^s+{\displaystyle \frac{2^{2s}a^s}{s1}}\zeta (s1,{\displaystyle \frac{3}{2}}+{\displaystyle \frac{c}{2a}})`$
$`i\mathrm{\hspace{0.17em}2}a^ss{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{(k\pi )^{1s}}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}[i^{s+2}e^{ik\pi (2l+1+\frac{c}{a})}\mathrm{\Gamma }(s,ik\pi (2l+1+{\displaystyle \frac{c}{a}}))`$
$$(i)^{s+2}e^{ik\pi (2l+1+\frac{c}{a})}\mathrm{\Gamma }(s,ik\pi (2l+1+\frac{c}{a}))]\}.$$
(A6)
We now use the integral representation for the incomplete gamma function
$`\mathrm{\Gamma }(\alpha ,x)={\displaystyle _x^{\mathrm{}}}𝑑te^tt^{\alpha 1}.`$
When doing so, and after interchanging the integral and the sum over $`l`$, the last term in equation(A6) can be written as
$`(2a)^ss{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{(k\pi )^2}}{\displaystyle _0^{\mathrm{}}}𝑑ue^u\left[\zeta (s+1,{\displaystyle \frac{3}{2}}+{\displaystyle \frac{c}{2a}}{\displaystyle \frac{iu}{2k\pi }})+\zeta (s+1,{\displaystyle \frac{3}{2}}+{\displaystyle \frac{c}{2a}}+{\displaystyle \frac{iu}{2k\pi }})\right]=`$
$`2(2a)^s{\displaystyle \frac{1}{\mathrm{\Gamma }(s)}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(1)^k{\displaystyle _0^{\mathrm{}}}𝑑tt^s{\displaystyle \frac{e^{(\frac{3}{2}+\frac{c}{2a})t}}{1e^t}}{\displaystyle \frac{1}{(k\pi )^2+(\frac{t}{2})^2}}`$
where we have used the integral form for the Hurwitz’s zeta functions, interchanged the integrals and performed the interior one.
Interchanging now the integral with the sum, and using equation (5), we obtain
$`2^{2s}a^s{\displaystyle \frac{1}{\mathrm{\Gamma }(s)}}{\displaystyle _0^{\mathrm{}}}𝑑tt^{s1}{\displaystyle \frac{e^{(\frac{3}{2}+\frac{c}{2a})t}}{1e^t}}\left[{\displaystyle \frac{e^{\frac{t}{2}}}{1e^t}}{\displaystyle \frac{1}{t}}\right]=`$
$`2^{2s}a^s\left[\zeta (s1,{\displaystyle \frac{c}{2a}}+1)({\displaystyle \frac{c}{2a}}+1)\zeta (s,{\displaystyle \frac{c}{2a}}+1){\displaystyle \frac{1}{s1}}\zeta (s1,{\displaystyle \frac{3}{2}}+{\displaystyle \frac{c}{2a}})\right].`$
When replaced in (A6), the final result is
$`\zeta (s;\overline{)}D^{}\overline{)}D)_{a=b}=\mu ^4\mathrm{\Omega }{\displaystyle \frac{a^2}{4\pi ^2}}\{c^s+2^{2s}a^s[\zeta (s1,{\displaystyle \frac{c}{2a}}+1){\displaystyle \frac{c}{2a}}\zeta (s,{\displaystyle \frac{c}{2a}}+1)]\}=`$
$$\mu ^4\mathrm{\Omega }\frac{a^2}{4\pi ^2}\left\{c^s+4(2a)^s\left(\zeta (s1,\frac{c}{2a})\frac{c}{2a}\zeta (s,\frac{c}{2a})\right)\right\}.$$
(A7)
This expression coincides whith the result obtained in (see equations (5.2.6) and (5.2.4) in that reference).
## B The weak-field limit
An unavoidable test our effective action must resist is its coincidence whith the well known result for weak fields . In order to check this is the case, we will develop the different contributions to the effective action (equations (13) to (18)) in powers of the fields over the squared mass. In the cases of equations (13) to (16), such development can be obtained by making use of the well known asymptotic expansions for $`\mathrm{log}\mathrm{\Gamma }(x)`$, $`\psi (x)`$, and $`\zeta (s,x)`$ (see also our equation (A1)). When doing so, and retaining terms up to the order of squared fields over mass to the fourth, one gets, after a straightforward though tedious calculation,
$$\frac{1}{2}\frac{}{s}\mathrm{A}(s)_{s=0}\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\{\frac{1}{6}ac^1+\frac{1}{2}\mathrm{log}(c)+\frac{1}{2a}(\mathrm{log}(c)1)c\}.$$
(B1)
$$\frac{1}{2}\frac{}{s}\mathrm{B}(s)_{s=0}\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\{\frac{1}{6}bc^1+\frac{1}{2}\mathrm{log}(c)+\frac{1}{2b}(\mathrm{log}(c)1)c\}.$$
(B2)
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{s}}\mathrm{C}_1(s)_{s=0}\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}{\displaystyle \frac{1}{ab}}\{({\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{4}}\mathrm{log}c+{\displaystyle \frac{1}{8}})c^2+({\displaystyle \frac{1}{2}}(a+b){\displaystyle \frac{1}{2}}(a+b)\mathrm{log}c)c{\displaystyle \frac{5}{24}}(a^2+b^2)`$
$$\frac{1}{2}ab\mathrm{log}c\frac{1}{24}(5ba^2+5ab^2+a^3+b^3)c^1+(\frac{1}{24}b^3a+\frac{1}{24}a^3b+\frac{1}{12}b^2a^2+\frac{7}{1440}a^4+\frac{7}{1440}b^4)c^2\}.$$
(B3)
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{s}}\mathrm{CD}_2(s)_{s=0}\mu ^4\mathrm{\Omega }{\displaystyle \frac{ab}{4\pi ^2}}{\displaystyle \frac{1}{24}}\{({\displaystyle \frac{a}{b}}+{\displaystyle \frac{b}{a}})\mathrm{log}c+(a+b+{\displaystyle \frac{a^2}{b}}+{\displaystyle \frac{b^2}{a}})c^1`$
$$\frac{1}{2}(2a^2+2b^2+\frac{a^3}{b}+\frac{b^3}{a}+\frac{4}{3}ba)c^2\}.$$
(B4)
As regards $`\frac{1}{2}\frac{}{s}\mathrm{CF}_2(s)_{s=0}`$, one has to use the asymptotic expansions for the incomplete $`\mathrm{\Gamma }`$ function and for the Hurwitz’ zeta functions (equation (A1)). After doing so, one obtains
$$\frac{1}{2}\frac{}{s}\mathrm{CF}_2(s)_{s=0}\mu ^4\mathrm{\Omega }\frac{ab}{4\pi ^2}\frac{7}{1440}(\frac{a^3}{b}+\frac{b^3}{a})c^2.$$
(B5)
By summing up the contributions in equations (B1) to (B5), plus the one coming from $`\frac{1}{2}\frac{}{s}\mathrm{D}(s)_{s=0}`$, the one-loop correction to the effective action is seen to reduce, in this weak-field limit, to
$$S^{(1)}=\mu ^4\mathrm{\Omega }\frac{1}{4\pi ^2}\left\{\left[\frac{3}{8}\frac{1}{4}\mathrm{log}(c)\right]c^2\frac{1}{6}(b^2+a^2)\mathrm{log}(c)+\left[\frac{7}{90}(ab)^2\frac{1}{90}(a^2+b^2)^2\right]c^2\right\}.$$
(B6)
Now, renormalizing according to the criterium discussed in Section IV, one is left with
$$S_{eff}=\frac{\mathrm{\Omega }}{2}(B^2+E^2)+\frac{\mathrm{\Omega }e^4}{8\pi ^2m^4}\left[\frac{7}{45}(EB)^2\frac{1}{45}(E^2+B^2)^2\right],$$
(B7)
where the definitions of $`a`$, $`b`$ and $`c`$ given in the paragraph following equation (3) were used.
The expression in (B7) is precisely the Euclidean version of the Euler-Heisenberg effective action for weak fields .
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# A Dual Gauge Model with Confinement
## 1 Introduction
There is a general statement that the color confinement is supported by the idea that the vacuum of quantum Yang-Mills (Y-M) theory is realized by a condensate of monopole-antimonople pairs . In such a vacuum the interacting field between two colored sources located in $`\stackrel{}{x}_1`$ and $`\stackrel{}{x}_2`$ is squeezed into a tube whose energy $`E_{tube}|\stackrel{}{x}_1\stackrel{}{x}_2|`$. This is a complete dual analogy to the magnetic monopole confinement in the type II superconductor. Since there is no monopoles as classical solutions with finite energy in a pure Y-M theory it has been suggested by ’t Hooft to go into the Abelian projection where the gauge group SU(2) is broken by a suitable gauge condition to its (may be maximal) Abelian subgroup U(1). It is proposed that the interplay between a quark-antiquark pair is analagous to the interaction between a monopole-antimonopole pair in a superconductor.
In fact, the topology of Y-M SU(N) manifold and that of its Abelian subgroup $`[U(1)]^{N1}`$ are different, and since any such gauge is singular, one might introduce the string by performing the singular gauge transformation with an Abelian gauge field $`A_\mu `$
$`A_\mu (x)A_\mu (x)+{\displaystyle \frac{g}{4\pi }}_\mu \mathrm{\Omega }(x),`$ (1.1)
where $`\mathrm{\Omega }(x)`$ is the solid angle subtended by the closed space-like curve described by the string at any point $`x=(x^0,x^1)`$, and $`g=2\pi /e`$ is responsible for the magnetic flux inside the string, $`e`$ being the Y-M coupling constant. Here, we choose a single string in the two-dimensional (2d) world sheet $`y_x(\tau ,\sigma )`$, for simplicity. Obviously, the Abelian field-strength tensor $`F_{\mu \nu }^A=_\mu A_\nu _\nu A_\mu `$ transforms as
$$F_{\mu \nu }^A(x)F_{\mu \nu }^A(x)+\stackrel{~}{G}_{\mu \nu }(x),$$
where a new term
$$\stackrel{~}{G}_{\mu \nu }(x)=\frac{g}{4\pi }[_\mu ,_\nu ]\mathrm{\Omega }(x),$$
is valid on the world sheet only
$$\stackrel{~}{G}_{\mu \nu }(x)=\frac{g}{2}ϵ_{\mu \nu \alpha \beta }𝑑\sigma 𝑑\tau \frac{(y^\alpha ,y^\beta )}{(\sigma ,\tau )}\delta _4[xy(\sigma ,\tau )].$$
Formally, a gauge group element, which transforms a generic SU(N) connection onto the gauge fixing surface in the space of connections, is not regular everywhere in spacetime. The projected (or transformed) connections contain topological singularities (or defects). Such a singular transformation (1.1) may form the worldline(s) of magnetic monopoles. Hence, this singularity leads to the monopole current $`J_\mu ^{mon}`$. This is a natural way of the transformation from the Y-M theory to a model dealing with Abelian fields. A dual string is nothing but a formal idealization of a magnetic flux tube in the equilibrium against the pressure of surrounding superfluid (Higgs-like field) which it displaces .
Recent lattice results give the promised picture that the monopole degrees of freedom can indeed form a condensate responsible for the confinement. The expression for the static heavy quark potential, using an effective dual Ginzburg-Landau model , has been presented in . In the paper , an analytic approximation to the dual field propagator without sources and in the presence of quark sources, and an expression for the static quark-antiquark potential were established.
The aim of this paper is to consider the model in 4d based on the dual description of a long-distance Y-M (LDY-M) theory which shows some kind of confinement. We study the model of Lagrangian where the fundamental variables are an octet of dual potentials coupled minimally to three octets of monopole (Higgs) fields. The dual gauge model is studied at the lowest order of the perturbative series using the canonical quantization. The basic manifestation of the model is that it generates the equations of motion where one of them for the scalar Higgs field looks like as a dipole-like field equation. The monopole fields obeying such an equation are classified by their two-point Wightman functions (TPWF). In the classical level there is some intersection with the Froissart model containing the scalar field satisfying the equation of the fourth order. In the scheme presented in this work, the flux distribution in the tubes formed between two heavy color charges is understood via the following statement: the Abelian monopoles are excluded from the string region while the Abelian electric flux is squeezed into the string region.
In Sec. 2, we introduce the essence of the dual gauge Higgs model and the classical solutions. In our model there are the dual gauge field $`\widehat{C}_\mu ^a(x)`$ and the monopole field $`\widehat{B}_i^a(x)`$ ($`i=1,\mathrm{},N_c(N_c1)/2`$; $`a`$=1,…,8 is a color index ) which are relevant modes for infrared behaviour. The local coupling of the $`\widehat{B}_i`$-field to the $`\widehat{C}_\mu `$-field provides the mass of the dual field and, hence, a dual Meissner effect. Although $`\widehat{C}_\mu (x)`$ is invariant under the local transformation of $`U(1)^{N_c1}SU(N_c)`$, $`\widehat{C}_\mu =\stackrel{}{C}_\mu \stackrel{}{H}`$ is an $`SU(N_c)`$-gauge dependent object and does not appear in the real world alone ($`N_c`$ is the number of colors and $`\stackrel{}{H}`$ stands for the Cartan superalgebra). The commutation relations, TPWF and Green’s functions as well-defined distributions in the space $`S(\mathrm{}^d)`$ of complex Schwartz test functions on $`\mathrm{}^d`$ will be defined in Sec. 3. In Sec. 4, we study the monopole- and dual gauge- field propagations. In Sec. 5, we obtain the asymptotic transverse behaviour of both the dual gauge field and the color-electric field. The analytic expression for the static potential is obtained in Sec. 6. Sec. 7 contains the discussion and conclusions.
## 2 A dual Higgs gauge model and classical solutions
The dual description of the LDY-M theory is simply understood by switching on the dual gauge field $`\widehat{C}_\mu (x)`$ and the three scalar octets $`\widehat{B}_i(x)`$ (necessary to give mass to all $`C_\mu ^a`$ and carrying color magnetic charge) in the Lagrangian density (LD) $`L`$
$`L=2Tr\left[{\displaystyle \frac{1}{4}}\widehat{F}^{\mu \nu }\widehat{F}_{\mu \nu }+{\displaystyle \frac{1}{2}}\left(D_\mu \widehat{B}_i\right)^2\right]W\left(\widehat{B}_i\right),`$ (2.1)
where
$$\widehat{F}_{\mu \nu }=_\mu \widehat{C}_\nu _\nu \widehat{C}_\mu ig[\widehat{C}_\mu ,\widehat{C}_\nu ],$$
$$D_\mu \widehat{B}_i=_\mu \widehat{B}_iig[\widehat{C}_\mu ,\widehat{B}_i],$$
$`\widehat{C}_\mu `$ and $`\widehat{B}_i`$ are the SU(3) matrices, g is the gauge coupling constant of the dual theory. The Higgs fields develop their vacuum expectation values (v.e.v.) $`\widehat{B}_{0_i}`$ and the Higgs potential $`W(\widehat{B}_i)`$ has a minimum at $`\widehat{B}_{0_i}`$. The v.e.v. $`\widehat{B}_{0_i}`$ produce a color monopole generating current confining the electric color flux. It is known, the LD (2.1) can generate classical equations of motion carrying a unit of the $`z_3`$ flux confined in a narrow tube along the $`z`$-axis (corresponding to quark sources at $`z=\pm \mathrm{}`$). This is a dual analogy to the Abrikosov magnetic vortex solution.
As the next step we introduce the color ansatz
$`\widehat{C}_\mu ={\displaystyle \underset{a}{}}C_\mu ^a{\displaystyle \frac{1}{2}}\lambda _a,`$ (2.2)
where the vector potential $`C_\mu ^a`$ is dual to an ordinary vector potential in the Y-M theory, (1/2)$`\lambda ^a`$ are generators of SU(3). Following paper in the sense of representing the quark sources by the Dirac string tensor $`\stackrel{~}{G}_{\mu \nu }(x)`$ having the same color structure as in (2.2), one can arrive at a more suitable form of the LD (2.1)
$`L(\stackrel{~}{G}_{\mu \nu })={\displaystyle \frac{1}{3}}G_{\mu \nu }^2+4|(_\mu igC_\mu )\varphi |^2+2(_\mu \varphi _3)^2W(\varphi ,\varphi _3),`$ (2.3)
where
$$G_{\mu \nu }=_\mu C_\nu _\nu C_\mu +\stackrel{~}{G}_{\mu \nu },$$
and $`\varphi (x)`$ and $`\varphi _3(x)`$ denote the complex scalar monopole fields. We choose the color structure for the QCD-monopole field $`\widehat{B}_i`$ (belonging to the fundamental representation of $`SU_c(3)`$) like in , and the effective potential stands
$$W(B,\overline{B},B_3)=\frac{2}{3}\lambda \{11[2(B^2+\overline{B}^2B_0^2)^2+(B_3^2B_0^2)^2]$$
$$+7[2(B^2+\overline{B}^2)+B_3^23B_0^2]^2\},$$
where
$$\varphi \varphi _1=\varphi _2=B_{1,2}i\overline{B}_{1,2},\varphi _3=B_3,$$
while $`\lambda `$ provides the weak couplings between the scalar fields.
The dual gauge field $`C_\mu `$ satisfies the relation
$$_\mu C_\nu _\nu C_\mu =^{}(_\mu A_\nu _\nu A_\mu )$$
in the absence of charges, and the duality transformation is realized by interchanging the gluon field $`A_\mu (x)`$ and $`C_\mu (x)`$. In the maximally Abelian gauge $`A_\mu (x)=A_\mu ^a(x)(\tau ^a/2)`$ behaves as the Abelian gauge field $`C_\mu (x)=A_\mu ^3(x)(\tau ^3/2)`$ approximately because the off-diagonal fields are suppressed by the gauge transformation. The LD ( 2.3 ) is invariant under the local gauge transformation of the dual gauge field $`C_\mu `$
$$C_\mu (x)C_\mu (x)_\mu \mathrm{\Lambda }(x)$$
and the phase transformation of the QCD-monopole field
$$\varphi _{1,2}(x)\mathrm{exp}[ig\mathrm{\Lambda }(x)]\varphi _{1,2}(x),$$
where $`\mathrm{\Lambda }(x)`$ is the real field in $`S(\mathrm{}^3)`$ at any fixed $`x^0`$. The local gauge symmetry is spontaneously broken because of $`\varphi _00`$ in (2.3). The generating current of (2.3) is nothing but the monopole current confining the electric color flux
$$J_\mu ^{mon}=4ig[\varphi ^{}(_\mu igC_\mu )\varphi \varphi (_\mu +igC_\mu )\varphi ^{}],$$
which enters into the equation of motion in the form
$$^\nu G_{\mu \nu }(x)=\frac{3}{2}J_\mu ^{mon}(x).$$
The formal consequence of the $`J_\mu ^{mon}`$ conservation, $`^\mu J_\mu ^{mon}=0`$, means that monopole currents form closed loops.
To find a solution of this model, one can consider the monopole-field as a solution where $`B(x)_00`$, $`\overline{B}(x)_00`$, $`B_3(x)_00`$. We choose
$$B(x)=b(x)+B_0,\overline{B}(x)=\overline{b}(x),B_3(x)=b_3(x)+B_0$$
with the boundary conditions at large distances $`\rho `$ from the center of the flux tube with $`\vartheta `$ as an angle in the cylindrical coordinates
$`\stackrel{}{C}{\displaystyle \frac{e}{g\rho }}\stackrel{}{e}_\vartheta ,\varphi B_0\mathrm{exp}(i\vartheta ),B_3B_0as\rho \mathrm{}.`$ (2.4)
In terms of the fields $`b(x),\overline{b}(x),b_3(x)`$ and $`C_\mu (x)`$ the LD (2.3) is divided into two parts
$`L=L_1+L_2,`$ (2.5)
where $`L_1`$ in the lowest order of $`g`$ and $`\lambda `$ and with the minimal weak interaction looks like
$`L_1={\displaystyle \frac{1}{3}}G_{\mu \nu }^2+\mathrm{\hspace{0.17em}4}\left[(_\mu b)^2+(_\mu \overline{b})^2+{\displaystyle \frac{1}{2}}(_\mu b_3)^2\right]`$ (2.6)
$`+m^2C_\mu ^2{\displaystyle \frac{4}{3}}\mu ^2(50b^2+18b_3^2)+8m_\mu \overline{b}C_\mu .`$ (2.7)
Here, $`mgB_0`$ and $`\mu \sqrt{2\lambda }B_0`$ are masses of the dual gauge field and the monopole field, respectively. The remaining part of (2.5) turns out to be
$$L_2=8g\left[(_\mu \overline{b})C_\mu b(_\mu b)C_\mu \overline{b}\right]+4g^2(_\mu C_\nu _\nu C_\mu )^2(\overline{b}^2+b^2+2B_0b)$$
$$\frac{4}{3}\lambda [25(b^4+\overline{b}^4)+9b_3^4+100B_0b(\overline{b}^2+b^2+B_0^2)+28B_0(\overline{b}^2+b^2+bb_3)b_3$$
$$+36B_0b_3(b_3^2+B_0^2)+2b^2(25\overline{b}^2+7b_3^2)2B_0^3(50b+18b_3)+56B_0^2bb_3+14\overline{b}^2b_3^2].$$
Let us consider the canonical quantization of (2.6). The equations of motion are
$$(\mathrm{\Delta }^2+\mu _1^2)b(x)=0;$$
$`\mathrm{\Delta }^2\overline{b}(x)+m(C)=0;`$ (2.8)
$$(\mathrm{\Delta }^2+\mu _2^2)b_3(x)=0;$$
$`(\mathrm{\Delta }^2+m_1^2)C_\mu (x)_\mu (C)+12m_\mu \overline{b}(x)^\nu \stackrel{~}{G}_{\mu \nu }(x)=0,`$ (2.9)
where $`\mu _1^2=(50/3)\mu ^2,\mu _2^2=12\mu ^2,m_1^2=3m^2`$. By taking the divergence of (2.9) one can get
$`\mathrm{\Delta }^2\overline{b}(x){\displaystyle \frac{1}{9m}}^\mu ^\nu \stackrel{~}{G}_{\mu \nu }(x)=0,`$ (2.10)
while the formal solution of equation (2.9) looks like
$$C_\mu (x)=\alpha ^\nu \stackrel{~}{G}_{\mu \nu }(x)\beta _\mu \overline{b}(x),$$
where $`\alpha (3m^2)^1`$, $`\beta 4/m`$. We see that the dual gauge field is defined via the divergence of the scalar field $`\overline{b}(x)`$ shifted by the divergence of the Dirac string tensor $`\stackrel{~}{G}_{\mu \nu }(x)`$. At the same time, we propose in the standard manner the Dirac string which can be understood as a straight line connecting two objects with opposite charges. For large enough $`\stackrel{}{x}`$, approaching such a string, the monopole field is going to its v.e.v. while $`C_\mu (\stackrel{}{x}\mathrm{})0`$. Hence,
$$J_\mu ^{mon}(\stackrel{}{x}\mathrm{})8m^2C_\mu .$$
For a very weak $`C_\mu `$-field one can say that in the $`d=2h`$ dimensions
$`\mathrm{\Delta }^{2h}\overline{b}(x)0,h=2,3,\mathrm{},`$ (2.11)
but
$$\mathrm{\Delta }^2\overline{b}(x)0.$$
Here, the solutions of equation (2.11) obey locality, Poincare covariance and spectral conditions, and look like the dipole ”ghosts” at h=2. In this model, the role of the dipole field at d=4 is held in the d=2 dimensions by the simple pole field, and the analogy of the behaviour between the d=2 and d=4 dimensions can be found at the level of Wightman functions in the free case, at least . This is related to the model proposed in in a study of the Higgs phenomenon, from which the present model is distinguished by the coupling to the dual gauge field $`C_\mu (x)`$, and the gauge field strengh tensor is shifted by the Dirac string tensor $`\stackrel{~}{G}_{\mu \nu }(x)`$. Thus, we obtain that the massless scalar field $`\overline{b}(x)`$ occurs in the model since the symmetry realizes such a way that the LD (2.3) is invariant under the local gauge transformations above mentioned but $`\overline{b}(x)_00`$.
Our aim is to find Green’s function of a scalar field $`\overline{b}(x)`$ obeying Eq. (2.11). The propagator $`\tau _h(x)`$ is defined via TPWF $`W_h(x)`$ in the d=2 h-dimensions
$`\tau _h(x)=T\overline{b}(x)\overline{b}(0)_0=\theta (x^0)W_h(x)+\theta (x^0)W_h(x),`$ (2.12)
where
$`W_h(x)=\overline{b}(x)\overline{b}(0)_0`$ (2.13)
is the distribution in the Schwartz space $`S^{}(\mathrm{}^{2h})`$ of temperate distributions on $`\mathrm{}^{2h}`$ and obeys the equation
$`\mathrm{\Delta }^{2h}W_h(x)=0.`$ (2.14)
The general solution of eq. (2.14) should be Lorentz invariant and is given in the form at h=2
$`W_2(x)=a_1\mathrm{ln}{\displaystyle \frac{l^2}{x_\mu ^2+iϵx^0}}+{\displaystyle \frac{a_2}{x_\mu ^2iϵx^0}}+a_3,`$ (2.15)
where $`a_i`$ (i=1,2,3) are the coefficients to be defined later, while $`l`$ is an arbitrary parameter with dimension minus one in mass units. Its origin becomes more transparent from . In fact, the first and the second terms in (2.15) are related to the scalar dipole field and the scalar pole field, respectively . In the first case, the solution is
$`W_2(x)=iE_2^{}(x)={\displaystyle \frac{1}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{l^2}{x_\mu ^2+iϵx^0}},`$ (2.16)
while for the second term the solution looks like
$`W_2(x)=iD_2^{}(x)={\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle \frac{1}{x_\mu ^2+iϵx^0}}.`$ (2.17)
## 3 TPWF as classical distributions
Before going into the quantization procedure, let us briefly consider the classical distributions (in the sense of generalized functions ) of TPWF and T-ordered TPWF for the $`\overline{b}(x)`$-field at large $`x_\mu ^2`$. The TPWF (2.13) at h=2 is provided by the distribution $`\theta (p^0)\delta ^{}(p^2)`$ as
$`W_2(x){\displaystyle d_4p\theta (p^0)\delta ^{}(p^2)\mathrm{exp}(ipx)}\mathrm{ln}[\stackrel{~}{\mu }^2x^2+i\theta (x^0)]`$ (3.1)
$`=\left[\mathrm{ln}|\stackrel{~}{\mu }^2x^2|+i\pi sgn(x^0)\theta (x^2)\right],`$ (3.2)
where $`\stackrel{~}{\mu }`$ is an arbitrary parameter (the infrared regularization parameter). The distribution $`\theta (p^0)\delta ^{}(p^2)`$ in (3.1) is uniquely defined only on the test functions $`f(p)S_0(\mathrm{}_4)`$ $`(S_0(\mathrm{}_4)=\{f(p)S(\mathrm{}_4),f(p=0)=0\})`$
$$d_4p\theta (p^0)\delta ^{}(p^2)f(p)=\underset{\nu ^20}{lim}\frac{}{\nu ^2}d_4p\theta (p^0)\delta ^{}(p^2\nu ^2)f(p),$$
$$=_{\mathrm{\Gamma }_0^+}\frac{d_3\stackrel{}{p}}{2p^0}\frac{1}{2(np)}\left[\frac{1}{(np)}(n)\right]f(p)$$
and for an arbitrary fixed time-like unit vector $`n_\mu `$ in the case we choose $`n_\mu =(1,\stackrel{}{0})`$ from $`V^+`$ where
$$V^+=\{x\mathrm{}:x^0>|x|\left[\mathrm{\Sigma }_{j=1}^3(x_j)^2\right]^{1/2}\}$$
is an open upper light cone in the M-space. Under the dilatation transformation $`x\alpha x`$ ($`\alpha >0`$) the TPWF (3.1) acquires the additional term
$`W_2(x)W_2(\alpha x)=W_2(x){\displaystyle \frac{1}{2(2\pi )^2}}\mathrm{ln}\alpha .`$ (3.3)
It could be interpreted as a spontaneous symmetry breaking of the dilatation invariance of (2.11). This is an important point in the special role of the field $`\overline{b}(x)`$.
In general, the $`\tau _h(x)`$-function (2.12) is a well-defined distribution in $`S^{}(\mathrm{}^{2h})`$ and obeys the $`h`$-ordered quadratic differential equation in the d-dimension of space-time
$$\mathrm{\Delta }^{2h}\tau _{h}^{}{}_{}{}^{d}(x)=\delta _d(x),h=2,3,\mathrm{},$$
$$\mathrm{\Delta }^2\frac{^2}{x_{}^{2}{}_{1}{}^{}}+\mathrm{}+\frac{^2}{x_{}^{2}{}_{m}{}^{}}\frac{^2}{x_{}^{2}{}_{m+1}{}^{}}\mathrm{}\frac{^2}{x_{}^{2}{}_{m+n}{}^{}};$$
in the following cases
i) even $`d`$ and $`hd/2`$
$`\tau _{h}^{}{}_{}{}^{d}(x)=(1)^{d/21}{\displaystyle \frac{\mathrm{exp}[(\pi /2)in]}{4^h(hd/2)!(h1)!}}\left[\stackrel{~}{\mu }^2P(x)+iϵ\right]^{d/2+h}`$ (3.4)
$`\mathrm{ln}\left[\stackrel{~}{\mu }^2P(x)+iϵ\right],`$ (3.5)
where $`P(x)=x_1^2+\mathrm{}+x_m^2x_{m+1}^2\mathrm{}x_{m+n}^2;`$
ii) even $`d`$ and $`h<d/2`$
$`\tau _{h}^{}{}_{}{}^{d}(x)=(1)^h{\displaystyle \frac{\mathrm{exp}[(\pi /2)in]}{4^h\pi ^{d/2}(h1)!}}\left[\stackrel{~}{\mu }^2P(x)+iϵ\right]^{d/2+h}\mathrm{\Gamma }\left({\displaystyle \frac{d}{2}}h\right).`$ (3.6)
In the case of $`d=2h`$-dimensions, expression (3.4) can be more simplified as follows:
$$\tau _{h}^{}{}_{}{}^{d=2h}(x)=(1)^{h1}\frac{\mathrm{exp}[(\pi /2)in]}{4^h(h1)!}\mathrm{ln}\left[\stackrel{~}{\mu }^2P(x)+iϵ\right].$$
For the following nearly realistic consideration of the confinement occurrence, we are interested in the case i) (3.4) which allows one to study highly singular objects for the confinement-like picture. In view of that, for $`h=2`$ in the M-space following Zwanziger the Fourier inversion $`F[\tau _{2}^{}{}_{}{}^{M}(x)]`$ of the distribution $`\tau _{2}^{}{}_{}{}^{M}(x)`$ looks like
$`F\left[\tau _{2}^{}{}_{}{}^{M}(x)\right]={\displaystyle \frac{\pi ^2}{2}}{\displaystyle \frac{}{p^\mu }}\left[{\displaystyle \frac{p^\mu \mathrm{ln}(l^2p^2iϵ)}{(p^2+iϵ)^2}}\right]`$ (3.7)
and obeying the equation
$$(p^2)^2F\left[\tau _{2}^{}{}_{}{}^{M}(x)\right]=1,p^2=(p^0)^2\mathrm{\Sigma }_{j=1}^3p_j^2,$$
where $`l=\mathrm{exp}(\gamma _E1/2)(2\stackrel{~}{\mu })^1`$ ($`\gamma _E`$ stands for the Euler constant). Note that in terms of weak derivatives the distribution (3.7) can be rewritten as
$`F\left[\tau _{h}^{}{}_{}{}^{d}(x)\right]=\underset{\kappa ^20}{lim}\left[{\displaystyle \frac{(1)^{d/2}}{(p^2\kappa ^2+iϵ)^{d/2}}}+i\pi ^{d/2}\delta _d(p)\mathrm{ln}\left({\displaystyle \frac{\kappa ^2}{4\stackrel{~}{\mu }^2}}\right)\right].`$ (3.8)
for any $`h`$ allowed in the dimension $`d`$. The remaining case ii) (3.6) leads to the following singular object:
$$F\left[\tau _{h}^{}{}_{}{}^{d}(x)\right]=\underset{\kappa 0}{lim}\stackrel{~}{\mu }^d\left(X^2+i\kappa \right)^h,$$
where $`\stackrel{~}{\mu }^2X^2=p_1^2+\mathrm{}+p_m^2p_{m+1}^2\mathrm{}p_{m+n}^2`$. In the space M one has
$$F\left[\tau _{h}^{}{}_{}{}^{M}(x)\right]=\underset{\kappa 0}{lim}\stackrel{~}{\mu }^4\left(\frac{\stackrel{~}{\mu }^2}{p^2+iϵ}\right)^2.$$
## 4 A dual field propagator
For calculation of the coefficients $`a_1`$ and $`a_2`$ in (2.15) one has to introduce invariant functions $`E_2(x)=E_2^{}(x)E_2^{}(x)`$ and $`D_2(x)=D_2^{}(x)D_2^{}(x)`$ (see also formulae (2.16) and (2.17)) which define the commutator written in the general form
$`[\overline{b}(x),\overline{b}(0)]=(2\pi )^2i\left[4a_1E_2(x)+a_2D_2(x)\right],`$ (4.1)
$$E_2(x)=(8\pi )^1sgn(x^0)\theta (x^2),$$
$$D_2(x)=(2\pi )^1sgn(x^0)\delta (x^2).$$
In the space $`S^{}(\mathrm{}^4)`$ on $`\mathrm{}^4`$ the propagator for the $`\overline{b}(x)`$-field looks like
$`\tau _2(x)=a_1\left[\mathrm{ln}|\stackrel{~}{\mu }^2x_\mu ^2|+i\pi \theta (x_\mu ^2)\right]+a_2\left[x_\mu ^2+i\pi \delta (x_\mu ^2)\right]+a_3.`$ (4.2)
The coefficients $`a_1`$ and $`a_2`$ in (4.2) can be fixed using the canonical commutation relations (CCR)
$$[C_\mu (x),\pi _{C_\nu }(0)]_{|_{x^0=0}}=ig_{\mu \nu }\delta ^3(\stackrel{}{x}),$$
$$[\overline{b}(x),\pi _{\overline{b}}(0)]_{|_{x^0=0}}=i\delta ^3(\stackrel{}{x}),$$
respectively. Here, the conjugate momenta $`\pi _{C_\mu }(x)`$ and $`\pi _{\overline{b}}(x)`$ look like
$$\pi _{\overline{b}}(x)=8\left[^0\overline{b}(x)+mC^0(x)\right],$$
$`\pi _{C_\mu }(x)={\displaystyle \frac{4}{3}}G_{0\mu }(x).`$ (4.3)
The direct calculation leads to (see also Appendix):
$$a_1=\frac{m^2}{12(2\pi )^2},$$
$$a_2=\frac{1}{6(2\pi )^2}.$$
Restricting the $`W_2(x)`$ function (2.15) to only the first term, one can obtain that $`\tau _2(x)`$ satisfies the equation
$`(\mathrm{\Delta }^2)^2\stackrel{~}{\tau }_2(x)=i\delta _4(x),`$ (4.4)
where $`\stackrel{~}{\tau }_2(x)=3(2/m)^2\tau _2(x)`$. The formal Fourier transformation in $`S^{}(\mathrm{}_4)`$ gives
$`\widehat{\tau }_2(p)=weak\underset{\stackrel{~}{\kappa }^2<<1}{lim}{\displaystyle \frac{i}{3(2\pi )^4}}\{m^2[{\displaystyle \frac{1}{(p^2\kappa ^2+iϵ)^2}}+i\pi ^2\mathrm{ln}{\displaystyle \frac{\kappa ^2}{\stackrel{~}{\mu }^2}}\delta _4(p)]`$ (4.5)
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{p^2\kappa ^2+iϵ}}\}.`$ (4.6)
Here, $`\kappa `$ is a parameter of representation and not the analogue of the infrared mass, $`\stackrel{~}{\kappa }^2\kappa ^2/p^2`$. To derive (4.5), we used the well-known mathematical trick with the generalized functions
$$weak\underset{\stackrel{~}{\kappa }^2<<1}{lim}d_{2h}p\mathrm{exp}(ipx)\frac{(1)^h}{(p^2\kappa ^2+iϵ)^h}=$$
$$=\frac{2i}{\mathrm{\Gamma }(h)(4\pi )^h}\{\mathrm{ln}2\gamma _E\mathrm{ln}(\kappa \sqrt{x_\mu ^2+iϵ})$$
$$+O[\kappa ^2x_\mu ^2,\kappa ^2x_\mu ^2\mathrm{ln}(\kappa \sqrt{x_\mu ^2+iϵ})]\}.$$
To go into the structure of the dual $`C_\mu `$-field propagator, we need to define the general form of the commutation relation $`[C_\mu (x),C_\nu (y)]`$. To do this procedure let us consider the canonical conjugate pair $`\{C_\mu ,\pi _{C_\nu }\}`$ (for the LD (2.6)) where $`\pi _{C_\mu }(x)`$ results from (4.3). The consequent CCR looks like
$`[{\displaystyle \frac{4}{3}}C_\mu (x),_\nu C_0(0)_0C_\nu (0)g_{0\nu }(C(0))+\mathrm{\Delta }_{0\nu }(0)]_{|_{x^0=0}}=ig_{\mu \nu }\delta ^3(\stackrel{}{x}),`$ (4.7)
where $`\mathrm{\Delta }_{\mu \nu }(x)=g_{\mu \nu }(C(x))\stackrel{~}{G}_{\mu \nu }(x)`$ tends to zero as $`x0`$ and the Dirac string tensor $`\stackrel{~}{G}_{\mu \nu }(x)`$ obeys the equation
$$\left(\mathrm{\Delta }^2+M^2\right)\stackrel{~}{G}_{\mu \nu }(x)=0,asx0.$$
Here, the equations of motion (2.8) and (2.10) were used, and $`M=3m`$. Obviously, the following form of the free $`C_\mu `$-field commutator (see Appendix)
$`[C_\mu (x),C_\nu (0)]=ig_{\mu \nu }\left[\xi m_1^2E_2(x)+cD_2(x)\right],`$ (4.8)
ensures the CCR (4.7) at large $`x_\mu ^2`$. Here, both $`\xi `$ and $`c`$ are real arbitrary numbers. The TPWF for $`C_\mu (x)`$ stands as (see also )
$`w_{\mu \nu }(x)=C_\mu (x)C_\nu (0)_0=ig_{\mu \nu }\left[\xi m_1^2E_2^{}(x)+cD_2^{}(x)+c_1F_2^{}(x)+c_2\right],`$ (4.9)
where $`F_2^{}(x)=\alpha (x)\alpha (0)_0=ix^2`$ is TPWF for the harmonic field $`\alpha (x)`$, $`c_1`$ and $`c_2`$ are arbitrary real numbers. In fact, the third and the fourth terms in (4.9) did not appear in (4.8) due to the properties of the $`F_2^{}(x)`$-function, namely, $`F_2(x)=F_2^{}(x)F_2^{}(x)=0`$ and the triviality, respectively. One can easily verify that (4.7) and (4.8) can be used to derive the following requirement for $`\xi `$
$`\xi ={\displaystyle \frac{3}{4}}4c.`$ (4.10)
The free dual gauge field propagator in $`S^{}(\mathrm{}_4)`$ in any local covariant gauge can be assumed as
$`\widehat{\tau }_{\mu \nu }(p)={\displaystyle d^4x\mathrm{exp}(ipx)\tau _{\mu \nu }(x)}`$ (4.11)
$`=i\left[g_{\mu \nu }\left(1{\displaystyle \frac{1}{\zeta }}\right){\displaystyle \frac{p_\mu p_\nu }{p^2+iϵ}}\right]\left[\xi m_1^2\widehat{t}_1(p)+c\widehat{t}_2(p)\right],`$ (4.12)
where
$$\tau _{\mu \nu }(x)=\frac{ig_{\mu \nu }}{(4\pi )^2}\left[\xi m_1^2\mathrm{ln}(\stackrel{~}{\mu }^2x_\mu ^2+iϵ)+\frac{c}{x_\mu ^2+iϵ}\right];$$
$$\widehat{t}_1(p)=weak\underset{\stackrel{~}{\kappa }^2<<1}{lim}\left[\frac{1}{(p^2\kappa ^2+iϵ)^2}+i\pi ^2\mathrm{ln}\left(\frac{\kappa ^2}{\stackrel{~}{\mu }^2}\right)\delta _4(p)\right];$$
$$\widehat{t}_2(p)=weak\underset{\stackrel{~}{\kappa }^2<<1}{lim}\frac{1}{2}\frac{1}{(p^2\kappa {}_{}{}^{2}+iϵ)}.$$
The gauge parameter $`\zeta `$ in (4.11) is a real number. According to Eq. (4.4), the following requirement on Green’s function $`\tau _{\mu \nu }(x)`$
$$(\mathrm{\Delta }^2)^2\stackrel{~}{\tau }_{\mu \nu }(x)=i\delta _4(x),$$
leads to that a constant $`c`$ has to be equal to zero where $`\stackrel{~}{\tau }_{\mu \nu }(x)=\tau _{\mu \nu }(x)/(\xi m_1^2)`$.
At the end of this section, let us consider the v.e.v.
$`V_\mu (x)=[J_\mu ^{mon}(x),\overline{b}(0)]_0,`$ (4.13)
where the monopole current expressed in terms of the $`\overline{b}`$-field and $`\stackrel{~}{G}_{\mu \nu }`$ Dirac tensor is
$$J_\mu ^{mon}(x)=\frac{2}{3}\left[\frac{3}{m}_\mu \mathrm{\Delta }^2\overline{b}(x)\alpha \mathrm{\Delta }^2^\nu \stackrel{~}{G}_{\mu \nu }(x)^\nu \stackrel{~}{G}_{\mu \nu }(x)\right].$$
Using the commutation relation (4.1) and the well-known relation between the following distributions like $`\mathrm{\Delta }^2sgn(x^0)\theta (x^2)=4sgn(x^0)\delta (x^2)`$, we arrive at the formal expression for $`V_\mu (x)`$
$`V_\mu (x)={\displaystyle \frac{im}{12\pi }}_\mu \left[sgn(x^0)\delta (x^2)\right]{\displaystyle \frac{2}{3}}(\alpha \mathrm{\Delta }^21)[^\nu \stackrel{~}{G}_{\mu \nu }(x),\overline{b}(0)]_0.`$ (4.14)
It is easily seen that the requirement of the Goldstone theorem for occurrence of a $`\delta (p_\mu ^2)`$ term in the Fourier transformation of $`V_\mu `$ is satisfied due to the presence of the first term in (4.14) which gives the Fourier transformed term $`p_\mu sgn(p^0)\delta (p^2)`$. The second term in (4.14) comes from the dual string effect.
## 5 The LDY-M solution for a gauge field
In this section, we consider the approximate topological solution for the dual gauge field in the LDY-M theory. The equations of motion are
$`^\nu G_{\mu \nu }=6ig[\varphi ^{}(_\mu igC_\mu )\varphi \varphi (_\mu +igC_\mu )\varphi ^{}];`$ (5.1)
$`(_\mu igC_\mu )^2\varphi ={\displaystyle \frac{2}{3}}\lambda (32B_0^225|\varphi |^27\varphi _3^2)\varphi .`$ (5.2)
At large distances Eq. (5.2) transforms into the following one:
$$(_\mu igC_\mu )^2\varphi =\frac{50}{3}\lambda (B_0^2|\varphi |^2)\varphi .$$
To get the solution, let us do the polar decomposition of the monopole field $`\varphi (x)`$ using new scalar variables $`\chi (x)`$ and $`f(x)`$
$$\varphi (x)=\frac{1}{\sqrt{2}}\mathrm{exp}(if(x))[\chi (x)+B_0].$$
Then, the equation of motion (5.1) transforms into the following one:
$`^\nu G_{\mu \nu }=6g(\chi +B_0)^2(gC_\mu _\mu f).`$ (5.3)
This means that the $`\overline{b}(x)`$-field is nothing but a mathematical realization of the ”massive” phase $`B_0f(x)`$ at large enough distancese:
$`\overline{b}(x)={\displaystyle \frac{B_0}{2}}S(x)f(x)+`$ (5.4)
$`m{\displaystyle [\mathrm{\hspace{0.17em}2}S(x)1]C_\mu (x)𝑑x^\mu }+{\displaystyle \frac{1}{12m}}{\displaystyle [\mathrm{\Delta }^2C_\mu (x)^\nu _\mu C_\nu (x)]𝑑x^\mu },`$ (5.5)
where $`S(x)(1+\chi (x)/B_0)^2`$. Now we have to define the flux as
$`\mathrm{\Pi }={\displaystyle G_{\mu \nu }(x)𝑑\sigma ^{\mu \nu }},`$ (5.6)
where $`\sigma ^{\mu \nu }`$ is the 2d surface element in the M-space. A substitution of $`C_\mu (x)`$ from (5.3) into (5.6) gives
$`\mathrm{\Pi }={\displaystyle \frac{\alpha }{2}}{\displaystyle _\mathrm{\Gamma }}S^1(x)^\nu \stackrel{~}{G}_{\mu \nu }(x)dx^\mu +{\displaystyle \stackrel{~}{G}_{\mu \nu }(x)𝑑\sigma ^{\mu \nu }}+{\displaystyle \frac{1}{g}}{\displaystyle _\mathrm{\Gamma }}_\mu f(x)dx^\mu ,`$ (5.7)
where $`\mathrm{\Gamma }`$ means the large closed loop where the current $`_\mu C_\nu _\nu C_\mu `$ is canceled. Integrating out over the loop $`\mathrm{\Gamma }`$ in the third term in (5.7) is nothing but the requirement that the phase $`f(x)`$ is varied by $`2\pi n`$ for any integer number $`n`$ associated with the topological charge inside the flux tube. One can present the fields $`C_\mu (x)`$ and $`\varphi (x)`$ in the cylindrical symmetry case using the radial coordinate $`r`$ (see also )
$$\stackrel{}{C}\frac{\stackrel{~}{C}(r)}{r}\stackrel{}{e}_\theta ,\varphi \varphi (r),$$
and $`f=2\pi n=n\theta `$ with $`\theta `$ being the azimuth around the z-axis. Thus, the field equation looks like
$`{\displaystyle \frac{d^2\stackrel{~}{C}(r)}{dr^2}}{\displaystyle \frac{1}{r}}{\displaystyle \frac{d\stackrel{~}{C}(r)}{dr}}3m^2[\mathrm{\hspace{0.17em}3}+2S(r)]\stackrel{~}{C}(r)+6nmB_0S(r)=0.`$ (5.8)
The following boundary conditions:
$$\stackrel{~}{C}(r)=\frac{4n}{7g},\chi (r)=(\sqrt{2}1)B_0,asr\mathrm{};$$
$$\stackrel{~}{C}(r)=0,\chi (r)=0,asr0$$
are conjugate with those in (2.4). At large enough $`r>>\mu ^1`$ $`((\mu =\sqrt{2\lambda }B_0)^1`$ defines the transverse dimension(s) of a monopole field around the tube) Eq. (5.8) transforms into the following one:
$$\frac{d^2\stackrel{~}{C}(r)}{dr^2}\frac{1}{r}\frac{d\stackrel{~}{C}(r)}{dr}+3g(4n7g\stackrel{~}{C})B_0^2=0$$
with the asymptotic transverse behaviour of its solution
$$\stackrel{~}{C}(r)\frac{4n}{7g}\sqrt{\frac{\pi mr}{2k}}e^{kmr}\left(\mathrm{\hspace{0.17em}1}+\frac{3}{8kmr}\right),k\sqrt{21}.$$
The field $`\stackrel{~}{C}(r)`$ grows rapidly when the radial distance from the center of the flux tube $`r<r_03fm`$ and approaches $`4n/7g`$ as soon as $`rr_0`$. Now it will be very instructive to clarify the singular properties of the field $`\overline{b}(x)`$ related to the phase $`f(x)`$ by means of (5.4). Since the phase is provided by the asimuth $`\theta `$ around the z-axis, $`\stackrel{}{}f=(n/r)\stackrel{}{e}_\theta `$. In the strong limit of the mass $`m`$ of the $`C_\mu `$ field (as well as the mass of the monopole field) at large distances
$$\stackrel{}{}\times \stackrel{}{}\overline{b}=2\pi nB_0\delta (x)\delta (y)\stackrel{}{e}_z,$$
where $`\stackrel{}{e}_z`$ is the unit vector along the axis $`z`$ and the $`\delta `$-functions stand for the center of the flux tube . Hence, the real singular character of the field $`\overline{b}(x)`$ is confirmed by its singular behaviour at the center of the flux tube for nonzero monopole condensate. At the end of this section we give the transverse distribution in $`r`$ of the color electric field $`E_z(r)`$ (see also )
$$\widehat{E}_c=\stackrel{}{}\times \stackrel{}{C}=\frac{1}{r}\frac{d\stackrel{~}{C}(r)}{dr}\stackrel{}{e}_zE_z(r)\stackrel{}{e}_z$$
which has the following profile in the flux tube at large $`r`$
$$E_z(r)=\sqrt{\frac{\pi m}{2kr}}\left(km\frac{1}{2r}\right)e^{kmr}.$$
## 6 Static potential
In this section, we intend to obtain the confinement potential in an analytic form for the system of interacting static test charges of quark and antiquark. Our statement is based on the dual character of the field $`C_\mu (x)`$ to a gluon field where $`C_\mu (x)`$ is just the interacting field provided by the monopole field $`\overline{b}(x)`$ and the divergence of $`\stackrel{~}{G}_{\mu \nu }(x)`$. We have found that $`\overline{b}(x)`$ plays the role of the dipole-type field at $`h=2`$ (see Eq. (2.11)). The mass of the dual field $`C_\mu (x)`$ is nonzero and equal to $`m=gB_0`$. It is assumed that the mass $`m`$ results from the Higgs-like mechanism when the dual field $`C_\mu (x)`$ interacts with the $`\varphi (x)`$ field, namely an octet of dual potentials $`C_\mu `$ coupled weakly with three octets of scalar fields $`\varphi _\alpha (x)(\alpha =1,2,3)`$.
According to the distribution (3.8), the first term in the expansion for the static potential
$$P_{stat}(r)=d_3\stackrel{}{p}e^{i\stackrel{}{p}\stackrel{}{r}}F\{\tau _{h}^{}{}_{}{}^{d}(x)\}_{|_{p^0=0}}$$
in $`\mathrm{}^3`$ is a rising function with $`r=|\stackrel{}{x}|`$
$`P_{stat}(r){\displaystyle \frac{1}{2^{2h}\pi ^{3/2}}}{\displaystyle \frac{1}{(h1)!}}\mathrm{\Gamma }(3/2h)r^{2h3}.`$ (6.1)
It is obvious that at $`h2`$ the distribution (6.1) increases with $`r`$ linearly ($`h=2`$) or faster ($`h>2`$). Let us represent the distribution $`r^\sigma `$ as the Taylor series around some regularization point $`\sigma _0`$
$`r^\sigma =\stackrel{~}{\mu }^\omega r^{\sigma _0}\left[1+\omega \mathrm{ln}(\stackrel{~}{\mu }r)+{\displaystyle \frac{1}{2}}\omega ^2\mathrm{ln}^2(\stackrel{~}{\mu }r)+\mathrm{}\right],`$ (6.2)
where $`\omega =\sigma \sigma _0`$ is an infinitesimal positive interval at $`\sigma d,d2,\mathrm{}`$ The potential (6.1) is simplified to
$$P_{stat}(r)\frac{1}{8\pi }r[1+\mathrm{ln}(\stackrel{~}{\mu }r)+\mathrm{}].$$
The Fourier transform $`F\{r^\sigma \}`$ of (6.2) into $`S^{}(\mathrm{}_d)`$ for the limit $`\omega 0`$ leads to the following singular distribution in the whole region of the existence of the analytic function $`r^\sigma `$ at $`\sigma d,d2,\mathrm{}`$:
$$F\{r^\sigma \}=\left(\frac{4\pi }{p^2}\right)^{(\sigma +d)/2}\pi ^{\sigma /2}\frac{\mathrm{\Gamma }[(\sigma +d)/2]}{\mathrm{\Gamma }(\sigma /2)}.$$
In general, the static potential is defined as
$`P_{stat}(r)=\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{T}}A(r),`$ (6.3)
where the action $`A(r)`$ is given by the colour source-current part of LD
$$L(p)=\stackrel{}{j}_\alpha ^\mu (p)\widehat{\tau }_{\mu \nu }(p)\stackrel{}{j}_\alpha ^\nu (p).$$
It is known that for such a system of heavy particles the sources are given by a c-number current
$$\stackrel{}{j}_\alpha ^\mu (x)=\stackrel{}{Q}_\alpha g^{\mu 0}[\delta _3(\stackrel{}{x}\stackrel{}{x_1})\delta _3(\stackrel{}{x}\stackrel{}{x_2})]$$
with $`\stackrel{}{Q}_\alpha =e\stackrel{}{\rho }_\alpha `$ being the Abelian color-electric charge of a quark while $`\stackrel{}{\rho }_\alpha `$ is the weight vector of the SU(3) algebra: $`\rho _1=(1/2,\sqrt{3}/6)`$, $`\rho _2=(1/2,\sqrt{3}/6)`$, $`\rho _3=(0,1/\sqrt{3})`$ . Here, $`\stackrel{}{x}_1`$ and $`\stackrel{}{x}_2`$ are the position vectors of a quark and an antiquark, respectively; the label $`\alpha `$=(1,2,3) corresponds to the color electric charge. The calculation of the potential (6.3) is most easily performed by taking into account the Fourier transformed quark current
$`\stackrel{}{j}_{\mu _\alpha }(p)=2\pi \stackrel{}{Q}_\alpha g_{\mu 0}\delta (p^0)\left(e^{i\stackrel{}{p}\stackrel{}{x}_1}e^{i\stackrel{}{p}\stackrel{}{x}_2}\right),`$ (6.4)
and by making use of the representation in the sense of generalized functions
$`weak\underset{\stackrel{~}{\kappa }^2<<1}{lim}\left[{\displaystyle \frac{1}{(p^2\kappa ^2+iϵ)^2}}+i\pi ^2\mathrm{ln}{\displaystyle \frac{\kappa ^2}{\stackrel{~}{\mu }^2}}\delta _4(p)\right]=`$ (6.5)
$`={\displaystyle \frac{1}{4}}{\displaystyle \frac{^2}{p^2}}{\displaystyle \frac{1}{p^2iϵ}}\mathrm{ln}{\displaystyle \frac{p^2iϵ}{\stackrel{~}{\mu }^2}}={\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{(p^2+iϵ)^2}}\left(53\mathrm{ln}{\displaystyle \frac{p^2iϵ}{\stackrel{~}{\mu }^2}}\right).`$ (6.6)
Due to the presence of the $`\delta (p^0)`$-function (see (6.4) in evaluating of the action $`A(r)`$, the remaining 3-dimensional integral over $`d_3p`$ will be easily calculated out using the instructive formula
$`{\displaystyle d_3\stackrel{}{p}e^{i\stackrel{}{p}\stackrel{}{x}}p^m\mathrm{ln}^\beta p}={\displaystyle \frac{1}{r^{m+3}}}{\displaystyle \underset{i=1}{\overset{\beta }{}}}{\displaystyle \frac{\mathrm{\Gamma }(\beta +1)(1)^i}{\mathrm{\Gamma }(\beta i+1)\mathrm{\Gamma }(i+1)}}\left({\displaystyle \frac{d}{dm}}\right)^{\beta i}H_m\mathrm{ln}^ir,`$ (6.7)
where
$$H_m=2^{m+3}\pi ^{3/2}\frac{\mathrm{\Gamma }[(m+3)/2]}{\mathrm{\Gamma }(m/2)},m3,5,\mathrm{},$$
$`r|\stackrel{}{x}|`$, $`p|\stackrel{}{p}|`$, $`\stackrel{}{x}\mathrm{}^3`$, $`\stackrel{}{p}\mathrm{}_3`$.
As a consequence of the dual field propagator (4.11) the static potential (6.3) at large distances looks like
$`P_{stat}(r)={\displaystyle \frac{\stackrel{}{Q}^2}{16\pi }}\left\{\xi m^2r[5+6(A+\mathrm{ln}\stackrel{~}{\mu }r)]+O\left({\displaystyle \frac{c}{r}}\right)\right\},`$ (6.8)
where $`A1+\sqrt{\pi }(1/2+2\sqrt{\pi }\gamma _E)(2\sqrt{\pi }+1)\mathrm{ln}2<0`$ and the last term in (6.8) is just the positive correction and not the analogue of the Coulomb part of the potential due to the one-gluon exchange. Neglecting the last term in (6.8) and taking into account formula (4.10), one can conclude
$`P_{stat}(r){\displaystyle \frac{3\stackrel{}{Q}^2}{64\pi }}m^2r(12.4+6\mathrm{ln}\stackrel{~}{\mu }r).`$ (6.9)
Hence, the string tension $`a`$ in $`P_{stat}(r)=ar`$ emerges as
$`a{\displaystyle \frac{\alpha _s}{16}}m^2\left(12.4+3\mathrm{ln}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m^2}}\right),\stackrel{~}{\mu }>9m,`$ (6.10)
where $`r`$ in the logarithmic function in (6.9) has been changed by the characteristic length $`r_c1/m`$ which determines the transverse dimension of the dual field concentration, while $`\stackrel{~}{\mu }`$ is associated with the ”coherent length” inverse and the dual field mass $`m`$ defines the ”penetration depth” in the type II superconductor. For a typical value of the electroweak scale $`\stackrel{~}{\mu }250GeV`$ we get $`a0.20GeV^2`$ for the mass of the dual $`C_\mu `$-field $`m=0.6GeV`$ and $`\alpha _s=e^2/(4\pi )`$=0.37 obtained from fitting the heavy quark-antiquark pair spectrum . The value of the string tension (6.10) is quite close to a phenomenological one (eg., coming from the Regge slope of the hadrons). Making the formal comparison of the result obtained here in the analytic form let us remind the analogue with the well-known expression of the energy per unit length of the vortex in the type II superconductor
$`ϵ_1={\displaystyle \frac{\varphi _{0}^{}{}_{}{}^{2}m_A^2}{32\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{m_\varphi }{m_A}}\right)^2,`$ (6.11)
where $`\varphi _0`$ is the magnetic flux of the vortex, $`m_A`$ and $`m_\varphi `$ are penetration depth and the coherent length inverse, respectively. On the other hand, the string tension in Nambu’s paper (see the first ref. in ) is given by
$`ϵ_2={\displaystyle \frac{g_m^2m_v^2}{8\pi }}\mathrm{ln}\left(1+{\displaystyle \frac{m_{s}^{}{}_{}{}^{2}}{m_{v}^{}{}_{}{}^{2}}}\right),`$ (6.12)
with $`m_s`$ and $`m_v`$ being the masses of scalar and vector fields and $`g_m`$ is a magnetic-type charge. It is clear from formula (6.8) that for a sufficiently long string $`r>>m^1`$ the $`r`$-behaviour of the static potential is dominant; for a short string $`r<<m^1`$ the singular interaction provided by the second term in (6.8) becomes important if the average size of the monopole is even smaller.
## 7 Summary and discussion
We studied the dual gauge model of the long-distance Yang-Mills theory in terms of two-point Wightman functions. The quantization of the model has been provided by using CCR, thus avoiding other methods (e.g., path integral use). We intended to give our understanding of the confinement by making use of nothing else but the well-known tools of quantum field theory based on LD given in as well as the renormalization model and symmetry properties. Among the physicists dealing with the models of interplay of a scalar (monopole, Higgs) field with a dual vector (gauge) boson field, where the vacuum state of the quantum Y-M theory is realized by a condensate of the monopole-antimonopole pairs, there is a strong belief that the flux-tube solution explains the scenarios of color confinement. Based on the flux-tube scheme approach of Abelian dominance and monopole condensation, we have obtained the analytic expressions for both the monopole and dual gauge boson field propagators (4.5) and (4.11), respectively, in $`S^{}(\mathrm{}_4)`$. These propagators lead to a consistent perturbative expansion of Green’s functions. However, the Fourier transformation of the first term in TPWF (2.15) gives the occurrence of the $`\delta ^{}(p_\mu ^2)`$-function. This is a consequence of the nonunitarity of the translations, and the spectral function with such a term gives an indefinite metric . We have found that the Goldstone theorem is valid in our model in the form taking into account the Dirac’s string effect (4.14). In fact, we obtained that the characteristic $`\delta (p_\mu ^2)`$ term naturally appears in the Fourier transformation of the v.e.v. of the commutator (4.13) of the monopole current and the scalar local field $`\overline{b}`$. In principle, a similar result has to be expected if one replaces the $`\overline{b}`$-field in (4.13) by any product of the gauge field $`C_\mu `$ and $`\overline{b}`$. We see that the fields $`b(x)`$ and $`b_3(x)`$ receive their masses and the $`\overline{b}(x)`$ field in combination with $`^\nu \stackrel{~}{G}_{\mu \nu }(x)`$ form the vector field $`C_\mu (x)`$ obeying the equation of motion for the massive vector field with the mass $`m=gB_0`$. The solution of the $`\overline{b}(x)`$-field can be identified as a ”ghost”-like particle in the substitute manner.
The occurrence of an arbitrary parameter $`l`$ in (2.15) and (2.16) leads to breaking their covariance under the dilatation transformation (3.3) and provides spontaneous symmetry breaking of the dilatation invariance of Eq. (2.11). The monopole condensation, formulated in the framework of LDY-M model, causes the strong and long-range interplay between heavy quark and antiquark, which gives the confining force, through the dual Higgs mechanism. We have obtained the analytic expression for the static potential (6.8) at large distances. The form of this potential grows linearly with the distance $`r`$ apart from logarithmic correction. The latter comes from the second term in the expression (6.5) (see also (6.7) ).
Making an analytic comparison of $`ϵ_1`$ (6.11) and $`ϵ_2`$ (6.12) with $`a`$ in (6.10), one can conclude that we have obtained a similar behaviour of the string tension $`a`$ to those in the magnetic flux picture of the vortex and in the Nambu scheme, respectively, as well as in the dual Ginzburg-Landau model .
Finally, it is to be noted that we have played the game with the choice of the gauge group where the Abelian group appears as a subgroup of the full Y-M gauge group. This is a very instructive method of calculating the confinement potential in the static limit in the analytic form. However, we understand that no real physics can depend on such a choice. Now, there is a next step in more formal consideration of the Y-M theory where it seems to be a new mechanism of confinement .
## 8 Acknowledgements
G.A.K. is grateful to G.M. Prosperi for the kind hospitality at the University of Milan where this work has partly been done.
## 9 Appendix
The following relations were used for calculation of the coefficients $`a_1`$ and $`a_2`$ in (4.1):
$$D_2(x^0=0,\stackrel{}{x})=0,_\mu D_2(x)_{|_{x^0=0}}=g_{o\mu }\delta _3(\stackrel{}{x});$$
$$(\mathrm{\Delta }^2)^2E_2(x)=_0^2E_2(x)_{|_{x^0=0}}=_0E_2(x)_{|_{x^0=0}}=E_2(0,\stackrel{}{x})=0;$$
$$_0^3E_2(x)_{|_{x^0=0}}=8\pi g_{0\mu }\delta _3(\stackrel{}{x}),\mathrm{\Delta }^2E_2(x)=D_2(x).$$
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# Gravitational and dilaton radiation from a relativistic membrane
## I Introduction
Now it is well understood that topological defects which could have been created at the early stages of cosmological expansion may leave their footprints in the spectrum of the background gravitational waves ViSh92 ; BaCaSh97 which can be detected in the future experiments. Until recently main attention was paid to the effects due to cosmic strings. Cosmic strings, presumably serving as seeds for the formation of galaxies, were thought to be produced at the phase transition of the grand unification scale Ki76 ; HiKi95 . Gravitational, dilaton and axion radiation emitted by oscillating strings was thoroughly studied in a number of papers (see e.g. Vi85 ; Bu85 ; VaVi85 ; GaVa87 ; DaVi97 and references therein). Although a recent analysis of the available astrophysical data seems to rule out the cosmic string scenario of the formation of structures in the Universe, these calculations may still be of interest in the context of the fundamental string/M theory.
Recently it was argued that the mass scale of the superstring/M theory may be as low as 10 TeV ArDiDv98 ; AnArDiDv98 at the expense of existence of some non-small extra dimensions. If true, this may drastically modify the full scenario of cosmological evolution. Based on this assumption various cosmological models were envisaged in which the observed Universe is treated as a three-brane embedded in a higher-dimensional world. According to this brane cosmology the standard model particles are confined to the brane while gravity lives in the full space. Such scenarios appeal to reconsider the problem of topological defects which should be associated rather with solitons of the fundamental string theory than with those of the field-theoretical models of grand unification.
Among the string theory solitons which may live on a three-brane there are one-dimensional strings and two-dimensional membranes. In the zero-thickness limit the dynamics of solitonic strings is governed by the Nambu-Goto action, while the membranes are described by a similar three-dimensional geometrical action (see e.g. CoTu76 ; Su83 ). Relativistic membrane in the four-dimensional spacetime geometrically are similar to domain walls encountered in the usual field-theoretical models. However their formation and dynamics in the string/brane cosmology are likely to be different from dynamics of the domain walls in the standard cosmological scenario. Here we do not discuss details of creation and evolution of membranes in the brane cosmology but concentrate on the problem of radiation.
Radiation produced by one-dimensional defects – the relativistic strings – is likely to be similar to radiation from the cosmic strings, apart from a different attribution of scalar and two-form fields involved. Radiation from an oscillating two-dimensional relativistic membrane seems not to have been studied in detail before, although some earlier estimates of gravitational radiation emitted by the cosmic domain walls are worth to be mentioned GlRo98 ; KoTuWa92 . Massless boson modes to which the membranes are coupled in the four-dimensional space-time include a dilaton, a graviton and a Ramond–Ramond three-form. The latter is non-dynamical in four dimensions, so oscillating membranes may emit dilatons and gravitons.
In the classical limit one can describe radiation from a vibrating membrane in a standard way using the retarded solutions of the corresponding wave equations. Unfortunately, contrary to the case of strings, a general solution of the (highly non-linear) membrane equation of motion can hardly be obtained in a closed form. To simplify the problem one can restrict oneself by some particular excitation mode. Here we study the dilaton and gravitational radiation emitted by a toroidal oscillating membrane excited along a homology cycle with the radius much smaller than the second radius of the torus. Under such an assumption the membrane equation of motion can be easily solved, and a closed analytic expression for the spectral power of radiation can be obtained. We analyse in detail the spectral-angular distributions with an emphasis on the effects due to two-dimensional nature of a membrane. Hopefully this simple model correctly describes basic features of the membrane radiation in more general case, although radiation due to more complicated excitation modes is certainly worth to be considered in the future.
The plan of the paper is as follows. In Sec. II we derive the linearized equations for the scalar and gravitational radiation from a membrane source. In Sec. III the solution of the equations of motion for a free toroidal membrane is obtained under an assumption of the relative smallness of the radius of an excited homology cycle. Sec. IV is devoted to an analysis of the dilaton radiation; we derive the spectral-angular distribution in terms of the Bessel’s functions and investigate the role of higher harmonics of the basic oscillation frequency. A simplified expression valid in the high-frequency range of the spectrum is obtained in terms of the MacDonald’s functions. In Sec. V the gravitational radiation is studied along similar lines. The results are summarized and discussed in Sec. VI. Some useful integrals involving Bessel’s functions used in the calculation are given in the Appendix.
## II Basic equations
Consider a relativistic membrane $`x^\mu =z^\mu (\sigma ^a),\mu =0,1,2,3,`$ propagating in the four-dimensional space-time with the metric $`g_{\mu \nu }`$ and interacting with the dilaton field $`\varphi `$ (our signature convention is $`(+)`$). The corresponding action can be written as a sum
$$S=S_m+S_f$$
(1)
of the purely geometrical membrane action CoTu76 ; Su83
$$S_m=\mu \mathrm{e}^{\alpha \varphi }\sqrt{\gamma }d^3\sigma ,$$
(2)
and the field term
$$S_f=\frac{1}{16\pi G}\left\{R+2_\mu \varphi _\nu \varphi g^{\mu \nu }\right\}\sqrt{g}d^4x.$$
(3)
Here $`\mu `$ is the membrane tension parameter of a dimension of the energy per unit area, $`\alpha `$ is the dilaton coupling constant. In our normalization the dilaton field $`\varphi `$ and the coupling constant $`\alpha `$ both are dimensionless. The membrane world-volume is parameterized by the internal coordinates $`\sigma ^a=(\tau ,\sigma ,\rho ),a=0,1,2`$. The induced metric on the world-volume reads
$$\gamma _{ab}=g_{\mu \nu }\frac{z^\mu }{\sigma ^a}\frac{z^\nu }{\sigma ^b},$$
(4)
where $`\gamma =det\gamma _{ab}`$.
Variation of the action (2) over the space-time coordinates $`z^\mu `$ gives the equation of motion
$$\frac{1}{\sqrt{\gamma }}_a\left\{\mathrm{e}^{\alpha \varphi }\sqrt{\gamma }\gamma ^{ab}_b\left(z^\nu g_{\mu \nu }\right)\right\}=0,$$
(5)
where $`\gamma ^{ab}`$ is the inverse metric on the world-volume. Up to the dilaton factor, this is a three-dimensional covariant D’Alembert equation. It is worth comparing it with the two-dimensional string equation of motion. In the latter case one has three gauge degrees of freedom corresponding to two diffeomorphisms and a Weyl transformation on the world-sheet. Consequently, three independent components of the world-sheet metric can be chosen flat, so that the equation of motion becomes linear. A general solution can be written as a sum of two arbitrary functions of null coordinates on the world-sheet: right and left movers. This is no more possible in the case of a membrane for which the number of gauge degrees of freedom is again three (diffeomorphisms of the world-volume), while the number of the independent world-volume metric components is six. Therefore the membrane equation of motion remains essentially non-linear in any gauge. To make it explicitly solvable one needs additional assumptions to be made about excitation modes.
Variation of the full action (1) with respect to the dilaton $`\varphi `$ gives the four-dimensional D’Alembert equation
$$\frac{1}{\sqrt{g}}_\mu \left(\sqrt{g}g^{\mu \nu }_\nu \varphi \right)=4\pi j(x)$$
(6)
with the source term
$$j(x)=G\mu \alpha \mathrm{e}^{\alpha \varphi }\frac{\delta ^4\left(xz(\sigma )\right)}{\sqrt{g}}\sqrt{\gamma }d^3\sigma .$$
(7)
Note that the effective coupling of the dilaton to the membrane is proportional to the dilaton exponential. In absence of the dilaton background this factor is of course negligible in the linearized limit relevant to the radiation problem.
Finally, varying the action (1) with respect to the space-time metric $`g_{\mu \nu }`$, one obtains the Einstein equations
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=8\pi G\left(T_{\mu \nu }+T_{\mu \nu }^{dil}\right),$$
(8)
with the source term consisting of the membrane stress–tensor
$$T^{\mu \nu }=\mu \mathrm{e}^{\alpha \varphi }\gamma ^{ab}_az^\mu _bz^\nu \frac{\delta ^4\left(xz(\sigma )\right)}{\sqrt{g}}\sqrt{\gamma }d^3\sigma ,$$
(9)
and the dilaton term. In what follows we neglect the interaction of the dilaton with gravity and omit the second term. The effective coupling of the membrane to gravity contains again the dilaton exponential factor which may be neglected in absence of the dilaton background.
The action (1) is invariant under the reparameterization of the world-volume $`\sigma ^a\sigma ^a^{}(\sigma ^b)`$, so three gauge conditions may be imposed on the world-volume metric $`\gamma _{ab}`$. We will treat the gravitational field perturbatively on the flat background
$$g_{\mu \nu }=\eta _{\mu \nu }+h_{\mu \nu },$$
(10)
where $`\eta _{\mu \nu }`$ is the Minkowski metric, and choose the parameterization of the world-volume ensuring diagonality of the metric $`\gamma _{ab}`$:
$`\dot{z}^\mu z_\mu ^{}=0,\dot{z}^\mu \overline{z}_\mu =0,z^\mu \overline{z}_\mu =0.`$ (11)
Here and below all the contractions are performed using the flat metric $`\eta _{\mu \nu }`$, and the following abbreviations are used $`\dot{z}^\mu =_\tau z^\mu ,z^\mu =_\sigma z^\mu ,\overline{z}^\mu =_\rho z^\mu `$. The induced metric then takes the form
$`\gamma _{ab}=\left(\begin{array}{ccc}\dot{z}^2& 0& 0\\ 0& z^2& 0\\ 0& 0& \overline{z}^2\end{array}\right),\gamma =\sqrt{\dot{z}^2z^2\overline{z}^2},`$ (15)
where $`\dot{z}^2=\dot{z}^\mu \dot{z}_\mu `$ etc. Note that $`\dot{z}^\mu `$ is a timelike vector, while $`z^\mu `$ and $`\overline{z}^\mu `$ are spacelike, so the signature of $`\gamma _{ab}`$ is $`+`$ (hence $`det\gamma _{ab}`$ is positive). In the diagonal gauge the flat-space membrane equation of motion (5) can be presented as the continuity equation (see, e.g., CoTu76 ):
$$\frac{P_\tau ^\mu }{\tau }+\frac{P_\sigma ^\mu }{\sigma }+\frac{P_\rho ^\mu }{\rho }=0,$$
(16)
where $`P_a^\mu `$ are generalized momenta obtained by varying the membrane action over $`x_\mu /\sigma ^a`$:
$`P_\tau ^\mu `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\gamma }}}\dot{z}^\mu \mu z^2\overline{z}^2,`$
$`P_\sigma ^\mu `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\gamma }}}z^\mu \dot{z}^2\overline{z}^2,`$ (17)
$`P_\rho ^\mu `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\gamma }}}\overline{z}^\mu \dot{z}^2z^2.`$
These equations are still highly non-linear. They can be solved, however, if one assumes certain symmetric configurations of the membrane.
In absence of the gravitational background field (10), the dilaton equation (6) can be considered in the Minkowski space-time. Denoting the flat wave operator by a box we have
$$\mathrm{}\varphi =4\pi j(x),$$
(18)
where the source term is
$$j(x)=\mu \alpha G\sqrt{\dot{z}^2z^2\overline{z}^2}\delta ^4\left(xz(\sigma )\right)d^3\sigma .$$
(19)
Finally, the linearized Einstein equations for the metric perturbations can be written in the Fock-de Donder gauge
$$\mathrm{}\psi ^{\mu \nu }=16\pi GT^{\mu \nu },$$
(20)
where as usual
$$\psi _{\mu \nu }=h_{\mu \nu }\frac{1}{2}h_\alpha ^\alpha \eta _{\mu \nu },_\mu \psi ^{\mu \nu }=0.$$
(21)
In the linearized theory the membrane stress-energy tensor simplifies to the following expression
$`T^{\mu \nu }=\mu {\displaystyle \left(\dot{z}^\mu \dot{z}^\nu z^2\overline{z}^2+z^\mu z^\nu \dot{z}^2\overline{z}^2+\overline{z}^\mu \overline{z}^\nu \dot{z}^2z^2\right)}`$
$`\times \delta ^4\left(xz(\sigma )\right)\left(\dot{z}^2z^2\overline{z}^2\right)^{1/2}d^3\sigma `$ . (22)
Note that the integral in the right hand side contains one integration more as compared with the case of a string and two integrations more compared with the point particle. The extended nature of the membrane source gives rise to the interference effects in the angular distribution of the radiation power.
## III Toroidal oscillating membrane
Our choice of the membrane model is motivated as follows. In view of complexity of the general situation, we intend to calculate radiation coming from the lowest excitation mode of the closed membrane. First the topology of the membrane 2-dimensional surface has to be specified. The simplest spherical membrane does not radiate in the main (spherical) excitation mode. So we choose a toroidal membrane in which case excitations naturally produce a non-vanishing quadrupole moment. We also assume that only one homology cycle of the torus is excited, this considerably simplifies the equations of motion.
Specifying $`z`$-direction of the coordinate system to be along the symmetry axis of the torus we arrive at the following parameterization of the membrane world-volume
$`z^0`$ $`=`$ $`\tau ,z^1=\left[R+r(\tau )\mathrm{cos}\rho \right]\mathrm{cos}\sigma ,`$
$`z^2`$ $`=`$ $`\left[R+r(\tau )\mathrm{cos}\rho \right]\mathrm{sin}\sigma ,z^3=r(\tau )\mathrm{sin}\rho ,`$ (23)
where $`R=`$const is the large radius and $`r(\tau )`$ is the variable small radius. Two parameters on the world volume
$$\sigma [0,2\pi ],\rho [0,2\pi ]$$
(24)
measure the angular distance along large and small homology cycles respectively, and $`\tau `$ is a time coordinate on the world-volume equal to the global time in the chosen gauge. This ansatz satisfies the equations of motion (16) only if at any moment $`\tau `$ one has
$$R^2r^2.$$
(25)
In this approximation the world-volume element is equal to
$$\sqrt{\gamma }=R\sqrt{r^2\left(1\dot{r}^2\right)},$$
(26)
substituting this into the equations of motion (16), one finds that under the assumption (25) the full system reduces to a single equation
$$\ddot{r}r+(1\dot{r}^2)=0.$$
(27)
This equation has the following solution
$$r=\mathrm{\Omega }^1\mathrm{sin}(\mathrm{\Omega }\tau ),$$
(28)
where $`\mathrm{\Omega }`$ is a constant, and it is assumed for simplicity that $`r(0)=0`$. The parameter $`\mathrm{\Omega }`$ is the oscillation frequency which is equal to the inverse radius of the smaller circle $`r_0`$ of the torus at the moments of its maximal extension.
In our approximation each loop $`\sigma =`$const belonging to the membrane moves independently. Actually the solution (28) satisfies exactly the equations of motion for an infinite cylindrical membrane. For a toroidal membrane this solution is only the main approximation valid if $`Rr_0`$.
Note that the function $`r`$ in (III) may take both positive and negative values. Indeed, two points on the loop $`\sigma =`$const which correspond to opposite ends of the diameter $`\rho =0,\pi `$ move to each other while $`r(\tau )`$ is decreased, merge when $`r=0`$, and then exchange their positions while $`r(\tau )`$ is negative. Therefore a half-period of $`r`$-oscillations corresponds to the full period of the true oscillations of the membrane. In other words, the period of the membrane motion is $`\pi /\mathrm{\Omega }`$ and not $`2\pi /\mathrm{\Omega }`$.
In our units the inverse angular frequency is equal to the small radius of the torus $`r_0=\mathrm{\Omega }^1`$ at the moments of its maximal extension. Therefore the quantity
$$v=\mathrm{\Omega }R=\frac{R}{r_0},$$
(29)
equal to the ratio of two radii of the toroidal membrane at the maximal extension, may be regarded as the shape parameter of the membrane.
## IV Dilaton radiation
The computation of the scalar radiation starting with the D’Alembert equation (18) is standard and amounts to the evaluation of the Fourier-transform of the current (19). Since the membrane motion is periodic with the period $`\pi /\mathrm{\Omega }`$, the radiation power can be presented as a sum over the even harmonics of the frequency $`\mathrm{\Omega }`$
$$P=\underset{n=1}{\overset{\mathrm{}}{}}P_n,\omega _n=2n\mathrm{\Omega }.$$
(30)
The angular distribution of radiation with the frequency $`\omega _n`$ is given by
$$\frac{dP_n}{d\mathrm{\Omega }}=\frac{\omega _n^2}{2\pi G}\left|j(𝐤,\omega _n)\right|^2$$
(31)
where $`𝐤`$ is the wave vector, $`𝐤^2=\omega _n^2`$, and the Fourier- -transform of the source current is defined as follows
$$j(𝐤,\omega _n)=\frac{\mathrm{\Omega }}{\pi }\underset{0}{\overset{\pi }{}}𝑑t\mathrm{e}^{i(\omega _nt\mathrm{𝐤𝐫})}j(𝐫,t)d^3x.$$
(32)
Parameterizing the wave vector $`𝐤`$ by the spherical angles $`\theta ,\phi `$:
$$𝐤=\omega _n(\mathrm{sin}\theta \mathrm{cos}\phi ,\mathrm{sin}\theta \mathrm{sin}\phi ,\mathrm{cos}\theta ),$$
(33)
and, using in (32) explicit expressions (19) and (III), one obtains after integration over the space-time variables $`t,𝐫`$:
$$j(𝐤,\omega _n)=\frac{\mu \alpha GR}{\pi \mathrm{\Omega }}\underset{0}{\overset{\pi }{}}𝑑\chi \underset{\pi }{\overset{\pi }{}}𝑑\rho \underset{0}{\overset{2\pi }{}}𝑑\sigma \mathrm{e}^{2ni\psi }\mathrm{sin}^2\chi ,$$
(34)
where
$$\psi =\left(\chi \mathrm{sin}\rho \mathrm{cos}\theta \mathrm{sin}\chi v\mathrm{cos}(\sigma \phi )\mathrm{sin}\theta \right).$$
(35)
Note that the azimuthal angle $`\phi `$ enters into this expression (as well as to all other formulas below) only in the combination $`\sigma \phi `$, and thus $`\phi `$ can be eliminated by the shift of the integration variable $`\sigma \sigma \phi `$. This means that all quantities are actually $`\phi `$-independent, as it could be expected in view of the axial symmetry of the membrane. Assuming without loss of generality $`\phi =0`$, one has
$$\psi =\left(\chi \mathrm{sin}\rho \mathrm{cos}\theta \mathrm{sin}\chi v\mathrm{cos}\sigma \mathrm{sin}\theta \right).$$
(36)
The integration in (35) can be performed in terms of the Bessel’s functions as follows. First one integrates over $`\chi `$ using the integral representation (73):
$`{\displaystyle \underset{0}{\overset{\pi }{}}}\mathrm{e}^{2ni\left(\chi z\mathrm{sin}\chi \right)}𝑑\chi ={\displaystyle \underset{0}{\overset{\pi }{}}}\mathrm{cos}\left[2n\left(\chi z\mathrm{sin}\chi \right)\right]𝑑\chi `$
$`+i{\displaystyle \underset{0}{\overset{\pi }{}}}\mathrm{sin}\left[2n\left(\chi z\mathrm{sin}\chi \right)\right]𝑑\chi =\pi J_{2n}(2nz)+iI,`$ (37)
where $`z=\mathrm{sin}\rho \mathrm{cos}\theta `$. It can shown that the contribution proportional to the imaginary part $`I`$ vanishes after integration over $`\rho `$. Indeed, a constant shift of the integration variable $`\chi \chi +\pi /2`$ leads to the following integral
$$\underset{\pi /2}{\overset{\pi /2}{}}𝑑\chi \underset{\pi }{\overset{\pi }{}}𝑑\rho \mathrm{sin}\left[2ni\left(\chi z\mathrm{cos}\chi \right)\right].$$
(38)
Here the integrand changes a sign under reflection of the domain of integration $`\chi \chi ,\rho \rho `$. Writing in (35) $`\mathrm{sin}^2\chi =(1\mathrm{cos}2\chi )/2`$ one can show that the $`iI`$ term gives no contribution in presence of $`\mathrm{cos}2\chi `$ as well.
The second integration over $`\rho `$ is performed using the formula (83) of the Appendix. Finally, integrating over $`\sigma `$ and using the recurrent relations (75) one obtains
$$j(𝐤,\omega _n)=2\pi ^2\mu \alpha G\frac{R}{\mathrm{\Omega }}J_0(y)\left(J_n^2(x)\mathrm{tan}^2\theta +J_{}^{}{}_{n}{}^{2}(x)\right),$$
(39)
where $`y=2nv\mathrm{sin}\theta `$ and $`x=n\mathrm{cos}\theta `$.
In this expression the Bessel’s functions $`J_n(x)`$ are typical for radiation from relativistic sources like the synchrotron radiation from a relativistic charge. Another Bessel’s function factor $`J_0(y)^2`$ has a different origin: it describes interference effects due to coherence of radiation emitted by different segments of the membrane. This factor reflects an extended nature of the membrane.
One can see that the angular distribution of the total power diverges in both directions along the symmetry axis $`\theta =0,\pi `$. For $`\mathrm{sin}\theta =0`$, the interference factor $`J_0(0)=1`$, so using an asymptotic formula (82) we obtain for $`n1`$:
$$\frac{dP_n(0)}{d\mathrm{\Omega }}=\frac{32}{3}\left(\frac{4}{3}\right)^{1/3}\mathrm{\Gamma }^4\left(\frac{1}{3}\right)G\mu ^2\alpha ^2R^2\pi ^3n^{2/3}.$$
(40)
The sum over $`n`$ diverges, but one can show that the solid angle into which the high $`n`$ radiation is peaked decreases with growing $`n`$, so that the total radiation power remains finite.
The angular distribution of the $`n`$-th harmonic is given by
$$\frac{dP_n}{d\mathrm{\Omega }}=8G\pi ^3(\mu \alpha nR)^2J_0^2(y)\left(\mathrm{tan}^2\theta J_n^2+J_{}^{}{}_{n}{}^{2}\right)^2,$$
(41)
where the argument of the Bessel’s functions $`J_n`$ and $`J_n^{}`$ is $`x=n\mathrm{cos}\theta `$. This is shown at Figs. 1,2 for $`n=1,2`$ and $`v=10`$. Angular distribution of the first harmonic has a maximum near the equatorial plane. Oscillations are due to the interference factor $`J_0^2(y)`$. Recall that our assumption (25) is equivalent to $`v1`$, so for non-small $`\theta `$ one can replace $`J_0(y)`$ by the asymptotic expression (80):
$$J_0(y)\sqrt{\frac{2}{2\pi vn\mathrm{sin}\theta }}\mathrm{cos}\left(2\pi vn\mathrm{sin}\theta \frac{\pi }{4}\right).$$
(42)
The distribution of the second harmonic exhibits maxima displaced from the equatorial plane to some intermediate region. It is modulated by the interference factor with doubled frequency.
Performing the angular integration one obtains the distribution over the harmonic number as shown at the Fig.3. One can see that about $`88\%`$ of total energy loss coresponds to the first, and about $`7\%`$ to the second harmonic, while the contribution from higher harmonics is relatively small.
The high frequency tail may be described using approximations for the Bessel’s functions which follow from the asymptotic formulas (74,77) valid for large $`n`$
$`J_n(na)`$ $``$ $`{\displaystyle \frac{\mathrm{sin}\theta }{\pi \sqrt{3}}}K_{1/3}\left({\displaystyle \frac{n}{3}}\mathrm{sin}^3\theta \right),`$
$`J_n^{}(na)`$ $``$ $`{\displaystyle \frac{\mathrm{sin}^2\theta }{\pi \sqrt{3}}}K_{2/3}\left({\displaystyle \frac{n}{3}}\mathrm{sin}^3\theta \right),`$ (43)
where $`K_{1/3}`$ and $`K_{2/3}`$ are MacDonald’s functions. Since $`K_\nu `$ decrease exponentially for large values of the argument (79), it is clear that the high-frequency part of the radiation is concentrated in a narrow angular regions around $`\theta =0,\pi `$
$$\mathrm{sin}\theta <n^{1/3}.$$
(44)
This beaming is a typical relativistic effect. The main contribution to the total power comes from the angular region in which the argument of the Bessel’s function $`J_0(y)`$ is large
$$y2\mathrm{\Omega }Rn^{1/3}1,$$
(45)
so the interference factor is rapidly oscillating. Therefore one can further simplify the angular distribution for higher harmonics using the asymptotic formula (42) and averaging over the interference ripples as follows
$$\overline{J_0^2(y)}\frac{1}{\pi y}.$$
(46)
As a result, the averaged over ripples angular distribution of radiation for $`n1`$ will be given by the formula
$$\frac{d\overline{P}_n}{d\mathrm{\Omega }}=\frac{4G\mu ^2\alpha ^2Rn}{9\pi ^2\mathrm{\Omega }}\mathrm{sin}^7\theta \left(K_{1/3}^2(z)+K_{2/3}^2(z)\right)^2,$$
(47)
where
$$z=\frac{n}{3}\mathrm{sin}^3\theta .$$
(48)
Note, that this representation remains qualitatively good even for relatively small $`n`$, see Fig.2.
In view of a rapid convergence of the integrals from the MacDonald’s functions, we can integrate this averaged distribution over angles. Passing to the argument of the MacDonald’s functions as the new integration variable and using the relation
$$\mathrm{sin}^8\theta d\theta =\frac{9}{n^3}z^2dz,$$
(49)
where an additional factor $`\mathrm{sin}\theta `$ comes from $`d\mathrm{\Omega }`$, one obtains for $`n1`$
$$\overline{P}_n=\frac{32I_1G\mu ^2\alpha ^2R}{\pi \mathrm{\Omega }}\frac{1}{n^2},$$
(50)
where
$$I_1=\left(K_{1/3}^2+K_{2/3}^2\right)^2𝑑z=4.33,$$
(51)
Therefore the contribution from high harmonics falls off as $`n^2`$.
One can also investigate the dependence of the radiated power over the shape parameter $`v`$ (the ratio of two radii of the torus). Numerical probes for different values of $`v`$ show a rather weak dependence of the total radiation power on this shape factor (see Fig.4)
## V Gravitational Radiation
Gravitational radiation power can be computed along similar lines. In addition one has to distinguish between the polarisation states. The angular distribution of radiation at the frequency $`\omega _n`$ can be presented as a sum over two circular polarisations
$$\frac{dP_n}{d\mathrm{\Omega }}=\frac{G\omega _n^2}{\pi }\underset{\pm }{}|T^\pm (𝐤,\omega _n)|^2,$$
(52)
where
$$T^\pm (𝐤,\omega _n)=\epsilon _{\mu \nu }^\pm T^{\mu \nu }(𝐤,\omega _n),$$
(53)
$$T^{\mu \nu }(𝐤,\omega _n)=\frac{\mathrm{\Omega }}{\pi }\underset{0}{\overset{\pi }{}}𝑑t\mathrm{e}^{i(\omega _nt\mathrm{𝐤𝐫})}T^{\mu \nu }(𝐫,t)d^3x.$$
(54)
Here the standard chiral graviton projectors are used
$`\epsilon _{\mu \nu }^\pm `$ $`=`$ $`\epsilon _\mu ^\pm \epsilon _\nu ^\pm ,`$
$`\epsilon _\mu ^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(e_\mu ^{(1)}\pm ie_\mu ^{(2)}),`$ (55)
where $`e_\mu ^{(1,2)}`$ are two linearly independent space-like unit four-vectors orthogonal to the wave four-vector $`k^\mu =(\omega _n,𝐤)`$ and its spatially reflected dual $`\overline{k}^\mu =(\omega _n,𝐤)`$,
$`e_\mu ^{(1,2)}k^\mu `$ $`=`$ $`e_\mu ^{(1,2)}\overline{k}^\mu =e_\mu ^{(1)}e^{(2)\mu }=0,`$
$`e^{(1)\mu }e_\mu ^{(1)}`$ $`=`$ $`e^{(2)\mu }e_\mu ^{(2)}=1.`$ (56)
Alternatively one could choose the linear polarisation states $`e_{\mu \nu }^{},e_{\mu \nu }^{}`$ corresponding to real and imaginary parts of (V):
$$\epsilon _{\mu \nu }^\pm =\frac{1}{\sqrt{2}}\left(e_{\mu \nu }^{}\pm ie_{\mu \nu }^{}\right),$$
(57)
or, explicitly,
$`e_{\mu \nu }^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(e_\mu ^{(1)}e_\nu ^{(1)}e_\mu ^{(2)}e_\nu ^{(2)}\right),`$
$`e_{\mu \nu }^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(e_\mu ^{(1)}e_\nu ^{(2)}+e_\mu ^{(2)}e_\nu ^{(1)}\right).`$ (58)
In the Lorentz frame in which $`e_0^{(1,2)}=0`$ two unit vectors orthogonal to $`𝐤`$ can be chosen along $`\theta `$ and $`\phi `$ directions:
$`𝐞^{(1)}`$ $`=`$ $`(\mathrm{cos}\theta \mathrm{cos}\phi ,\mathrm{cos}\theta \mathrm{sin}\phi ,\mathrm{sin}\theta ),`$
$`𝐞^{(2)}`$ $`=`$ $`(\mathrm{sin}\phi ,\mathrm{cos}\phi ,0).`$ (59)
As in the dilaton case, the chiral projections of the Fourier-transform of the stress tensor depend only on the difference $`\sigma \phi `$, so without loss of generality we can set $`\phi =0`$. After some simplifications one finds
$$T^\pm (𝐤,\omega _n)=\frac{\mu R}{2\pi \mathrm{\Omega }}\underset{0}{\overset{\pi }{}}𝑑\chi \underset{\pi }{\overset{\pi }{}}𝑑\rho \underset{0}{\overset{2\pi }{}}𝑑\sigma \mathrm{e}^{2ni\psi }\mathrm{\Phi }^\pm ,$$
(60)
where
$$\mathrm{\Phi }^\pm =\left(\mathrm{cos}\sigma \mathrm{cos}\theta \mathrm{cos}\rho \mathrm{sin}\theta \mathrm{sin}\rho \pm i\mathrm{sin}\sigma \mathrm{cos}\rho \right)^2,$$
(61)
and $`\psi `$ is given again by the Eq.(36).
After the integration over $`\sigma ,\rho `$ both chiral amplitudes lead to the same result. Indeed, purely imaginary terms in $`\mathrm{\Phi }^\pm `$ are proportional either to $`\mathrm{sin}2\sigma `$ or to $`\mathrm{sin}2\rho `$; taking into account an explicit form of the phase factor $`\psi `$, one finds that both these terms do not contribute to the integral because of the antisymmetry of the integrand in appropriate variables. In view of the relation (57) between chiral and linear polarisation projectors this means that the gravitational radiation in any direction has only one linear polarisation component $`e_{\mu \nu }^{}`$. By similar reasoning one can show that the product of the first two terms in (61) also vanishes after an integration over $`\rho `$, so $`\mathrm{\Phi }^\pm `$ can be replaced by the sum of the squares of each of the three terms. After some rearrangements, we find the following equivalent expression suitable for further integration
$$\mathrm{\Phi }_1=\frac{\mathrm{sin}^2\theta }{4}(13\mathrm{cos}2\rho )+\frac{1+\mathrm{cos}^2\theta }{4}\mathrm{cos}2\sigma (1+\mathrm{cos}2\rho ).$$
(62)
Note that, contrary to the dilaton case, this function is independent on $`\chi `$, (i.e. on time). Therefore the integration over $`\chi `$ gives again the expression (IV), and the imaginary part of the integral will vanish after the integration over $`\rho `$; this remains true for terms in (62) containing $`\mathrm{cos}2\rho `$. The subsequent integration over $`\rho `$ is performed using the integral (83). The integration over $`\sigma `$ directly amounts to the integral representation for the Bessel’s functions (79), and finally one obtains
$`T^+(𝐤,\omega _n)`$ $`={\displaystyle \frac{\pi ^2\mu R}{2\mathrm{\Omega }}}[(1+\mathrm{cos}^2\theta )J_2(y)(J_{}^{}{}_{n}{}^{2}\mathrm{tan}^2\theta J_n^2)`$
$`+\mathrm{sin}^2\theta `$ $`J_0(y)(J_n^2(1+3/\mathrm{cos}^2\theta )3J_{}^{}{}_{n}{}^{2})],`$ (63)
where the argument of the Bessel’s functions of the order $`n,n\pm 1`$ is $`x=n\mathrm{cos}\theta `$. Now an interference of radiation produced by different parts of the membrane is accounted by two different Bessel’s functions $`J_0(y)`$ and $`J_2(y)`$; apparently this is related to the graviton spin. Since $`J_2(0)=0`$, the first term in (V) vanishes at $`\theta =0,\pi `$, so does the second term containing an explicit factor $`\mathrm{sin}^2\theta `$. Thus, contrary to the dilaton case, the amplitude of the gravitational radiation strictly along the symmetry axis is zero. Still the main part of radiation is concentrated near this direction, see Figs. 5,6.
The maximal contribution to the total radiation power comes from the first harmonic, but its relative contribution is smaller than in the dilaton case (Fig. 7). The high frequency tail can be described in terms of the MacDonald’s functions. For large $`y`$ corresponding to the main contribution to the totat power one has $`J_0J_2`$ (see (80), this can be used before averaging over the ripples. An averaged distribution in the high frequency region is then obtained as follows
$$\frac{d\overline{P}_n}{d\mathrm{\Omega }}=\frac{4\mu ^2Rn}{9\pi ^2\mathrm{\Omega }}\mathrm{sin}^7\theta \left(3K_{1/3}^2(z)K_{2/3}^2(z)\right)^2.$$
(64)
After the angular integration one gets
$$\overline{P}_n=\frac{8\mu ^2RI_2}{\pi \mathrm{\Omega }}\frac{1}{n^2},$$
(65)
where
$$I_2=\left(3K_{1/3}^2K_{2/3}^2\right)^2𝑑z=1.24.$$
(66)
This is similar to the previous result for the dilaton (47).
Numerical integration for different $`v`$ shows that the total radiation power depends on the ratio of the torus radii even weaker than in the dilaton case (Fig. 8).
## VI Concluding remarks
Let us summarize our results. We have shown that the dilaton and gravitational radiation from a vibrating relativistic membrane exhibits features typical for extended relativistic sources: the radiation power contains high harmonics of the basic frequency beamed along the symmetry axis of the torus and its angular distribution is substantially modulated due to interference. The relative contribution of the high frequency tail depends on the spin and it is more pronounced in the gravitational case. For high $`n`$ the averaged over ripples radiation power falls off as $`\overline{P}_n1/n^2`$ both in the dilaton and gravitational cases. The high frequency part of radiation can be conveniently described in terms of the MacDonald’s functions like in the case of ultrarelativistic point particles.
The total power of the radiation from a toroidal membrane can be presented in the form
$$P=G\mu ^2Af(v),$$
(67)
where $`ARr_0`$ is the membrane area at the maximal extension, and the shape factor $`f(v)`$ is a function of the ratio of two radii of the torus (actually we observed rather weak variation of $`f(v)`$). On dimensional grounds it can be expected that this formula will be valid for other excitation modes up to some factor, which we set to unity $`f(v)1`$ for a qualitative estimate. Then in terms of the total mass of the membrane $`M=\mu A`$ we get
$$P\frac{GM^2}{A},$$
(68)
so the radiative lifetime may be estimated as
$$\tau _{rad}=\frac{M}{P}\frac{A}{GM}.$$
(69)
This quantity is proportional to the ratio of the area at the maximal extension to the gravitational radius of the membrane as a whole (an inverse speed of light factor has to be inserted in the ordinary units). In our model the membrane area at the moments of a maximal extension is $`4\pi ^2R/\mathrm{\Omega }`$ (recall that $`\mathrm{\Omega }`$ is equal to an inverse radius of the small circle of the torus, and we assumed that $`R1/\mathrm{\Omega }`$), therefore the quantity $`\mu R/\mathrm{\Omega }`$ is proportional to the total mass. Thus the membrane lifetime
$$\tau \frac{\mu R/\mathrm{\Omega }}{P}\frac{1}{G\mu }$$
(70)
depends only on the membrane tension. In contrary, the lifetime the radiating string loop is proportional to the loop size.
In the string case the radiative power in terms of a total mass has the order
$$P^{st}\frac{GM^2}{L^2},$$
(71)
where $`L`$ is the string length. Correspondingly, the radiative lifetime is
$$\tau _{rad}^{st}=\frac{M}{P^{st}}\frac{L^2}{GM}.$$
(72)
Actually our toroidal membrane has a shape of a ’thick’ string, following this similarity one can identify $`LR,ARr_0`$. Then it is easy to see that the membrane lifetime is $`R/r_01`$ times shorter than that of a string of an equal length.
It can be expected that gravitational radiation from membranes will depend on the membrane topology. For a toroidal membrane the radiation coming from the main excitation mode (as we have assumed here) in non-zero, while in the case of the spherical membrane the main sprerically symmetric excitation mode does not contribute to the gravitational radiation at all.
This work was supported in part by the RFBR grant 00-02-16306.
## VII Appendix
Evaluating the integral in the representation of the Bessel’s functions
$$J_n(y)=\frac{1}{\pi }\underset{0}{\overset{\pi }{}}\mathrm{cos}(n\sigma y\mathrm{sin}\sigma )𝑑\sigma $$
(73)
for $`y=nx`$ and $`n1`$ in the stationary phase approximation one finds
$$J_n(nx)\frac{\sqrt{2(1x)}}{\pi \sqrt{3}}K_{1/3}\left(\frac{2\sqrt{2}}{3}n(1x)^{3/2}\right),$$
(74)
where $`K`$ is the MacDonald’s function. For the Bessel’s functions one has the following reccurent relations
$$J_{n\pm 1}(nx)=\frac{1}{x}J_nJ_n^{},$$
(75)
while for the MacDonald’s functions with non-integer $`\nu `$
$$zK_\nu ^{}+\nu K_\nu =zK_{1\nu },$$
(76)
(note that $`K_\nu =K_\nu `$). Differentiating (74) and using (76) one obtains in the leading approximation for large $`n`$ the corresponding formula for the derivative of the Bessel’s function
$$J_n^{}(nx)\frac{2(1x)}{\pi \sqrt{3}}K_{2/3}\left(\frac{2\sqrt{2}}{3}n(1x)^{3/2}\right).$$
(77)
For small arguments the MacDonald’s function behaves as
$$K_\nu (z)\frac{\pi }{2\mathrm{sin}\pi \nu }\mathrm{\Gamma }^1(1\nu )\left(\frac{2}{z}\right)^\nu ,$$
(78)
while for large ones
$$K_\nu \sqrt{\frac{\pi }{2z}}\mathrm{e}^z.$$
(79)
For Bessel’s functions at large $`y`$ one has
$$J_p(y)\sqrt{\frac{2}{\pi y}}\mathrm{cos}\left(y\frac{p\pi }{2}\frac{\pi }{4}\right)$$
(80)
Passing in the Eqs.(74,77) to the limit $`x1`$, in which case the asymptotic behavior (78) holds, one obtains for $`n1`$
$$J_n(n)3^{2/3}\mathrm{\Gamma }^1(2/3)\left(\frac{2}{n}\right)^{1/3},$$
(81)
$$J_{}^{}{}_{n}{}^{}(n)3^{1/3}\mathrm{\Gamma }^1(1/3)\left(\frac{2}{n}\right)^{2/3}.$$
(82)
In the main text we have used the following integral:
$$\underset{0}{\overset{\pi }{}}\mathrm{cos}2\mu \rho J_{2\nu }(2a\mathrm{sin}\rho )𝑑\rho =\pi \mathrm{cos}\pi \mu J_{\nu \mu }(a)J_{\nu +\mu }(a).$$
(83)
In particular, in view of the Eq.(75) one has
$$\underset{0}{\overset{\pi }{}}J_{2n}(2nx\mathrm{sin}\rho )\mathrm{cos}2\rho d\rho =\pi \left(J_{}^{}{}_{n}{}^{2}(nx)\frac{1}{x^2}J_n^2(nx)\right).$$
(84)
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# 1 Introduction
## 1 Introduction
Employment of quantum states as the information carriers opened new prospects for the realization of various cryptographic protocols \[1–5\]. An important result in this field was achieved when it was shown that within the framework of the non-relativistic quantum mechanics some basic classical cryptographic protocols (whose security in classical physics is based on the computational complexity only) become unconditionally secure, i.e. their security is based on the laws of nature (non-relativistic quantum mechanics) and cease to depend on the available computational resources. However, for a number of basic tasks, e.g. bit commitment, all the numerous efforts to develop an unconditionally secure quantum protocol failed \[6–9\]. Moreover, it was shown later that no ideal unconditionally secure bit commitment protocol can be constructed in the framework of the non-relativistic quantum mechanics \[10-11\].
The only information carriers suitable for the realistic large-scale quantum cryptographic systems are the photons which are essentially relativistic particles. Since the relativistic quantum mechanics does not allow any sensible physical interpretation, the theory describing the relativistic quantum systems arises as already the quantum field theory . The relativistic quantum field theory imposes additional fundamental restrictions on the speed of information transfer and the processes of the quantum field state measurements. The latter circumstance has the consequence that the secret protocols which cannot be realized in the non-relativistic quantum mechanics become realizable in the classical relativistic theory and in the quantum field theory \[13–15\].
Below we shall show that in the relativistic quantum field theory the reliable (with unit probability) distinguishability of two orthogonal states requires an infinite time (which actually is a consequence of the microcausality principle), in contrast to the non-relativistic quantum mechanics where this can be done instantly at any time. This circumstance allows to construct new cryptographic protocols in the quantum field theory.
All quantum cryptographic protocols actually employ the following two features of quantum theory. The first one is the no cloning theorem , i.e. the impossibility of copying of an arbitrary quantum state which is not known beforehand or, in other words, the impossibility of the following process:
$$|A|\psi U(|A|\psi )=|B_\psi |\psi |\psi ,$$
where $`|A`$ and $`|B_\psi `$ are the apparatus states before and copying act, respectively, and $`U`$ is a unitary operator. Such a process is prohibited by the linearity and unitary nature of quantum evolution. Actually, even a weaker process of obtaining any information about one of the two non-orthogonal states without disturbing it is impossible, i.e. the final states of the apparatus $`|A_{\psi _1}`$ and $`|A_{\psi _2}`$ corresponding to the initial input states $`|\psi _1`$ and $`|\psi _2`$, respectively, after the unitary evolution $`U`$,
$$|A|\psi _1U(|A|\psi _1)=|A_{\psi _1}|\psi _1,$$
$$|A|\psi _2U(|A|\psi _2)=|A_{\psi _2}|\psi _2,$$
can only be different, $`|A_{\psi _1}|A_{\psi _2}`$, if $`\psi _1|\psi _2=0`$ , which means the impossibility of reliable distinguishing between non-orthogonal states. There is no such a restriction for orthogonal states. Moreover, in the non-relativistic quantum mechanics generally there is no restriction on the instant (arbitrarily fast) reliable distinguishability of any two orthogonal states at any time without disturbing them.
That is why the orthogonal states are not even discussed in the context of non-relativistic quantum cryptographic protocols.
On the contrary, in the quantum field theory it turns out that the distinguishability of two orthogonal states with the probability arbitrarily close to unity requires finite time, and the reliable distinguishability (with unit probability) can only be achieved in infinite time.
To clarify the difference between the non-relativistic quantum mechanics and the relativistic quantum field theory, consider first the non-relativistic case.
Suppose we have a pair of orthogonal states in the Hilbert space: $`|\psi _{1,2}`$ and $`\psi _1|\psi _2=0`$. To reliably distinguish these states, one can for example use the following orthogonal identity resolution in $``$:
$$𝒫_1+𝒫_2+𝒫_{}=I,𝒫_{1,2}=|\psi _{1,2}\psi _{1,2}|,𝒫_{}=I𝒫_1𝒫_2,$$
(1)
where $`𝒫_{1,2}`$ are the projectors on the subspaces $`_{1,2}`$ spanned by the states $`|\psi _1`$ and $`|\psi _2`$, while $`𝒫_{}`$ is the projector on $`_{12}^{}=(_1_2)^{}`$. The probability of obtaining, for example, the result in the channel $`𝒫_1`$ if in input state was $`|\psi _1`$, is
$$\text{Pr}_1\{|\psi _1\}=\text{Tr}\{|\psi _1\psi _1|𝒫_1\}1,$$
(2)
while it is identically equal to zero in the channels $`𝒫_{2,}`$:
$$\text{Pr}_{2,}\{|\psi _1\}=\text{Tr}\{|\psi _1\psi _1|𝒫_{2,}\}0,$$
(3)
and similarly for the input state $`|\psi _2`$. The equations (1–3) mean that the orthogonal states can be reliably (with unit probability) distinguished. To answer the question of whether or not this can be done instantly and without disturbing the measured states, one should consider the measurement procedure in more detail. Indeed, there are three natural levels of description of the measurement process in quantum mechanics which differ in the amount of information they provide \[17–19\]. The simplest description of the measurement procedure lists only the possible measurement outcomes (i.e., specifies the space of possible measurement outcomes) and specifies the relative frequencies (probabilities) of occurrence of a particular outcome for a given input state of the measured quantum system. In that sense the measurements are in one-to-one correspondence with the positive identity resolutions in the Hilbert space of the system states \[17–19\]. However, this approach completely ignores the problem of finding the state of the system after the measurement which gave a particular result. At the next level of description of quantum mechanical measurement each measurement procedure is associated with the so-called instrument $`T`$ which actually is, speaking informally, a rule allowing to ascribe to each pair consisting of the input density matrix $`\rho _0`$ of the measured system and the measurement result $`i`$ a positive matrix $`\stackrel{~}{\rho }_i`$ with the trace $`\mathrm{Tr}\stackrel{~}{\rho }_\mathrm{i}=\mathrm{p}_\mathrm{i}<1`$, which is interpreted as an unnormalized density matrix of the ensemble of the systems selected by the condition that the performed measurement procedure gave the outcome $`i`$ (the probability of obtaining the $`i`$-th outcome being $`p_i`$ and the state of the system after the outcome $`i`$ is obtained being $`\rho _i=\stackrel{~}{\rho }_i/p_i`$). In this approach the von Neumann – Lüders projection postulate is equivalent to the statement that the instrument $`T`$ corresponding to the measurement described by Eqs.(1–3) is given by the formula
$$\stackrel{~}{\rho }_i=𝒫_i\rho 𝒫_i,$$
so that for example for the input state $`|\psi _1`$ we have
$$\stackrel{~}{\rho }_1=\frac{𝒫_1|\psi _1\psi _1|𝒫_1}{\text{Tr}\{|\psi _1\psi _1|𝒫_1\}}=|\psi _1\psi _1|.$$
(4)
Equation (4) means that under the assumption of the possibility of the realization of von Neumann measurement satisfying the projection postulate the orthogonal states can be reliably distinguished without any disturbance. One should also note that the reliable distinguishing of orthogonal states without any disturbance can be performed at arbitrarily chosen moment of time, i.e. the measurement procedure can be started at any moment although the procedure itself generally depends on the chosen starting moment. Indeed, since the states $`|\psi _1`$ and $`|\psi _2`$ are assumed to be known and one should only determine which of these two states is actually given for a particular measurement run, the temporal evolution of these states ($`|\psi _i(t)=U(t)|\psi _i`$, where $`|\psi _i`$ is the system state at $`t=0`$) should also be considered as a known function of time. If the measurement is performed at time $`t`$, the identity resolution should employ $`𝒫_i(t)=U(t)𝒫_iU^1(t)`$ as the corresponding projectors.
An important point is that the duration of the measurement procedure $`\tau `$ has not yet been actually mentioned. To introduce $`\tau `$, one should consider the third, most detailed, level description of the measurement . In the quantum mechanical measurement theory it is shown that any instrument can be represented in the following way: first the measured quantum system $`S`$ begins to interact with an auxiliary system $`A`$ (ancilla) and they perform a joint evolution during time $`\tau `$. Then the system $`A`$ is subjected to a von Neumann measurement (assumed to be completed arbitrarily fast) whose outcome determines the state of system $`S`$ immediately after the measurement. Therefore, the speed with which a particular measurement procedure can actually be realized depends on the possibility of choosing a suitable auxiliary system $`A`$ and the realization of required interaction between $`A`$ and $`S`$.
## 2 Non-relativistic case
Since the statement of the possibility of reliable distinguishability at any time of the two orthogonal states is quite general, we shall make the transition to the relativistic case smoother by considering in the non-relativistic case two states $`|\psi _1`$ and $`|\psi _2`$ of a one-dimensional free particle.
It is intuitively clear that to identify a state one should have access to the entire spatial domain where the wavefunction of the particle is present. In the non-relativistic case there are no restrictions on the maximum speed of the information transfer so that an entire extended spatial domain can be accessed instantly at any time, i.e. there are no restrictions on instant non-local measurements.
It is natural to assume that the identification of the state of a quantum field also requires access to the entire domain of the Minkowski space-time where the field is present. However, the existence of the maximum speed of information transfer no extended domain can be accessed instantly, i.e. non-local any measurement requires a finite time.
In the non-relativistic quantum mechanics the above considerations are obvious because the states of a non-relativistic particle can be described by a wavefunction. In the relativistic quantum field theory these intuitively appealing considerations should be analyzed more rigorously, since the states of a quantum filed cannot be described by a wavefunction (or even by an operator field).
The set of observables of a quantum system is represented by an operator algebra with a unit and a conjugation (involution) which are defined on a dense subspace $`\mathrm{\Omega }`$ so that $`A|\phi \mathrm{\Omega }`$ if $`|\phi \mathrm{\Omega }`$. The Hilbert state $``$ itself is a completion of $`\mathrm{\Omega }`$ with respect to the convergence defined by the norm induced by the scalar product in $``$. The convergence conditions in $`\mathrm{\Omega }`$ are stronger than in $``$ and can be chosen in a suitable way for a particular physical system. Continuous functionals defined on elements $`|\phi \mathrm{\Omega }`$ constitute a linear space $`\mathrm{\Omega }^{}`$ dual to $`\mathrm{\Omega }`$. The space $`\mathrm{\Omega }^{}`$ contains more elements than $`^{}`$ dual to $``$. The space $`^{}`$ consists of continuous linear functionals defined on elements $`|\phi `$, frequently written as $`f|^{}`$ (the value of a functional is a $`c`$-number $`f|\phi `$). The space $`^{}`$ is actually known to be isomorphic to $``$ itself.
The elements of $`\mathrm{\Omega }^{}`$ frequently arise as the generalized eigenvectors $`|\psi _\alpha `$ of the operator $`A`$ with a continuous spectrum
$$A|\psi _\alpha =\lambda _\alpha |\psi _\alpha ,|\psi _\alpha \mathrm{\Omega }^{},$$
(5)
$$\psi _\alpha |\psi _\alpha ^{}=\delta (\alpha \alpha ^{}).$$
Any state vector $`|\phi \mathrm{\Omega }`$ can be expanded in a complete set of generalized eigenvectors
$$|\phi =\phi (\alpha )|\psi _\alpha 𝑑\alpha ,$$
(6)
where the coefficients (wavefunctions) $`\phi (\alpha )`$ are given by the values of a functional from $`\mathrm{\Omega }^{}`$
$$\phi (\alpha )=\psi _\alpha |\phi .$$
(7)
The construction $`\mathrm{\Omega }\mathrm{\Omega }^{}`$ is called rigged Hilbert space (Gel’fand triplet) \[22–24\].
The self-adjoint operator with a continuous spectrum possesses a spectral resolution built of the spectral projectors associated with the generalized eigenvectors
$$A=\lambda _\alpha |\psi _\alpha \psi _\alpha |𝑑\alpha ,$$
(8)
and the corresponding identity resolution
$$I=|\psi _\alpha \psi _\alpha |𝑑\alpha .$$
(9)
A particular functional realization of the rigged space can be chosen in the form of $`𝒥(x)_2(x,dx)𝒥^{}(x)`$ (where $`𝒥(x)`$ is the space of smooth rapidly decreasing functions, $`_2(x,dx)`$ is the space of square integrable functions, and $`𝒥^{}(x)`$ is the space of tempered distributions \[23–25\]).
Consider two orthogonal states of a free non-relativistic one-dimensional particle:
$$|\phi _{1,2}=_{\mathrm{}}^{\mathrm{}}\phi _{1,2}(x)|x𝑑x,|\phi _{1,2}𝒥(x)_2(x,dx),$$
(10)
$$\phi _1|\phi _2=0,|x𝒥^{}(x),x|x^{}=\delta (xx^{}),$$
where $`|x`$ is the generalized eigenvector of the position operator. The identity resolution based on the generalized eigenvectors is familiarly written as
$$I=_{\mathrm{}}^{\mathrm{}}|xx|𝑑x.$$
(11)
The measurement allowing to reliably distinguish between the two states $`|\phi _1`$ and $`|\phi _2`$ without disturbing them is given by the identity resolution defined by Eq.(1). For the measurement outcomes we have
$$\text{Pr}_i\{|\phi _j\}=\text{Tr}\{|\phi _j\phi _j|P_j\}=_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\phi _j^{}(x)\delta (xx^{})\phi _i(x^{})𝑑x𝑑x^{}^2=\delta _{ij}.$$
(12)
This measurement is non-local in the sense that the reliable state identification requires the access to the entire spatial domain where the wavefunction is different from zero. In the non-relativistic quantum mechanics there is no restriction on the maximum speed of information transfer so that the non-local measurements can be made instantly (arbitrarily fast).
Thus, the physical apparatus implementing the measurement (1–3,12) should be non-local in space. The spatial position of the observer which is the ultimate element in any measurement procedure can be chosen completely arbitrarily since due to the unlimited speed of information transfer the data read off with any non-local device can be instantly gathered to any spatial point where the observed is located.
## 3 Relativistic case
Consider now the relativistic case. The states of a relativistic quantum system (field) are described by the rays in the physical Hilbert space $``$ where a unitary representation of the Poincaré group is realized . The local quantum field $`\phi (\widehat{x})`$ (here $`\widehat{x}=(t,𝐱)`$ is a point in the Minkowski space) is defined as a tensor (if the field has more than one component) operator valued distribution. To be more precise, any function (or a set of function if the field has several components) $`f(\widehat{x})𝒥(\widehat{x})`$ is associated with an operator $`\phi (f)`$ which can formally be written as
$$\phi (f)=\underset{j=1}{\overset{r}{}}\phi _j(f_j)=\underset{j=1}{\overset{r}{}}\phi _j(\widehat{x})f_j(\widehat{x})𝑑\widehat{x},$$
(13)
All the operators $`\phi (f)`$ and $`\phi ^{}(f)`$ have a common domain $`\mathrm{\Omega }`$ which does not depend on $`f(\widehat{x})`$ and is a dense subspace $``$ mapped by the operators into itself, $`\phi (f)\mathrm{\Omega }\mathrm{\Omega }`$ ($`\phi ^{}(f)\mathrm{\Omega }\mathrm{\Omega }`$), so that for the vectors $`|\varphi ,|\psi \mathrm{\Omega }`$, the quantity $`\varphi |\phi (f)|\psi `$ is a distribution from $`𝒥^{}(\widehat{x})`$.
The subspace $`\mathrm{\Omega }`$ contains a cyclic vector, called the vacuum state, $`|0\mathrm{\Omega }`$, such that the set of all polynomials $`P(\phi ,f)`$ (field operator algebra) generate the entire $`\mathrm{\Omega }`$. The field operator algebra elements are defined as
$$P(\phi ,f)=f_0+\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{}\phi (\widehat{x}_1)\phi (\widehat{x}_2)\mathrm{}\phi (\widehat{x}_n)f(\widehat{x}_1,\widehat{x}_2,\mathrm{},\widehat{x}_n)𝑑\widehat{x}_1𝑑\widehat{x}_2\mathrm{}𝑑\widehat{x}_n.$$
(14)
The field operators $`\phi (\widehat{x})`$ map the regular states from $`\mathrm{\Omega }`$ to the generalized states $`P(\phi (\widehat{x}))\mathrm{\Omega }\mathrm{\Omega }^{}`$. The microcausality principle is assumed to be satisfied or, to be more precise, if the functions $`f(\widehat{x}),g(\widehat{y})`$ have supports separated by a space-like interval ($`\text{supp}f(\widehat{x})g(\widehat{y})(\widehat{x}\widehat{y})^2<0`$), the field operators are assumed to either commute or anticommute, i.e. for any $`|\psi \mathrm{\Omega }`$ the following microcausality relation holds:
$$[\phi (f),\phi (g)]_\pm |\psi =0,(\widehat{x}\widehat{y})^2<0.$$
(15)
At the infinity, the test functions from $`𝒥(\widehat{x})`$ vanish faster than any polynomial. However, the space of test functions $`𝒥(\widehat{x})`$ contains a dense set of compactly supported functions $`𝒟(\widehat{x})𝒥(\widehat{x})`$ which are zero outside a certain compact domain; therefore, any function from $`𝒥(\widehat{x})`$ can be approximated by a compactly supported function. The equation (15) should be interpreted as the statement that for the field $`\phi (\widehat{x})`$ the measurements performed in the domains separated by a space-like interval do not affect each other since no interaction can propagate faster than light.
It is known that if one requires that (i) the system states are described by the rays in a Hilbert space where a unitary representation of the Poincaré group is realized and (ii) the spectrum of the group generators lies in the front part of the light cone in the momentum representation then the Lorentz-invariant quantum field can only be realized as an operator-valued distribution rather than the field of operators $`\phi (\widehat{x})`$ acting in $``$. Therefore, one can only make sense of $`\phi (f)`$ as an unbounded operator generating a state in $`\mathrm{\Omega }`$. The function $`f(\widehat{x})`$ can be interpreted (with some reservations) as the amplitude (“shape”) of a one-particle packet.
The interpretation of a quantized field as a field of operators leads to the trivial two-point function $`0|\phi ^{}(\widehat{x})\phi ^+(\widehat{y})|0=const`$ and violation of the microcausality principle.
In the relativistic case the entire physical Hilbert space $``$ is a direct sum of the coherent subspaces where different representations of the Poincaré group are realized. Roughly speaking, different subspaces correspond to different types of particles. Considered in the rest of the paper are the one-particle states of a free field; we shall restrict ourselves to the neutral scalar field, spinor field of Dirac electrons and a gauge field (photons).
Suppose we are given two orthogonal states of the neutral scalar field, $`|\phi _{1,2}\mathrm{\Omega }`$
$$|\phi _{1,2}=\phi ^+(f_{1,2})|0=\phi ^+(\widehat{x})f_{1,2}(\widehat{x})𝑑\widehat{x}|0,$$
(16)
$$\phi ^\pm (\widehat{x})=\frac{1}{(2\pi )^{3/2}}_{V_m^+}\text{e}^{\pm i\widehat{p}\widehat{x}}a^\pm (𝐩)\frac{d𝐩}{\sqrt{2p_0}},$$
(17)
$$|𝐩=a^+(𝐩)|0\mathrm{\Omega }^{},\phi ^+(\widehat{x})|0\mathrm{\Omega }^{},$$
where $`a^\pm (𝐩)`$ are the field operators in the momentum representation which generate the generalized eigenvectors from $`\mathrm{\Omega }^{}`$. The integration in Eq.(16) is performed over the front part of the mass shell $`V_m^+`$ ($`p_0^2𝐩^2=m^2`$, $`p_0>0`$). The operator valued distribution $`\phi ^\pm (\widehat{x})`$ satisfies the Klein-Gordon equation
$$(\mathrm{}+m^2)\phi (\widehat{x})=0,\phi (\widehat{x})=\phi ^+(\widehat{x})+\phi ^{}(\widehat{x}).$$
(18)
Since in the relativistic case, just in the non-relativistic quantum mechanics, the field states $`|\phi _{1,2}=\phi ^+(f_{1,2})|0\mathrm{\Omega }`$ are described by the rays in the Hilbert space, the appropriate measurement has the form analogous to Eq.(1) because the orthogonal identity resolution of that kind is only based on the geometrical properties of $``$ (projection on the rays corresponding to the states $`|\phi _{1,2}`$). The measurement allowing to reliably distinguish between the two orthogonal states is given by the identity resolution
$$𝒫_1+𝒫_2+𝒫_{}=I,𝒫_{}=I𝒫_1𝒫_2,$$
(19)
$$𝒫_{1,2}=|\phi _{1,2}\phi _{1,2}|=\phi ^+(f_{1,2})|00|\phi ^{}(f_{1,2}),$$
$$I=\phi ^+(\widehat{x})|00|\phi ^{}(\widehat{x})d𝐱=_{V_m^+}|𝐩𝐩|\frac{d𝐩}{2p_0}.$$
The above identity resolution is written in terms of the generalized eigenstates $`\phi ^+(\widehat{x})|0\mathrm{\Omega }^{}`$.
The probabilities of obtaining the outcomes in different possible channels are
$$\text{Pr}_i\{|\phi _j\}=\text{Tr}\{|\phi _i\phi _i|𝒫_j\}=|\phi _j|\phi _i|^2=|f_j^{}(\widehat{x})D_m^+(\widehat{x}\widehat{x}^{})f_i(\widehat{x}^{})𝑑\widehat{x}𝑑\widehat{x}^{}|^2=$$
(20)
$$=|f_j^{}(\widehat{p})f_i(\widehat{p})\theta (p_0)\delta (\widehat{p}^2m^2)𝑑\widehat{p}|^2=|_{V_m^+}f_j^{}(𝐩)f_i(𝐩)\frac{d𝐩}{2p_0}|^2=\delta _{ij},$$
where the generalized commutator function for the field with mass $`m`$ is
$$D_m^\pm (\widehat{x})=\pm \frac{1}{i(2\pi )^{3/2}}\text{e}^{i\widehat{p}\widehat{x}}\theta (\pm p_0)\delta (\widehat{p}^2m^2)𝑑\widehat{p}=$$
(21)
$$\frac{1}{4\pi }\epsilon (x_0)\delta (\widehat{x}^2)\frac{im}{8\pi \sqrt{\widehat{x}^2}}\theta (\widehat{x}^2)\left[N_1(m\sqrt{\widehat{x}^2})i\epsilon (x_0)J_1(m\sqrt{\widehat{x}^2})\right]\pm \frac{im}{4\pi ^2\sqrt{\widehat{x}^2}}\theta (\widehat{x}^2)K_1(m\sqrt{\widehat{x}^2}),$$
$$\epsilon (x_0)\delta (\widehat{x}^2)\frac{\delta (x_0|𝐱|)\delta (x_0+|𝐱|)}{2|𝐱|}.$$
To within the exponential tales, the commutator function is only different from zero inside the light cone and has a singularity at the cone itself $`\lambda ^2=(\widehat{x}\widehat{x}^{})^2=0`$; outside the light cone the functions $`D^\pm (\lambda )`$ vanish exponentially at the Compton length as $`|\lambda |^{3/4}\mathrm{exp}(m\sqrt{|\lambda |})`$ . For a fixed point $`\widehat{x}`$ the contribution to the integral is only given by the points $`\widehat{x}^{}`$ lying inside the light cone issued from the point $`\widehat{x}`$, which is actually the consequence of the microcausality principle and reflects the impossibility of the faster-than-light field propagation.
The non-zero tails of the commutator function for the massive particles $`m0`$ at the Compton length outside the light cone do not result in any contradictions with the macroscopic causality .
The physical states $`\phi ^+(f_i)|0`$ corresponding to two different functions $`f_i`$ can generally be identical, so that there is no one-to-one correspondence between the functions $`f(\widehat{x})𝒥(\widehat{x})`$ and the states they generate (the generating functions $`f_i`$ are only recoverable from the state $`|\phi _i`$ to within the equivalence class). To be more precise, the states are only determined (as it is seen from Eqs.(16,17)) by the values of the generating function in the momentum representation $`f(\widehat{p})`$ on the mass shell $`p_0=\sqrt{𝐩^2+m^2}`$. The functions $`f_{1,2}(𝐩)`$ should be interpreted as the restriction of the functions $`f_{1,2}(\widehat{p})`$ to the mass shell, so that the states corresponding to any functions $`f(\widehat{x})`$ coinciding in the $`\widehat{p}`$-representation on the mass shell are physically identical.
To eliminate this ambiguity it is convenient to rewrite Eq.(20) in the following equivalent way:
$$\text{Pr}_i\{|\phi _j\}=\text{Tr}\{|\phi _i\phi _i|𝒫_j\}=|\phi _j|\phi _i|^2=|f_j^{}(𝐱)D_m^+(𝐱𝐱^{})f_i(𝐱^{})𝑑𝐱𝑑𝐱^{}|^2=\delta _{ij},$$
(22)
where
$$f_i(𝐱)=\text{e}^{i(t_0p_0(𝐩)\mathrm{𝐩𝐱})}f_i(𝐩)𝑑𝐩.$$
(23)
The commutator function $`D_m^+(𝐱𝐱^{})`$ is obtained from Eq.(21) if one sets $`t=t^{}`$ in $`D_m^+(\widehat{x}\widehat{x}^{})`$ ($`\widehat{x}=(t,𝐱)`$, $`\widehat{x}^{}=(t^{},𝐱^{})`$), i.e. the non-zero contributions to the measurement at a given time slice $`t=t^{}`$ are only given by spatially coinciding points. Here $`f_i(𝐱)`$ is understood as the amplitude (playing the role similar to the one-particle wavefunction in the position representation in the non-relativistic case) taken at all points at the same moment of time corresponding to the beginning of the measurement procedure, i.e. to the moment of time starting from which the entire spatial domain embracing the considered state ($`f_i(𝐱)`$) becomes accessible to the measuring apparatus.
The measurement defined by Eqs.(20,22) is non-local in the sense that it requires the access to the entire domain of the coordinate space $`𝐱`$ where the functions $`f_i(𝐱)`$ are different from zero. Formally, this measurement can be interpreted as a non-local one in the coordinate space performed at a particular moment of time. The device implementing such a measurement should occupy an extended (formally even infinite) domain in the $`𝐱`$-space, i.e. it should have simultaneous access to the entire state. The non-local measurements of that kind are not forbidden in the relativistic case. However, the outcome of the measurement performed with a non-local device cannot be obtained at the same moment as the measurement starts if the measurement encompasses the entire spatial domain where the state is present since the information relevant to the measurement outcome cannot be gathered instantly to the observer located at a certain spatial point from all the points of an extended spatial domain. The information can only be gathered in finite time since actually the indicated spatial domain should be covered by the past part of the light cone issued from the point where the observer is located (Fig.1). The entire domain where the state is present should reside in the interior part of the light cone. The minimum time required for the measurement (for fixed input states) can be determined by examining all the observer positions corresponding to the situation where the entire domain where the state is present at the time $`t_0`$ and which is accessible to the measuring apparatus is completely covered by the past part of the light cone.
Since the required measurement time is only determined by the condition of covering the domain by the light cone, it does not depend on the choice of the reference frame because the light cone is a Lorentz-invariant object. The latter can be graphically illustrated for a one-dimensional state. Suppose that the prepared state in a certain reference frame has a characteristic extent $`x_0x_1=L`$ along the $`x`$-axis. The time required in that system is $`t=L/2c`$ (the optimal observer position is at the center of the domain). At first site, the observer in a moving reference frame would see a Lorentz-contracted domain where the state is defined. Indeed, transformation from the initial reference frame $`(x,y,z,t)`$ to the moving one $`(x^{},y^{},z^{},t^{})`$ through the hyperbolic rotation (Lorentz transformation)
$$\left(\begin{array}{c}x^{}\\ t^{}\end{array}\right)=\left(\begin{array}{cc}\text{ch}\psi & \text{sh}\psi \\ \text{sh}\psi & \text{ch}\psi \end{array}\right)\left(\begin{array}{c}x\\ t\end{array}\right),y^{}=y,z^{}=z,\beta =\text{th}\psi ,$$
(24)
makes the domain size equal to $`x_0^{}x_1^{}=(x_0x_1)\sqrt{1\beta ^2}`$. The time necessary for obtaining the information on the state in the new reference frame is $`t_0^{}t_1^{}=(x_0^{}x_1^{})/2c=(x_0x_1)\sqrt{1\beta ^2}/2c`$ and can be done arbitrarily small in the moving reference frame. However, the time elapsed in the initial reference frame $`t_0t_1=(t_0^{}t_1^{})/\sqrt{1\beta ^2}=(x_0x_1)/2c`$ remains the same<sup>1</sup><sup>1</sup>1This point is important for construction of cryptographic protocols, for example, for the quantum cryptography based on orthogonal states, since it prohibits eavesdropping employing the twin paradox..
For the three-dimensional case the minimal required time will be determined by the size of the maximum cross-section of the spatial domain.
The functions $`f_i(𝐱)`$ rapidly decrease at the infinity, but do not become identical zero outside of any compact domain actually because it is impossible to obtain a function with compact support in the position space taking a Fourier transform of a function which is only defined on the mass shell in the momentum space (the proof based on the Wiener–Paley theorem can be found, e.g. in Refs.). Formally, this means that the time necessary for the reliable identification (with unit probability) of one of the two orthogonal filed states is infinite since because of the infinite support of the functions generating the field states one should have access to the entire coordinate space. However, the states still can be identified in a finite time (which, of course, depends on the structure of the chosen states) with the probability arbitrarily close to unity.
Consider now the case of a multicomponent spinor field of Dirac electrons. The operator valued spinor field distribution has the form
$$\psi (\widehat{x})=\frac{\sqrt{m}}{(2\pi )^{3/2}}\frac{d𝐩}{\sqrt{\epsilon _𝐩}}\underset{\zeta =\pm 1/2}{}\left\{a_\zeta ^{(+)}(𝐩)u_\zeta ^{(+)}(𝐩)\text{e}^{i\widehat{p}\widehat{x}}+a_\zeta ^{()}(𝐩)u_\zeta ^{()}(𝐩)\text{e}^{i\widehat{p}\widehat{x}}\right\},$$
(25)
and the Dirac conjugate operator is
$$\stackrel{~}{\psi }(\widehat{x})=\frac{\sqrt{m}}{(2\pi )^{3/2}}\frac{d𝐩}{\sqrt{\epsilon _𝐩}}\underset{\zeta =\pm 1/2}{}\left\{a_\zeta ^{(+)}(𝐩)\stackrel{~}{u}_\zeta ^{(+)}(𝐩)\text{e}^{i\widehat{p}\widehat{x}}+a_\zeta ^{()}(𝐩)\stackrel{~}{u}_\zeta ^{()}(𝐩)\text{e}^{i\widehat{p}\widehat{x}}\right\}$$
(26)
with the normalization conditions of the Dirac spinors
$$2\underset{\zeta =\pm 1/2}{}u_\zeta ^{(\pm )\alpha }(𝐩)\stackrel{~}{u}_{\zeta \beta }^{(\pm )}(𝐩)=\left(\frac{\widehat{p}}{m}\right)_\beta ^\alpha \pm \delta _\beta ^\alpha ,\epsilon _𝐩=\sqrt{𝐩^2+m^2},\alpha ,\beta =1\mathrm{}4.$$
(27)
The operator valued distribution (25) satisfies the Dirac equation
$$(\widehat{p}+m)\psi (\widehat{x})=0,\widehat{p}=i\gamma ^\mu _\mu ,$$
(28)
and the anticommutation relations
$$[\psi ^\alpha (\widehat{x}),\psi ^\beta (\widehat{x}^{})]_+=[\stackrel{~}{\psi }_\alpha (\widehat{x}),\stackrel{~}{\psi }_\beta (\widehat{x}^{})]_+=0,$$
$$[\psi ^\alpha (\{\widehat{x}),\stackrel{~}{\psi }_\beta (\widehat{x}^{})]_+=iS_\beta ^\alpha (\widehat{x}\widehat{x}^{})=(i\gamma ^\mu _\mu +m)D_m(\widehat{x}\widehat{x}^{}),D_m(\widehat{x})=D_m^+(\widehat{x})+D_m^{}(\widehat{x}).$$
(29)
The smeared operator functions can be written in the form
$$|\psi (𝐟)=\psi ^\alpha (\widehat{x})f^\alpha (\widehat{x})𝑑\widehat{x},f^\alpha (\widehat{x})𝒥(\widehat{x}).$$
(30)
The two orthogonal states $`|\phi _{1,2}`$ ($`\phi _1|\phi _2=0`$) of the spinor field can be written
$$|\phi _{1,2}=\left(\psi ^\alpha (\widehat{x})f^\alpha (\widehat{x})𝑑\widehat{x}\right)|0\mathrm{\Omega },$$
(31)
with the corresponding identity resolution employing the generalized states
$$I=\left(\psi (\widehat{x})\right)|00|\left(\stackrel{~}{\psi }(\widehat{x})\right)d𝐱=\underset{\zeta =\pm \frac{1}{2},s=\pm }{}|𝐩,\zeta ,s𝐩,\zeta ,s|\frac{md𝐩}{2\epsilon (𝐩)},$$
(32)
$$|𝐩,\zeta ,s=a_\zeta ^{(s)}(𝐩)|0,$$
and similarly for the projection operators on the states $`|\phi _1`$ and $`|\phi _2`$.
The probabilities of obtaining different measurement outcomes are
$$\text{Pr}_i\{|\phi _j\}=\text{Tr}\{|\phi _i\phi _i|𝒫_j\}=|f_j^\alpha (\widehat{x})S_{m\beta }^{+\alpha }(\widehat{x}\widehat{x}^{})f_{i\beta }(\widehat{x}^{})𝑑\widehat{x}𝑑\widehat{x}^{}|^2=$$
(33)
$$|f_j^\alpha (\widehat{p})(\widehat{p}m)_\beta ^\alpha f_{i\beta }(\widehat{p})\theta (p_0)\delta (\widehat{p}^2m^2)𝑑\widehat{p}|^2=|_{V_m^+}f_j^\alpha (𝐩)(\widehat{p}m)_\beta ^\alpha f_{i\beta }(𝐩)\frac{d𝐩}{2\epsilon (𝐩)}|^2=\delta _{ij},$$
Here $`f_{i\beta ,j\alpha }(𝐩)`$ are the values of the amplitude on the mass shell with positive energy. In the present case, just as for the scalar field, the answer is expressed through the derivative of the commutator function $`D^\pm (\widehat{x}\widehat{x}^{})`$ so that everything said above on the finiteness of the time required for the distinguishing between two states applies as well to the fermionic field.
Consider now the case of a gauge field which is most interesting from the viewpoint of applications, i.e. the photon field. The electromagnetic field operators are written as
$$A_\mu ^\pm (\widehat{x})=\frac{1}{(2\pi )^{3/2}}\frac{d𝐤}{\sqrt{2k_0}}\text{e}^{\pm i\widehat{k}\widehat{x}}e_\mu ^m(𝐤)a_m^\pm (𝐤)$$
(34)
and satisfy the commutation relations
$$[A_\mu ^{}(\widehat{x}),A_\nu ^{}(\widehat{x}^{})]_{}=ig_{\mu \nu }D_0^{}(\widehat{x}\widehat{x}^{}),$$
(35)
where $`D_0^{}(\widehat{x}\widehat{x}^{})`$ is the commutator function for the massless field (21). There are four types of photons: two transverse, one longitudinal, and one temporal. The two latter types are actually fictitious particles and can be eliminated at the expense of introducing an indefinite metric . For our purposes the shortest way to the required result consists in employing a specific gauge. We shall further work in the subspace of physical space using the Coulomb gauge $`A_\mu =(𝐀,\phi =0)`$ dealing with the two physical transverse states of the electromagnetic field.
The operator valued distribution is a vector in the three-dimensional space:
$$\stackrel{}{𝝍}(\widehat{x})=\frac{1}{(2\pi )^{3/2}}_{V_0^+}\frac{d𝐤}{\sqrt{2k_0}}\underset{s=\pm 1}{}𝐰(𝐤,s)\{a(𝐤,s)\text{e}^{i\widehat{k}\widehat{x}}+a^+(𝐤,s)\text{e}^{i\widehat{k}\widehat{x}}\};$$
(36)
here $`𝐰(𝐤,s)`$ is a three-dimensional vector describing the polarization state $`s=\pm 1`$,
$$𝐰(𝐤,\pm )=\frac{1}{\sqrt{2}}[𝐞_1(𝐤)\pm i𝐞_2(𝐤)],𝐞_1(𝐤)𝐞_2(𝐤),|𝐰(𝐤,s)|^2=1,$$
(37)
where $`𝐞_{1,2}(𝐤)`$ are the orthogonal vectors normal to $`𝐤`$. The field operator satisfies the Maxwell equations
$$\times \stackrel{}{𝝍}(\widehat{x})=i\frac{}{t}\stackrel{}{𝝍}(\widehat{x}),$$
(38)
$$\stackrel{}{𝝍}(\widehat{x})=0.$$
The smeared field operators can be written as
$$\stackrel{}{𝝍}(f_{1,2})=\underset{s=\pm 1}{}\stackrel{}{𝝍}(\widehat{x},s)f_{1,2}(\widehat{x},s)𝑑\widehat{x}=$$
(39)
$$\frac{1}{(2\pi )^{3/2}}_{V_0^+}\frac{d𝐤}{\sqrt{2k_0}}\underset{s=\pm 1}{}𝐰(𝐤,s)\{f_{1,2}(𝐤,s)a(𝐤,s)\text{e}^{i\widehat{k}\widehat{x}}+f_{1,2}(𝐤,s)a^+(𝐤,s)\text{e}^{i\widehat{k}\widehat{x}}\},$$
where the values of the functions $`f_{1,2}(𝐤,s)`$ are taken at the mass shell (light cone $`V_0^+`$ in the momentum representation). Two orthogonal states of the photon field can be written in the form
$$\stackrel{}{|𝝍}_{1,2}=\left(\stackrel{}{𝝍}(f_{1,2})\right)|0.$$
(40)
The corresponding orthogonal projectors on the states $`|\stackrel{}{𝝍}_{1,2}`$ are
$$𝓟_{1,2}=|\stackrel{}{𝝍}_{1,2}\stackrel{}{𝝍}_{1,2}|,$$
(41)
$$I=\left(\stackrel{}{𝝍}^+(\widehat{x})\right)|00|\left(\stackrel{}{𝝍}^{}(\widehat{x})\right)d𝐱=\underset{s=\pm 1}{}_{V_0^+}\left(𝐰(𝐤,s)a^+(𝐤,s)\right)|00|\left(𝐰(𝐤,s)a(𝐤,s)\right)\frac{d𝐤}{2|𝐤|}.$$
(42)
The probabilities of obtaining different measurement outcomes are
$$\text{Pr}_i\{|\stackrel{}{𝝍}_j\}=\text{Tr}\{|\stackrel{}{𝝍}_i\stackrel{}{𝝍}_i|𝓟_j\}=|\stackrel{}{𝝍}_j|\stackrel{}{𝝍}_i|^2=$$
(43)
$$|f_j^{}(𝐱)D_0^+(𝐱𝐱^{})f_i(𝐱^{})𝑑𝐱𝑑𝐱^{}|^2=|_{V_0^+}f_i^{}(𝐤)f_j(𝐤)\frac{d𝐤}{2|𝐤|}|^2=\delta _{ij},$$
where $`D_0^+(𝐱𝐱^{})`$ is the commutator function for the massless field at the time slice $`t=t^{}`$ has the form
$$D_0^+(𝐱𝐱^{})=\frac{1}{4\pi }\frac{\delta (|𝐱𝐱^{}|)}{2|𝐱𝐱^{}|}.$$
It should be noted that arising in the non-relativistic case instead of the commutator function $`D_0^+`$ is the usual $`\delta (𝐱𝐱^{})`$-function (see Eq.(12)). The latter is actually related to the fact that in the non-relativistic case the integration in the scalar product in the momentum representation is performed with the Galilean-invariant measure $`d\mu (𝐩)=d𝐩`$ instead of the Lorentz-invariant measure $`d\mu (\widehat{p})=\theta (p_0)\delta (\widehat{p}^2)d\widehat{p}`$ residing at the mass shell in the relativistic case.
The time necessary for the reliable identification of one of the pair of orthogonal states is determined by the spatial localization of the generating functions $`f_{1,2}(𝐤)`$ in the coordinate (position) representation or, to be more precise, by the localization of $`|𝐟_{1,2}(𝐱,t,s)|^2`$ $`|𝐟_{1,2}(𝐱,t,s)|^2`$
$$𝐟_{1,2}(𝐱,t,s)=_{V_0^+}𝐰(𝐤,s)f_{1,2}(𝐤,s)\text{e}^{i(|𝐤|t\mathrm{𝐤𝐱})}𝑑𝐤,$$
(44)
where the integration is performed over the mass shell $`V_0^+`$ (surface of the light cone in the momentum representation). It has long been known (e.g., see Ref. and references therein) that it is impossible to obtain a strictly localizable (with compact support) function in the coordinate space from a normalized function defined on the mass shell. The latter holds for both massive and massless particles. However, there exist the functions with arbitrarily close to the exponential fall off at the infinity. Recently, the one-particle photon states of that kind were explicitly constructed . The restrictions on the arbitrarily close to the exponential fall off stem from the Wiener-Paley theorem for square integrable functions and the impossibility of mixing of the positive and negative frequency states on the mass shell for one-particle states.
The possibility of the existence of the states whose localization is arbitrarily close to the exponential one is actually related to the choice of the test function space $`𝒥(\widehat{x})`$ which contains the dense subspace of compactly supported functions $`𝒟(\widehat{x})`$. The latter circumstance is closely related to the local nature of the theory since the existence of a dense set of compactly supported functions defined in the Minkowski $`\widehat{x}`$-space allows to achieve the local properties of the distributions, including the commutator functions, appearing in the microcausality principle.
Reduction of the set of the test functions can result in the non-local nature of the theory (e.g., see Refs.\[23,29–32\]). However, non-locality of the theory does not imply violation of the causality principle at the macroscopic level (at the level of the observer) .
## 4 Conclusion
Thus, the reliable (with unit probability) identification of one-particle field states generally requires an infinite time because of the requirement of having access to the entire spatial domain where the state “is present”. Formally, the construction of the measurement at the level of the corresponding identity resolution in the relativistic case is completely analogous to the non-relativistic quantum mechanics since in both cases the states are described by the rays in the Hilbert state. At this level, the description of measurement involves only the geometrical properties of the state space (the scalar product and projections on the corresponding orthogonal states). The difference between the non-relativistic and relativistic cases arises at the level of the internal structure of the scalar product and conveying the measurement outcome to the observer. In the non-relativistic case, there are no restrictions on the instantaneous spatially non-local measurements and conveying their outcomes to the observer (which can be located at an arbitrary point) since there are no limits on the maximum speed of information transfer. In the relativistic case the Lorentz-invariant scalar product already contains (through the commutator functions reflecting the microcausality principle) the information on the maximum speed of the field propagation. The latter should only be interpreted as the impossibility for the field to propagate into an extended device (fill it) faster than light. A similar situation occurs if an extended device moves into the domain where the field is present. However, the presence of the field in the device alone does not yet mean obtaining of the measurement outcome. The measurement outcome should be conveyed to a single point (location of the observer). The latter can only be done in a finite time depending on the spatial localization of the states since the spatial domain should be “covered” by the past part of the light cone issued from the point where the observer is situated (the observer can be interpreted as a classical device registering the final measurement outcome).
In the above approach all the information on the non-local (extended) device is contained in the projectors on the corresponding rays in the Hilbert space. Formally, these projectors should be considered as an (non-local) observable which, from the viewpoint of the measurement theory, is quite similar to a non-local in the coordinate space observable (e.g., momentum). Therefore, associated with this observable should be a physical apparatus implementing the corresponding measurement.
The conclusion on the finiteness of the time interval required for the state identification does not depend on a particular measurement procedure. It is often experimentally easier to realize a non-local measurement as a local one employing a localized auxiliary system $`A`$ interacting with the measured system (it is known that any measurement can be realized in this way ). In the course of the joint evolution the state of the auxiliary system is gradually changed and finally the measurement is performed over that system. For the orthogonal states there are no restrictions on the existence of non-disturbing measurements (see Eqs.(1–3)) in the sense that just after the measurement only the state of the auxiliary system $`A`$ is changed while state of the studied system $`S`$ coincides with its initial state just before the interaction between $`S`$ and $`A`$ was turned on. In that case it is also intuitively clear that a finite time is required (if the field state is extended) for the field to pass through a point where the system $`A`$ is localized and the interaction between $`A`$ and $`S`$ occurs which cannot be done faster than $`tL/c`$ because of the existence of a maximum speed of the field propagation. It is also clear that this time cannot be reduced by enhancement of the local coupling between $`A`$ and $`S`$ because this time is limited by the field propagation speed. The required time cannot also be reduced employing the Lorentz contraction of the domain where the field is present occurring in a moving reference frame (see discussion above). Naturally, the auxiliary system $`A`$ should also be described by the quantum field theory and one should assume that it is so strongly spatially localized that the time necessary to perform a measurement on this system can be done arbitrarily small. However, it is much more difficult to obtain any general results in this approach since it inevitably requires making certain assumptions concerning the system $`S`$ itself and its interaction with the auxiliary system $`A`$.
Formally, the von Neumann measurement described by the orthogonal resolution of identity given by Eq.(43) in the subspace of one-particle states in the Hilbert space does not disturb the orthogonal states since the situation here is analogous to the non-relativistic case described by Eqs.(1–4). However, because of the non-local nature of the projectors, the reliable and non-disturbing identification of the orthogonal states achieved by this measurement formally requires infinite time.
With respect to the time necessary for the measurement, the situation is quite similar for the non-orthogonal states, the only difference being that the latter cannot be reliably distinguished even in an infinite time.
In conclusion, it should be noted that the non-relativistic quantum cryptographic protocols implicitly assume the possibility of performing a measurement distinguishing between the two states in a finite time (ideally, an instantaneous measurement). Briefly, all the non-relativistic protocols can be described in the following way. User A sends either the state $`|\psi _1`$ or $`|\psi _2`$ to user B who performs quantum-mechanical measurements. Part of information on the measurement outcomes is discussed through the open communication channel. It is assumed that the states $`|\psi _1`$ or $`|\psi _2`$ are prepared and measured in certain agreed moments of time so that the preparation and measurement procedures can be performed at arbitrary moments of time. Thus, it is implicitly assumed that the protocol can be implemented in a finite time interval.
The relativistic quantum field theory does not allow access to the entire state in any finite time because of the non-localizable nature of the state. Therefore, all the exchange protocols with finite duration should inevitably involve an error in the preparation and detection of both orthogonal and non-orthogonal one-particle states even for the ideal communication channel. In other words, the non-localizability of the states results in the effective “noise” because of the impossibility of the reliable detection of orthogonal states in a finite time although for a fixed measurement duration the error can be made arbitrarily (exponentially) small by choosing more and more strongly localized state. This circumstance should be taken into account in the realistic cryptographic protocols employing photons as the information carriers.
The authors are grateful to W.L.Golo, S.V.Iordanskii and V.M.Edelstein for the interest and discussion of the results of this work.
This work was supported by the Russian Foundation for Basic Research (project No 99-02-18127), the project “Physical Principles of the Quantum Computer”, and the program “Advanced Devices and Technologies in Micro- and Nanoelectronics” (project No 02.04.5.2.40.T.50).
This work was also supported by the Wihuri Foundation, Finland.
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# Effects of the Fourth Generation on Δ𝑀_𝐵_{𝑑,𝑠} in 𝐵⁰-𝐵̄⁰ Mixing
## 1 Introduction
The Standard Model (SM) is a very successful theory of the elementary particles known today. But it must be incomplete because it has too many unpredicted parameters ($`ninteen!`$) to be put by hand. Most of these parameters are in the fermion part of the theory. We don’t know the source of the quarks and leptons, as well as how to determinate their mass and number theoretically. We have to get their information all from experiment. There is still no successful theory which can be descripted them with a unified point, even if the Grand Unified Theory and Supersymmetry. Perhaps elementary particles have substructure ( like $`preon`$) and we need to progress more elementary theories. But this is beyond our current experimental level.
From the point of phenomenology, for fermions, there is a realistic question is number of the fermions generation or weather there are other additional quarks or leptons. The present experiments can tell us there are only three generation fermions with $`light`$ neutrinos which mass are less smaller than $`M_Z/2`$ but the experiments don’t exclude the existence of other additional generation, such as the fourth generation, with a $`heavy`$ neutrino, i.e. $`m_{\nu _4}M_Z/2`$. Many refs. have studied models which extend the fermions part, such as vector-like quark models, sterile neutrino models and the sequential four generation standard model (SM4) which we talk in this note. We consider a sequential fourth generation non -SUSY model, which is added an up-like quark $`t^{^{}}`$, a down-like quark $`b^{^{}}`$, a lepton $`\tau ^{^{}}`$, and a heavy neutrino $`\nu ^{^{}}`$ in the SM. The properties of these new fermions are all the same as their corresponding counterparts of other three generations except their masses and CKM mixing, see tab.1,
There are a lot of refs. about the fourth generation. Some refs. devoted to the mass spectrum of the fourth generation particles, as well as some discussed the mass bounds of the fourth leptons. There are many papers talked about other problems about the fourth generation and the experimental search of the fourth generation paticles. In our previous papers, we investigated the rare $`B`$ meson decays with the fourth generation and $`ϵ^{^{}}/ϵ`$ in $`K^0`$ systems in SM4. We got some interesting results, such as the new effects of the 4th generation particle on the meson decays and CP violation. We also got the constraints of the fourth generation CKM matrix factors, like $`V_{t^{^{}}s}^{}V_t^{^{}}b`$ from $`BX_s\gamma `$ and $`V_{t^{^{}}s}^{}V_t^{^{}}d`$ from $`ϵ^{^{}}/ϵ`$. In other words, these rare decays provided possible test of the fourth generation existence, as well as CP violation. In this note, we talk about the mass difference $`\mathrm{\Delta }M_{B_{d,s}}`$ in $`B^0\overline{B}^0`$ with the fourth generation. We will give the prediction of $`\mathrm{\Delta }M_{B_s}`$ in SM4 and the constraints of a new fourth generation CKM matrix factor $`V_{t^{^{}}b}^{}V_t^{^{}}d`$ from $`\mathrm{\Delta }M_{B_d}`$ in SM4. Like refs.huo1,huo2, this note can also provide a test of 4th generation.
Particle-antiparticle mixing is responsible for the small mass differences between the mass eigenstates of neutral mesons, such as $`\mathrm{\Delta }M_K`$ in $`K_\mathrm{L}K_\mathrm{S}`$ mixing and $`\mathrm{\Delta }M_{B_{d,s}}`$ in $`B^0\overline{B}^0`$ mixing. Being an FCNC process it involves heavy quarks in loops and consequently it is a perfect testing ground for heavy flavor physics. For example, $`B^0\overline{B}^0`$ mixing gave the first indication of a large top quark mass. $`K_\mathrm{L}K_\mathrm{S}`$ mixing is also closely related to the violation of CP symmetry which is experimentally known since 1964. They are sensitive measures of the top quark $`t`$ couplings $`V_{ti}(i=d,s,b)`$ and of the top quark mass $`m_t`$. The experimental measurements of $`\mathrm{\Delta }B_d`$ is used to determine the CKM matrix elements $`V_{td}`$. It offer an improved determination of the unitarity triangle with the future accurate measurement of $`\mathrm{\Delta }M_{B_s}`$. For physics beyond the SM, there are a number of studies of the new physics effects in $`B_d`$ decays. But $`B_s`$ system has received somewhat less attention from new physics point of view. Experimentally, $`\mathrm{\Delta }M_{B_d}`$ has been accurately measured, $`\mathrm{\Delta }M_{B_d}=0.473\pm 0.016(ps)^1`$. But $`\mathrm{\Delta }M_{B_s}`$ has only lower bound, $`\mathrm{\Delta }M_{B_s}>14.3(ps)^1`$.
In this note, We want to investigate $`\mathrm{\Delta }M_{B_{d,s}}`$ in $`B^0\overline{B}^0`$ mixing in SM4. First, if we add a sequential fourth up-like quark $`t^{^{}}`$, there produce new prediction of the the mass difference $`\mathrm{\Delta }M_{B_s}`$ through the new Wilson coefficients and a fourth generation CKM matrix factor $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$, which constrained for the rare decay $`BX_s\gamma `$ in . We found, as like the analysis of $`BX_sl^+l^{}`$ in , our results of the prediction of $`\mathrm{\Delta }M_{B_s}`$ in SM4 are quite different from that of SM and can satisfy the lower experimental bound in one case of the values $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$ taken. In another case, it is almost the same as the SM one. The new effects of the fourth generation show clearly in the first case. Second, we can get the constraint of a fourth generation CKM matrix factor, $`V_{t^{^{}}b}^{}V_{t^{^{}}d}`$ from the experimental measurement of $`\mathrm{\Delta }M_{B_d}`$. We get one kind of reasonable analytical solutions of $`V_{t^{^{}}b}^{}V_{t^{^{}}d}`$. They are very small, $`1.0\times 10^4V_{t^{^{}}b}^{}V_{t^{^{}}d}0.5\times 10^4`$. These result don’t contradicted the unitarity constraints for quark $`d,b`$. Moreover, This small absolute value, $`\lambda ^4\lambda ^5`$ order, is in agreement with the hierarchy in the CKM matrix elements. It seems to give the possible test of this hierarchy and the existence of the fourth generation. We give the analysis of the hierarchy in the four generation CKM matrix.
In sec. 2, we give the basic formulae for the mass difference $`\mathrm{\Delta }M_{B_{d,s}}`$ in $`B^0\overline{B}^0`$ with the sequential fourth generation up-like quark $`t^{^{}}`$ in SM4 model. In sec. 3, we give the prediction of mass difference $`\mathrm{\Delta }M_{B_s}`$ in SM4 and the numerical analysis. Sec. 4 is devoted to the numerical analysis of one fourth generation CKM matrix factor $`V_{t^{^{}}b}^{}V_{t^{^{}}d}`$ from the experimental measurements of the mass difference $`\mathrm{\Delta }M_{B_d}`$ in SM4. We also analyze the hierarchy of the four generation CKM matrix in this section. Finally, in sec. 5, we give our conclusion.
## 2 Basic formulae for $`\mathrm{\Delta }M_{B_{d,s}}`$ with $`t^{^{}}`$
$`B_{d,s}^0\overline{B}_{d,s}^0`$ mixing proceeds to an excellent approximation only through box diagrams with internal top quark exchanges in SM. In SM, the effective Hamiltonian $`_{\mathrm{eff}}(\mathrm{\Delta }B=2)`$ for $`B_{d,s}^0\overline{B}_{d,s}^0`$ mixing, relevant for scales $`\mu _b=𝒪(m_b)`$ is given by
$`_{\mathrm{eff}}^{\mathrm{\Delta }B=2}={\displaystyle \frac{G_\mathrm{F}^2}{16\pi ^2}}M_W^2(V_{tb}^{}V_{tq})^2S_0(x_t)Q(\mathrm{\Delta }B=2)+h.c.`$ (1)
where $`Q(\mathrm{\Delta }B=2)=(\overline{b}_\alpha q_\alpha )_{VA}(\overline{b}_\beta q_\beta )_{VA}`$, with $`q=d,s`$ for $`B_{d,s}^0\overline{B}_{d,s}^0`$ respectively and $`S_0(x_t)`$ is the Wilson coefficient which is taken the form
$`S_0(x_t)={\displaystyle \frac{4x_t11x_t^2+x_t^3}{4(1x_t)^2}}{\displaystyle \frac{3}{2}}{\displaystyle \frac{x_t^3}{(1x_t)^3}}\mathrm{ln}x,`$ (2)
where $`x_t=m_t^2/M_W^2`$. The mass differences $`\mathrm{\Delta }M_{d,s}`$ can be expressed in terms of the off-diagonal element in the neutral $`B`$-meson mass matrix
$`\mathrm{\Delta }M_{d,s}`$ $`=`$ $`2|M_{12}^{d,s}|`$ (3)
$`2m_{B_{d,s}}|M_{12}^{d,s}|`$ $`=`$ $`|\overline{B}_{d,s}^0|_{\mathrm{eff}}(\mathrm{\Delta }B=2)|B_{d,s}^0.`$
Now, we turn to the case of SM4. If we add a fourth sequential fourth generation up-like quark $`t^{}`$, the above equations would have some modification. There exist other box diagrams contributed by $`t^{^{}}`$ (see fig. 1), similar to the leading box diagrams in MSSM.
The effective Hamiltonian in Standard Model, eq.(1), chang into the form,
$`_{\mathrm{eff}}^{\mathrm{\Delta }B=2}`$ $`=`$ $`{\displaystyle \frac{G_\mathrm{F}^2}{16\pi ^2}}M_W^2[\eta _t(V_{tb}^{}V_{tq})^2S_0(x_t)+`$ (4)
$`+`$ $`\eta _t^{^{}}(V_{t^{^{}}b}^{}V_{t^{^{}}q})^2S_0(x_t^{^{}})+\eta _{tt^{^{}}}(V_{t^{^{}}b}^{}V_{t^{^{}}q})(V_{tb}^{}V_{tq})S_0(x_t,x_t^{^{}})]Q(\mathrm{\Delta }B=2)+h.c.`$
The mass differences $`\mathrm{\Delta }M_{d,s}`$ in SM4 can be expressed
$`\mathrm{\Delta }M_d`$ $`=`$ $`{\displaystyle \frac{G_\mathrm{F}^2}{6\pi ^2}}M_W^2m_{B_d}(\widehat{B}_{B_d}\widehat{F}_{B_d}^2)[\eta _t(V_{tb}^{}V_{td})^2S_0(x_t)+`$ (5)
$`+`$ $`\eta _t^{^{}}(V_{t^{^{}}b}^{}V_{t^{^{}}d})^2S_0(x_t^{^{}})+\eta _{tt^{^{}}}(V_{t^{^{}}b}^{}V_{t^{^{}}d})(V_{tb}^{}V_{td})S_0(x_t,x_t^{^{}})]`$
$`\mathrm{\Delta }M_s`$ $`=`$ $`{\displaystyle \frac{G_\mathrm{F}^2}{6\pi ^2}}M_W^2m_{B_s}(\widehat{B}_{B_s}\widehat{F}_{B_s}^2)[\eta _t(V_{tb}^{}V_{ts})^2S_0(x_t)+`$ (6)
$`+`$ $`\eta _t^{^{}}(V_{t^{^{}}b}^{}V_{t^{^{}}s})^2S_0(x_t^{^{}})+\eta _{tt^{^{}}}(V_{t^{^{}}b}^{}V_{t^{^{}}s})(V_{tb}^{}V_{ts})S_0(x_t,x_t^{^{}})]`$
where $`(\widehat{B}_{B_s}\widehat{F}_{B_s}^2)=\xi _s^2(\widehat{B}_{B_d}\widehat{F}_{B_d}^2)`$. The new Wilson coefficients $`S_0(x_t^{^{}})`$ present the contribution of $`t^{}`$, which like $`S_0(x_t)`$ SM in eq. (5) except exchanging $`t^{}`$ quark not $`t`$ quark. $`S_0(x_t,x_t^{^{}})`$ present the contribution of a mixed $`tt^{}`$, which is taken the form
$`S_0(x,y)`$ $`=`$ $`xy[{\displaystyle \frac{1}{yx}}({\displaystyle \frac{1}{4}}+{\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{1x}}{\displaystyle \frac{3}{4}}{\displaystyle \frac{1}{(1x)^2}}\mathrm{ln}x+`$ (7)
$`+`$ $`(yx){\displaystyle \frac{3}{4}}{\displaystyle \frac{1}{(1x)(1y)}}]`$
where $`x=x_t=m_t^2/M_W^2`$, $`y=x_t^{^{}}=m_t^{^{}}^2/M_W^2`$. The numerical results of $`S_0(x_t^{})`$ and $`S_0(x_t,x_t^{})`$ is shown on the tab. 2.
The short-distance QCD correction factors $`\eta _t^{^{}}`$ and $`\eta _{tt^{^{}}}`$ can be calculated like $`\eta _c`$ and $`\eta _{ct}`$ in the mixing of $`K^0\overline{K}^0`$, which the NLO values are given in refs, relevant for scale not $`𝒪(\mu _c)`$ but $`𝒪(\mu _b)`$. In leading-order, $`\eta _t`$ is calculated by
$`\eta _t^0=[\alpha _s(\mu _t)]^{(6/23)},\alpha _s(\mu _t)=\alpha _s(M_Z)[1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(\beta _0{\displaystyle \frac{\alpha _s(M_Z)}{2\pi }}\mathrm{In}{\displaystyle \frac{M_Z}{\mu _t}})^n],`$ (8)
with its numerical value in tab. 3. The formulae of factor $`\eta _t^{^{}}`$ is similar to the above equation except for exchanging $`t`$ by $`t^{^{}}`$.In ref., to leading order, $`\eta _{tt^{^{}}}`$ was taken as
$`\eta _t^0{}_{}{}^{}=[\alpha _s(\mu _t)]^{(6/23)}[{\displaystyle \frac{\alpha _s(\mu _b^{^{}})}{\alpha _s(\mu _t)}}]^{(6/21)}[{\displaystyle \frac{\alpha _s(\mu _t^{^{}})}{\alpha _s(\mu _b^{^{}})}}]^{(6/19)},`$ (9)
with value 0.58 for the same as for $`\eta _{tt^{^{}}}^K`$ in $`K^0\overline{K}^0`$. For simplicity, we take $`\eta _{tt^{^{}}}=\eta _t^{^{}}`$. We give the numerical results in tab.4.
In the last of this section, we give other input parameters necessary in this note. (See the following tab.).
## 3 Prediction of $`\mathrm{\Delta }M_{B_s}`$ with $`t^{^{}}`$
Experimentally, the mass difference $`\mathrm{\Delta }M_{B_s}`$ of the $`B^0\overline{B}^0`$ mixing is unclear. It has only low bound, $`\mathrm{\Delta }M_{B_s}^{\mathrm{exp}}>14.3(ps)^1`$. We have given the calculation formula of $`\mathrm{\Delta }M_{B_s}`$ in eq. (6) and the numerical results of Wilson coefficients $`S_0`$ and QCD correction coefficients $`\eta `$. Now, if we constrain the fourth generation CKM factor $`V_{t^{^{}}b}^{}V_{t^{^{}}s}`$, we can predict $`\mathrm{\Delta }M_{B_s}`$ in our four generations model. Fortunately, from our previous paper, we have given the constraints of $`V_{t^{^{}}b}^{}V_{t^{^{}}s}`$ from experimental measurements of $`BX_s\gamma `$. Here, we give only the basic scheme and the final numerical results.
The leading logarithmic calculations can be summarized in a compact form as follows :
$$R_{\mathrm{quark}}=\frac{Br(BX_s\gamma )}{Br(BX_ce\overline{\nu }_e)}=\frac{|V_{ts}^{}V_{tb}|^2}{|V_{cb}|^2}\frac{6\alpha }{\pi f(z)}|C_7^{\mathrm{eff}}(\mu _b)|^2.$$
(10)
In the case of four generation there is an additional contribution to $`BX_s\gamma `$ from the virtual exchange of the fourth generation up quark $`t^{^{}}`$. The Wilson coefficients of the dipole operators are given by
$$C_{7,8}^{\mathrm{eff}}(\mu _b)=C_{7,8}^{(\mathrm{SM})\mathrm{eff}}(\mu _b)+\frac{V_{t^{^{}}s}^{}V_{t^{^{}}b}}{V_{ts}^{}V_{tb}}C_{7,8}^{(4)\mathrm{eff}}(\mu _b),$$
(11)
where $`C_{7,8}^{(4)\mathrm{eff}}(\mu _b)`$ present the contributions of $`t^{^{}}`$ to the Wilson coefficients, and $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$ are the fourth generation CKM matrix factor which we need now. With these Wilson coefficients and the experiment results of the decays of $`BX_s\gamma `$ and $`Br(BX_ce\overline{\nu }_e)`$ , we obtain the results of the fourth generation CKM factor $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$. There exist two cases, a positive factor and a negative one:
$`V_{t^{^{}}s}^{}V_{t^{^{}}b}^{(+)}`$ $`=`$ $`[C_7^{(0)\mathrm{eff}}(\mu _b)C_7^{(\mathrm{SM})\mathrm{eff}}(\mu _b)]{\displaystyle \frac{V_{ts}^{}V_{tb}}{C_7^{(4)\mathrm{eff}}(\mu _b)}}`$ (12)
$`=`$ $`[\sqrt{{\displaystyle \frac{R_{\mathrm{quark}}|V_{cb}|^2\pi f(z)}{|V_{ts}^{}V_{tb}|^26\alpha }}}C_7^{(\mathrm{SM})\mathrm{eff}}(\mu _b)]{\displaystyle \frac{V_{ts}^{}V_{tb}}{C_7^{(4)\mathrm{eff}}(\mu _b)}}`$
$$V_{t^{^{}}s}^{}V_{t^{^{}}b}^{()}=[\sqrt{\frac{R_{\mathrm{quark}}|V_{cb}|^2\pi f(z)}{|V_{ts}^{}V_{tb}|^26\alpha }}C_7^{(\mathrm{SM})\mathrm{eff}}(\mu _b)]\frac{V_{ts}^{}V_{tb}}{C_7^{(4)\mathrm{eff}}(\mu _b)}$$
(13)
as in table. 3. With these values, we can give the prediction of $`\mathrm{\Delta }M_{B_s}`$ in SM4 by the figs. 2. It is very interesting that the fanal analytic result is same as that in decay of $`BX_sl^+l^{}`$.
The mass difference $`\mathrm{\Delta }M_{B_s}`$ in the two cases of $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$ are shown in the figs. 2(a) and 2(b) respectively. In the first case, which the value of $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$ takes positive, i.e. $`(V_{t^{^{}}s}^{}V_{t^{^{}}b}^{})`$, the curve of $`\mathrm{\Delta }M_{B_s}`$ to $`m_t^{^{}}`$ is almost overlap with that of SM. That is, the results in SM4 are the same as that in SM, except a peak in the curve when $`m_t^{^{}}`$ takes value about 170GeV. The reason is not because there is new prediction deviation from the SM but only because there is a term of $`(xy)`$ in denominator of the formulae of eq. (7). In this case, it does not show the new effects of $`t^{^{}}`$. The mass difference $`\mathrm{\Delta }M_{B_s}`$ is nathless unclear. Also, we can not obtain the information of existence of the fourth generation from $`\mathrm{\Delta }M_{B_s}`$, although we can not exclude them either. This is because, from tab. 5, the values of $`V_{t^{^{}}s}^{}V_{t^{^{}}b}^{()}`$ are positive. But they are of order $`10^3`$ and is very small. The values of $`V_{ts}^{}V_{tb}`$ are about ten times larger than them ($`V_{ts}^{}=0.038`$, $`V_{tb}=0.9995`$, see ref. ). Furthermore, the last two terms about $`m_t^{^{}}`$ in eq. (6) are approximately same order. The contribution of them counteract each other.
But in the second case, when the values of $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$ are negative, i.e. $`(V_{t^{^{}}s}^{}V_{t^{^{}}b}^{()})`$. The curve of $`\mathrm{\Delta }M_{B_s}`$ is quite different from that of the SM. This can be clearly seen from fig. 2(b). The enhancement of $`\mathrm{\Delta }M_{B_s}`$ increases rapidly with increasing of $`t^{^{}}`$ quark mass. In this case, the fourth generation effects are shown clearly. The reason is that $`V_{t^{^{}}s}^{}V_{t^{^{}}b}^{(+)}`$ is 2-3 times larger than $`V_{ts}^{}\dot{V}_{tb}`$ so that the last two terms about $`m_t^{^{}}`$ in eq. (6) becomes important and it depends on the $`t^{^{}}`$ mass strongly. Thus, the effect of the fourth generation is significant. Meanwhile, The prediction of $`\mathrm{\Delta }M_{B_s}`$ in SM4 can satisfy the experimental low bound of $`\mathrm{\Delta }M_{B_s}14.3(ps^1)`$. So, the sequential fourth generation model could be one of the ways of searching new physics about $`\mathrm{\Delta }M_{B_s}`$. If $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$ choose this case, the mass difference $`\mathrm{\Delta }M_{B_s}`$ in $`B^0\overline{B}^0`$ mixing could be a good probe to the existence of the fourth generation.
## 4 Constrains of the fourth generation CKM factor $`V_{t^{^{}}b}^{}V_{t^{^{}}d}`$ from experimental measurements of $`\mathrm{\Delta }M_{B_d}`$
Unlike $`\mathrm{\Delta }M_{B_s}`$, the mass difference $`\mathrm{\Delta }M_{B_d}`$ of $`B_d^0\overline{B}_d^0`$ mixing is experimental clear, $`\mathrm{\Delta }M_{B_d}^{\mathrm{exp}}=0.473\pm 0.016(ps)^1`$. We can get the constraints of the fourth generation CKM factor $`V_{t^{^{}}b}^{}V_{t^{^{}}d}`$ from the present experimental value of $`\mathrm{\Delta }M_{B_d}`$.
We change the form of eq. (5) as a quadratic equation about $`V_{t^{^{}}b}^{}V_{t^{^{}}d}`$. By solving it , we can get two analytical solution $`V_{t^{^{}}d}^{}V_{t^{^{}}b}^{(1)}`$ (absolute value is the large one) and $`V_{t^{^{}}d}^{}V_{t^{^{}}b}^{(2)}`$ (absolute value is the small one), just like the other 4th generation CKM matrix factor $`V_{t^{^{}}s}^{}V_{t^{^{}}b}^{(\pm )}`$ in last section.
However, experimentally, it is not accurate for the measurement of CKM matrix element $`V_{td}`$. So, we have to search other ways to solve this difficulty. Fortunately, the CKM unitarity triangle, i.e. the graphic representation of the unitarity relation for $`d,b`$ quarks, which come from the orthogonality condition on the first and third row of $`V_{\mathrm{CKM}}`$,
$`V_{ud}V_{ub}^{}+V_{cd}V_{cb}^{}+V_{td}V_{tb}^{}=0,`$ (14)
can be conveniently depicted as a triangle relation in the complex plane, as shown in the fig. 3(a).
From the above equation, we can give the constraints of $`V_{td}V_{tb}^{}`$,
$$0.005|V_{td}V_{tb}^{}|0.013$$
(15)
Then, we give the final results as shown in the figs. 4(a) and 4(b).
We must announce that figs. 4 only show the curves with $`V_{t^{^{}}d}^{}V_{t^{^{}}b}^{(2)}`$ (absolute value is the small one) firstly. Because the absolute value of $`V_{t^{^{}}d}^{}V_{t^{^{}}b}^{(1)}`$ is generally larger than 1. This is contradict to the unitarity of CKM matrix. So, we don’t think about this solution. From the figs. 4, we found all curves are in the range from $`1\times 10^4`$ to $`0.5\times 10^4`$ when we considering the constraint of $`V_{td}V_{tb}^{}`$. That is to say, the absolute value of $`V_{t^{^{}}d}^{}V_{t^{^{}}b}`$ is about $`10^4`$ order. This is a very interesting result.
First, these CKM matrix elements obey unitarity constraints. With the fourth generation quark $`t^{^{}}`$, eq. (12) change to ,
$$V_{ud}^{}V_{ub}+V_{cd}^{}V_{cb}+V_{td}^{}V_{tb}+V_{t^{^{}}d}^{}V_{t^{^{}}b}=0.$$
(16)
This a quadrilatral, (see fig. 3(b)). We take the average values of the SM CKM matrix elements from Ref. . The sum of the first three terms in eq. (12) is about $`10^2`$ order. If we take the value of $`V_{t^{^{}}s}^{}V_{t^{^{}}b}^{(2)}`$ the result of the left of (14) is better and more close to $`0`$ than that in SM, when $`V_{t^{^{}}s}^{}V_{t^{^{}}b}^{(2)}`$ takes negative values. Even if $`V_{t^{^{}}s}^{}V_{t^{^{}}b}^{(2)}`$ takes positive values, the sum of (16) would change very little because the values of $`V_{t^{^{}}d}^{}V_{t^{^{}}b}^{(2)}`$ are about $`10^4`$ order, two orders smaller than the sum of the first three ones in the left of (14). Considering that the data of CKM matrix is not very accurate, we can get the error range of the sum of these first three terms. It is much larger than $`V_{t^{^{}}d}^{}V_{t^{^{}}b}^{(2)}`$. Thus,in the case the values of $`V_{t^{^{}}d}^{}V_{t^{^{}}b}`$ satisfy the CKM matrix unitarity constraints.
Second, this small order of $`V_{t^{^{}}d}^{}V_{t^{^{}}b}`$ doesn’t contradict to the hierarchy of the CKM matrix elements or the quarks mixing angles. Moreover, it seem to prove the hierarchy. The hierarchy in the quarks mixing angles is clearly presented in the Wolfenstein parameterization of the CKM matrix. Let’s see CKM matrix firstly,
$$V_{\mathrm{CKM}}=\left(\begin{array}{cccc}V_{ud}\hfill & V_{us}& \hfill V_{ub}& \hfill \mathrm{}\\ V_{cd}\hfill & V_{cs}& \hfill V_{cb}& \hfill \mathrm{}\\ V_{td}\hfill & V_{ts}& \hfill V_{tb}& \hfill \mathrm{}\\ \mathrm{}\hfill & \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}\end{array}\right)\left(\begin{array}{cccc}1\hfill & \lambda & \hfill \lambda ^3& \hfill \mathrm{}\\ \lambda \hfill & 1& \hfill \lambda ^2& \hfill \mathrm{}\\ \lambda ^3\hfill & \lambda ^2& \hfill 1& \hfill \mathrm{}\\ \mathrm{}\hfill & \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}\end{array}\right)$$
(17)
with $`\lambda =\mathrm{sin}^2\theta =0.23`$. Now, the hierarchy can be expressed in powers of $`\lambda `$. We found, the magnitudes of the mixing angles are about 1 among the $`same`$ generations, $`V_{ud}`$, $`V_{cs}`$ and $`V_{tb}`$. For different generations, the magnitudes are about $`\lambda `$ order between $`1st`$ and $`2nd`$ generation, $`V_{us}`$ and $`V_{cd}`$, as well as about $`\lambda ^2`$ order between $`2nd`$ and $`3rd`$ generation, $`V_{cb}`$ and $`V_{ts}`$. The magnitudes are about $`\lambda ^3`$ order between the $`1st`$ and $`third`$ generation, $`V_{ub}`$ and $`V_{td}`$. Then, there should be an interesting problem: If the fourth generation quarks exist, how to choose the order do the magnitude of the mixing angles concern the fourth generation quarks? Because there is not direct experimental measurement of the fourth generation quark mixing angles, one have to look for other indirect methods to solve the problem. Many refs. have already talked about these additional CKM mixing angles, like the vector-like quark models, the four neutrinos models and the sequential four generations models. For simple, we give a guess for the magnitude of the fourth generation mixing angles. Similar to the general CKM matrix elements magnitude order, the fourth generation ones are about $`\lambda ^4\lambda ^5`$ order between the $`1st`$ and $`4th`$ generation, such as $`V_{t^{^{}}d}`$, as well as $`\lambda ^2\lambda ^3`$ between the $`2nd`$ and $`4th`$ generation, such as $`V_{t^{^{}}s}`$. For the mixing between the $`3rd`$ and $`4th`$ generation quarks, such as $`V_{t^{^{}}b}`$, we take the magnitude as 1 because the mass of the fourth generation quark $`t^{^{}}`$ is the same order, $`10^2`$, as the top quark $`t`$. So $`V_{t^{^{}}b}`$ should take the order of $`V_{tb}`$. Then, the magnitude order of the fourth generation CKM factor $`V_{t^{^{}}d}^{}V_{t^{^{}}b}`$ is about $`\lambda ^4\lambda ^5`$, i.e. $`<\lambda ^4`$. From figs. 4, we found that the numerical results, $`V_{t^{^{}}d}^{}V_{t^{^{}}b}^{(2)}`$, satisfy this guess.
At last, the factor $`V_{t^{^{}}d}^{}V_{t^{^{}}b}`$ constrained from $`\mathrm{\Delta }M_{B_d}`$ does not contradict to the CKM matrix texture. Moreover, it seem to support the existence of the fourth generation.
## 5 Conclusion
In this note, we have investigated the mass differences $`\mathrm{\Delta }M_{B_{d,s}}`$ in the mixing $`B_{d,s}^0\overline{B}_{d,s}^0`$ with a new up-like quark $`t^{}`$ in a sequential fourth generation model. We give the basic formulae for $`\mathrm{\Delta }M_{B_{d,s}}`$ in this model and calculated the new Wilson coefficients $`S_0(m_t^{^{}}/m_W)`$ and $`S_0(m_t^{^{}}/m_W)`$ in the effective $`\mathrm{\Delta }S=2`$ Hamiltonian We also calculated the short distance QCD factors $`\eta _t^{^{}}`$ and $`\eta _{tt^{^{}}}`$. With these values, first, we analysis the final numerical results of $`\mathrm{\Delta }M_{B_s}`$, which is the function of $`m_t^{}`$ through a fourth generation CKM factor $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$. This factor was constrained from the rare decay $`BX_s\gamma `$ and had two kinds of numerical results. We found the new results in SM4 were almost same as in SM with one case of $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$ numerical values and satisfied with the present experimental low-bound of $`\mathrm{\Delta }M_{B_s}`$ in other case. We gave the figs of $`\mathrm{\Delta }M_{B_s}`$ to $`m_t^{^{}}`$ and the numerical analysis. Second, we investigated the mass difference $`\mathrm{\Delta }M_{B_d}`$ to get the constraint of the other fourth generation CKM factor $`V_{t^{^{}}b}^{}V_{t^{^{}}d}`$, (which is the function of $`t^{^{}}`$), from the experimental measurements of $`\mathrm{\Delta }M_{B_d}`$. We first got the constraints of the general CKM factor $`V_{tb}^{}V_{td}`$ from the CKM unitarity triangle. Then we give the figs. of $`V_{t^{^{}}b}^{}V_{t^{^{}}d}`$ to the mass of $`t^{^{}}`$ and to the general factor $`V_{td}^{}V_{tb}`$ and the numerical analysis too. We also talked about the texture of the fourth generation CKM matrix.
The fourth generation quark $`t^{^{}}`$ will give obviously new effects on the mass difference $`\mathrm{\Delta }M_{B_{d,s}}`$ if it really exists. At least, the present experimental statue of $`\mathrm{\Delta }M_{B_{d,s}}`$ could not exclude the space of the fourth generation, Furthmore, the progress of theoretical calculation and experimental measurement $`\mathrm{\Delta }M_{B_{d,s}}`$ could provide the strong test of the existence of the fourth generation. In other words, as one of the directions beyond the SM, $`\mathrm{\Delta }M_{B_{d,s}}`$ could provide a possible test of the fourth generation or perhaps a signal of the new physics.
## Acknowledgments
This research was supported in part by the National Nature Science Foundation of China. I am grateful to prof. C.S. Huang and Prof. Y.L. Wu for useful discussions and valuable modification andcomments on the manusript.
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# Proposition
## References.
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James N. Brawner, The Gaussian-Wahl Map for Tetragonal Curves. Preprint 1996.
Ciro Ciliberto and Rick Miranda, Gaussian maps for certain families of canonical curves. In: Complex Projective Geometry. London Math. Soc. Lecture Note Ser. 179, Cambridge Univ. Press, Cambridge, 1992, pp. 106–127.
Ciro Ciliberto and Guiseppe Pareschi, Pencils of minimal degree on curves on a $`K3`$ surface. J. Reine Angew. Math. 460 (1995), 15–36.
Ron Donagi and David R. Morrison, Linear systems on $`K3`$-sections. J. Differential Geom. 29 (1989), 49–64.
R. Drewes and J. Stevens, Deformations of Cones over Canonical Trigonal Curves. Abh. Math. Sem. Univ. Hamburg 66 (1996), 289–315.
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D. H. J. Epema, Surfaces with Canonical Hyperplane Sections. Thesis, Leiden 1983. Also: CWI Tract 1, Centrum voor Wiskunde en Informatica, Amsterdam, 1984.
Miles Reid, Surfaces with $`p_g=3,K^2=4`$ according to E. Horikawa and D. Dicks. Text of a lecture, Univ. of Utah and Univ. of Tokyo 1989.
Miles Reid, Chapters on Algebraic Surfaces. In: Complex algebraic geometry (Park City 1993), IAS/Park City Math. Ser. 3, Amer. Math. Soc., 1997, pp. 3–159.
Karl Rohn, Ueber die Flächen vierter Ordnung mit dreifachem Punkte. Math. Ann. 24 (1884), 55–151.
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Jan Stevens, Deformations of cones over hyperelliptic curves. J. Reine Angew. Math. 473 (1996), 87–120.
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Address of the author:
Jan Stevens
Matematik
Göteborgs universitet
Chalmers tekniska högskola
SE 412 96 Göteborg, Sweden
e-mail: stevens@math.chalmers.se
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# Appendix A: Mathematical Preliminaries
## Appendix A: Mathematical Preliminaries
In this appendix we list certain mathematical relations that we make use of later in the paper.
The following Gaussian integrals play a central role in our technique for studying mixed systems:
For a $`2n\times 2n`$, antisymmetric matrix $`M_{}`$, we have the relation
$$D\phi _{}e^{(\phi _{},M_{}\phi _{})}=\sqrt{\text{det}M_{}},$$
(A.1)
where $`\phi _{}`$ denotes a vector whose elements $`\phi _j`$ are Grassmannian variables, and
$$D\phi _{}:=\underset{j}{}\frac{d\phi _j}{\sqrt{2}}.$$
We use the convention that
$$𝑑\phi _j\phi _j=1.$$
For a symmetric matrix $`M_+`$, with a positive definite real part, we have the relation
$$D\phi _+e^{(\phi _+,M_+\phi _+)}=\left[\text{det}M_+\right]^{1/2},$$
(A.2)
where $`\phi _+`$ denotes a vector whose elements $`\phi _{+j}`$ are real variables and $`D\phi _+:=_j(1/\sqrt{\pi })d\phi _{+j}`$.
A similar identity for integrals over pairs of complex variables $`\phi _{+j},\overline{\phi }_{+j}`$, which is valid for any matrix $`M`$ such that $`(M+M^{})`$ is positive definite, is
$$𝒟\phi _+e^{(\overline{\phi }_+,M\phi _+)}=\left[\text{det}M\right]^1,$$
(A.3)
where
$$𝒟\phi _+=\underset{j}{}\frac{d\overline{\phi }_{+j}d\phi _{+j}}{2\pi i}.$$
The determinant of any matrix $`M`$ can be expressed as an integral over Grassmannian variables as follows
$$𝒟\phi _{}e^{(\overline{\phi }_{},M\phi _{})}=\text{det}M_{},$$
(A.4)
where $`\overline{\phi }_{}`$ and $`\phi _{}`$ denote vectors whose elements, denoted by $`\overline{\phi }_j`$ and $`\phi _j`$, are Grassmannian variables; the measure is given by $`𝒟\phi _{}=D\overline{\phi }_{}D\phi _{}`$, with $`D\overline{\phi }_{}:=_jd\psi _j`$ and $`D\phi _{}:=_jd\phi _j`$.
The following identities involving determinants and traces of matrices are also used frequently: Let $`A`$ be an $`(n\times n)`$ matrix and $`B`$ be an $`(m\times m)`$ matrix. Then
* $$\text{det}(AB)=(\text{det}A)^m(\text{det}B)^n$$
(A.5)
* $$\text{det}(A1\mathrm{I}_m+1\mathrm{I}_nB)=\underset{\alpha =1}{\overset{n}{}}\underset{j=1}{\overset{m}{}}\left(a_\alpha +b_j\right)=\underset{\alpha =1}{\overset{n}{}}\text{det}\left(a_\alpha 1\mathrm{I}_m+B\right)=\underset{j=1}{\overset{m}{}}\text{det}\left(A+b_j1\mathrm{I}_n\right),$$
(A.6)
where $`a_\alpha `$ and $`b_j`$ denote the eigenvalues of the matrices $`A`$ and $`B`$ respectively. The symbol $`1\mathrm{I}_j`$ is used to denote the $`j\times j`$ identity matrix.
* If $`\mathrm{tr}A^j=\mathrm{tr}B^j`$ for any arbitrary integer $`j`$, then
$$\text{det}(1\mathrm{I}_n+A)=\text{det}(1\mathrm{I}_m+B).$$
(A.7)
## Appendix B: Properties of the matrix $`Q_{}`$
The $`2n\times 2n`$ matrix $`Q_{}`$ is self-adjoint and satisfies the property
$$Q_{}^t=\mathrm{{\rm Y}}_{}Q_{}\mathrm{{\rm Y}}_{}$$
(B.1)
where
$$\mathrm{{\rm Y}}=\left(\begin{array}{cc}0\hfill & \hfill 1\mathrm{I}_n\\ 1\mathrm{I}_n\hfill & \hfill 0\end{array}\right)$$
We will show that such matrices have doubly degenerate eigenvalues and can be diagonalised by unitary matrices, which also belong to the symplectic group. More precisely
$$Q_{}=Uq_{}U^{},$$
(B.2)
where
$$UU^{}=U^{}U=1$$
(B.3)
$$U^t\mathrm{{\rm Y}}_{}U=\mathrm{{\rm Y}}_{}$$
(B.4)
and $`q_{}`$ is the diagonal matrix of eigenvalues of $`Q_{}`$. Consider the spectral decomposition of $`Q_{}`$:
$$Q_{}=\underset{j=1}{\overset{n}{}}\lambda _jP_j.$$
(B.5)
Then the projectors $`P_j`$ also have the property (B.1), which implies that
$$P_j(1\alpha |1\beta )=P_j(2\beta |2\alpha ).$$
hence,
$$\mathrm{tr}P_j=2\underset{\alpha =1}{\overset{n}{}}P_j(1\alpha |1\alpha ),$$
(B.6)
which implies that for almost all $`Q_{}`$ (with respect to the Lebesgue measure), $`TrP_j=2`$, and, therefore, the eigenvalues of $`Q_{}`$ are doubly degenerate.
On the other hand, (B.1) and (B.2) can be expressed as
$$Mq_{}=q_{}M,$$
(B.7)
where
$$M=U^{}\mathrm{{\rm Y}}_{}(U^{})^t.$$
(B.8)
However,
$$M^t=M\mathrm{and}M^{}M=1\mathrm{I}_{2n},$$
(B.9)
and from (B.7) it follows that
$$M(p\alpha |q\beta )=\delta _{\alpha ,\beta }e^{i\psi _\alpha }[\delta _{p,1}\delta _{q,2}\delta _{p,2}\delta _{q,1}]$$
(B.10)
Noting, however, that $`U`$ can be replaced by $`Ue^{i\phi }`$ in (B.2), where $`e^{i\phi }`$ is a diagonal matrix, we see from (B.8) that we can choose $`\psi _\alpha =0`$ in (B.10). In other words
$$M=\mathrm{{\rm Y}}_{}$$
which gives the desired property (B.4) for $`U`$.
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# Can Rotational Properties of General Relativistic Compact Objects Be Predicted from Static Ones?
## I INTRODUCTION
Rotation is a basic physical property of astrophysical objects. Study of rotational properties of stars can lead to some restrictions imposed on the nuclear equation of state of dense matter at densities larger than nuclear density.Oscillations of rapidly rotating stars can become unstable hence producing detectable gravitational wave emissions.
In the past decades there were numerous attempts to construct analytic models for rotating perfect fluid bodies in general relativity. But obtaining an exact interior solution for a rotating body proved to be a formidable task.The first step in this direction was the fundamental result of Kerr who obtained the solution for the vacuum domain, outside the rotating star. It took no less than three decades of investigations to obtain the first, highly idealized model of a general relativistic, thin rotating disk of dust, by Neugebauer and Meinel . Various schemes have been developed for obtaining stationary and axisymmetric perfect fluid solutions of the gravitational equations like Petrov type D, local rotational symmetry, fluid kinematics, non-trivial Killing tensor, vanishing Simon tensor, electric magnetic Weyl curvature, lagrangian or static- stationary symmetry, geodesic eigenrays etc. (for a recent review of rotating perfect fluid models in general relativity see ).
If in the case of rapidly rotating stellar configurations there are still many unsolved problems, a remarkable progress has been made in the study of slowly and rigidly rotating perfect fluid configurations. By casting the metric in the form
$$ds^2=e^\nu dt^2e^\lambda dr^2r^2\left[d\theta ^2+\mathrm{sin}^2\theta \left(d\varphi \omega dt\right)^2\right]+O\left(\mathrm{\Omega }^2\right),$$
(1)
Hartle obtained a formalism that proved to be very useful in many investigations of the rotational properties of stars. However, this model, which considers only first order corrections in $`\mathrm{\Omega }`$, can not be used to compute models of rapidly rotating relativistic stars with sufficient accuracy.
On the other hand there has been recently a considerable advance in the numerical understanding of rotating stars. Several high precision numerical codes are now avalaible and it has been shown that they agree with each other to remarkable accuracy (see for a review of recent developments in numerical study of rotatation, nonaxisymmetric oscillations and instabilities of general relativistic stars).
The non-sphericity of rapidly rotating stationary stellar configurations and the complicated character of the interplay of the effects of rotation and of those of general relativity seem not to permit a simple universal description of rotating compact objects. However, Haensel and Zdunik and Friedman, Ipser and Parker have found a simple relation connecting the maximum rotation frequency $`\mathrm{\Omega }_{max}`$ with the maximum mass $`M_{max}^{stat}`$ and radius $`R_{max}^{stat}`$ of the static configuration:
$$\mathrm{\Omega }_{max}=C_S\left(\frac{M_{max}^{stat}}{M_{}}\right)^{\frac{1}{2}}\left(\frac{R_{max}^{stat}}{10km}\right)^{\frac{3}{2}},$$
(2)
with $`C_S`$ a constant which does not depend on the equation of state of the dense matter. The value of the constant $`C_S`$ has been obtained by fitting equation (2) with data obtained by numerically integrating the gravitational field equations. It is given by $`C_S=7200s^1`$ , $`C_S=7700s^1`$ or by $`C_S=7730s^1`$ . The empirical relation (2) has been checked for many realistic equations of state for neutron stars , - and for the case of strange stars where the empirical formula also holds with a very good precision, the relative deviations do not exceeding 2% . Equation (2), obtained on the basis of analyzing numerical solutions of the gravitational field equations, provides an enormous simplification of the problem of dynamical effects of rotation because the solutions of the complicate general relativistic equations for a rotating star can be replaced with the solutions of the much simpler TOV equation. Some attempts to explain this empirical relation were not concluded with satisfactory result. Weber and Glendenning , used numerical models of slowly rotating relativistic stars to show that the formula still hold, but with $`C_S=7500s^1`$. Also in the slow rotation limit Glendenning and Weber derived a formula relating $`\mathrm{\Omega }_{max}`$ to $`M_{max}/R_{max}^3`$, in terms of the mass, equatorial radius, moment of inertia, angular momentum and quadrupole moment of the maximally rotating configuration only. But it is not clear how the formula (2) follows from their results.Up to now a clear physical understanding of this relation is still missing.
On the other hand a universal relation between the maximum mass and radius of non-rotating neutron star configuration and the mass $`M_{max}^{rot}`$ and radius $`R_{max}^{rot}`$ of the configuration rotating with $`\mathrm{\Omega }_{max}`$ of the form
$$M_{max}^{rot}=C_MM_{max}^{stat},R_{max}^{rot}=C_RR_{max}^{stat},$$
(3)
has also been found . In (3) $`C_M`$ and $`C_R`$ are specific equation of state dependent constants, whose values have been calculated, for a broad set of realistic EOS, in . Their mean values are $`C_M=1.18`$ and $`C_R=1.34`$ . The empirical constant $`C_S`$ can be obtained, within an approximation better than $`5\%`$ , from the formula
$$C_S\left(\frac{C_M}{C_R^3}\right)^{1/2}.$$
(4)
The general validity of (2) and (3) suggests the possibility that all basic physical parameters of general relativistic rotating stellar objects (like mass and radius) can be somehow related to the similar parameters of the static configuration. It is the purpose of the present paper to propose a general description of basic physical parameters (mass, radius and oblateness) of compact general relativistic objects in terms of the physical properties of the static configuration and of the angular velocity only. To obtain the representation in terms of static physical parameters we use only the general relativistic conditions of equilibrium for static and rotating stars and some phenomenological assumptions.The resulted approximate mass and radius formulae are compared with data obtained from numerical integration of gravitational field equations in case of neutron stars described by realistic equations of state and strange stars described by the bag model equation of state.
The present paper is organized as follows.The general formalism allowing to obtain mass and radius formulae for rotating general relativistic stars is presented in Section 2. In Section 3 we apply our results to the case of neutron stars described by realistic equations of state. In Section 4 we consider the case of strange stars. We discuss and conclude our results in Section 5.
## II THE GENERAL FORMALISM
For a static equilibrium stellar type configuration, with interior described by the metric
$$ds^2=e^{2\nu (r)}dt^2e^{\lambda (r)}dr^2r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right),$$
(5)
the condition of the hydrostatic equilibrium, which follows from the Bianchi identities, can be written as
$$2\nu +2\frac{dp}{\rho +p}=const.,$$
(6)
with $`\rho `$ and $`p`$ the energy density and the pressure of the matter respectively ( in the present paper we shall use units so that $`c=G=1`$).
We shall also make the assumption that the matter EOS is a one parameter dependent function,
$$\rho =\rho (n),p=p(n),$$
(7)
with $`n`$ the proper baryon density. For such an equation of state the heat function defined by
$$H(n)=\frac{1}{\rho (n)+p(n)}\frac{dp(n)}{dn},$$
(8)
is a regular function of $`n`$.It can always be written in the form
$$H(n)=\mathrm{ln}f(n),$$
(9)
where
$$f(n)=\frac{\rho (n)+p(n)}{n},$$
(10)
is the enthalpy per baryon.
As applied at the center and at the surface of the star respectively, the hydrostatic equilibrium condition yields
$$\left(2\nu +2\mathrm{ln}f(n)\right)_{center}=\left(2\nu +2\mathrm{ln}f(n)\right)_{surface}.$$
(11)
At the vacuum boundary of the static star the Schwarzschild exterior solution gives the metric and we have
$$2\nu _S=\mathrm{ln}\left(1\frac{2M_{stat}}{R_{stat}}\right),$$
(12)
with $`M_{stat}`$ and $`R_{stat}`$ the mass and radius of the static stellar configuration. We denote by $`n_C`$ and $`n_S`$ the baryon density at the center of the star and at the surface, respectively. Therefore from (11) and (12) we obtain the following general and exact expression for the mass-radius ratio of the static star:
$$\frac{M_{stat}}{R_{stat}}=\frac{1}{2}\left[1C\left(n_C\right)F(n_C,n_S)\right],$$
(13)
where we have denoted
$$F(n_C,n_S)=\frac{f^2\left(n_C\right)}{f^2\left(n_S\right)},$$
(14)
and
$$C\left(n_C\right)=\mathrm{exp}\left(\nu _C\right)=g_{00}_{r=0}.$$
(15)
For a given equation of state $`C\left(n_C\right)`$ is a function of the central density only. From a physical point of view $`C\left(n_C\right)`$ can be related via the relation $`\frac{1}{C}1=z_c`$ to the redshift $`z_c`$ of a photon emitted from the center of the star.
For the mass of the star we can obtain another general representation by assuming
$$M_{stat}=\frac{4\pi \rho _{spec}}{3}R_{stat}^3A\left(n_C\right),$$
(16)
where $`\rho _{spec}`$ is a specific density that can be arbitrarily chosen (for example it is the central density of the minimum mass configuration) and $`A\left(n_C\right)`$ is a function describing the effects of the variation of the central density on the basic parameters of the stellar configuration.
The two unknown functions $`C\left(n_C\right)`$ and $`A\left(n_C\right)`$ can be determined by fitting equations (13) and (16) with the exact values of the mass and radius obtained by numerically integrating the TOV equation for a given equation of state. Their knowledge allows us to construct exact mass and radius formulae for sequences of static general relativistic stars having different central densities. From equations (13) and (16) we obtain the radius and the mass of the star in the form
$$R_{stat}=\left(\frac{3}{16\pi }\frac{1}{A\left(n_C\right)\rho _{spec}}\right)^{1/2}\left[1C\left(n_C\right)F(n_C,n_S)\right]^{1/2},$$
(17)
$$M_{stat}=\frac{4\pi }{3}\rho _{spec}A\left(n_C\right)\left(\frac{3}{16\pi }\frac{1}{A\left(n_C\right)\rho _{spec}}\right)^{3/2}\left[1C\left(n_C\right)F(n_C,n_S)\right]^{3/2}.$$
(18)
Solving the equations $`\frac{dR_{stat}}{dn_C}=0`$ and $`\frac{dM_{stat}}{dn_C}=0`$ for the value of the central density $`n_C`$ will give, with the use of (17) and (18) the maximum values of the radius and mass of the static star, respectively.
To describe the interior of the rotating general relativistic star we shall adopt the formalism presented in and . Under the hypothesis of stationarity, axial symmetry and purely azimuthal motion a coordinate system $`(t,r,\theta ,\varphi )`$ can be chosen so that inside the star the line element takes the form
$$ds^2=N^2dt^2B^2r^2\mathrm{sin}^2\theta \left(d\varphi N^\varphi dt\right)^2A^2\left(dr^2+r^2d\theta ^2\right),$$
(19)
where $`N`$,$`N^\varphi `$,$`A`$ and $`B`$ are functions of $`r`$ and $`\theta `$ only. As measured by the locally non-rotating observer the fluid 3-velocity is given by $`U=Br\mathrm{sin}\theta \frac{\mathrm{\Omega }N^\varphi }{N}`$ , where $`\mathrm{\Omega }=\frac{d\varphi }{dt}`$ is the angular velocity of a fluid element moving in the $`\varphi `$-direction (physically it is the angular velocity as measured by an observer at spatial infinity) ,.
From the point of view of our phenomenological approach the most important result is the equation of the stationary motion, which results from the Bianchi identities and which for a rotating perfect fluid reduces to
$$\frac{1}{\rho +p}\frac{p}{x^i}+\frac{\nu }{x^i}\frac{}{x^i}\left(\mathrm{ln}\mathrm{\Gamma }\right)=S\frac{\mathrm{\Omega }}{x^i},$$
(20)
where $`S=\mathrm{\Gamma }^2\frac{A^2B}{N}Ur\mathrm{sin}\theta `$ , $`\nu =\mathrm{ln}N`$ and $`\mathrm{\Gamma }=\left(1U^2\right)^{1/2}`$. If $`\mathrm{\Omega }=const.`$ (case called uniform or rigid rotation) equation (20) can be integrated to give the following fundamental result describing the stationary equilibrium of a rotating general relativistic star:
$$\left(2\nu +2\mathrm{ln}f(n)+\mathrm{ln}\left(1U^2\right)\right)_{center}=\left(2\nu +2\mathrm{ln}f(n)+\mathrm{ln}\left(1U^2\right)\right)_{surface}.$$
(21)
Equation (21) is just the generalization to the case of rotation of the well-known static Bianchi identity $`2\nu +2\frac{dp}{\rho +p}=const.`$ we have already used to describe static stellar configurations.
We assume that the vacuum boundary of the rotating star is described by the Kerr metric in the Boyer-Lindquist coordinates, which has the form :
$`ds^2`$ $`=`$ $`\left(1{\displaystyle \frac{2M_{rot}r}{r^2+a^2\mathrm{cos}^2\theta }}\right)dt^2\left({\displaystyle \frac{r^2+a^2\mathrm{cos}^2\theta }{r^2+a^22GM_{rot}r}}+r^2+a^2\mathrm{cos}^2\theta \right)dr^2`$ (23)
$`(r^2+a^2\mathrm{sin}^2\theta )d\theta ^2\left[(r^2+a^2)\mathrm{sin}^2\theta +{\displaystyle \frac{2M_{rot}ra^2\mathrm{sin}^4\theta }{r^2+a^2\mathrm{cos}^2\theta }}\right]d\varphi ^22{\displaystyle \frac{2M_{rot}ra\mathrm{sin}^2\theta }{r^2+a^2\mathrm{cos}^2\theta }}dtd\varphi .`$
In this form the Kerr metric is manifestly axially symmetric and closely resembles the Schwarzschild solution in its standard form.$`M_{rot}`$ is the mass of the source and the parameter $`a=\frac{J}{M_{rot}}`$ is the ratio between the angular momentum $`J`$ and the mass$`M_{rot}`$ of the rotating star.
Let’s apply equation (21) at two points: at the center of the dense core and at the pole of the rotating star. We denote by $`C_{rot}^{(1)}`$ the value of the metric tensor component $`g_{00}`$ at the center of the star. At the polar point $`\theta =0`$ and $`r=R_p=const.`$, where $`R_p`$ is the polar radius of the star. We also have $`U=0`$ . Therefore at this point the line element is given by
$$ds_{pol}^2=\left[1\frac{2M_{rot}}{R_p}\left(1+a^2/R_p^2\right)^1\right]dt^2.$$
(24)
Consequently from equations (21) and (24) we obtain the following exact mass-polar radius relation for the rotating general relativistic configuration:
$$\frac{M_{rot}}{R_p}=\frac{1}{2}\left(1+a^2/R_p^2\right)^1\left[1C_{rot}^{(1)}(n_C,\mathrm{\Omega },\mathrm{})F(n_C,n_S)\right].$$
(25)
The function $`F(n_C,n_S)`$ is identical to that of the static case. Let’s apply now equation (21) for two points situated in the equatorial plan of the rotating star: at the center of the star and at the equator respectively. At the equator $`r=R_e=const.`$ and $`\theta =\frac{\pi }{2}`$ ($`R_e`$ is the equatorial radius of the star). For a uniform rotation the rotation angle of the source/observer at the equator is $`\varphi =\mathrm{\Omega }t`$ . Taking into account these results we obtain the Kerr metric at the equator of the rotating star in the form:
$$ds_{eq}^2=\left[1\frac{2M_{rot}}{R_e}\left(1+a\mathrm{\Omega }\right)^2R_e^2\mathrm{\Omega }^2\left(1+\frac{a^2}{R_e^2}\right)\right]dt^2.$$
(26)
With the use of (21) and (26) we obtain the following mass-equatorial radius relation:
$$\frac{M_{rot}}{R_e}=\frac{1}{2}\left[\frac{1C_{rot}^{(2)}(n_C,\mathrm{\Omega },R_e)F(n_C,n_S)R_e^2\mathrm{\Omega }^2\left(1+\frac{a^2}{R_e^2}\right)}{\left(1+a\mathrm{\Omega }\right)^2}\right],$$
(27)
where a new function
$$C_{rot}^{(2)}(n_C,\mathrm{\Omega },R_e)=C_{rot}^{(1)}(n_C,\mathrm{\Omega },R_e)/\left(1U^2\right)_{r=R_e},$$
(28)
has also been defined.
From Eqs. (25) and (27) we obtain the ratio $`e`$ of the polar and equatorial radius (the oblateness) of the star in the form:
$$e=\frac{R_p}{R_e}=\frac{1+\frac{a^2}{R_p^2}}{(1+a\mathrm{\Omega })^2}\frac{1C_{rot}^{(2)}(n_C,\mathrm{\Omega },R_e)F(n_C,n_S)R_e^2\mathrm{\Omega }^2\left(1+\frac{a^2}{R_e^2}\right)}{1C_{rot}^{(1)}(n_C,\mathrm{\Omega },R_e)F(n_C,n_S)}.$$
(29)
In the case of the static star we have also proposed the alternative Eq.(16) for providing another mass-radius relation. We shall generalize this equation to the rotating case by assuming that the following formula relates the mass of the rotating star to its equatorial radius:
$$M_{rot}=\frac{4\pi \rho _{spec}}{3}R_e^3C_{rot}^{(3)}(n_C,\mathrm{\Omega },R_e),$$
(30)
with $`C_{rot}^{(3)}(n_C,\mathrm{\Omega },R_e)`$ a function describing the combined general relativistic effects of rotation and central density on the mass of the star and with $`\rho _{spec}`$ the same specific density as used in the static case. For $`\mathrm{\Omega }=0`$, $`R=R_e`$ and Eq.(30) must reduce to the static case equation (16).
Equations (27),(29) and (30) give a complete and exact description of the mass and radius of the rotating general relativistic star. But unfortunately the present approach, which is basically thermodynamic in its essence, can not predict the exact form and the values of the three unknown functions entering in the formalism. The only thing we can do is to assume, also based on the static case, some empirical forms for the functions $`C_{rot}^{(i)}(n_C,\mathrm{\Omega },R_e),i=1,2,3`$ and to check whether the resulting formulae can give a satisfactory description of rotating star configurations. Therefore in the following we shall use the following five approximations:
i) As a first approximation we shall define the moment of inertia $`I_{rot}`$ of the rotating compact relativistic object via the Newtonian expression
$$I_{rot}=\frac{2}{5}M_{rot}R_e^2.$$
(31)
In fact over a wide range of $`M_{rot}`$ and $`R_e`$ the corrections added by general relativistic effects to the moment of inertia can be approximated by $`I_{rot}=0.21\frac{M_{rot}R_e^2}{1\frac{2M_{rot}}{R_e}}`$ , but we shall not use this representation. By adopting the Newtonian formula we obtain
$$a=\frac{J}{M_{rot}}=\frac{2}{5}R_e^2\mathrm{\Omega }.$$
(32)
ii) We assume that the function $`C_{rot}^{(1)}(n_C,\mathrm{\Omega },R_e)`$, describing the metric tensor component $`e^\nu `$ at the center of the rotating star is given by
$$C_{rot}^{(1)}(n_C,\mathrm{\Omega },R_e)=\left[\left(1\frac{1}{2}g\left(\frac{M_{stat}}{R_{stat}}\right)R_e^2\mathrm{\Omega }^2\right)\right]C\left(n_C\right),$$
(33)
with $`g\left(\frac{M_{stat}}{R_{stat}}\right)`$ an EOS dependent function given by
$$g\left(\frac{M_{stat}}{R_{stat}}\right)=\frac{\alpha }{1\frac{2M_{stat}}{R_{stat}}},$$
(34)
and $`\alpha `$ a non-negative constant. $`C\left(n_C\right)`$ is also the function corresponding to the static case.
iii) We assume that inside the rotating star $`C_{rot}^{(2)}(n_C,\mathrm{\Omega },R_e)`$ is independent of the angular velocity $`\mathrm{\Omega }`$ of the rotating compact object and can be represented by the function corresponding to the static case:
$$C_{rot}^{(2)}(n_C,\mathrm{\Omega },R_e)=C\left(n_C\right).$$
(35)
iv) We suppose that
$$C_{rot}^{(3)}(n_C,\mathrm{\Omega },R_e)=A\left(n_C\right)e.$$
(36)
$`A\left(n_C\right)`$ is again the function corresponding to the static case.
With these four phenomenological assumptions Eqs. (27)-(30) lead to the following representation of the basic physical parameters of the rotating general relativistic star:
$$R_e=\left(\frac{M_{stat}}{R_{stat}}\right)^{1/2}\left[\frac{M_{stat}}{R_{stat}^3}\frac{\alpha }{4}\mathrm{\Omega }^2\right]^{1/2},$$
(37)
$$e=\frac{\left(1+\frac{4R_e^2\mathrm{\Omega }^2}{25}\frac{1}{e^2}\right)}{\left(1\frac{2}{5}R_e^2\mathrm{\Omega }^2\right)^2}\frac{\frac{2M_{stat}}{R_{stat}}R_e^2\mathrm{\Omega }^2\left(1+\frac{4}{25}R_e^2\mathrm{\Omega }^2\right)}{1\left(1\frac{2M_{stat}}{R_{stat}}\right)\left(1\frac{1}{2}\frac{\alpha }{1\frac{2M_{stat}}{R_{stat}}}R_e^2\mathrm{\Omega }^2\right)},$$
(38)
$$M_{rot}=\frac{4\pi \rho _{spec}}{3}A\left(n_C\right)R_e^3e=M_{stat}\left(\frac{R_e}{R_{stat}}\right)^3e.$$
(39)
We shall also suppose that $`\alpha `$ is a universal constants and we shall choose $`\alpha 5`$.In this formulation of the general relativistic problem of the rotation the oblateness parameter $`e`$ of the star is given by the roots of the third order algebraic equation (38).
The equatorial radius is defined only for values of the angular velocity satisfying the condition $`\frac{\alpha }{4}\mathrm{\Omega }^2\frac{M_{stat}}{R_{stat}^3}`$.Therefore for the maximum admissible constant angular velocity of the maximally rotating star in uniform rotation we obtain the relation
$$\mathrm{\Omega }C\left(\frac{M_{max}^{stat}}{M_{}}\right)^{1/2}\left(\frac{R_{max}^{stat}}{10km}\right)^{3/2},$$
(40)
where
$$C=\frac{2}{\sqrt{\alpha }}\sqrt{G}\frac{\left(M_{}\right)^{1/2}}{10^9}=10330.489s^1,$$
(41)
for $`\alpha =5`$.
Equation (40) is very similar to the ”empirical” formula discussed in -. The coefficient of proportionality in (40) is independent of the equation of state of the dense matter, but its numerical value does not fit the calculated value.On the other hand for $`\mathrm{\Omega }=C\left(\frac{M_{max}^{stat}}{M_{}}\right)^{1/2}\left(\frac{R_{max}^{stat}}{10km}\right)^{3/2}`$ the radius of the star tends to infinity.
An alternative expression can be obtained by imposing the restriction
$$R_e^{max}\mathrm{\Omega }_{max}\beta \sqrt{\frac{M_{max}^{rot}}{R_{max}^{rot}}}\sqrt{\frac{C_MM_{max}^{stat}}{C_RR_{max}^{stat}}}\sqrt{\frac{M_{max}^{stat}}{M_{max}^{stat}}},$$
(42)
where $`\beta `$ is the maxium allowable equatorial speed of the star and we have also used the Newtonian force balance equation between the gravitational and centrifugal force. Therefore we obtain
$$\mathrm{\Omega }_{max}C\left(\frac{M_{max}^{stat}}{M_{}}\right)^{1/2}\left(\frac{R_{max}^{stat}}{10km}\right)^{3/2}\left[1+\frac{4M_{max}^{stat}}{\alpha \beta ^2R_{max}^{stat}}\right]^{1/2},$$
(43)
It is interesting to note that the values of $`2M_{max}^{stat}/R_{max}^{stat}`$ are in a narrow range of (0.467, 0.667) , . Taking for $`2M_{max}^{stat}/R_{max}^{stat}`$ a mean value of 0.58 it follows that $`\beta 0.54`$ and we find
$$\mathrm{\Omega }_{max}C_S\left(\frac{M_{max}^{stat}}{M_{}}\right)^{1/2}\left(\frac{R_{max}^{stat}}{10km}\right)^{3/2},$$
(44)
with $`C_S=7708s^1`$, value which coincides with that proposed in and differs only within $`1\%`$ from the value $`C_S=7840s^1`$ obtained in .
By taking for twice of the ratio of the maximum static mass and radius a mean value of 0.58 and considering $`\beta 0.54`$, we obtain
$$\mathrm{\Omega }_{max}C_S\left(\frac{M_{max}^{stat}}{M_{}}\right)^{1/2}\left(\frac{R_{max}^{stat}}{10km}\right)^{3/2},$$
(45)
with $`C_S=7708s^1`$,value wich coincides with the value obtained in that differs within $`5\%`$ from the value $`C_S=7840s^1`$ obtained in .
The maximum radius of the maximally rotating configuration can be obtained from $`R_{max}^{rot}\frac{\beta }{\mathrm{\Omega }_{max}}`$, and is given by
$$R_{max}^{rot}C_R\left(\frac{M_{max}^{stat}}{R_{max}^{stat}}\right)R_{max}^{stat},$$
(46)
where
$$C_R\left(\frac{M_{\mathrm{max}}^{stat}}{R_{\mathrm{max}}^{stat}}\right)=\frac{\beta }{10km}\frac{1}{C_S}\sqrt{\frac{M_{}}{10km}}\left(\frac{M_{\mathrm{max}}^{stat}}{R_{\mathrm{max}}^{stat}}\right)^{1/2},$$
(47)
The mass $`M_{maxrad}^{rot}`$ of the maximal radius rotating neutron star follows from (39) and is given by
$$M_{maxrad}^{rot}=C_{RM}\left(\frac{M_{max}^{stat}}{R_{max}^{stat}}\right)M_{max}^{stat},$$
(48)
where
$$C_{RM}\left(\frac{M_{max}^{stat}}{R_{max}^{stat}}\right)=C_R^3\left(\frac{M_{max}^{stat}}{R_{max}^{stat}}\right)e_{max}^{rad}\left(\frac{M_{max}^{stat}}{R_{max}^{stat}}\right),$$
(49)
$$e_{max}^{rad}\left(\frac{M_{max}^{stat}}{R_{max}^{stat}}\right)=\frac{1}{\left(1+\frac{2}{5}\beta ^2\right)^2}\frac{1\left(1\beta ^2\right)^\gamma \left(1\frac{2M_{max}^{stat}}{R_{max}^{stat}}\right)\beta ^2\left(1+\frac{4}{25}\beta ^2\right)}{1\left(1\frac{2M_{max}^{stat}}{R_{max}^{stat}}\right)\left[1\frac{\alpha \beta ^2}{2}/\left(1\frac{2M_{max}^{stat}}{R_{max}^{stat}}\right)\right]}.$$
(50)
The maximum mass $`M_{max}^{rot}`$ of the rotating star can be obtained from the equation
$$\frac{M}{\mathrm{\Omega }}_{\mathrm{\Omega }=\mathrm{\Omega }_{max}}=\frac{M_{stat}}{R_{stat}}\frac{}{\mathrm{\Omega }}\left(R_e^3e\right)_{\mathrm{\Omega }=\mathrm{\Omega }_{max}}=\frac{M_{stat}}{R_{stat}}\left(3R_e^2\frac{R_e}{\mathrm{\Omega }}e+R_e^3\frac{e}{\mathrm{\Omega }}\right)_{\mathrm{\Omega }=\mathrm{\Omega }_{max}}=0,$$
(51)
leading to
$$e_{max}\left(\mathrm{\Omega }_{max}\right)=\frac{C(M_{stat},R_{stat})}{R_e^3\left(\mathrm{\Omega }_{max}\right)}=\frac{C_M\left(\frac{M_{stat}}{R_{stat}}\right)R_{stat}^3}{R_e^3\left(\mathrm{\Omega }_{max}\right)},$$
(52)
where $`C_M\left(\frac{M_{stat}}{R_{stat}}\right)`$ is a dimensionless EOS dependent function. Therefore for the maximum mass of the maximally rotating configuration we obtain
$$M_{max}^{rot}=C_M\left(\frac{M_{stat}}{R_{stat}}\right)M_{max}^{stat}.$$
(53)
Equations (46) and (53) show the existence of a proportionality relation between maximum mass and radius of the rotating and non-rotating configuration, respectively, as has already been suggested, on an empirical basis,in .
An investigation of 12 EOS performed in for 12 realistic EOS of nuclear matter led to a (mean) value of $`C_M=1.1807`$ , while the same calculation performed by us with the use of data presented in for other 14 different EOSs gives $`C_M=1.177`$ , leading to a mean value of $`C_M=1.1790`$ for the 26 considered EOSs.
Eqs. (37)-(39) lead, for $`\beta 0.54`$ and $`2M_{max}^{stat}/R_{max}^{stat}0.58`$ to values of $`e_{max}=0.365`$ and $`C_M=1.233`$.
Therefore we may conclude that the proportionality between the maximum mass of the rotating star and the maximum mass of the static configuration is universal, being with a very good approximation independent of the equation of state of dense matter. On the other hand the relation between the radius of the rotating and non-rotating configuration is EOS dependent,the coefficient of proportionality slightly decreasing with the increase of the mass-radius ratio of the static star.
## III APPLICATIONS TO NEUTRON STARS
The correct mathematical and physical modelling of millisecond pulsars can be done only in the framework of general relativistic equilibrium models for rapidly rotating neutron stars.Such models are solutions of the Einstein’s equations for the axisymmetric stationary gravitational field and they must be constructed numerically. Recently several independent numerical codes have been developed by different groups of researchers and have been used to obtain rapidly rotating neutron star models based on a variety of realistic equations of state.Hence a large amount of numerical data is now available.
In the present Section we shall apply the results of the phenomenological formalism presented in the previous Section to the case of neutron stars described by realistic equations of state. The data are selected from the paper by Cook, Shapiro and Teukolsky ,who constructed general relativistic rotating star sequences for 14 nuclear matter equations of state. Detailed data are presented only for 5 equations of state. For the sake of comparison we have used the equations of state denoted A , AU , FPS and L . At low densities all these equations of state employ the Feynman,Metropolis and Teller EOS and then join onto the Baym,Bethe and Pethick EOS up to neutron drip.The equations of state are given in a tabular form and small changes in the way the tabulated equation of state is constructed do have a small effect on the resulting neutron star model.
In Tables 1-4 we present the comparison of the basic physical parameters of rotating stars obtained, for these four equations of state, with the help of Eqs. (37)-(39) and by numerically integrating the gravitational field equations.
For $`1.4M_{}`$ sequences described by EOS A and FPS the maximum error of our prediction is around $`9\%`$. For the maximum mass normal sequences of EOS AU and L the maximum error in the predicted value of the mass does not exceed $`14\%`$ but for large angular speeds it is around $`16\%`$ for the radius of the rotating compact object.
## IV APPLICATIONS TO STRANGE STARS
It is generally believed today that strange quark matter,consisting of u-,d- and s quarks is energetically the most favorable state of quark matter. Witten suggested that there are two ways of formation of the strange matter: the quark-hadron phase transition in the early universe and conversion of neutron stars into strange ones at ultrahigh densities.In the theories of strong interactions quark bag models suppose that the breaking of physical vacuum takes place inside hadrons. As a result the vacuum energy densities inside and outside a hadron become essentially different and the vacuum pressure on a bag wall equilibrates the pressure of quarks thus stabilizing the system.
If the hypothesis of the quark matter is true, then some of neutrons stars could actually be strange stars, built entirely of strange matter ,. Caldwell and Friedman have presented arguments against the existence of strange stars. For a recent review of strange star properties see .
There are several proposed mechanisms for the formation of quark stars after galaxy formation. Strange stars are expected to form during the collapse of the core of a massive star after the supernova explosion as a result of a first or second order phase transition, resulting in deconfined quark matter .Another possibility for strange star formation is that some rapidly spinning neutron stars in low-mass X-ray binaries (LXMBs) can accrete sufficient mass to undergo a phase transition to become strange stars . In this scenario it is supposed that at the beginning of accretion the mass of the neutron star is 1.4. It has been shown that the amounts of matter accreted by 18 millisecond pulsars in binary systems exceed $`0.5M_{}`$. Hence some of the millisecond pulsars may be strange stars.Strange stars have also been proposed as sources of unusual astrophysical phenomena, e.g. soft $`\gamma `$-ray repeaters ,pulsating X-ray burster , cosmological $`\gamma `$-ray bursts , \[37-38\] etc.The mechanism of the phase transition from neutron to quark stars in low LXMBs also results in the excitation of stellar radial oscillations that can be damped by gravitational wave radiation instead of internal viscosity .The discovery of kHz quasi-periodic oscillation in LXMBs implies that the compact stellar object must have very soft equation of state, which is consistent with that of strange stars \[41-42\].
Assuming that interactions of quarks and gluons are sufficiently small the energy density $`\rho `$ and pressure $`p`$ of a quark-gluon plasma at temperature $`T`$ and chemical potential $`\mu _f`$ (the subscript f denotes the various quark flavors u,d,s etc.) can be calculated by thermal theory. Neglecting quark masses in first order perturbation theory and supposing that quarks are confined to the bag volume (in the case of a bare strange star, the boundary of the bag coincides with stellar surface), the equation of state is
$$p=\frac{(\rho 4B)}{3},$$
(54)
where $`B`$ is the difference between the energy density of the perturbative and non-perturbative QCD vacuum (the bag constant). Equation (54) is essentially the equation of state of a gas of massive particles with corrections due to the QCD trace anomaly and perturbative interactions.These are always negative , reducing the energy density at given temperature by about a factor two . In the limit $`p0`$ (at the star’s surface) we have $`\rho 4B`$. The equation of state (54) does not depend upon quark flavor number,hence it will be correct either for strange quark matter ($`m_s0`$) or for normal quark matter ($`m_s\mathrm{}`$). For any intermediate values of $`m_s`$ the state equation (54) gives the pressure with error less than 4% .Thus the equation of state of strange matter is mainly determined by the vacuum energy density $`B`$.
The bag model equation of state (54) has been the basis for the study of most of the static relativistic models of strange stars ,. Based on the numerical integration of the mass continuity and hydrostatic equilibrium TOV (Tolman-Oppenheimer-Volkoff) equations for different values of the bag constant these authors obtained a complete description of static strange stars. Using numerical methods Witten and Haensel et al. obtained the maximum gravitational mass $`M_{max}`$, the maximum baryon mass $`M_{B,max}1.66\times 10^{27}kg\times N_B`$ ($`N_B`$-the total baryon number of the stellar configuration) and the maximum radius $`R_{max}`$ of the strange star ,as a function of the bag constant, in the form ,,:
$$M_{max}=\frac{1.9638M_{}}{\sqrt{B_{60}}},M_{B,max}=\frac{2.6252M_{}}{\sqrt{B_{60}}},R_{max}=\frac{10.172km}{\sqrt{B_{60}}},$$
(55)
where $`B_{60}B/(60MeVfm^3)`$.
Colpi and Miller and Glendenning and Weber have investigated the rotational properties of strange stars in the slow rotation approximation. As far as rotational deformations are concerned, there are a number of detailed differences between the strange star models and standard neutron stars. Exact numerical calculations of rapidly rotating strange stars were done by Lattimer et al. , Gourgoulhon et al. (by using a multi-domain spectral method that enable to treat exactly the density discontinuity at the surface of strange stars) and by Stergioulas,Kluzniak and Bulik . Rotation increases maximum allowable mass of strange stars and the equatorial radius of the maximum mass configuration. Gourgoulhon et al. obtained for the maximum mass and radius of quark stars the following two exact formulae
$$M_{max}^{rot}=\frac{2.831M_{}}{\sqrt{B_{60}}},R_{eq,M_{max}}^{rot}=\frac{16.54km}{\sqrt{B_{60}}}.$$
(56)
In the present section we shall derive, by using the formalism presented in Section I, analytic mass and radius formulae for general relativistic static and rotating equilibrium strange matter configurations described by the bag model equation of state (54).We shall begin with the study of the static strange star, but presenting and alternative and physically more involved discussion of this case.
The changes caused by the general theory of relativity in the conditions of thermal equilibrium,taking into account the gravitational field of the body,are of fundamental importance.In a constant gravitational field we must distinguish the conserved energy $`E_0`$ of any small part of the stellar object from the energy $`E`$ measured by an observer situated at a given point. These two quantities are related by $`E_0=E\sqrt{g_{00}}`$ , where $`g_{00}`$ is the time component of the metric tensor. A similar change occurs in the condition of the constancy of the chemical potential throughout the star. The chemical potential is defined as the derivative of the energy with respect to the number of particles $`N`$, $`\mu =(\frac{E}{N})_{S,V}`$.Since this number is a constant for the stellar object, $`N=constant`$,for the chemical potential measured at any point inside the gravitating body we have the relation :
$$\mu \sqrt{g_{00}}=constant.$$
(57)
A similar relation also holds for the temperature $`T`$, $`T\sqrt{g_{00}}=constant.`$, since we suppose that the strange star is in thermal equilibrium .Consequently, $`\frac{\mu }{T}=constant`$ inside the compact object. Hence $`\frac{d\mu }{\mu }=\frac{dT}{T}`$.At constant volume (equal to unity) we have $`dp=sdT+nd\mu `$,where $`s`$ and $`n`$ are the entropy and number of particles in unit volume of the body,respectively. With the use of $`dT=T\frac{d\mu }{\mu }`$ and taking into account that $`\mu n+sT=\rho +p`$ we obtain the following equation relating the equilibrium chemical potential to the energy density and pressure of the star :
$$\frac{d\mu }{\mu }=\frac{dp}{\rho +p}.$$
(58)
Consider now a static equilibrium quark matter configuration satisfying the bag model equation of state (54). Let us compare the values of the chemical potential $`\mu =\mu _f`$ at two points:at the center of the star and at the vacuum boundary. From equation (57) we obtain
$$\left[\mu \sqrt{g_{00}}\right]_{center}=\left[\mu _S\sqrt{g_{00}}\right]_{boundary},$$
(59)
where the indices $`C`$ and $`S`$ refer to the center and to surface of quark star respectively.At the vacuum boundary the gravitational field of the strange star is described by the Schwarzschild solution, which gives :
$$g_{00}=1\frac{2M_{stat}}{R_{stat}},$$
(60)
where $`M_{stat}`$ and $`R_{stat}`$ are the total mass and radius of the static strange star,respectively.At the center of the star the time component of the metric tensor has a constant value (this also follows from the Bianchi identity $`\frac{d}{dr}\mathrm{ln}(g_{00})=\frac{2}{ϵ+p}\frac{dp}{dr}`$ and we denote
$$g_{00}_{center}=C(\rho _C,B).$$
(61)
From a physical point of view $`C(\rho _C,B)`$ can be related via the relation $`\frac{1}{C}1=z_c`$ to the redshift $`z_c`$ of a photon emitted from the center of the quark star. For a given static strange matter configuration the value of $`C`$ depends only on the central density of the quark star $`\rho _C`$ and on the bag constant. Therefore from equation (59) we obtain
$$\frac{\mu _C^2}{\mu _S^2}=\frac{1}{C(\rho _C,B)}\left(1\frac{2M_{stat}}{R_{stat}}\right).$$
(62)
With the use of the bag model equation of state (54) we can integrate equation (58) to obtain
$$\mu =C_0(\rho B)^{1/4}.$$
(63)
The integration constant $`C_0`$ can be determined by calculating the chemical potential at the center of the quark star. Hence we obtain
$$C_0=\frac{\mu _C}{(\rho _CB)^{1/4}}.$$
(64)
The variation of the chemical potential inside the quark star can be represented as
$$\mu =\frac{\mu _C}{(\rho _CB)^{1/4}}(\rho B)^{1/4}.$$
(65)
At the surface of the star $`\rho _S4B`$.Therefore from equation (65) it immediately follows that
$$\mu _S=(3B)^{1/4}\frac{\mu _C}{(\rho _CB)^{1/4}}.$$
(66)
In order to simplify notation we shall introduce a dimensionless parameter $`\eta =\frac{\rho _C}{B}`$, so that $`C(\rho _C,B)=C(\eta )`$.By eliminating $`\frac{\mu _C}{\mu _S}`$ from equations (62) and (66) we obtain the following exact formula for the mass-radius ratio of a strange star:
$$\frac{M_{stat}}{R_{stat}}=[1C(\eta )\left(\frac{\eta 1}{3}\right)^{\frac{1}{2}})].$$
(67)
For a given equation of state the mass-radius ratio of the star depends on the values of the metric tensor component at the center of the star,$`C(\eta )`$, only.A possible representation for the function giving the values of $`g_{00}`$ at the center of the quark star is in the form of a power series $`C(\eta )=const.+\frac{a_i}{\eta ^i}`$,with $`a_i`$ constants.
As applied on the star surface the mass continuity equation leads to a rough approximation of the quark star mass of the form $`\frac{dM_{stat}}{dr}\frac{M_{stat}}{R_{stat}}16\pi BR_{stat}^2`$. A mass-radius relation of this form could also describe zero pressure quark matter,with $`ϵ=4B`$,$`p0`$ and $`M_{stat}(R_{stat})=\frac{16\pi BR_{stat}^3}{3}`$. But for densities greater than $`4B`$ the effects determined by the large central density become important. Hence for strange quark stars we propose the following mass-radius relation:
$$M_{stat}=\frac{16\pi B}{3}a(\eta )R_{stat}^3,$$
(68)
with $`a(\eta )`$ a function describing the variation in the quark star mass due to the increase of the central density.
The exact form and the values of the functions $`C(\eta )`$ and $`a(\eta )`$ can be determined only by numerically integrating the gravitational field equations.By fitting the numerical data given in for the mass and radius of the strange star with the expressions (67) and (68) we obtain the following representations for these functions (in the present paper we shall consider $`B=10^{14}g/cm^3=56Mev/fm^3`$):
$$C(\eta )=44.005\frac{1}{\eta ^3}6.68158\frac{1}{\eta ^2}+2.7403\frac{1}{\eta }+0.0554667,$$
(69)
$$a(\eta )=0.0000521833\eta ^30.00378523\eta ^2+0.114564\eta +0.624094.$$
(70)
The numerical constants in equations (69) and (70) depend on $`B`$ because the numerical data have been calculated at a given $`B`$. For the polynomial fittings (69) and (70) the correlation coefficient $`r=0.9997`$ and the probability $`P<0.001`$. Therefore for a given value of the bag constant $`B`$ we obtain the following exact representations for the radius and mass of the static strange matter configuration obeying the MIT bag model equation of state:
$$R_{stat}(\eta )=\left(\frac{3c^2}{32\pi GB}\right)^{1/2}\left(\frac{1}{a(\eta )}\left(1C(\eta )\left(\frac{\eta 1}{3}\right)^{\frac{1}{2}}\right)\right)^{\frac{1}{2}},$$
(71)
$$M_{stat}(\eta )=\frac{\sqrt{3}}{2}\left(\frac{c^2}{32\pi GB}\right)^{\frac{1}{2}}\left(a(\eta )^{\frac{1}{2}}\left(1C(\eta )\left(\frac{\eta 1}{3}\right)^{\frac{1}{2}}\right)\right)^{\frac{3}{2}}.$$
(72)
The variations of the radius and mass for a strange star ($`B=10^14g/cm^3`$) as a function of the parameter $`\eta `$ are represented in Figures 1 and 2. For the sake of comparison we have also presented the data obtained by numerically integrating the TOV and hydrostatic equilibrium equations .Using (69)-(72) we can reproduce the values of the mass and radius of the quark star obtained by numerical integration with an error smaller than 1%. The maximum radius of the strange star is obtained from the condition $`\frac{dR}{d\eta }=0`$.The corresponding algebraic equation has the solution $`\eta _{max}^R=9.99012`$ (this value depends of course on the value of $`B`$),giving the value of the ratio of the central pressure and bag constant for the maximum allowable radius $`R_{max}`$ of the static strange star.This can be expressed as
$$R_{max}^{stat}=0.569906\times \sqrt{\frac{3c^2}{32\pi BG}},$$
(73)
and its numerical value for $`B=10^{14}g/cm^3`$ is $`R_{max}=1.1436\times 10^6cm`$.From the condition $`\frac{dM}{d\eta }=0`$ it follows that $`\eta _{max}^{(M)}=22.41173`$ and the maximum mass of the static quark star is given by
$$M_{max}^{stat}=0.297866\times \frac{\sqrt{3}}{2(32\pi B)^{\frac{1}{2}}}\left(\frac{c^2}{G}\right)^{\frac{3}{2}}.$$
(74)
From equation (74) and for the chosen value of the bag constant,we obtain a value of $`M_{max}=2.016M_{}`$.These results are in good agreement with the previously proposed (Witten ,Haensel,Zdunik and Schaeffer ) maximum radius and mass values,given by equations (55) (from equations (55) and for $`B=56MeV/fm^3`$ we obtain $`M_{max}=2.03M_{}`$).For values of $`\eta >\eta _{max}^{(M)}`$ static quark star models would be unstable to radial perturbations.
As an application of the mass and radius formulae obtained for the static strange stars we shall derive an explicit expression for the total energy of the quark star. The total energy (including the gravitational field contribution) inside an equipotential surface $`S`$ can be defined, according to Lynden-Bell and Katz and Gron and Johannesen to be
$$E=E_M+E_F=\frac{1}{8\pi }\xi _s_S[K]𝑑s,$$
(75)
where $`\xi ^i`$ is a Killing vector field of time translation, $`\xi _s`$ its value at $`S`$ and $`[K]`$ is the jump across the shell of the trace of the extrinsic curvature of $`S`$, considered as embedded in 2-space $`t=constant`$.$`E_M=_ST_i^k\xi ^i\sqrt{g}𝑑S_k`$ and $`E_F`$ are the energy of the matter and of the gravitational field,respectively. This definition is manifestly coordinate invariant. In the case of the static strange star with the use of equation (71)and (72) we obtain for the total energy (also including the gravitational contribution) the following exact expression:
$$E=E_{SQM}+E_F=\sqrt{\frac{3}{32\pi B}}C(\eta )\left[\frac{1}{a(\eta )}\frac{\eta 1}{3}\left(1C(\eta )\sqrt{\frac{\eta 1}{3}}\right)\right]^{1/2},$$
(76)
where $`E_{SQM}`$ is the total energy of the quark matter.
The variation of the total energy of the strange star as a function of the parameter $`\eta `$ is represented in Figure 3.
The minimum value of the total-matter plus gravitational-energy of the strange matter configuration is obtained for $`\eta _{min}=5.68171`$.The most stable static stellar configuration made of strange matter is given by quark stars with radius $`R_{stab}=9.97179\times 10^5cm`$ and with mass $`M_{stab}=0.96558M_{}`$, corresponding to values of the central density of the order of $`\rho _C=5.681\times B=5.681\times 10^{14}g/cm^3`$.
We shall consider now the study of the rotating strange star configurations. We shall compare our results obtained with the use of equations (37)-(39) and (71)-(72) with the results provided by Stergioulas,Kluzniak and Bulik and obtained by numerically integrating the gravitational field equations for maximally rotating (”Keplerien”) models of strange stars.The results of Stergioulas et al. are also in very good agreement with the results of the exact numerical models of rotating strange stars built of self bounded quark matter of Gourgoulhon et al. , the difference between these two works being smaller than $`1\%`$. In order to improve the accuracy of the expressions (37)-(39) we shall consider that the function $`C_{rot}^{(2)}(\eta ,R_e,\mathrm{\Omega })`$ can be expressed in a more general form as
$$C_{rot}^{(2)}(\eta ,R_e,\mathrm{\Omega })=\left(1R_e^2\mathrm{\Omega }^2\right)^\gamma C\left(\eta \right),$$
(77)
and we will assume that the parameters $`\alpha `$ and $`\gamma `$ are not constants, but some angular velocity dependent functions given by
$$\alpha =\frac{1.23188\times 10^9}{\mathrm{\Omega }^2}\frac{2.15445\times 10^5}{\mathrm{\Omega }}+13.1399,$$
(78)
$$\gamma =9.66251\times 10^{11}\mathrm{\Omega }^34.56952\times 10^8\mathrm{\Omega }^2+6.99642\times 10^4\mathrm{\Omega }^12.54401.$$
(79)
Equations (78)-(79) take into account the variation of the central density of the maximally rotating strange star due to the increase of the angular velocity.
In Fig.4 we have represented the variation of the radius of the strange stars,given by equation (80) together with the angular velocity dependent functions $`\alpha `$ and $`\gamma `$ and the values given in Stergioulas et al., calculated for the same values of the central density and angular velocity.The mean of the difference between these two sets of values is smaller than $`1\%`$. The variation of the radius of the maximally rotating strange star as a function of both central density and angular velocity is represented in Fig.5.Figs.6 and 7 present the variation of the mass of the rapidly rotating strange star as a function of the angular velocity $`\mathrm{\Omega }`$ and of $`\mathrm{\Omega }`$ and central density, respectively.
From equation (37) and equations (71) and (72) it follows that the radius of the rotating strange star can be expressed, as a function of the central density and angular velocity only,in the following form:
$$R_e=\sqrt{\frac{12\left(1C(\eta )\sqrt{\frac{\eta 1}{3}}\right)}{64\pi Ba(\eta )3\alpha \mathrm{\Omega }^2}}.$$
(80)
Hence in this approximation the basic rotational parameters of maximally rotating strange star can be represented in terms of the static configuration and of the angular velocity only.
## V DISCUSSIONS AND FINAL REMARKS
In the present paper we have suggested the possibility of the existence of a universal pattern expressing the basic properties of rotating compact object as simple functions of the parameters of the static object and of the angular velocity only. We have obtained exact formulae which give the dependence of the radius and mass of the static and rotating stars on the central density of the stellar object and of its angular velocity.In the static case this is made possible due to the constancy of the chemical potential. The two unknown functions involved in the model must be obtained by fitting the exact formulae with data obtained from the numerical integration of the structure equations of the neutron or quark star. The resulting analytical expressions can reproduce the radius and mass of the strange star with an error smaller than $`1\%`$ and they also provide a simple way to obtain the maximum mass and radius of the static configuration.
In the rotating case,with the use of the hydrostatic equilibrium condition,which is the consequence of the Bianchi identities, we have also obtained exact mass-radius relations,depending on three functions describing the effect of rotation on star structure.These relations are exact in the sense that they have been obtained without any special assumptions.By assuming some appropriate forms for the unknown functions we have obtained a general description of the mass and radius of the rotating neutron or strange stars, which generally and for a broad class of equations of state can reproduce the values obtained by numerical integration of the gravitational field equations with a mean error of around $`5\%`$.The expressions of the unknown parameters have been chosen following a close analogy with the static case,whose relevance for the study of rotating general relativistic configurations seems to be more important than previously believed.These functions also incorporate some other general relativistic effects not explicitly taken into account,like the variation of the moment of inertia of the star with the angular velocity. Ravenhall and Pethick have presented a formula valid for a broad range of realistic equations of state of dense matter expressing the moment of inertia in terms of stellar mass and radius. We have not used these results,obtained in the slow rotation limit because,at least in the case of strange stars, the Newtonian expression of the moment of inertia also leads to a quite accurate physical description of the rotating objects.As an application of the obtained formulae we had given a derivation of the ”empirical” formula relating the maximum angular velocity to the mass and radius of the static maximal stellar configurations.
The possibility of obtaining the basic parameters of general relativistic rotating objects in terms of static parameters could lead to a major computational simplification in the study of rotation.The relation between the presented formalism and the Einstein gravitational field equations will be the subject of a future publication.
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# 1 Introduction
## 1 Introduction
The results that this paper presents is about gravitational couplings for y deformed and generalized orientifold planes (yGOp-planes) . The usual orientifold planes do not have gauge fields on their worldvolumes and do no have any kind of topological y-deformation over their worldvolumes . The y deformed and generalized orientifold planes that this paper consider have SO(2k) Yang-Mills gauge fields-bundles over their corresponding worldvolumes and have topological deformations of the fields-bundles over their corresponding worldvolumes. The aim of the present paper is to display the Wess-Zumino part of the effective action for such y-deformed and generalized orientifold planes.
For the usual orientifold planes the Wess-Zumino action has the following form,which can be derived both from anomaly cancellation arguments and from direct computation on string scattering amplitudes:
$`𝑺_{𝑾𝒁\mathbf{(}𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}`$
Where the Mukai vector of RR charges for the usual orientifold p-plane is given by:
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}\mathbf{=}\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}`$
In this formula C is the vector of the RR potential forms. L is the Hirzebruch genus that generates the Hirzebruch polynomials which are given in terms of Pontryaguin classes for real bundles. The Pontryaguin classes are given in terms of the 2-form curvature of the corresponding real bundle. The formula for Q involves two real bundles over the worldvolume of the usual orientifold plane. These two bundles are the tangent bundle for the worldvolume and the normal bundle by respect to space-time for such worldvolume. Q is given then in terms of the curvatures for the tangent and normal bundles and does not have contributions from the others real bundles such as SO(2k) Yang-Mills gauge bundles and does not have any kind of topological deformation.
In a recent work was presented the Mukay vector of RR charges for the generalized orientifold planes which have two SO(2k) Yang-Mills gauge bundles on their worldvolumes. Such vector of RR charges is given by the following formula:
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}}}`$
For the generalized orientifold planes the Wess-Zumino action has the following form:
$`𝑺_{𝑾𝒁\mathbf{(}𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}}}`$
The formula for the vector of RR charges corresponding to a generalized orientifold plane involves now four real bundles over the worldvolume: the tangent bundle, the normal bundle and two new SO(2k) YM gauge bundles. When one of these new SO(2k) bundles is the tangent bundle and the other is the normal bundle, one obtain the usual formula for Q corresponding to the usual orientifold planes using the following identity:
$`𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}`$
Then, one has:
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}}}`$
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{=}𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}`$
In these formulas, A denotes the roof-Dirac genus and Mayer denotes the Mayer class for one SO(2k) YM gauge bundle.
In other recent work,also, was presented the Mukay vector of RR charges for the y-deformed orientifold planes which have topological y-deformations on their worldvolumes. Such vector of RR charges is given by the following formula:
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}`$
For the y-deformed orientifold planes the Wess-Zumino action has the following form:
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}`$
The formula for the vector of RR charges corresponding to a y-deformed orientifold plane involves now two real bundles over the worldvolume: the tangent bundle and the normal bundle, but in this case one has a topological y-deformation over the worldvolume. When the parameter y of the topological y-deformation is 1, then one obtain the usual formula for Q corresponding to the usual orientifold planes using the following identity:
$`𝑪𝑯𝑰_\mathrm{𝟏}\mathbf{(}𝑹\mathbf{)}\mathbf{=}𝑳\mathbf{(}𝑹\mathbf{)}`$
Then, one has:
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_\mathrm{𝟏}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_\mathrm{𝟏}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}`$
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{,}\mathrm{𝟏}\mathbf{,}\mathbf{)}\mathbf{=}\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{=}𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}`$
In these formulas, CHI sub y denotes the chi-y- genus which when y=1 is the Hirzebruch-genus and when y=0, is the Todd genus.
In this paper is presented the Mukay vector of RR charges for the y- deformed and generalized orientifold planes which have two SO(2k) Yang-Mills gauge bundles on their worldvolumes and have topological y-deformations of the all bundles that are living on their worldvolumes . Such vector of RR charges is given by the following formula:
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}}`$
For the y-deformed and generalized orientifold planes the Wess-Zumino action has the following form:
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}}`$
The formula for the vector of RR charges corresponding to a y-deformed and generalized orientifold plane involves now four y-deformed real bundles over the worldvolume: the tangent bundle, the normal bundle and two new SO(2k) YM gauge bundles. When one of these new SO(2k) bundles is the tangent bundle and the other is the normal bundle, one obtain the formula for Q corresponding to the y-deformed orientifold plane using the following identity:
$`𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{,}𝒚\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{,}𝒚\mathbf{)}\mathbf{=}𝑪𝑯𝑰_𝒚\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}`$
Then, one has:
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}}`$
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{=}𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{,}𝒚\mathbf{)}`$
When the parameter y of the topological deformation is 1 , one obtains the formula for Q corresponding to the generalized orientifold plane using the following identity:
$`𝑨\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}𝑨\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}`$
Then, one has:
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}}}`$
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{=}𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}`$
When one of these new SO(2k) bundles is the tangent bundle and the other is the normal bundle,and the parameter y of the topological deformation is 1, then, one obtains the formula for Q corresponding to the usual orientifold plane using the following identity:
$`𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}`$
Then, one has:
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}}}`$
$`𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{=}𝑸\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{,}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}`$
In these formulas, appears the y-deformed roof-Dirac genus and the y-deformed Mayer class for one SO(2k) YM gauge bundle.
In the following section the Mukay vector of RR charges for a such y-deformed and generalized orientifold p-plane (yGOp-plane) , will be given in terms of the powers of the curvatures for the four y-deformed real bundles involved over the worldvolume.
In the third section are presented the elementary processes corresponding to the power expansion for the three Q’s. In the final four section some conclutions are presented about other yGOp-planes, about non-BPS yGOp-planes and non commutative yGOp-planes.
## 2 The Power Expantions for Q’s
In this section are obtained three series power-curvature expantions corresponding to the three Q associated respectively to the GOp-planes, yOp-planes and yGOp-planes. The yGOp-planes are an unification of the GOp-planes and yOp-planes. The yGOp-planes contains the usual Op-planes, as limiting cases.
### 2.1 The Power Expantion for GOp-plane
Let E be a SO(2k)-bundle over the worldvolume of a generalized orientifold plane and consider a formal factorisation for the total Pontryaguin classs of the real bundle E, which has the following form:
$`𝒑\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^𝒌\mathbf{(}\mathrm{𝟏}\mathbf{+}𝒚_𝒊^\mathrm{𝟐}\mathbf{)}`$
The total Pontryaguin classs of the real bundle E,has the following formal sumarisation in terms of the corresponding Pontryaguin classes:
$`𝒑\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝒑_𝒋\mathbf{(}𝑬\mathbf{)}`$
The total Mayer class for the real bundle E has the following formal factorisation:
$`𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^𝒌𝒄𝒐𝒔𝒉\mathbf{(}\frac{𝐲_𝐢}{\mathrm{𝟐}}\mathbf{)}`$
The total Mayer class for the real bundle E has the following formal sumarisation in terms of the Mayer polynomials which are formed from the corresponding Pontryaguin classes :
$`𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝑴𝒂𝒚𝒆𝒓_𝒋\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}\mathbf{}\mathbf{,}𝒑_𝒋\mathbf{(}𝑬\mathbf{)}\mathbf{)}`$
The Mayer polynomials are given by:
$`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟎}\mathbf{(}𝒑_\mathrm{𝟎}\mathbf{(}𝑬\mathbf{)}\mathbf{)}\mathbf{=}𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟎}\mathbf{(}\mathrm{𝟏}\mathbf{)}\mathbf{=}\mathrm{𝟏}`$
$`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟏}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{)}\mathbf{=}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}}{\mathrm{𝟖}}`$
$`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟐}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟐}\mathbf{(}𝑬\mathbf{)}\mathbf{)}\mathbf{=}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟒}𝐩_\mathrm{𝟐}\mathbf{(}𝐄\mathbf{)}}{\mathrm{𝟑𝟖𝟒}}`$
$`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟑}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟐}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟑}\mathbf{(}𝑬\mathbf{)}\mathbf{)}\mathbf{=}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}^\mathrm{𝟑}\mathbf{+}\mathrm{𝟏𝟐}𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}𝐩_\mathrm{𝟐}\mathbf{(}𝐄\mathbf{)}\mathbf{+}\mathrm{𝟒𝟖}𝐩_\mathrm{𝟑}\mathbf{(}𝐄\mathbf{)}}{\mathrm{𝟒𝟔𝟎𝟖𝟎}}`$
The pontryaguin classes of the real bundle E have the following realizations in terms of the powers of the 2-form curvature for such bundle. For this curvature the y’s are the eigenvalues:
$`𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{=}𝒑_\mathrm{𝟏}\mathbf{(}𝑹_𝑬\mathbf{)}\mathbf{=}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟖}𝐩𝐢^\mathrm{𝟐}}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}`$
$`𝒑_\mathrm{𝟐}\mathbf{(}𝑬\mathbf{)}\mathbf{=}𝒑_\mathrm{𝟐}\mathbf{(}𝑹_𝑬\mathbf{)}\mathbf{=}\frac{\mathrm{𝟏}}{\mathrm{𝟏𝟔}𝐩𝐢^\mathrm{𝟒}}\mathbf{[}\frac{\mathrm{𝟏}}{\mathrm{𝟖}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟒}}𝒕𝒓𝑹_𝑬^\mathrm{𝟒}\mathbf{]}`$
$`𝒑_\mathrm{𝟑}\mathbf{(}𝑬\mathbf{)}\mathbf{=}𝒑_\mathrm{𝟑}\mathbf{(}𝑹_𝑬\mathbf{)}\mathbf{=}\frac{\mathrm{𝟏}}{\mathrm{𝟔𝟒}𝐩𝐢^\mathrm{𝟔}}\mathbf{[}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟒𝟖}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟑}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟔}}𝒕𝒓𝑹_𝑬^\mathrm{𝟔}\mathbf{+}\frac{\mathrm{𝟏}}{\mathrm{𝟖}}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}𝒕𝒓𝑹_𝑬^\mathrm{𝟒}\mathbf{]}`$
Using all these expretions one can to obtain the following expantion:
$`𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐄\mathbf{)}}{\mathrm{𝟑𝟐}}\mathbf{+}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐄\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟒}𝐩_\mathrm{𝟐}\mathbf{(}𝐑_𝐄\mathbf{)}}{\mathrm{𝟔𝟏𝟒𝟒}}\mathbf{+}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐄\mathbf{)}^\mathrm{𝟑}\mathbf{+}\mathrm{𝟏𝟐}𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐄\mathbf{)}𝐩_\mathrm{𝟐}\mathbf{(}𝐑_𝐄\mathbf{)}\mathbf{+}\mathrm{𝟒𝟖}𝐩_\mathrm{𝟑}\mathbf{(}𝐑_𝐄\mathbf{)}}{\mathrm{𝟐𝟗𝟒𝟗𝟏𝟐𝟎}}\mathbf{+}\mathbf{}`$
Now one has the following expantions:
$`𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\mathrm{𝟏}\mathbf{}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑\mathbf{)}}{\mathrm{𝟗𝟔}}\mathbf{+}\frac{\mathrm{𝟕}𝐩_\mathrm{𝟏}\mathbf{(}𝐑\mathbf{)}^\mathrm{𝟐}\mathbf{}\mathrm{𝟒}𝐩_\mathrm{𝟐}\mathbf{(}𝐑\mathbf{)}}{\mathrm{𝟗𝟐𝟏𝟔𝟎}}\mathbf{+}\mathbf{}`$
$`𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑\mathbf{)}}{\mathrm{𝟒𝟖}}\mathbf{+}\frac{\mathbf{}𝐩_\mathrm{𝟏}\mathbf{(}𝐑\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟕}𝐩_\mathrm{𝟐}\mathbf{(}𝐑\mathbf{)}}{\mathrm{𝟏𝟏𝟓𝟐𝟎}}\mathbf{+}\mathbf{}`$
Using these three expantions it is easy to obtain the following identities:
$`𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}`$
$`𝑨\mathbf{(}𝑹\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑹\mathbf{)}\mathbf{=}𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}`$
$`𝑨\mathbf{(}\mathrm{𝟐}𝑹\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\mathrm{𝟐}𝑹\mathbf{)}\mathbf{=}𝑳\mathbf{(}𝑹\mathbf{)}`$
$`𝑨\mathbf{(}\mathrm{𝟐}^𝒒𝑹\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\mathrm{𝟐}^𝒒𝑹\mathbf{)}\mathbf{=}𝑳\mathbf{(}\mathrm{𝟐}^{𝒒\mathbf{}\mathrm{𝟏}}𝑹\mathbf{)}`$
$`\mathbf{[}𝑨\mathbf{(}𝑹\mathbf{)}\mathrm{𝟐}^𝒌𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑹\mathbf{)}\mathbf{]}_{𝒕𝒐𝒑𝒇𝒐𝒓𝒎}\mathbf{=}𝑳\mathbf{(}𝑹\mathbf{)}_{𝒕𝒐𝒑𝒇𝒐𝒓𝒎}`$
With the help from these identities one has that:
$`\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{=}\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}`$
Using all these equations it is easy to obtain the following power expantion for Q:
$`\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟏𝟓𝟑𝟔}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟓𝟏𝟐}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟒𝟕𝟏𝟖𝟓𝟗𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟐𝟗𝟒𝟗𝟏𝟐𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟓𝟐𝟒𝟐𝟖𝟖}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟗𝟔𝟔𝟎𝟖}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟒}\mathbf{)}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟕𝟖𝟔𝟒𝟑𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}`$
When the bundle E is the tangent bundle and the bundle F is the normal bundle one obtain the usual power expantion for Q corresponding to the usual orientifold plane:
$`\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟏𝟓𝟑𝟔}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟓𝟏𝟐}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟒𝟕𝟏𝟖𝟓𝟗𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟐𝟗𝟒𝟗𝟏𝟐𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟓𝟐𝟒𝟐𝟖𝟖}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟗𝟔𝟔𝟎𝟖}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟕𝟖𝟔𝟒𝟑𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}`$
$`\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{=}\mathrm{𝟏}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟕𝟔𝟖}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟏𝟕𝟗𝟔𝟒𝟖}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathrm{𝟕}\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟒𝟕𝟒𝟓𝟔𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}`$
$`\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{=}\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}`$
### 2.2 The Power Expantion for yOp-plane
For the other hand in the case of the y-Op-plane, the total Chern Class for a complex n-dimensional bundle V over the worldvolume has the following sumarization:
$`𝒄\mathbf{(}𝑽\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝒄_𝒋\mathbf{(}𝑻\mathbf{)}`$
also, the total Chern Class for the such bundle has the following factorization:
$`𝒄\mathbf{(}𝑽\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^𝒏\mathbf{(}\mathrm{𝟏}\mathbf{+}𝒙_𝒊\mathbf{)}`$
The CHI-y- genus for the complex bundle V has the following formal factorisation:
$`𝑪𝑯𝑰_𝒚\mathbf{(}𝑽\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^𝒏\frac{\mathbf{(}\mathrm{𝟏}\mathbf{+}\mathrm{𝐲𝐞𝐱𝐩}\mathbf{(}\mathbf{}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}𝐱_𝐢\mathbf{)}\mathbf{)}𝐱_𝐢}{\mathrm{𝟏}\mathbf{}\mathrm{𝐞𝐱𝐩}\mathbf{(}\mathbf{}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}𝐱_𝐢\mathbf{)}}`$
The CHI-y- genus for the complex bundle V has the following formal sumarisation in terms of the y-deformed Todd polynomials which are formed from the corresponding Chern classes and from the polynomials on y :
$`𝑪𝑯𝑰_𝒚\mathbf{(}𝑽\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝑻_𝒋\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}\mathbf{}\mathbf{,}𝒄_𝒋\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}`$
The y-Todd polynomials are given by:
$`𝑻_\mathrm{𝟎}\mathbf{(}𝒄_\mathrm{𝟎}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}𝑻_\mathrm{𝟎}\mathbf{(}\mathrm{𝟏}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\mathrm{𝟏}`$
$`𝑻_\mathrm{𝟏}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathbf{(}\mathrm{𝟏}\mathbf{}𝐲\mathbf{)}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟐}}`$
$`𝑻_\mathrm{𝟐}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathbf{(}𝐲^\mathrm{𝟐}\mathbf{}\mathrm{𝟏𝟎}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟏𝟐}}`$
$`𝑻_\mathrm{𝟑}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟑}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathbf{}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}𝐜_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathrm{𝟏𝟐}𝐲\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟑}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟐𝟒}}`$
$`𝑻_\mathrm{𝟒}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟑}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟒}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathbf{(}\mathbf{}𝐲^\mathrm{𝟒}\mathbf{+}\mathrm{𝟒𝟕𝟒}𝐲^\mathrm{𝟐}\mathbf{}\mathrm{𝟏𝟐𝟒}𝐲\mathbf{}\mathrm{𝟏}\mathbf{}\mathrm{𝟏𝟐𝟒}𝐲^\mathrm{𝟑}\mathbf{)}𝐜_\mathrm{𝟒}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathbf{(}𝐲^\mathrm{𝟐}\mathbf{}\mathrm{𝟓𝟖}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}𝐜_\mathrm{𝟑}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{(}\mathrm{𝟑}𝐜_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟒}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}\mathbf{}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}^\mathrm{𝟒}\mathbf{)}}{\mathrm{𝟕𝟐𝟎}}`$
Now the relations between the Pontryaguin classes and the Chern Classes for the bundle V are given by the following formulas:
$`𝒑_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{=}\mathbf{}\mathrm{𝟐}𝒄_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{+}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}^\mathrm{𝟐}`$
$`𝒑_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{=}\mathrm{𝟐}𝒄_\mathrm{𝟒}\mathbf{(}𝑽\mathbf{)}\mathbf{}\mathrm{𝟐}𝒄_\mathrm{𝟑}\mathbf{(}𝑽\mathbf{)}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{+}𝒄_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}^\mathrm{𝟐}`$
Using these relations the y-deformed Todd polynomials can be written as follows:
$`𝑻_\mathrm{𝟎}\mathbf{(}𝒄_\mathrm{𝟎}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}𝑻_\mathrm{𝟎}\mathbf{(}\mathrm{𝟏}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\mathrm{𝟏}`$
$`𝑻_\mathrm{𝟏}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathbf{(}\mathrm{𝟏}\mathbf{}𝐲\mathbf{)}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟐}}`$
$`𝑻_\mathrm{𝟐}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐩_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathrm{𝟑}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟏𝟐}}`$
$`𝑻_\mathrm{𝟑}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟑}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathbf{}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}𝐜_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathrm{𝟏𝟐}𝐲\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟑}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟐𝟒}}`$
$`𝑻_\mathrm{𝟒}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟑}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟒}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathbf{}\mathrm{𝟏𝟓}\mathbf{(}𝐲^\mathrm{𝟐}\mathbf{+}\mathrm{𝟏𝟒}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟒}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathrm{𝟏𝟓}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}𝐜_\mathrm{𝟑}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{(}\mathrm{𝟕}𝐩_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}\mathbf{}𝐩_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}^\mathrm{𝟐}\mathbf{)}}{\mathrm{𝟕𝟐𝟎}}`$
When y=1 the y-deformed Todd polynomials are the same Hirzebruch polynomials:
$`𝑻_\mathrm{𝟎}\mathbf{(}𝒄_\mathrm{𝟎}\mathbf{(}𝑽\mathbf{)}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}𝑻_\mathrm{𝟎}\mathbf{(}\mathrm{𝟏}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\mathrm{𝟏}\mathbf{=}𝑳_\mathrm{𝟎}`$
$`𝑻_\mathrm{𝟏}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\frac{\mathbf{(}\mathrm{𝟏}\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟐}}\mathbf{=}\mathrm{𝟎}`$
$`𝑻_\mathrm{𝟐}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\frac{\mathbf{(}\mathrm{𝟏}\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐩_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathrm{𝟑}\mathbf{(}\mathrm{𝟏}\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟏𝟐}}\mathbf{=}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟑}}\mathbf{=}𝑳_\mathrm{𝟏}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{)}`$
$`𝑻_\mathrm{𝟑}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟑}\mathbf{(}𝑽\mathbf{)}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\frac{\mathbf{}\mathbf{(}\mathrm{𝟏}\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}\mathrm{𝟏}\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}𝐜_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathrm{𝟏𝟐}\mathbf{(}\mathrm{𝟏}\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟑}\mathbf{(}𝐕\mathbf{)}}{\mathrm{𝟐𝟒}}\mathbf{=}\mathrm{𝟎}`$
$`𝑻_\mathrm{𝟒}\mathbf{(}𝒄_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟑}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒄_\mathrm{𝟒}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{,}\mathrm{𝟏}\mathbf{)}\mathbf{=}\frac{\mathbf{}\mathrm{𝟏𝟓}\mathbf{(}\mathrm{𝟏}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟏𝟒}\mathbf{+}\mathrm{𝟏}\mathbf{)}\mathbf{(}\mathrm{𝟏}\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟒}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathrm{𝟏𝟓}\mathbf{(}\mathrm{𝟏}\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}\mathrm{𝟏}\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}𝐜_\mathrm{𝟑}\mathbf{(}𝐕\mathbf{)}\mathbf{+}\mathbf{(}\mathrm{𝟏}\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{(}\mathrm{𝟕}𝐩_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}\mathbf{}𝐩_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}^\mathrm{𝟐}\mathbf{)}}{\mathrm{𝟕𝟐𝟎}}\mathbf{=}\frac{\mathrm{𝟕}𝐩_\mathrm{𝟐}\mathbf{(}𝐕\mathbf{)}\mathbf{}𝐩_\mathrm{𝟏}\mathbf{(}𝐕\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟒𝟓}}\mathbf{=}𝑳_\mathrm{𝟐}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑽\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟐}\mathbf{(}𝑽\mathbf{)}\mathbf{)}`$
Using all these expretions one can to obtain the following expantion:
$`𝑪𝑯𝑰_𝒚\mathbf{(}\frac{𝐑_𝐕}{\mathrm{𝟒}}\mathbf{)}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟏}\mathbf{}𝐲\mathbf{)}𝐜_\mathrm{𝟏}\mathbf{(}𝐑_𝐕\mathbf{)}}{\mathrm{𝟖}}\mathbf{+}\frac{\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐕\mathbf{)}\mathbf{+}\mathrm{𝟑}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟐}\mathbf{(}𝐑_𝐕\mathbf{)}}{\mathrm{𝟏𝟗𝟐}}\mathbf{+}\frac{\mathbf{}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟏}\mathbf{(}𝐑_𝐕\mathbf{)}𝐜_\mathrm{𝟐}\mathbf{(}𝐑_𝐕\mathbf{)}\mathbf{+}\mathrm{𝟏𝟐}𝐲\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟑}\mathbf{(}𝐑_𝐕\mathbf{)}}{\mathrm{𝟏𝟓𝟑𝟔}}\mathbf{+}\frac{\mathbf{}\mathrm{𝟏𝟓}\mathbf{(}𝐲^\mathrm{𝟐}\mathbf{+}\mathrm{𝟏𝟒}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟒}\mathbf{(}𝐑_𝐕\mathbf{)}\mathbf{+}\mathrm{𝟏𝟓}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}𝐜_\mathrm{𝟏}\mathbf{(}𝐑_𝐕\mathbf{)}𝐜_\mathrm{𝟑}\mathbf{(}𝐑_𝐕\mathbf{)}\mathbf{+}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{(}\mathrm{𝟕}𝐩_\mathrm{𝟐}\mathbf{(}𝐑_𝐕\mathbf{)}\mathbf{}𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐕\mathbf{)}^\mathrm{𝟐}\mathbf{)}}{\mathrm{𝟏𝟖𝟒𝟑𝟐𝟎}}\mathbf{+}\mathbf{}`$
When the first chern class of V is trivial, one obtain, using again the relations between pontryaguin classes and Chern classes, the following result:
$`𝑪𝑯𝑰_𝒚\mathbf{(}\frac{𝐑_𝐕}{\mathrm{𝟒}}\mathbf{)}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟐}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{}\mathrm{𝟑}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{)}𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐕\mathbf{)}}{\mathrm{𝟑𝟖𝟒}}\mathbf{+}\frac{\mathrm{𝟏𝟐}𝐲\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}𝐜_\mathrm{𝟑}\mathbf{(}𝐑_𝐕\mathbf{)}}{\mathrm{𝟏𝟓𝟑𝟔}}\mathbf{+}\frac{\mathbf{(}\mathbf{}\mathrm{𝟔𝟎}\mathbf{(}𝐲^\mathrm{𝟐}\mathbf{+}\mathrm{𝟏𝟒}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟓𝟔}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{)}𝐩_\mathrm{𝟐}\mathbf{(}𝐑_𝐕\mathbf{)}\mathbf{}\mathbf{(}\mathbf{}\mathrm{𝟏𝟓}\mathbf{(}𝐲^\mathrm{𝟐}\mathbf{+}\mathrm{𝟏𝟒}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟖}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{)}𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐕\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟏𝟒𝟕𝟒𝟓𝟔𝟎}}\mathbf{+}\mathbf{}`$
Finally using this last expansion and the relations between the Pontryaguin classes and the 2-form curvature, one can to obtain the following development for the Q of the yOp-planes:
$`\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{\mathbf{(}𝐲^\mathrm{𝟐}\mathbf{}\mathrm{𝟏𝟎}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟔𝟏𝟒𝟒}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟑}}{\mathrm{𝟐𝟓𝟔}}𝒚\mathbf{(}𝒚\mathbf{}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑻\mathbf{)}\mathbf{}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑵\mathbf{)}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟕𝟓𝟒𝟗𝟕𝟒𝟕𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒚^\mathrm{𝟐}\mathbf{}\mathrm{𝟏𝟎}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟖𝟖𝟕𝟒𝟑𝟔𝟖𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}\mathbf{}\mathrm{𝟒}𝒚^\mathrm{𝟒}\mathbf{}\mathrm{𝟒𝟗𝟔}𝒚^\mathrm{𝟑}\mathbf{+}\mathrm{𝟏𝟖𝟗𝟔}𝒚^\mathrm{𝟐}\mathbf{}\mathrm{𝟒𝟗𝟔}𝒚\mathbf{}\mathrm{𝟒}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}`$
When y=1, then one obtain the development for the Q of the usual Op-plane:
$`\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_\mathrm{𝟏}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_\mathrm{𝟏}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟏}^\mathrm{𝟐}\mathbf{}\mathrm{𝟏𝟎}\mathbf{+}\mathrm{𝟏}\mathbf{)}\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟔𝟏𝟒𝟒}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟑}}{\mathrm{𝟐𝟓𝟔}}\mathrm{𝟏}\mathbf{(}\mathrm{𝟏}\mathbf{}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑻\mathbf{)}\mathbf{}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑵\mathbf{)}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟕𝟓𝟒𝟗𝟕𝟒𝟕𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}\mathrm{𝟏}^\mathrm{𝟐}\mathbf{}\mathrm{𝟏𝟎}\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟖𝟖𝟕𝟒𝟑𝟔𝟖𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}\mathbf{}\mathrm{𝟒}\mathbf{}\mathrm{𝟒𝟗𝟔}\mathbf{+}\mathrm{𝟏𝟖𝟗𝟔}\mathbf{}\mathrm{𝟒𝟗𝟔}\mathbf{}\mathrm{𝟒}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}`$
$`\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_\mathrm{𝟏}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_\mathrm{𝟏}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{=}\mathrm{𝟏}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟕𝟔𝟖}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟏𝟕𝟗𝟔𝟒𝟖}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathrm{𝟕}\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟒𝟕𝟒𝟓𝟔𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}`$
$`\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_\mathrm{𝟏}\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_\mathrm{𝟏}\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{=}\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}`$
### 2.3 The Power Expantion for yGOp-plane
Let E be a y-deformed SO(2k)-bundle over the worldvolume of a y-deformed and generalized orientifold plane and consider a formal factorisation for the total Pontryaguin classs of the y-deformed real bundle E, which has the following form:
$`𝒑\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^𝒌\mathbf{(}\mathrm{𝟏}\mathbf{+}𝒚_𝒊^\mathrm{𝟐}\mathbf{)}`$
The total Pontryaguin classs of the real SO(2k)- bundle E,has the following formal sumarisation in terms of the corresponding Pontryaguin classes:
$`𝒑\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝒑_𝒋\mathbf{(}𝑬\mathbf{)}`$
From the other hand the total Chern class of the complex SU(k)-bundle E, has the following formal factorisation:
$`𝒄\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^𝒌\mathbf{(}\mathrm{𝟏}\mathbf{+}𝒚_𝒊\mathbf{)}`$
also, the total Chern Class for the such bundle has the following sumarisation:
$`𝒄\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝒄_𝒋\mathbf{(}𝑬\mathbf{)}`$
The total y-deformed Mayer class for the real-complex bundle E has the following formal factorisation:
$`𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑬\mathbf{,}𝒚\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^𝒌\frac{\mathrm{𝐞𝐱𝐩}\mathbf{(}\frac{𝐲_𝐢\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}}{\mathrm{𝟒}}\mathbf{)}\mathbf{+}\mathrm{𝐲𝐞𝐱𝐩}\mathbf{(}\mathbf{}\frac{𝐲_𝐢\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}}{\mathrm{𝟒}}\mathbf{)}}{\mathrm{𝟐}}`$
The total y-deformed Mayer class for the real-complex bundle E has the following formal sumarisation in terms of the y-deformed Mayer polynomials which are formed from the corresponding Pontryaguin classes,from the corresponding Chern classes and from polynomials for the parameter y :
$`𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑬\mathbf{,}𝒚\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝑴𝒂𝒚𝒆𝒓_𝒋\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}\mathbf{}\mathbf{,}𝒄_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}\mathbf{}\mathbf{,}𝒚\mathbf{)}`$
The y-deformed Mayer polynomials are given by:
$`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟎}\mathbf{(}𝒑_\mathrm{𝟎}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟎}\mathbf{(}\mathrm{𝟏}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟔}}`$
$`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟏}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{𝐲\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}}{\mathrm{𝟏𝟐𝟖}}`$
$`𝑴𝒂𝒚𝒆𝒓_{\frac{\mathrm{𝟑}}{\mathrm{𝟐}}}\mathbf{(}𝒄_\mathrm{𝟑}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{𝐲\mathbf{(}𝐲\mathbf{}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}𝐜_\mathrm{𝟑}\mathbf{(}𝐄\mathbf{)}}{\mathrm{𝟐𝟓𝟔}}`$
$`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟐}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟐}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{𝐲\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{(}\mathbf{(}\mathrm{𝟏}\mathbf{}𝐲\mathbf{+}𝐲^\mathrm{𝟐}\mathbf{)}𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}^\mathrm{𝟐}\mathbf{}\mathrm{𝟐}\mathbf{(}\mathbf{}\mathrm{𝟒}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{+}𝐲^\mathrm{𝟐}\mathbf{)}𝐩_\mathrm{𝟐}\mathbf{(}𝐄\mathbf{)}\mathbf{)}}{\mathrm{𝟔𝟏𝟒𝟒}}`$
The total Pontryaguin classs of the real tangent bundle T of the worldvolume of the y-GOp-plane,has the following formal sumarisation in terms of the corresponding Pontryaguin classes:
$`𝒑\mathbf{(}𝑻\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝒑_𝒋\mathbf{(}𝑻\mathbf{)}`$
also, the formal factorisation for the total Pontryaguin classs of the y-deformed real tangent bundle T, has the following form:
$`𝒑\mathbf{(}𝑻\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^{\frac{𝐩\mathbf{+}\mathrm{𝟏}}{\mathrm{𝟐}}}\mathbf{(}\mathrm{𝟏}\mathbf{+}𝒙_𝒊^\mathrm{𝟐}\mathbf{)}`$
The total y-deformed Dirac-roof genus for the real bundle T has the following formal factorisation:
$`𝑨\mathbf{(}𝑻\mathbf{,}𝒚\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^{\frac{𝐩\mathbf{+}\mathrm{𝟏}}{\mathrm{𝟐}}}\frac{\frac{𝐱_𝐢}{\mathrm{𝟐}}}{\mathrm{𝐬𝐢𝐧𝐡}\mathbf{(}\frac{\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}𝐱_𝐢}{\mathrm{𝟒}}\mathbf{)}}`$
The total y-deformed Dirac-roof genus for the real bundle T has the following formal sumarisation in terms of the y-deformed Dirac polynomials which are formed from the corresponding Pontryaguin classes and from polynomials for the parameter y :
$`𝑨\mathbf{(}𝑻\mathbf{,}𝒚\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝑨_𝒋\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑻\mathbf{)}\mathbf{,}\mathbf{}\mathbf{,}𝒑_𝒋\mathbf{(}𝑻\mathbf{)}\mathbf{,}𝒚\mathbf{)}`$
The y-deformed Dirac polynomials are given by:
$`𝑨_\mathrm{𝟎}\mathbf{(}𝒑_\mathrm{𝟎}\mathbf{(}𝑻\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}𝑨_\mathrm{𝟎}\mathbf{(}\mathrm{𝟏}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathrm{𝟏𝟔}}{\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}}`$
$`𝑨_\mathrm{𝟏}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑻\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\mathbf{}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐓\mathbf{)}}{\mathrm{𝟔}\mathbf{(}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}}`$
$`𝑨_\mathrm{𝟐}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑻\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟐}\mathbf{(}𝑻\mathbf{)}\mathbf{,}𝒚\mathbf{)}\mathbf{=}\frac{\mathrm{𝟕}𝐩_\mathrm{𝟏}\mathbf{(}𝐓\mathbf{)}^\mathrm{𝟐}\mathbf{}\mathrm{𝟒}𝐩_\mathrm{𝟐}\mathbf{(}𝐓\mathbf{)}}{\mathrm{𝟓𝟕𝟔𝟎}}`$
It is easy to check that when y=1 the y-deformed Mayer polynomials and the y-deformed Dirac polynomials are the same usuales Mayer polynomials and Dirac polynomials.
Using all these y-deformed polynomials and the relations between the Pontryaguin classes and the 2-form curvatures, one can to obtain the following expantion for the Q of the yGOp-planes:
$`\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟔𝟏𝟒𝟒}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟓𝟏𝟐}𝐩𝐢^\mathrm{𝟐}}𝒚\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟑}}{\mathrm{𝟐𝟓𝟔}}𝒚\mathbf{(}𝒚\mathbf{}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑬\mathbf{)}\mathbf{}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑭\mathbf{)}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟕𝟓𝟒𝟗𝟕𝟒𝟕𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟒𝟕𝟏𝟖𝟓𝟗𝟐𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟓𝟐𝟒𝟐𝟖𝟖}𝐩𝐢^\mathrm{𝟒}}𝒚^\mathrm{𝟐}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟑𝟗𝟑𝟐𝟏𝟔}𝐩𝐢^\mathrm{𝟒}}𝒚\mathbf{(}\mathbf{}\mathrm{𝟒}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{+}𝒚^\mathrm{𝟐}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟒}\mathbf{)}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟑𝟏𝟒𝟓𝟕𝟐𝟖}𝐩𝐢^\mathrm{𝟒}}𝒚\mathbf{(}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}`$
When y=1 the yGOp-plane is reduced to the GOp-plane. When E=T and F=N, the yGOp-plane is reduced to the yOp-plane. When y=1 and E=T and F=N, then the yGOp-plane is reduced to the usual Op-plane.
## 3 The Elementary Processes
In this section are presented the elementary gravitational processes for Op-planes, GOp-planes, yOp-planes and yGOp-planes corresponding to the series power-curvature expantios of the three Q’s obtained in the section two.
### 3.1 The Elementary Processes for Op-planes
The WZ action for the usual orientifold p-plane can be writen as a sum of the WZ actions for three elementary processes:
$`𝑺_{𝑾𝒁\mathbf{(}𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟑}𝑺_{𝑾𝒁\mathbf{(}𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}𝒋}`$
The WZ actions for the three elementary processes are given by the following expretions:
$`𝑺_{𝑾𝒁\mathbf{(}𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟏}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{+}\mathrm{𝟏}}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟐}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟑}}\mathbf{[}\mathbf{}\mathbf{(}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟕𝟔𝟖}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{)}\mathbf{]}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟑}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟕}}\mathbf{(}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟏𝟕𝟗𝟔𝟒𝟖}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathrm{𝟕}\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟒𝟕𝟒𝟓𝟔𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}\mathbf{)}`$
The first WZ action describes an elementary process for which the usual orientifold p-plane emites one (p+1)-form RR potential. The second WZ action describes an elementary process for which the usual Op-plane absorbs two gravitons and emits one (p-3)-form RR potential. The third WZ action describes an elementary process for which the Op-plane absorbs four gravitons and emits one (p-7)-form RR potential.
### 3.2 The Elementary Processes for GOp-planes
From the result of the section two, the WZ action for a generalized orientifold p-plane can be writen as a sum of the WZ actions for some elementary processes:
$`𝑺_{𝑾𝒁\mathbf{(}𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟔}𝑺_{𝑾𝒁\mathbf{(}𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}𝒋}`$
The WZ actions for the six elementary processes are given by the following expretions:
$`𝑺_{𝑾𝒁\mathbf{(}𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟏}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{+}\mathrm{𝟏}}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟐}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟑}}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟏𝟓𝟑𝟔}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟑}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟑}}\mathbf{(}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟓𝟏𝟐}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟒}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟕}}\mathbf{(}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟒𝟕𝟏𝟖𝟓𝟗𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟐𝟗𝟒𝟗𝟏𝟐𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟓}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟕}}\mathbf{(}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟓𝟐𝟒𝟐𝟖𝟖}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟗𝟔𝟔𝟎𝟖}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟒}\mathbf{)}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟔}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟕}}\mathbf{(}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟕𝟖𝟔𝟒𝟑𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}\mathbf{)}`$
The first WZ action describes an elementary process for which the generalized orientifold p-plane emites one (p+1)-form RR potential. The second WZ action describes an elementary process for which the generalized Op-plane absorbs two gravitons and emits one (p-3)-form RR potential. The third WZ actuib describes an elementary process for which the generalized Op-plane absorbs two gaugeons and emits one (p-3)-form RR potential. The fourth WZ action describes an elementary process for which the GOp-plane absorbs four gravitons and emits one (p-7)-form RR potential. The fifth WZ action describes an elementary process for which the GOp-plane absorbs four gaugeons and emits one (p-7)-form RR potential. The sixth WZ action describes an elementary process for which the GOp-planes absorbs two gravitons and two gaugeons and emits one (p-7)-form RR potential.
When the gaugeons corresponding to the bundles E and F are the same gravitons corresponding to the bundles T and N respectively, then the six elementary process for the GOp-plane are reduced to the usuals three elementary process for the usual Op-plane: Op-plane emites one (p+1)-form RR potential,Op-plane absorbs two gravitons and emits one (p-3)-form RR potential; and, Op-plane absorbs four gravitons and emits one (p-7)-form RR potential.
### 3.3 The Elementary Processes for yOp-planes
Of other hand, from the result of the section two, the WZ action for a y-deformed orientifold p-plane can be writen as a sum of the WZ actions for some elementary processes:
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟒}𝑺_{𝑾𝒁\mathbf{(}𝒚𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}𝒋}`$
The WZ actions for the four elementary processes are given by the following expretions:
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟏}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{+}\mathrm{𝟏}}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟐}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟑}}\mathbf{[}\mathbf{}\mathbf{(}\frac{\mathbf{(}𝐲^\mathrm{𝟐}\mathbf{}\mathrm{𝟏𝟎}𝐲\mathbf{+}\mathrm{𝟏}\mathbf{)}\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟔𝟏𝟒𝟒}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{)}\mathbf{]}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟑}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟓}}\mathbf{(}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟑}}{\mathrm{𝟐𝟓𝟔}}𝒚\mathbf{(}𝒚\mathbf{}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑻\mathbf{)}\mathbf{}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑵\mathbf{)}\mathbf{)}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟒}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟕}}\mathbf{(}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟕𝟓𝟒𝟗𝟕𝟒𝟕𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒚^\mathrm{𝟐}\mathbf{}\mathrm{𝟏𝟎}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟏𝟖𝟖𝟕𝟒𝟑𝟔𝟖𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}\mathbf{}\mathrm{𝟒}𝒚^\mathrm{𝟒}\mathbf{}\mathrm{𝟒𝟗𝟔}𝒚^\mathrm{𝟑}\mathbf{+}\mathrm{𝟏𝟖𝟗𝟔}𝒚^\mathrm{𝟐}\mathbf{}\mathrm{𝟒𝟗𝟔}𝒚\mathbf{}\mathrm{𝟒}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}\mathbf{)}`$
The first WZ action describes an elementary process on which the yOp-plane emites one (p+1)-form RR potential. The second WZ action describes an elementary process for which the y-deformed Op-plane absorbs two gravitons and emits one (p-3)-form RR potential. The third WZ action describes an elementary process for which the y-deformed Op-plane absorbs three gravitons and emits one (p-5)-form RR potential. The fourth WZ action describes an elementary process for which the yOp-plane absorbs four gravitons and emits one (p-7)-form RR potential. When y=1,then the four elementary process for the yOp-plane are reduced to the usuals three elementary process for the usual Op-plane: Op-plane emites one (p+1)-form RR potential,Op-plane absorbs two gravitons and emits one (p-3)-form RR potential; and, Op-plane absorbs four gravitons and emits one (p-7)-form RR potential.
### 3.4 The Elementary Processes for yGOp-planes
From the result of the section two, the WZ action for a y-deformed and generalized orientifold p-plane can be writen as a sum of the WZ actions for some elementary processes:
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟕}𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}𝒋}`$
The WZ actions for the seven elementary processes are given by the following expretions:
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟏}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{+}\mathrm{𝟏}}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟐}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟑}}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟔𝟏𝟒𝟒}𝐩𝐢^\mathrm{𝟐}}\mathbf{(}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟑}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟑}}\mathbf{(}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟐}}{\mathrm{𝟓𝟏𝟐}𝐩𝐢^\mathrm{𝟐}}𝒚\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟒}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟓}}\mathbf{(}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟑}}{\mathrm{𝟐𝟓𝟔}}𝒚\mathbf{(}𝒚\mathbf{}\mathrm{𝟏}\mathbf{)}\mathbf{(}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑬\mathbf{)}\mathbf{}𝒄_\mathrm{𝟑}\mathbf{(}𝑹_𝑭\mathbf{)}\mathbf{)}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟓}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟕}}\mathbf{(}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟕𝟓𝟒𝟗𝟕𝟒𝟕𝟐}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟒𝟕𝟏𝟖𝟓𝟗𝟐𝟎}𝐩𝐢^\mathrm{𝟒}}\mathbf{(}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟒}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟒}\mathbf{)}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟔}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟕}}\mathbf{(}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟓𝟐𝟒𝟐𝟖𝟖}𝐩𝐢^\mathrm{𝟒}}𝒚^\mathrm{𝟐}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{+}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟑𝟗𝟑𝟐𝟏𝟔}𝐩𝐢^\mathrm{𝟒}}𝒚\mathbf{(}\mathbf{}\mathrm{𝟒}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{+}𝒚^\mathrm{𝟐}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟒}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟒}\mathbf{)}\mathbf{)}`$
$`𝑺_{𝑾𝒁\mathbf{(}𝒚𝑮𝑶𝒑\mathbf{}𝒑𝒍𝒂𝒏𝒆\mathbf{)}\mathbf{,}\mathrm{𝟕}}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪_{𝒑\mathbf{}\mathrm{𝟕}}\mathbf{(}\mathbf{}\frac{\mathbf{(}\mathrm{𝟒}𝐩𝐢^\mathrm{𝟐}\mathrm{𝐚𝐥𝐟𝐚}\mathbf{)}^\mathrm{𝟒}}{\mathrm{𝟑𝟏𝟒𝟓𝟕𝟐𝟖}𝐩𝐢^\mathrm{𝟒}}𝒚\mathbf{(}𝒚\mathbf{+}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}\mathbf{(}𝒕𝒓𝑹_𝑻^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑵^\mathrm{𝟐}\mathbf{)}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{}𝒕𝒓𝑹_𝑭^\mathrm{𝟐}\mathbf{)}\mathbf{)}`$
The first WZ action describes an elementary process for which the y-deformed and generalized orientifold p-plane emites one (p+1)-form RR potential. The second WZ action describes an elementary process for which the y-deformed and generalized Op-plane absorbs two gravitons and emits one (p-3)-form RR potential. The third WZ action describes an elementary process for which the y-deformed and generalized Op-plane absorbs two gaugeons and emits one (p-3)-form RR potential. The fourth WZ action describes an elementary process for which the yGOp-plane absorbs three gaugeons and emits one (p-5)-form RR potential. The fifth WZ action describes an elementary process for which the yGOp-plane absorbs four gravitons and emits one (p-7)-form RR potential. The sixth WZ action describes an elementary process for which the yGOp-planes absorbs four gaugeons a and emits one (p-7)-form RR potential. The seventh WZ action describes an elementary process for which the yGOp-planes absorbs two gravitons and two gaugeons and emits one (p-7)-form RR potential. When y=1 the elementary processes for the yGOp-plane are reduced to the elementary processes for the GOp-plane. When E=T and F=N the elementary processes for the yGOp-plane are reduced to the elementary processes for the yOp-plane. When y=1 and E=T and F=N the elementary processes for the yGOp-plane are reduced to the elementary processes for the usual Op-plane.
## 4 Conclutions
The WZ action for the yGOp-planes can be modified or extended by various ways. When the bundles have non-trivial second Stiefel-Whitney classes one can to write the following WZ action which incorporates an effect of the magnetic monopoles:
$`𝑺_{𝑾𝒁}\mathbf{=}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}𝐞^{\frac{𝐝_\mathrm{𝟏}}{\mathrm{𝟐}}}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}𝐞^{\frac{𝐝_\mathrm{𝟐}}{\mathrm{𝟐}}}}}`$
where:
$`𝒅_\mathrm{𝟏}\mathbf{=}𝒓𝒆𝒅𝒖𝒄𝒕𝒊𝒐𝒏\mathbf{.}𝒎𝒐𝒅\mathbf{.2}\mathbf{(}𝒘_\mathrm{𝟐}\mathbf{(}𝑻\mathbf{)}\mathbf{+}𝒘_\mathrm{𝟐}\mathbf{(}𝑬\mathbf{)}\mathbf{)}`$
$`𝒅_\mathrm{𝟐}\mathbf{=}𝒓𝒆𝒅𝒖𝒄𝒕𝒊𝒐𝒏\mathbf{.}𝒎𝒐𝒅\mathbf{.2}\mathbf{(}𝒘_\mathrm{𝟐}\mathbf{(}𝑵\mathbf{)}\mathbf{+}𝒘_\mathrm{𝟐}\mathbf{(}𝑭\mathbf{)}\mathbf{)}`$
This action describes processes on which the yGOp-plane emites RR-forms and absorbs gravitons, gaugeons and magnetic monopoles.
From the other side one can to write the following actions for GOp-planes non standard:
$`𝑺_{𝑾𝒁}\mathbf{=}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\mathbf{(}\mathrm{𝟐}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{)}}}\mathbf{}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{)}`$
$`𝑺_{𝑾𝒁}\mathbf{=}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\mathbf{(}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{)}}}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{)}`$
In the same way, one can to write the following actions for yOp-planes non standard:
$`𝑺_{𝑾𝒁}\mathbf{=}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\mathbf{(}\mathrm{𝟐}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{)}}}\mathbf{}\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{)}`$
$`𝑺_{𝑾𝒁}\mathbf{=}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\mathbf{(}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{)}}}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{)}`$
These actions correspond respectively to the Sp-type yOp-planes and the yOp-planes that give rise to gauge symmetries of type SO(2n+1). Such non-standard yOp-planes are building from combinations of the D-p-branes and standard yOp-planes.
In the same way, one can to write the following actions for yGOp-planes non stantard:
$`𝑺_{𝑾𝒁}\mathbf{=}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\mathbf{(}\mathrm{𝟐}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{)}}}\mathbf{}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}}\mathbf{)}`$
$`𝑺_{𝑾𝒁}\mathbf{=}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\mathbf{(}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{)}}}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}}\mathbf{)}`$
These actions correspond respectively to the Sp-type yGOp-planes and the yGOp-planes that give rise to gauge symmetries of type SO(2n+1). Such non-standard yGOp-planes are building from combinations of the D-p-branes and standard yGOp-planes.
By combination of Dp-branes,yDp-branes, Op-planes, GOp-planes,yOp-planes and yGOp-planes one can to have gauge teories with symmetries Sp and SO-odd whose WZ actions are give respectively by:
$`𝑺_{𝑾𝒁}\mathbf{=}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\mathbf{(}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{)}}}\mathbf{+}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{,}𝐲\mathbf{)}}}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟒}}\mathbf{(}\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{+}\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{+}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{)}\mathbf{+}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}}\mathbf{)}\mathbf{)}`$
$`𝑺_{𝑾𝒁}\mathbf{=}\frac{𝐓_𝐩}{\mathrm{𝐤𝐚𝐩𝐩𝐚}}\mathbf{}_{𝒑\mathbf{+}\mathrm{𝟏}}𝑪\mathbf{(}\frac{\mathrm{𝟏}}{\mathrm{𝟐}}\mathbf{(}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{)}}}\mathbf{+}\sqrt{\frac{𝐀\mathbf{(}𝐑_𝐓\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}𝐑_𝐍\mathbf{,}𝐲\mathbf{)}}}\mathbf{)}\mathbf{}\mathrm{𝟐}^{𝒑\mathbf{}\mathrm{𝟒}}\frac{\mathrm{𝟏}}{\mathrm{𝟒}}\mathbf{(}\sqrt{\frac{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{\mathrm{𝐂𝐇𝐈}_𝐲\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{+}\sqrt{\frac{𝐋\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟒}}\mathbf{)}}{𝐋\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟒}}\mathbf{)}}}\mathbf{+}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{)}}}\mathbf{)}\mathbf{+}\sqrt{\frac{𝐀\mathbf{(}\frac{𝐑_𝐓}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}{𝐀\mathbf{(}\frac{𝐑_𝐍}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}\mathrm{𝐌𝐚𝐲𝐞𝐫}\mathbf{(}\frac{𝐑_𝐅}{\mathrm{𝟐}}\mathbf{,}𝐲\mathbf{)}}}\mathbf{)}\mathbf{)}`$
Finally one can to think about non-BPS GOp-planes, non-BPS yOp-planes and non-BPS yGOp-planes with the tachyon effect. One can also to think about noncommutative Op-planes,GOp-planes,yOp-planes and yGOp-planes
In conclution gauge theories with symmetries SO-even,Sp and SO-odd can be obtained from the combination of the Dp-branes,yDp-branes,Op-planes, GOp-planes, yOp-planes and yGOp-planes of the string theory.
## 5 References
About WZ action for usual orientifold planes:
K. Dasgupta, D. Jatkar and S. Mukhi, Gravitational couplings and Z2 orientifolds, Nucl. Phys. B523 (1998) 465, hep-th/9707224.
K. Dasgupta and S. Mukhi, Anomaly inflow on orientifold planes, J. High Energy Phys. 3 (1998) 4, hep-th/9709219.
J. Morales, C. Scrucca and M. Serone, Anomalous couplings for D-branes and O-planes, hep-th/9812071.
B.Stefanski,Jr., Gravitational Couplings of D-branes and O-planes, hep-th/9812088
Ben Craps and Frederik Roose, (Non-)Anomalous D-brane and O-plane couplings:the normal bundle, hep-th/9812149.
About WZ action for non-standard orientifold planes:
Sunil Mukhi and Nemani V. Suryanarayana, Gravitational Couplings, Orientifolds and M-Planes, hep-th/9907215
About Mayer class , Mayer integrality theorem and CHI-y-genus:
F. Hirzebruch, Topological Methods in Algebraic Geometry, 1978
Christian Bar, Elliptic Symbols, december 1995, Math. Nachr. 201, 7-35 (1999)
Taras E. Panov, Calculation of Hirzebruch Genera for manifolds acted on by the Group Z/p via invariantes of the action, math.AT/9909081
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# Non-Abelian Stokes Theorem and Computation of Wilson Loop
## Abstract
It is shown that the application of the non-Abelian Stokes theorem to the computation of the operators constructed with Wilson loop will lead to ambiguity, if the gauge field under consideration is a non-trivial one. This point is illustrated by the specific examples of the computation of a non-local operator.
The non-Abelian Stokes theorem<sup>1-5</sup> is widely applied to compute Wilson loop (closed non-Abelian phase factor), $`\mathrm{\Psi }(C)=Pexp\left(ig_C𝑑z^\mu A_\mu (z)\right)`$ ($`A_\mu `$ is the simplified notation for $`\underset{a=1}{\overset{dimG}{}}A_\mu ^aT^a`$), which is important for the construction of gauge invariant operators in the non-perturbative approaches to QCD. The power of the theorem lies in transforming the line integral in a Wilson loop to a more tractable surface integral over the surface $`S`$ enclosed by contour $`C`$:
$`Pexp\left(ig{\displaystyle _C}𝑑z^\mu A_\mu (z)\right)=\mathrm{\Psi }(x,y;\overline{C}_x)Pexp\left(ig/2{\displaystyle _S}𝑑\sigma ^{\mu \nu }\widehat{F}_{\mu \nu }(y,z;C_z)\right)\mathrm{\Psi }(y,x;C_x),`$ (1)
where $`y`$ is an arbitrary reference point on $`S`$, $`\mathrm{\Psi }(y,x;C_x)`$ the phase factor connecting the initial and final point $`x`$ of $`C`$ to $`y`$ through a path $`C_x`$ and its inverse path $`\overline{C}_x`$. The shifted field strength $`\widehat{F}_{\mu \nu }(y,z;C_z)`$ here is defined as $`\widehat{F}_{\mu \nu }(y,z;C_z)=\mathrm{\Psi }(y,z;C_z)F_{\mu \nu }(z)\mathrm{\Psi }^{}(y,z;C_z)`$, and its ordered integral over $`S`$ is given as follows:
$`Pexp\left(ig/2{\displaystyle _S}𝑑\sigma ^{\mu \nu }\widehat{F}_{\mu \nu }(y,z;C_z)\right)=\underset{n\mathrm{}}{lim}{\displaystyle \underset{i=n}{\overset{1}{}}}(I+ig/2\widehat{F}_{\mu \nu }(y,z_i;C_{z_i})\mathrm{\Delta }\sigma ^{\mu \nu })`$ (2)
$$=I+ig/2_0^1𝑑s_0^1𝑑t\frac{(z^\mu ,z^\nu )}{(s,t)}\widehat{F}_{\mu \nu }(y,z(s,t))+(ig/2)^2_0^1𝑑s_0^1𝑑t\frac{(z^\mu ,z^\nu )}{(s,t)}\widehat{F}_{\mu \nu }(y,z(s,t))$$
$`\times {\displaystyle _0^1}ds_1{\displaystyle _0^t}dt_1{\displaystyle \frac{(z_1^\mu ,z_1^\nu )}{(s_1,t_1)}}\widehat{F}_{\mu \nu }(y,z_1(s_1,t_1))+\mathrm{}.`$ (3)
Eq. (2) means that the ordered surface integral is equivalent to the infinite product of the phase factor on the net of the handled small plaquettes around each $`z_i`$ in $`S`$. Obviously the ‘handle’ is referred to as the phase factor connecting $`z_i`$ and $`y`$:
$$(I+ig/2\widehat{F}_{\mu \nu }(y,z_i;C_{z_i})\mathrm{\Delta }\sigma ^{\mu \nu })=\mathrm{\Psi }(y,z_i;C_{z_i})(I+ig/2F_{\mu \nu }(z_i)\mathrm{\Delta }\sigma ^{\mu \nu })\mathrm{\Psi }^{}(y,z_i;C_{z_i}).$$
If the shifted field strength $`\widehat{F}_{\mu \nu }(y,z;C_z)`$ is well defined as the function of $`\{s,t\}`$ in the equations, the ordered surface integral in Eq. (2) can be performed definitely, and it makes the calculation of the expectation value $`<Tr\mathrm{\Psi }(C)>`$ convenient by means of the cumulant expansion technique<sup>6</sup>. For an overview see Ref. . However, we find that the transformation from the line integral to the surface integral in Eq. (1) might also give rise to an ambiguity in the computation of Wilson loop $`\mathrm{\Psi }(C)`$, if it is applied in the situation of a non-trivial gauge field. In this letter we will present the specific examples to illustrate the point. For simplicity we specify that the gauge field under consideration is two-dimensional and the gauge group is SU(2).
Before our discussion we review some properties of the non-Abelian phase factor $`\mathrm{\Psi }(x,x_0;C_x)`$ that is defined as the solution of the following differential equation:
$`{\displaystyle \frac{d\mathrm{\Psi }}{dt}}(t)=A(t)\mathrm{\Psi }(t),`$ (4)
with the boundary condition
$$\mathrm{\Psi }(t_0)=I,$$
where $`A(t)=igA_\mu (x(t))dx^\mu (t)/dt`$. The path $`C_x`$ is parametrized by $`t`$, with $`x(t)=x`$ and $`x(t_0)=x_0`$. The solution of Eq. (4) is
$$\mathrm{\Psi }(t)=Pexp\left(_{t_0}^t𝑑sA(s)\right)=I+_{t_0}^t𝑑sA(s)$$
$$+_{t_0}^t𝑑sA(s)_{t_0}^s𝑑s_1A(s_1)+\mathrm{}+_{t_0}^t𝑑sA(s)\mathrm{}_{t_0}^{s_{n1}}𝑑s_nA(s_n)+\mathrm{}$$
$`=\underset{\mathrm{\Delta }t0}{lim}(I+A(t_{n1})\mathrm{\Delta }t)\mathrm{}(I+A(t_0)\mathrm{\Delta }t).`$ (5)
From it we immediately obtain the following two properties of $`\mathrm{\Psi }(x,x_0;C)`$:
$`\mathrm{\Psi }(x_3,x_1;C_{x_2}C_{x_1})=\mathrm{\Psi }(x_3,x_2;C_{x_2})\mathrm{\Psi }(x_2,x_1;C_{x_1});`$ (6)
$`\mathrm{\Psi }(x_1,x_2;\overline{C}_{x_1})=\mathrm{\Psi }^{}(x_2,x_1;C_{x_1}).`$ (7)
It is just through these relations that the infinite product of the phase factor on the net of the handled plaquettes (Eq. (2)) should be equivalent to the original Wilson loop $`\mathrm{\Psi }(C)`$. To study the property of $`\mathrm{\Psi }(x,x_0;C)`$ under the gauge transformation
$`A_\mu ^{}(x)=U(x)A_\mu (x)U^{}(x){\displaystyle \frac{i}{g}}_\mu U(x)U^{}(x),`$ (8)
we perform a transformation, $`\mathrm{\Psi }^{}(t)=U(x(t))\mathrm{\Psi }(t)`$, in Eq. (4). After the rearrangement of the terms we arrive at the gauge transformation of $`\mathrm{\Psi }(x,x_0,C)`$:
$`Pexp\left({\displaystyle _{t_0}^t}𝑑sA^{}(s)\right)=U(x(t))Pexp\left({\displaystyle _{t_0}^t}𝑑sA(s)\right)U^{}(x(t_0)).`$ (9)
Applying these results to the general case when $`A_\mu (x)=\underset{a=1}{\overset{3}{}}A_\mu ^a(x)\sigma ^a`$, with $`A_\mu ^a(x)0`$ for $`a=1,2,3`$, we first study the behavior of the shifted field strength $`\widehat{F}_{\mu \nu }(y,x;C_z)`$ in Eqs. (2) and (3). If it can be given as a function of the surface parameters $`\{s,t\}`$, there is a well-defined ordered surface integral $`Pexp(ig/2_S𝑑\sigma ^{\mu \nu }\widehat{F}_{\mu \nu }(y,z;C_z))`$, and the non-Abelian Stokes theorem (Eq. (1)) will be surely valid. From the definition of the shifted field strength it is true as long as the phase factor $`\mathrm{\Psi }(y,z;C_z)`$ can be expressed by the surface parameters $`\{s,t\}`$.
As a matter of fact, however, the phase factor cannot always be given as the function of $`\{s,t\}`$ if $`A_\mu (x)dx^\mu `$ is not an exact form, i.e. $`_\mu A_\nu (x)_\nu A_\mu (x)=0`$ identically. It is proved as follows. Since the reference point $`y`$ in Eq. (2) is arbitrarily chosen on $`S`$, we can set the origin of the surface coordinate $`\{s,t\}`$ at $`y`$. Then the line integral in Eq. (5) for the phase factor $`\mathrm{\Psi }(y,z;C_z)`$ is transformed in terms of the surface parameters $`\{s,t\}`$ to
$$_y^zA_\mu (\stackrel{~}{z})𝑑\stackrel{~}{z}^\mu =_{(0,0)}^{(s,t)}\left\{A_\mu (\stackrel{~}{z}(\stackrel{~}{s},\stackrel{~}{t}))\frac{\stackrel{~}{z}^\mu }{\stackrel{~}{s}}d\stackrel{~}{s}+A_\mu (\stackrel{~}{z}(\stackrel{~}{s},\stackrel{~}{t}))\frac{\stackrel{~}{z}^\mu }{\stackrel{~}{t}}d\stackrel{~}{t}\right\}.$$
If this line integral is not a path-independent one in the surface coordinate system $`\{s,t\}`$, it cannot be given as a function of $`\{s,t\}`$ and, therefore, the phase factor $`\mathrm{\Psi }(z,y;C_z)`$ fails to be a function of $`\{s,t\}`$ too. Of course, if we choose some special surface coordinates, e.g. that of a homotopic path family in Ref., the phase factor can be actually expressed as a function of the surface parameters, $`\mathrm{\Psi }(z,y;C_z)=\mathrm{\Psi }(s,t)`$. Applying the inverse function theorem in analysis to the transformation, $`x^1=x^1(s,t)`$ and $`x^2=x^2(s,t)`$, we have a locally defined relations $`s=s(x^1,x^2)`$ and $`t=t(x^1,x^2)`$, if $`det((x^1,x^2)/(s,t))0`$. Thus $`\mathrm{\Psi }(s,t)`$ can also be given as a function of $`\{x^1,x^2\}`$: $`\mathrm{\Psi }(s,t)=\stackrel{~}{\mathrm{\Psi }}(x^1,x^2)`$. $`\stackrel{~}{\mathrm{\Psi }}(x^1,x^2)`$, however, doesn’t satisfy the partial differential equation
$`_\mu \stackrel{~}{\mathrm{\Psi }}(x)=igA_\mu (x)\stackrel{~}{\mathrm{\Psi }}(x),`$ (10)
unless $`A_\mu (x)dx^\mu `$ is an exact form. For a non-trivial gauge field this condition doesn’t generally hold, and there is no equivalence between Eqs. (4) and (10).
In the general situation we can say that the flaw with the application of the non-Abelian Stokes theorem might be in the ‘handle’ phase factor $`\mathrm{\Psi }(y,z;C_z)`$, which makes the shifted field strength $`\widehat{F}_{\mu \nu }(y,z;C_z)`$ a path-dependent operator that can’t be integrated with respect to the surface parameters. Actually, even if there is no effect of the ‘handle’ phase factor, we can still find examples to demonstrate some ambiguity from the non-Abelian Stokes theorem. To remove the effect of the ‘handle’ phase factor, we just restrict the generators of the Lie algebra to its Cartan subalgebras, e.g. $`A_\mu (x)=A_\mu ^1(x)\sigma ^1`$ in SU(2) gauge theory, then the shifted field strength operator will reduce to field strength operator $`F_{\mu \nu }(x)`$. Our examples are the application of the non-Abelian Stokes theorem to a special case of the gauge invariant operator $`F_{\mu \nu \rho \sigma }(x,x^{},y;C_x,C_x^{})=Tr\{\widehat{F}_{\mu \nu }(y,x;C_x)\widehat{F}_{\rho \sigma }(y,x^{};C_x^{})\}`$, the expectation value of which (the correlator of two shifted field strength) is the basic object of the stochastic vacuum model (SVM)<sup>8-10</sup>. From Eq. (9) this operator should be gauge invariant under an arbitrary gauge transformation Eq. (8). When $`x=x^{}`$ and $`C=\overline{C}_x^{}C_x`$ forms a closed path with the initial and final point at the same point $`x`$, the operator reduces to
$$F_{\mu \nu \rho \sigma }(x,x;C)=Tr\{Pexp\left(ig_C𝑑z^\mu A_\mu (z)\right)F_{\mu \nu }(x)Pexp\left(ig_C𝑑z^\mu A_\mu (z)\right)F_{\rho \sigma }(x)\}.$$
Obviously, with the help of the non-Abelian Stokes theorem, $`F_{\mu \nu \rho \sigma }(x,x;C)`$ can be rewritten as
$`F_{\mu \nu \rho \sigma }(x,x;C)=Tr\{Pexp(ig/2{\displaystyle _S}d\sigma ^{\mu \nu }\widehat{F}_{\mu \nu }(y,z;C_z))\widehat{F}_{\mu \nu }(y,x;C_x)`$ (11)
$`\times Pexp(ig/2{\displaystyle _S}d\sigma ^{\mu \nu }\widehat{F}_{\mu \nu }(y,z;C_z))\widehat{F}_{\rho \sigma }(y,x;C_x)\},`$
where $`y`$ is an arbitrary reference point chosen on $`S`$.
A specific situation we will study with Eq. (11) is described in Fig. (1). The contour $`S`$ here is a rectangular one, and $`A_\mu (x)=A_\mu ^1(x)\sigma ^1`$ for each $`xS`$, then the operator $`F_{\mu \nu \rho \sigma }(x_0,x_0;S)`$ is independent of the gauge field on $`S`$, i.e. $`F_{\mu \nu \rho \sigma }(x_0,x_0;S)=2F_{\mu \nu }^1(x_0)F_{\rho \sigma }^1(x_0)`$.
We apply Eq. (11) to the computation of the operator after the following discontinuos gauge transformation on X-Y plane:
$`U(x,y)=exp(i{\displaystyle \frac{\pi }{4}}H(xa)\sigma ^3),`$ (12)
where $`H(xa)`$ is a Heaviside function,
$$H(xa)=\{\begin{array}{cc}0,& x<a,\\ 1,& x>a,\end{array}$$
which transforms the gauge field, $`A_\mu (x)=A_\mu ^1(x)\sigma ^1`$, to
$$A_\mu ^{}(x)=A_\mu ^1(x)\sigma ^1+H(xa)(A_\mu ^1(x)\sigma ^2A_\mu ^1(x)\sigma ^1)\frac{1}{g}\frac{\pi }{4}\delta (xa)\sigma ^3𝐞_𝐱.$$
Here $`𝐞_𝐱`$ is the unit vector in the direction of X. The difference will arise if we apply Eq. (11) to the transformed operator. The reference point here is chosen at $`y=(a,0)`$, and the whole $`S`$ is paved with the net shown in Fig. (1). From Eqs. (6) and (7), the operator is equivalent to the product of the handled phase factor on $`S`$:
$`F_{\mu \nu \rho \sigma }(x_0,x_0;S)=Tr\{Pexp(ig{\displaystyle _{L_1}}dz^\mu A_\mu (z))Pexp(ig{\displaystyle _{L_2}}dz^\mu A_\mu (z))\widehat{F}_{\mu \nu }(y,x_0)`$ (13)
$`\times Pexp\left(ig{\displaystyle _{L_2}}dz^\mu A_\mu (z)\right)Pexp\left(ig{\displaystyle _{L_1}}dz^\mu A_\mu (z)\right)\widehat{F}_{\rho \sigma }(y,x_0)\},`$
where $`Pexp\left(ig_{L_i}𝑑z^\mu A_\mu (z)\right)`$, $`i=1,2`$, is the product of the phase factor on the net of the handled plaquettes on $`S_i`$ ($`i=1,2`$), and is equal to $`Pexp\left(ig/2_{S_i}𝑑\sigma ^{\mu \nu }\widehat{F}_{\mu \nu }(y,z;C_z)\right)`$, respectively. After the gauge transformation Eq. (12), there are
$$\widehat{F^{}}_{\mu \nu }(y,x_0)=F_{\mu \nu }^1(x_0)\sigma ^1,$$
$$\widehat{F^{}}_{\mu \nu }(y,z)=F_{\mu \nu }^1(z)\sigma ^1$$
on $`S_1`$, and
$$\widehat{F^{}}_{\mu \nu }(y,z)=F_{\mu \nu }^1(z)\sigma ^2$$
on $`S_2`$. Substituting these results into Eq. (11), we obtain
$`F_{\mu \nu \rho \sigma }^{}(x_0,x_0;C)=Tr\{Pexp(ig/2{\displaystyle _{S_2}}d\sigma ^{\mu \nu }F_{\mu \nu }^1(z)\sigma ^2)F_{\mu \nu }^1(x_0)\sigma ^1`$ (14)
$`\times Pexp(ig/2{\displaystyle _{S_2}}d\sigma ^{\mu \nu }F_{\mu \nu }^1(z)\sigma ^2)F_{\rho \sigma }^1(x_0)\sigma ^1\}.`$
It is obviously not equal to the previous result because it involves the gauge field at other points than $`x_0`$ and, therefore, an ambiguity of the gauge invariance of the operator $`F_{\mu \nu \rho \sigma }(x_0,x_0;S)`$ will arise unless the field is a pure gauge one, i.e. the field strength $`F_{\mu \nu }(z)`$ vanishes identically on $`S`$.
Furthermore, we can find an example that directly shows the ambiguity of the non-Abelian Stokes theorem in the computation of the Wilson loop, when the generators of the Lie algebra are still confined in its Cartan subalgebras. Here we adopt polar coordinate (Fig. 2). In this case the gauge field on $`S`$ is given as a discontinuous one:
$$A_\mu (x)=\{\begin{array}{cc}A_\mu ^1(r,\theta )\sigma ^1,& r>r_2orr<r_1,\\ \frac{1}{g}\frac{\pi }{4}_rF(r)𝐞_𝐫\sigma ^3,& r_1<r<r_2.\end{array}$$
The smooth function $`F(r)`$ is defined as
$$F(r)=_r^{r_2}f(r)𝑑r/_{r_1}^{r_2}f(r)𝑑r,$$
where the function $`f(r)`$ is given as
$$f(r)=\{\begin{array}{cc}exp(\frac{1}{rr_1}\frac{1}{rr_2}),& r_1<r<r_2,\\ 0,& elsewhere.\end{array}$$
Such a smooth function has the following property:
$$F(r)=\{\begin{array}{cc}1,& rr_1,\\ 0,& rr_2,\\ decreasefrom1to0,& r_1<r<r_2.\end{array}$$
The whole surface $`S`$ is the union of the three parts: $`S=S_1S_2S_3`$. According to the original definition of the operator $`F_{\mu \nu \rho \sigma }(x_0,x_0;S)`$, it is given as follows:
$`F_{\mu \nu \rho \sigma }(x_0,x_0;S)=Tr\{F_{\mu \nu }^1(x_0)\sigma ^2Pexp\left(ig{\displaystyle _{C_3}}A_\mu (z)\sigma ^1dz^\mu \right)`$ (15)
$`\times F_{\rho \sigma }^1(x_0)\sigma ^2Pexp(ig{\displaystyle _{C_3}}A_\mu (z)\sigma ^2dz^\mu )\},`$
where the path $`C_3`$ is $`S_3`$ excluding the arc $`r_2S`$. The operator is determined by the gauge field $`A_\mu (x)`$ on $`C_3`$, since the phase factor between $`r=r_1`$ and $`r=r_2`$, $`Pexp(_{r_1}^{r_2}i\frac{\pi }{4}_rF(r)\sigma ^3dr)=exp(i\frac{\pi }{4}\sigma ^3)`$, rotates the field strength operator $`F_{\mu \nu }^1(x_0)\sigma ^1`$ at $`x_0`$ to the direction of $`\sigma ^2`$ in the internal space.
To apply the non-Abelian Stokes theorem to the operator, we choose the reference point at $`(r_2,\theta _0)`$. Then we have
$$\widehat{F}_{\mu \nu }(y,x_0)=exp(i\frac{\pi }{4}\sigma ^3)F_{\mu \nu }^1(x_0)\sigma ^1exp(i\frac{\pi }{4}\sigma ^3)=F_{\mu \nu }^1(x_0)\sigma ^2,$$
and
$$Pexp(ig/2_S\widehat{F}_{\mu \nu }(y,z)𝑑\sigma ^{\mu \nu })=Pexp(ig/2_{S_{12}}F_{\mu \nu }^1(z)\sigma ^2𝑑\sigma ^{\mu \nu })$$
$$\times Pexp(ig/2_{S_3}F_{\mu \nu }^1(z)\sigma ^1𝑑\sigma ^{\mu \nu })Pexp(ig/2_{S_{11}}F_{\mu \nu }^1(z)\sigma ^2𝑑\sigma ^{\mu \nu }),$$
where $`S_1=S_{11}S_{12}`$ with respect to the ordering of the plaquettes on the net (Fig. 2), since the contribution from $`S_2`$, a area with pure gauge, is zero. After substituting the contributions from the three different parts into Eq. (11), and considering the fact:
$$Pexp(ig/2_{S_3}F_{\mu \nu }^1(z)\sigma ^1𝑑\sigma ^{\mu \nu })=Pexp(ig_{S_3}A_\mu ^1(z)\sigma ^1𝑑z^\mu ),$$
we have
$`F_{\mu \nu \rho \sigma }(x_0,x_0;C)`$ $`=`$ $`Tr\{Pexp(ig{\displaystyle _{S_3}}A_\mu ^1(z)\sigma ^1dz^\mu )F_{\mu \nu }^1(x_0)\sigma ^2`$
$`\times `$ $`Pexp\left(ig{\displaystyle _{S_3}}A_\mu ^1(z)\sigma ^1dz^\mu \right)F_{\rho \sigma }^1(x_0)\sigma ^2\}.`$
The difference between the results given in Eq. (15) and Eq. (16) is an additional phase factor $`Pexp(ig_{r_2^{}}A_\mu (z)𝑑z^\mu )`$, where $`r_2^{}=r_2+ϵ`$ ($`ϵ`$ is an infinitesimal positive number), in the phase factor $`Pexp(ig_{_{S_3}}A_\mu (z)𝑑z^\mu )`$ in Eq.(16), because $`S_3=r_2^{}C_3`$.
It is clearer to see the difference if we choose a new gauge by performing a gauge transformation, $`U(r,\theta )=exp(i\frac{\pi }{4}F(r)\sigma ^3)`$, over $`S`$. The gauge field on $`S`$ will therefore be transformed to
$$A_\mu (x)=\{\begin{array}{cc}A_\mu (r,\theta )\sigma ^2,& r<r_1,\\ A_\mu (r,\theta )\sigma ^1,& r>r_2,\\ 0,& r_1<r<r_2.\end{array}$$
Then it is easier to compute the shifted field strength operators, $`\widehat{F}_{\mu \nu }(y,z)`$, on the three different parts of $`S`$ and reproduce the results in Eq. (15) and Eq. (16). The same operator $`F_{\mu \nu \rho \sigma }(x_0,x_0;S)`$, with the concerned Wilson loop treated differently as a ‘one-dimensional object’ or a ‘two-dimensional object’ through the non-Abelian Stokes theorem, will be given two different results. This example demonstrates the ambiguity from the transformation of the line integral to the surface integral in Eq. (1), even if the effect of the ‘handle’ phase factor is removed.
Finally we should mention a problem in lattice gauge field theory that is related to our discussion. In the literature available the non-Abelian Stokes theorem is widely used to the expansion of a plaqutte operator with respect to lattice spacing $`a`$, which is crucial in the construction of the lattice actions with the lattice artifact removed to a certain order of the lattice spacing $`a`$ (improved action approach). A plaquette operator $`U_{\mu \nu }`$ is expanded according to the theorem as follows<sup>11</sup>:
$`U_{\mu \nu }=U_\mu U_\nu U_\mu ^{}U_\nu ^{}=Pexp\left(a^2{\displaystyle _0^1}𝑑s{\displaystyle _0^1}𝑑t\widehat{F}_{\mu \nu }(x+as\widehat{\mu }+at\widehat{\nu })\right)`$ (17)
$`=I+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle _0^{s_{i1}}}𝑑s_i{\displaystyle _0^1}𝑑t_ia^2\widehat{F}_{\mu \nu }(x+as_i\widehat{\mu }+at_i\widehat{\nu })\times \mathrm{}a^2(x+as_n\widehat{\mu }+at_n\widehat{\nu }).`$
Then one needs to perform the Taylor expansion of $`\widehat{F}(x+as\widehat{\mu }+at\widehat{\nu })`$ around $`x`$ with respect to $`(s,t)`$ with the help of the relation
$`_\mu ^n_\nu ^m\widehat{F}_{\mu \nu }(x)=D_\mu ^nD_\nu ^mF_{\mu \nu }(x).`$ (18)
In doing so, it is taken for granted that the shifted field strength operator,
$$\widehat{F}_{\mu \nu }(x+as\widehat{\mu }+at\widehat{\nu })=\mathrm{\Psi }(s,t)F_{\mu \nu }(x+as\widehat{\mu }+at\widehat{\nu })\mathrm{\Psi }^{}(s,t),$$
is a function of $`(s,t)`$ as the field strength operator $`F_{\mu \nu }(x+as\widehat{\mu }+at\widehat{\nu })`$ itself. Our previous discussions prove that it is not generally valid for the situation of a non-trivial gauge field and, moreover, Eq. (18) is true only when $`A_\mu (x)dx^\mu `$ is an exact form, i.e. there is the relation Eq. (10).
Of course a plaquette operator can also be expanded by means of Baker-Hausdorff formula or the choice of axial gauge<sup>12-13</sup>, e.g. by imposing
$$A_1(x_1,x_2)=0,$$
$$A_2(0,x_2)=0,$$
on the field configuration, which simplifies the plaquette operator considerably. However, for a gauge field in the general situation, i.e. $`A_\mu (x)0`$ ($`\mu =1,2`$) in any of a gauge we choose, the non-Abelian Stokes theorem is the only effective tool for a convenient expansion of a plaquette operator. With the problems in its application we have discussed, it should be taken an approximate rather than an exact approach if we are dealing with a non-trivial gauge field.
ACKNOWLEDGMENTS. The work is partly supported by National Nature Science Foundation of China under Grant 19677205.
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# Right-handed Neutrinos in 𝝂_𝝁+𝒆⁻→𝝂_𝝁+𝒆⁻ Elastic Scattering
## Abstract
In this paper a scenario admitting the participation of the exotic right-handed scalar $`g_S^R`$ coupling in addition to the standard left-handed $`g_V^L,g_A^L`$ couplings in the low-energy $`\nu _\mu +e^{}\nu _\mu +e^{}`$ scattering is considered. The main goal is to show how the presence of the right-handed neutrinos in the $`(\nu _\mu e^{})`$ scattering changes the laboratory differential cross section in relation to the Standard Model prediction. The $`(\nu _\mu e^{})`$ scattering is studied at the level of the four-fermion point interaction. Muon-neutrinos are assumed to be polarized Dirac fermions and to be massive. In the laboratory differential cross section the new interference term between the standard $`g_V^L`$ coupling of the left-handed neutrino and exotic $`g_S^R`$ coupling of the right-handed neutrino appears which does not vanish in the limit of massless neutrino. This additional contribution, including information about the transverse components of neutrino polarization, generates the azimuthal asymmetry in the angular distribution of the recoil electrons. This regularity would be the proof of the participation of the right-handed neutrinos in the $`(\nu _\mu e^{})`$ scattering.
The $`(VA)`$ structure of weak interactions at low energies describes only what has been measured so far. We mean here the measurement of the electron helicity Bob , the indirect measurement of the neutrino helicity Gold , the asymmetry in the distribution of the electrons from $`\beta `$-decay CWu and the experiment with muon decay Gar which confirmed parity violation Lee . Feynman, Gell-Mann and independently Sudarshan, Marshak Gell established that only left-handed vector $`V`$, axial $`A`$ couplings can take part in weak interactions because this yields the maximum symmetry breaking under space inversion, under charge conjugation; the two-component neutrino theory of negative helicity; the conservation of the combined symmetry $`CP`$ and of the lepton number. In consequence it led to the conclusion that produced neutrinos in $`VA`$ interaction can only be left-handed. However Wu SWu indicated that both standard left-handed $`(V,A)_L`$ couplings and exotic right-handed $`(S,T,P)_R`$ couplings may be responsible for the negative electron helicity observed in $`\beta `$-decay. It would mean that generated neutrinos in $`(S,T,P)`$ interactions may also be right-handed. Recent tests do not provide a unique answer as to the presence of the exotic weak interactions. So Shimizu et al. Shimizu determined the ratio of the strengths of scalar and tensor couplings to the standard vector coupling in $`K^+\pi ^0+e^++\nu _e`$ decay at rest assuming the only left-handed neutrinos for all interactions. Their results indicated the compatibility with the Standard Model (SM) Glashow ; Wein ; Salam prediction. Bodek et al. at the PSI Bodek looked for the evidence of the violation of time reversal invariance measuring $`T`$-odd transverse components of the positron polarization in $`\mu ^+`$-decay. They also admitted the presence of the only left-handed neutrinos produced by the scalar interaction. The recent results presented by the DELPHI Collaboration Delphi concerning the measurement of the Michel parameters and the neutrino helicity in $`\tau `$ lepton decays indicated the consistency with the standard $`VA`$ structure of the charged current weak interaction. However on the other hand, the achieved precision of measurements still admits the deviation from the pure $`VA`$ interaction, i.e. the possible participation of the exotic interactions with the right-handed neutrinos beyond the SM. It is necessary to carry out the new high-precision tests of the Lorentz structure and of the handedness structure of the weak interactions at low energies in which the components of the neutrino polarization would be measured, because in the conventional observables the interference contributions coming from the right-handed neutrinos vanish in the limit of massless neutrino Wolf ; Miran ; Sobkow . Frauenfelder et al. spin pointed out that one has to measure either the neutrino polarization (spin) or the neutrino-electron correlations to determine the full Lorentz structure of the weak interactions.
The main our goal is to show how the presence of the right-handed neutrinos in the $`(\nu _\mu e^{})`$ scattering changes laboratory differential cross section in relation to the SM prediction with the left-handed neutrinos. We will consider the process of the $`(\nu _\mu e^{})`$ scattering at the level of the four-fermion point (contact) interaction. Muon-neutrinos are assumed to be massive Dirac fermions and to be polarized. In these considerations, the incoming neutrinos come from the muon-capture, where the production plane is spanned by the initial muon polarization and the outgoing neutrino momentum. However in practice, the suitable low-energy polarized neutrino source could be the strong chromium source $`(^{51}Cr)`$. Because there exist the models of the spontaneous symmetry breaking under time reversal TDLee in which the scalar weak interaction can appear, the possible participation of the exotic right-handed scalar $`S`$ coupling in addition to the standard left-handed vector $`V`$ and axial $`A`$ couplings is admitted. It means that in the $`VA`$ interaction, the incoming muon-neutrinos are always left-handed, but in the scalar $`S`$ interaction the initial neutrinos are right-handed (the outgoing neutrinos are left-handed). The couplings constants are denoted as $`g_V^L,g_A^L`$ and $`g_S^R`$ respectively to the incoming neutrino handedness:
$``$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}\{(\overline{u}_e^{}\gamma ^\alpha (g_V^Lg_A^L\gamma _5)u_e)(\overline{u}_{\nu _\mu ^{}}\gamma _\alpha (1\gamma _5)u_{\nu _\mu })`$ (1)
$`+{\displaystyle \frac{1}{2}}g_S^R(\overline{u}_e^{}u_e)(\overline{u}_{\nu _\mu ^{}}(1+\gamma _5)u_{\nu _\mu })\},`$
where $`u_e`$ and $`u_e^{}`$ $`(u_{\nu _\mu }`$ and $`u_{\nu _\mu ^{}})`$ are the Dirac bispinors of the initial and final electron (neutrino) respectively, $`G_F=1.16639(1)\times 10^5GeV^2`$ is the Fermi coupling constant Data . In our considerations the system of natural units with $`\mathrm{}=c=1`$, Dirac-Pauli representation of the $`\gamma `$-matrices and the $`(+,,,)`$ metric are used Mandl .
Because the carried out investigations for the $`\mu ^{}`$-capture Sobkow led to the conclusion that the production of the right-handed neutrinos by the exotic scalar interaction manifests the non-vanishing transverse components of the neutrino polarization in the limit of massless neutrino, one expects the similar regularity in the $`(\nu _\mu e^{})`$ scattering. The matrix element for the $`\mu ^{}`$-capture and formula on the magnitude of the transverse neutrino polarization vector $`|𝜼_𝝂^𝑻|`$, in the limit of vanishing neutrino mass, were as follows:
$`_\mu ^{}`$ $`=`$ $`C_V^L(\overline{\mathrm{\Psi }}_\nu \gamma _\lambda (1\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_n\gamma ^\lambda \mathrm{\Psi }_p)`$
$`+C_A^L(\overline{\mathrm{\Psi }}_\nu i\gamma _5\gamma _\lambda (1\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_ni\gamma ^5\gamma ^\lambda \mathrm{\Psi }_p)+C_S^R(\overline{\mathrm{\Psi }}_\nu (1\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_n\mathrm{\Psi }_p),`$
$`|𝜼_𝝂^𝑻|`$ $`=`$ $`{\displaystyle \frac{\sqrt{<𝐒_\nu (\widehat{𝐏}_\mu \times \widehat{𝐪})>_f^2+<𝐒_\nu \widehat{𝐏}_\mu >_f^2}}{<\mathrm{𝟏}>_𝐟}}=|𝐏_\mu ||{\displaystyle \frac{C_S^R}{C_V^L}}|(1+{\displaystyle \frac{q}{2M}})`$
$`\times \{(3+{\displaystyle \frac{q}{M}})|{\displaystyle \frac{C_A^L}{C_V^L}}|^2+(1+{\displaystyle \frac{q}{M}})+|{\displaystyle \frac{C_S^R}{C_V^L}}|^2{\displaystyle \frac{q}{M}}|{\displaystyle \frac{C_A^L}{C_V^L}}|cos(\alpha _{AV}^L)\}^1,`$
$`<𝐒_\nu (\widehat{𝐏}_\mu \times \widehat{𝐪})>_f`$ $`=`$ $`{\displaystyle \frac{\varphi _\mu (0)^2}{4\pi }}|𝐏_\mu |(1+{\displaystyle \frac{q}{2M}})Im(C_V^LC_S^R),`$ (4)
$`<𝐒_\nu \widehat{𝐏}_\mu >_f`$ $`=`$ $`{\displaystyle \frac{\varphi _\mu (0)^2}{4\pi }}|𝐏_\mu |(1+{\displaystyle \frac{q}{2M}})Re(C_V^LC_S^R),`$ (5)
where $`C_V^L,C_A^L,C_S^R`$ \- the complex coupling constants for the standard vector, axial and exotic scalar weak interactions denoted respectively to the outgoing neutrino handedness; $`<𝐒_\nu (\widehat{𝐏}_\mu \times \widehat{𝐪})>_f,<𝐒_\nu \widehat{𝐏}_\mu >_f,<\mathrm{𝟏}>_𝐟`$ \- the $`T`$odd and $`T`$even transverse components of neutrino polarization and the probability of muon capture, respectively; $`q,M`$ \- the value of the neutrino momentum and the nucleon mass; $`|𝐏_\mu |`$ \- the value of the muon polarization in $`1s`$ state; $`\varphi _\mu (0)`$ \- the value of the large radial component of the muon Dirac bispinor for $`r=0`$; $`\widehat{𝐪},\widehat{𝐏}_\mu `$ \- the direction of the neutrino momentum and of the muon polarization, respectively; $`\alpha _{AV}^L\alpha _A^L\alpha _V^L`$ \- the relative phase between the standard $`C_A^L`$ and $`C_V^L`$ couplings.
It can be seen that the neutrino observables consist exclusively of the interference term between the standard left-handed vector $`V_L`$ coupling and exotic right-handed scalar $`S_R`$ coupling. It can be understood as the interference between the neutrino waves of negative and positive handedness. If we admit also the presence of the left-handed scalar $`S_L`$ coupling in addition to the right-handed exotic scalar $`S_R`$ coupling, we get the interferences between the $`(V,A)_L`$ and $`S_L`$ couplings proportional to the neutrino mass, and they vanish in the limit of massless neutrino. Our coupling constants $`C_{V,A}^L,C_S^R`$ can be expressed by Fetscher’s couplings $`g_{ϵ\mu }^\gamma `$ for the normal and inverse muon decay Data , assuming the universality of weak interactions. The induced couplings generated by the dressing of hadrons are neglected as their presence does not change qualitatively the conclusions about transverse neutrino polarization. Here, $`\gamma =S,V,T`$ indicates a scalar, vector, tensor interaction; $`ϵ,\mu =L,R`$ indicate the chirality of the electron or muon and the neutrino chiralities are uniquely determined for given $`\gamma ,ϵ,\mu `$. We get the following relations:
$`C_V^L`$ $`=`$ $`A(g_{LL}^V+g_{RL}^V),`$ (6)
$`C_A^L`$ $`=`$ $`A(g_{LL}^Vg_{RL}^V),`$
$`C_S^R`$ $`=`$ $`A(g_{LR}^S+g_{RL}^S),`$
where $`A(4G_F/\sqrt{2})cos\theta _c`$, $`\theta _c`$ is the Cabbibo angle. In this way, the lower limits on the $`C_{V,A}^L`$ and upper limit on the $`C_S^R`$ can be calculated, using the current data Data ; $`|C_V^L|>0.850A,|C_A^L|>1.070A,|C_S^R|<0.974A`$. In consequence, one gives the upper bound on the magnitude of the transverse neutrino polarization vector proportional to the value of the muon polarization; $`|𝜼_𝝂^𝑻|0.334|𝐏_\mu |`$ (for $`\alpha _{AV}^L=0`$). The obtained limit has to be divided by the $`|𝐏_\mu |`$ to have the upper bound on the physical value of the transverse neutrino polarization vector generated by the exotic scalar interaction; $`|𝜼_𝝂^{}_{}{}^{\mathbf{}}𝑻|=|𝜼_𝝂^𝑻|/|𝐏_\mu |0.334`$.
To describe $`(\nu _\mu e^{})`$ scattering the following observables are used: $`𝜼_𝝂`$ \- the full vector of the initial neutrino polarization in the rest frame, $`𝐪`$ \- the incoming neutrino momentum, $`𝐩_𝐞^{}`$ \- the outgoing electron momentum. The laboratory differential cross section for the $`\nu _\mu e^{}`$ scattering, in the limit of vanishing neutrino mass, is of the form:
$`{\displaystyle \frac{d^2\sigma }{dyd\varphi _e^{}}}=({\displaystyle \frac{d^2\sigma }{dyd\varphi _e^{}}})_{(V,A)}+({\displaystyle \frac{d^2\sigma }{dyd\varphi _e^{}}})_{(S)}+({\displaystyle \frac{d^2\sigma }{dyd\varphi _e^{}}})_{(VS)},`$ (7)
$`({\displaystyle \frac{d^2\sigma }{dyd\varphi _e^{}}})_{(V,A)}`$ $`=`$ $`B\{(1𝜼_𝝂\widehat{𝐪})[(g_V^L+g_A^L)^2+(g_V^Lg_A^L)^2(1y)^2{\displaystyle \frac{m_ey}{E_\nu }}((g_V^L)^2(g_A^L)^2)]\},`$ (8)
$`({\displaystyle \frac{d^2\sigma }{dyd\varphi _e^{}}})_{(S)}`$ $`=`$ $`B\{{\displaystyle \frac{1}{8}}y(y+2{\displaystyle \frac{m_e}{E_\nu }})[|g_S^R|^2(1+𝜼_𝝂\widehat{𝐪})]\},`$ (9)
$`({\displaystyle \frac{d^2\sigma }{dyd\varphi _e^{}}})_{(VS)}`$ $`=`$ $`B\{\sqrt{y(y+2{\displaystyle \frac{m_e}{E_\nu }})}[𝜼_𝝂(\widehat{𝐩}_𝐞^{}\times \widehat{𝐪})Im(g_V^Lg_S^R)+(𝜼_𝝂\widehat{𝐩}_𝐞^{})Re(g_V^Lg_S^R)]`$
$`y(1+{\displaystyle \frac{m_e}{E_\nu }})(𝜼_𝝂\widehat{𝐪})Re(g_V^Lg_S^R)\},`$
$`B`$ $``$ $`{\displaystyle \frac{E_\nu m_e}{(2\pi )^2}}{\displaystyle \frac{G_F^2}{2}},`$ (11)
$`y`$ $``$ $`{\displaystyle \frac{E_e^{}^R}{E_\nu }}={\displaystyle \frac{m_e}{E_\nu }}{\displaystyle \frac{2cos^2\theta }{(1+\frac{m_e}{E_\nu })^2cos^2\theta }},`$ (12)
where $`y`$\- the ratio of the kinetic energy of the recoil electron $`E_e^{}^R`$ to the incoming neutrino energy $`E_\nu `$, $`\theta `$\- the angle between the direction of the outgoing electron momentum $`\widehat{𝐩}_e^{}`$ and the direction of the incoming neutrino momentum $`\widehat{𝐪}`$ (recoil electron scattering angle), $`m_e`$\- the electron mass,$`𝜼_𝝂\widehat{𝐪}`$ \- the longitudinal polarization of the incoming neutrino, $`\varphi _e^{}`$ \- the angle between the production plane and the reaction plane.
It can be noticed that the main non-standard contributions to the laboratory differential cross section come from the interference between the standard left-handed vector $`g_V^L`$ coupling and exotic right-handed scalar $`g_S^R`$ coupling, whose occurrence does not depend explicitly on the neutrino mass. This interference may be understood as the interference between the neutrino waves of the same handedness for the the final neutrinos. The similar regularity as for $`\mu ^{}`$-capture appeared here Sobkow . The correlation $`𝜼_𝝂\widehat{𝐩}_𝐞^{}`$ proportional to $`Re(g_V^Lg_S^R)`$ lies in the reaction plane spanned by the $`\widehat{𝐩}_e^{}`$ and $`\widehat{𝐪}`$, and it includes both $`T`$-even longitudinal and transverse components of the initial neutrino polarization:
$`(𝜼_𝝂\widehat{𝐩}_𝐞^{})`$ $`=`$ $`cos\theta (𝜼_𝝂\widehat{𝐪})+(𝜼_𝝂^{}_{}{}^{\mathbf{}}𝑻𝐩_𝐞^{}^𝐓),`$ (13)
where index ”T” describes transverse components of the outgoing electron momentum and of the incoming neutrino polarization, respectively. The other correlation $`𝜼_𝝂(\widehat{𝐩}_𝐞^{}\times \widehat{𝐪})`$ proportional to $`Im(g_V^Lg_S^R)`$ lies along the direction perpendicular to the reaction plane and it includes only $`T`$-odd transverse component of the initial neutrino polarization:
$`𝜼_𝝂(\widehat{𝐩}_𝐞^{}\times \widehat{𝐪})`$ $`=`$ $`𝜼_𝝂^{}_{}{}^{\mathbf{}}𝑻(\widehat{𝐩}_𝐞^{}\times \widehat{𝐪}).`$ (14)
It can be shown that in the full interference term, Eq. (10), the contributions from the longitudinal components of the neutrino polarization annihilate, and in consequence one gives the interference including only the transverse components of the initial neutrino polarization, both $`T`$-even and $`T`$-odd:
$`({\displaystyle \frac{d^2\sigma }{dyd\varphi _e^{}}})_{(VS)}`$ $`=`$ $`B\{\sqrt{{\displaystyle \frac{m_e}{E_\nu }}y[2(2+{\displaystyle \frac{m_e}{E_\nu }})y]}`$
$`\times |g_V^L||g_S^R||𝜼_𝝂^{}_{}{}^{\mathbf{}}𝑻|cos(\varphi \alpha )\},`$
where $`\alpha \alpha _V^L\alpha _S^R`$ \- the relative phase between the $`g_V^L`$ and $`g_S^R`$ couplings, $`\varphi `$ \- the angle between the reaction plane and the transverse neutrino polarization vector and it is connected with the $`\varphi _e^{}`$ in the following way; $`\varphi =\varphi _0\varphi _e^{}`$, where $`\varphi _0`$ \- the angle between the production plane and the transverse neutrino polarization vector. It can be noticed that the contribution from the interference between the $`g_V^L`$ and $`g_S^R`$ couplings, involving the transverse neutrino polarization components, will be substantial at lower neutrino energies $`E_\nu m_e`$ but negligibly small at large energies and vanishes for $`\theta =0`$ or $`\theta =\pi /2`$. The occurrence of the interference term in the cross section depends on the relative phase between the angle $`\varphi `$ and phase $`\alpha `$ and does not vanish for $`\varphi \alpha \pi /2`$. The situation is illustrated in the Fig.1, when $`m_e/E_\nu =1`$, $`y[0,0.66]`$, $`\varphi \alpha =0`$ (ESI1(y) - dashed line) and $`\varphi \alpha =\pi `$ (ESI2(y) - dashed line), respectively.
We use the present model-independent upper limits on the non-standard coupling constants Data obtained from the normal and inverse muon-decay, assuming the universality of weak interaction. We take the upper bound on the $`|𝜼_𝝂^{}_{}{}^{\mathbf{}}𝑻|0.334`$ for the neutrinos coming from the muon-capture (however the phase $`\alpha `$ is still unknown). The value $`|𝜼_𝝂^{}_{}{}^{\mathbf{}}𝑻|=0.334`$ is used to get the upper limit on the expected effect from the right-handed neutrinos in the cross section for the $`(\nu _\mu e^{})`$ scattering. It means that the value of the longitudinal neutrino polarization is equal to $`𝜼_𝝂^𝒍𝜼_𝝂\widehat{𝐪}=0.943`$. The plot for the SM is made with the use of the present experimental values for $`g_V^L=0.041\pm 0.015`$, $`g_A^L=0.507\pm 0.014`$ Data , $`𝜼_𝝂\widehat{𝐪}=1`$ and $`m_e/E_\nu =1`$, Fig.1 (SM(y) - solid line).
It is known that in the SM the angular distribution of the recoil electrons does not depend on the $`\varphi _e^{}`$. It is necessary to analyse all the possible reaction planes corresponding to the given recoil electron scattering angle, e.g. $`\theta =\pi /3`$, $`\theta =\pi /4`$, $`\theta =2\pi /9`$ (for $`E_\nu m_e)`$ to verify if the azimuthal asymmetry in the cross section appears. The regularity of this type would indicate the possible participation of the right-handed neutrinos in the $`(\nu _\mu e^{})`$ scattering.
The low-energy high-precision neutrino-electron scattering experiments using the beta-radioactive intense and polarized neutrino source could be used to search for new effects coming from the right-handed neutrinos, e. g. the Borexino neutrino experiment with the chromium source Miranda (to be published).
###### Acknowledgements.
I am greatly indebted to Prof. S. Ciechanowicz for many useful and helpful discussions.
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# Approximating Spectral invariants of Harper operators on graphs
## Introduction
Both the Harper operator and the discrete magnetic Laplacian (DML) on the Cayley graph of $`^2`$ have been extremely well studied in mathematical physics, as they arise as the Hamiltonian for the discrete model describing the quantum mechanics of free electrons in the presence of a magnetic field. In particular, the DML is the Hamiltonian of the discrete model for the integer quantum Hall effect, cf. \[Bel\]. These operators can be easily generalized to the Cayley graph of an arbitrary discrete group. This and a further generalization to general graphs with a free co-compact action of a discrete group with finite quotient, was defined by Sunada \[Sun\] and studied in the context of noncommutative Bloch theory and the fractional quantum Hall effect in \[CHMM\], \[CHM\], \[MM\], \[MM2\].
In this paper we are primarily concerned with the approximation of the spectral density function of the DML by examining restrictions of the DML over a sequence of finite subgraphs. Consider a graph $`X`$ which has a finite fundamental domain under the free action of an amenable group $`\mathrm{\Gamma }`$. By Følner’s characterization of amenability (see also \[Ad\]) there must exist a *regular exhaustion* of $`\mathrm{\Gamma }`$: a tower of finite subsets $`\mathrm{\Lambda }_m\mathrm{\Lambda }_{m+1}`$, $`_m\mathrm{\Lambda }_m=\mathrm{\Gamma }`$ satisfying
(0.1)
$$\underset{k\mathrm{}}{lim}\frac{\mathrm{\#}_\delta \mathrm{\Lambda }_k}{\mathrm{\#}\mathrm{\Lambda }_k}=0$$
where $`_\delta \mathrm{\Lambda }_k=\{\gamma \mathrm{\Gamma }:d_1(\gamma ,\mathrm{\Lambda }_k)<\delta \text{ and }d_1(\gamma ,\mathrm{\Gamma }\mathrm{\Lambda }_k)<\delta \}`$ is a $`\delta `$-neighborhood of the boundary of $`\mathrm{\Lambda }_k`$, and $`d_1`$ is the word metric on $`\mathrm{\Gamma }`$ with respect to some generating set. Given a choice $``$ of fundamental domain for the $`\mathrm{\Gamma }`$ action on $`X`$, a corresponding sequence of subgraphs $`X_m`$ is constructed, with $`X_m`$ being the largest subgraph of $`X`$ contained within $`\{\gamma |\gamma \mathrm{\Lambda }_m\}`$. These $`X_m`$ satisfy a similar property to (0.1):
(0.2)
$$\underset{m\mathrm{}}{lim}\frac{\mathrm{\#}\mathrm{Vert}_\delta X_m}{\mathrm{\#}\mathrm{Vert}X_m}=0\delta >0,$$
where $`_\delta X_m`$ refers to the subgraph which is the $`\delta `$-boundary of $`X_m`$ in $`X`$, that being the intersection of the $`\delta `$-neighbourhoods of $`X_m`$ and $`XX_m`$ in the simplicial metric.
The DML restricted to the space of functions over the vertices of $`X_m`$ is finite dimensional, and these restricted DMLs constitute a sequence of finite approximations to the DML itself. This leads us to the following conjecture.
###### Conjecture 0.1.
Let $`X`$ be a graph on which there is a free group action by an amenable group $`\mathrm{\Gamma }`$, with finite fundamental domain. Let $`\left\{X_m\right\}_{m=1}^{\mathrm{}}`$ be the regular exhaustion of $`X`$ corresponding to a regular exhaustion $`\mathrm{\Lambda }_m`$ of $`\mathrm{\Gamma }`$. Then
(0.3)
$$\underset{m\mathrm{}}{lim}\frac{E_m(\lambda )}{N_m}=F(\lambda ),\lambda ,$$
where $`N_m`$ is the size of $`\mathrm{\Lambda }_m`$, $`E_m(\lambda )`$ denotes the number of eigenvalues less than or equal to $`\lambda `$ of the DML restricted to $`X_m`$ with either Dirichlet or Neumann boundary conditions, and $`F`$ is the spectral density function of the DML on $`X`$, which is defined using the von Neumann trace.
We prove two approximation theorems which partially resolve this conjecture. The first states that the equality (0.3) holds at all but at most a countable set of points.
###### Theorem 0.2 (Rough approximation theorem for spectral density functions).
In the notation above, one has
(0.4)
$$\underset{m\mathrm{}}{lim}\frac{E_m(\lambda )}{N_m}=F(\lambda )$$
at every point of continuity of $`F`$. In particular, (0.4) holds for all but possibly a countable set of points.
Alternative proofs for this theorem exist in mathematical physics, using for instance the ‘Shubin formula’ (see \[Bel2\]). We thank J. Bellissard for bringing this reference to our attention.
Our main result is a refined approximation theorem, that holds only for DMLs associated with rational weight functions.
###### Theorem 0.3 (Refined approximation theorem; the rational case).
In the notation above, suppose that the DML is associated with a rational weight function. Then one has
(0.5)
$$\underset{m\mathrm{}}{lim}\frac{E_m(\lambda )}{N_m}=F(\lambda ),\lambda .$$
The strategy of proof of these approximation theorems is similar to that in \[DM\], but the details differ and are more involved in this case. The key lemma for the proof of Theorem 0.2 is the combinatorial analogue of the principle of not feeling the boundary.
Lemmas 2.1 and 2.2 in section 2 establish an approximation for $`F(\lambda )`$ in terms of the restricted DMLs $`\mathrm{\Delta }_\sigma ^{(m)}`$ and a sequence of polynomial approximations to the characteristic function $`\chi _{[0,\lambda ]}`$:
(0.6)
$$F(\lambda )=\underset{n\mathrm{}}{lim}\underset{m\mathrm{}}{lim}\frac{1}{N_m}\text{Tr}_{}\left(p_n\left(\mathrm{\Delta }_\sigma ^{(m)}\right)\right).$$
The proof of the refined approximation theorem 0.3 has at its core a variant of Lemma 2.8 in \[Lu\]. Such a result is called ‘Log Hölder continuity’ in the mathematical physics literature. By providing bounds on the growth of scaled spectral densities $`E_m(\lambda )/N_m`$, our variant of Lück’s lemma allows us to effectively interchange the limits in (0.6) and provide us with our result. Showing that the lemma is applicable is more involved in the case of the DML; we require a positive lower bound on the modulus of the product of non-zero eigenvalues of the restricted DMLs translated by $`\lambda `$. Such a bound for algebraic $`\lambda `$ is established in section 3, using some interesting facts about algebraic number fields and their rings of integers. Proposition 2.5, which states that the point spectrum of the DML with rational weight function must be contained within the algebraic numbers, completes the proof.
The conjecture is known to be true for all weight functions - not just rational ones - when $`X`$ is the Cayley graph of $`^2`$. In this case the spectral density function of the DML can be shown to be continuous (\[Sh\], \[DS\]), and the result then follows from Theorem 0.2. The proof that the spectral density function is continuous appears to utilize the special geometry of the Cayley graph of $`^2`$. We do not know if the spectral density function is continuous for general amenable graphs, and our proof of Theorem 0.3 is thus quite different to that of the $`^2`$ result.
In the final section of this paper we examine some consequences of these approximation theorems. Where the weight function of the DML is rational, Corollary 4.2 provides a criterion for the existence of spectral gaps which is in principle computable; the existence of spectral gaps for general DMLs is unknown and is a central open question in the area. Further in this section, we consider the Fuglede-Kadison determinant of DML$`\lambda `$, and show that it is greater than zero for most $`\lambda `$. This qualification on the spectrum of the DML may have applications in Mathematical Physics. In joint work with Dodziuk, Linnell and Schick \[DLMSY\], we have also worked out the analogs of the results of this paper for the combinatorial Laplacian on $`L^2`$ cochains on covering spaces of finite $`CW`$ complexes.
## 1. The Discrete magnetic Laplacian
Consider a combinatorial graph $`X`$ on which there is a free action by a group $`\mathrm{\Gamma }`$, with finite fundamental domain. The edge set $`\mathrm{Edge}X`$ is a collection of oriented edge; each combinatorial edge in $`X`$ has two corresponding oriented edges in $`\mathrm{Edge}X`$, one for each choice of orientation.
If $`e`$ is an oriented edge, $`\overline{e}`$ will denote the edge with opposite orientation, and $`𝔱(e)`$ and $`𝔬(e)`$ will denote the terminus and origin respectively.
It will be convenient to regard a subset $`E^+`$ of $`\mathrm{Edge}X`$ in which each combinatorial edge has exactly one oriented representative; $`E^+`$ corresponds to a choice of orientation for the graph.
If $`g`$ is a function defined on the edge set $`\mathrm{Edge}X`$ with values in some group $`H`$, then for compatibility with the orientation reversing operation on the edges, we demand that $`g(\overline{e})=g(e)^1`$. The graph structure gives a coboundary operator $`d`$ from the space of $`L^2`$ functions on the vertices of $`X`$ to the $`L^2`$ functions on the edges, giving the chain complex
$$0C_{(2)}^0(X)\stackrel{d}{}C_{(2)}^1(X)0$$
where
$$(df)(e)=f(𝔱(e))f(𝔬(e)).$$
Here $`𝔱(e)`$ denotes the terminus of the edge $`e`$, and $`𝔬(e)`$ the origin.
The adjoint $`d^{}`$ is given by
$$(d^{}g)(v)=\underset{\begin{array}{c}eE^+\\ 𝔱(e)=v\end{array}}{}g(e)\underset{\begin{array}{c}eE^+\\ 𝔬(e)=v\end{array}}{}g(e)$$
One can then define the discrete Laplacian on $`C_{(2)}^0(X)`$ by
$$\mathrm{\Delta }f=d^{}df$$
$$(\mathrm{\Delta }f)(v)=𝒪(v)f(v)\underset{\begin{array}{c}eE^+\\ 𝔱(e)=v\end{array}}{}f(𝔬(e))\underset{\begin{array}{c}eE^+\\ 𝔬(e)=v\end{array}}{}f(𝔱(e))$$
where $`𝒪(v)`$ is the valence of the vertex $`v`$.
Comparing this with the random walk operator on $`C_{(2)}^0(X)`$
$$(Rf)(v)=\underset{\begin{array}{c}eE^+\\ 𝔱(e)=v\end{array}}{}f(𝔬(e))+\underset{\begin{array}{c}eE^+\\ 𝔬(e)=v\end{array}}{}f(𝔱(e)),$$
we see that
$$(\mathrm{\Delta }f)(v)=𝒪(v)f(v)(Rf)(v).$$
The random walk operator is the basis of the generalized Harper operator of Sunada, where the value of $`f`$ at each of the neighbouring points is weighted with a complex number depending on the edge.
More precisely, consider the cochain of $`U(1)`$ valued functions on the graph $`X`$. Two functions in $`C^1(X;U(1))`$ are said to be equivalent if they belong in the same cohomology class. The weight function $`\sigma `$ that we choose on the graph must be weakly $`\mathrm{\Gamma }`$-invariant, meaning that $`\sigma `$ is equivalent to its left translation by any element of $`\mathrm{\Gamma }`$. Specifically, for each $`\gamma \mathrm{\Gamma }`$ there must exist a function $`s_\gamma :\mathrm{Vert}XU(1)`$ such that
$$\sigma (\gamma e)=\sigma (e)s_\gamma (𝔱(e))\overline{s_\gamma (𝔬(e))}.$$
The Harper operator is then given by
$$(H_\sigma f)(v)=\underset{\begin{array}{c}eE_v^+\\ 𝔱(e)=v\end{array}}{}\sigma (e)f(𝔬(e))+\underset{\begin{array}{c}eE_v^+\\ 𝔬(e)=v\end{array}}{}\overline{\sigma (e)}f(𝔱(e))$$
and one can define the discrete magnetic Laplacian (DML) in terms of $`H_\sigma `$,
(1.1)
$$(\mathrm{\Delta }_\sigma f)(v)=𝒪(v)f(v)(H_\sigma f)(v).$$
The discrete Laplacian arises naturally from the cochain complex on $`X`$. The DML can also be formed this way by examining a twisted complex with coboundary operator $`d_\tau `$
$$d_\tau :C_{(2)}^0(X)C_{(2)}^1(X)$$
$$(d_\tau f)(e)=\tau (e)f(𝔱(e))\overline{\tau (e)}f(𝔬(e)).$$
where $`\tau `$ is a weight function satisfying $`\tau ^2=\overline{\sigma }`$.
While the normal discrete Laplacian commutes with the the left $`\mathrm{\Gamma }`$-translations, the DML does not. In order to discuss its spectral properties we need to determine the appropriate von Neumann algebra in which to examine it.
The magnetic translation operators are the maps
$$(T_\gamma f)(x)=s_\gamma (\gamma ^1x)f(\gamma ^1x).$$
The form a projective representation of $`\mathrm{\Gamma }`$ with
(1.2)
$$T_\gamma T_\gamma ^{}=\mathrm{\Theta }(\gamma ,\gamma ^{})T_{\gamma \gamma ^{}}\gamma ,\gamma ^{}\mathrm{\Gamma },$$
where $`\mathrm{\Theta }`$ is a $`U(1)`$-valued group 2-cocycle determined by $`\sigma `$ up to cohomology.
The DML commutes with these operators, and the natural context for the examination of its spectral properties is the finite von Neumann algebra $`B(l^2(X))^{\mathrm{\Gamma },\sigma }`$ consisting of all the bounded operators that commute with the $`T_\gamma `$. Hereafter $`\text{Tr}_{\mathrm{\Gamma },\sigma }`$ will denote the unique trace on this algebra, with associated von Neumann dimension $`dim_{\mathrm{\Gamma },\sigma }`$.
More precisely, let $`AB(l^2(X))^{\mathrm{\Gamma },\sigma }`$ and $`k(v_1,v_2),v_i\mathrm{Vert}X`$ denote its kernel or matrix. It follows that
(1.3)
$$s_\gamma (v_1)k(\gamma v_1,\gamma v_2)s_\gamma (v_2)^1=k(v_1,v_2)$$
for all $`\gamma \mathrm{\Gamma }`$ and for all $`v_1,v_2\mathrm{Vert}X`$. The von Neumann trace of $`A`$ is given by the expression
(1.4)
$$\text{Tr}_{\mathrm{\Gamma },\sigma }(A)=\underset{v}{}k(v,v),$$
where $``$ denotes a fundamental domain for the action of $`\mathrm{\Gamma }`$ on $`X`$. The expression is well defined in view of (1.3).
To express the DML in terms of an equivariant coboundary operator, it is necessary to define first compatible magnetic translation operators on $`C_{(2)}^1(X)`$.
In order to define such translations that commute with the twisted coboundary operator, a further condition must be imposed upon $`s_\gamma `$. As the sole restriction on $`s_\gamma `$ is that $`ds_\gamma \sigma =\gamma ^{}\sigma `$, one can scale it by $`k_\gamma `$, constant on connected components of $`X`$. The expression $`s_\mu (x)s_\gamma (\mu x)\overline{s_{\gamma \mu }(x)}`$ is constant on connected components of $`X`$ (\[Sun\]) as is $`s_1(x)`$. So setting $`k_\gamma (x)^2=\overline{s_\gamma (x)s_{\gamma ^1}(\gamma x)}`$ and letting $`s_\gamma ^{}(x)=s_\gamma (x)k_\gamma (x)`$, one has $`ds_\gamma =ds_\gamma ^{}`$ and $`s_\gamma ^{}(x)^1=s_{\gamma ^1}^{}(\gamma x)`$. Hereafter it will be assumed that $`s_\gamma (x)`$ enjoys this normalization property.
Recall that the weakly $`\mathrm{\Gamma }`$-invariant weight function $`\tau `$ satisfies $`\tau ^2=\overline{\sigma }`$. Then one can choose $`t_\gamma `$ with $`t_\gamma ^2=\overline{s_\gamma }`$ and $`\tau (\gamma e)=\tau (e)t_\gamma (𝔱(e))\overline{t_\gamma (𝔬(e))}`$. $`t_\gamma `$ can be extended to edges by
$$t_\gamma (e)=\overline{t_\gamma (𝔱(e))t_\gamma (𝔬(e))}.$$
Extending the magnetic translation operators to $`C_{(2)}^1(X)`$ by
$$(T_\gamma g)(e)=t_\gamma (\gamma ^1e)g(\gamma ^1e)$$
then allows them to commute with the twisted coboundary operator $`d_\tau `$, while preserving the property 1.2.
In order to obtain the bounds required for the proofs of the approximation theorems in the next section, we need the following lemmas characterizing the near-diagonality of the DML and giving a common bound on the operator norm of the DML and its finite restrictions $`\mathrm{\Delta }_\sigma ^{(m)}`$.
Let $`d(v_1,v_2)`$ and $`d_m(v_1,v_2)`$ denote the distance between vertices $`v_1`$ and $`v_2`$ in the simplicial metric on $`X`$ and $`X_m`$ respectively. Recall that $`d(v_1,v_2)`$ is the length of the shortest sequence of edges $`\{e_1,\mathrm{},e_k\}`$ with $`𝔬(e_1)=v_1`$, $`𝔱(e_k)=v_2`$ and $`𝔱(e_i)=𝔬(e_{i+1})`$ for all $`i`$.
Let $`D(v_1,v_2)=\mathrm{\Delta }_\sigma \delta _{v_1},\delta _{v_2}`$ denote the matrix coefficient of the DML, $`\mathrm{\Delta }_\sigma `$, and $`D^{(m)}(v_1,v_2)=\mathrm{\Delta }_\sigma ^{(m)}\delta _{v_1},\delta _{v_2}`$ denote the matrix coefficient of the DML, $`\mathrm{\Delta }_\sigma ^{(m)}`$.
###### Lemma 1.1.
$`D(v_1,v_2)=0`$ whenever $`d(v_1,v_2)>1`$ and $`D^{(m)}(v_1,v_2)=0`$ whenever $`d_m(v_1,v_2)>1`$. There is also a positive constant $`C`$ independent of $`v_1,v_2`$ such that $`|D(v_1,v_2)|C`$ and $`|D^{(m)}(v_1,v_2)|C`$.
Let $`D^k(v_1,v_2)=\mathrm{\Delta }_\sigma ^k\delta _{v_1},\delta _{v_2}`$ denote the matrix coefficient of the $`k`$-th power of the DML, $`\mathrm{\Delta }_\sigma ^k`$, and $`D^{(m)k}(v_1,v_2)=\left(\mathrm{\Delta }_\sigma ^{(m)}\right)^k\delta _{v_1},\delta _{v_2}`$ denote the matrix coefficient of the $`k`$-th power of the DML, $`\mathrm{\Delta }_\sigma ^{(m)k}`$. Then
$$D^k(v_1,v_2)=\underset{u_1,\mathrm{}u_{k1}\mathrm{Vert}X}{}D(v_1,u_1)D(u_1,u_2)\mathrm{}D(u_{k1},v_2)$$
and
$$D^{(m)k}(v_1,v_2)=\underset{u_1,\mathrm{}u_{k1}\mathrm{Vert}X_m}{}D^{(m)}(v_1,u_1)D^{(m)}(u_1,u_2)\mathrm{}D^{(m)}(u_{k1},v_2)$$
Then the following lemma follows easily from Lemma 1.1.
###### Lemma 1.2.
Let $`k`$ be a positive integer. Then $`D^k(v_1,v_2)=0`$ whenever $`d(v_1,v_2)>k`$ and $`D^{(m)k}(v_1,v_2)=0`$ whenever $`d_m(v_1,v_2)>k`$. There is also a positive constant $`C`$ independent of $`v_1,v_2`$ such that $`|D^k(v_1,v_2)|C^k`$ and $`|D^{(m)k}(v_1,v_2)|C^k`$.
Since the DML $`\mathrm{\Delta }_\sigma ^k`$ commutes with magnetic translations, $`\mathrm{\Delta }_\sigma ^k`$ belongs to $`B(l^2(X))^{\mathrm{\Gamma },\sigma }`$ and the von Neumann trace is given as in (1.4),
(1.5)
$$\text{Tr}_{\mathrm{\Gamma },\sigma }(\mathrm{\Delta }_\sigma ^k)=\underset{v}{}D^k(v,v),$$
###### Lemma 1.3.
There is a positive number $`K`$ such that the operator norms of $`\mathrm{\Delta }_\sigma `$ and of $`\mathrm{\Delta }_\sigma ^{(m)}`$ for all $`m=1,2\mathrm{}`$ are smaller than $`K^2`$.
###### Proof.
The proof is similar to that in \[Lu\], Lemma 2.5 and uses Lemma 1.1 together with uniform local finiteness of $`X`$. More precisely we use the fact that the valency of any vertex in $`X`$ is uniformly bounded, say $`b`$. *A fortiori* the same is true (with the same constant $`b`$) for $`X_m`$ for all $`m`$. We now estimate the $`\mathrm{}^2`$ norm of $`\mathrm{\Delta }_\sigma \kappa `$ for a function $`\kappa =_va_v\delta _v`$ Now
$$\mathrm{\Delta }_\sigma \kappa =\underset{u}{}\left(\underset{v}{}D(u,v)a_v\right)\delta _u$$
so that
$`{\displaystyle \underset{u}{}}\left|{\displaystyle \underset{v}{}}D(u,v)a_v\right|^2`$ $`{\displaystyle \underset{u}{}}\left({\displaystyle \underset{d(u,v)1}{}}|D(u,v)|^2\right)\left({\displaystyle \underset{d(u,v)1}{}}|a_v|^2\right)`$
$`C^2b{\displaystyle \underset{u}{}}{\displaystyle \underset{d(u,v)1}{}}|a_v|^2,`$
where we have used Lemma 1.1 and Cauchy-Schwartz inequality. In the last sum above, for every vertex $`u`$, $`|a_u|^2`$ appears at most $`b`$ times. This proves that $`\mathrm{\Delta }_\sigma \kappa ^2C^2b^2\kappa ^2`$. Identical estimate holds (with the same proof) for $`\mathrm{\Delta }_\sigma ^{(m)}`$ which yields the lemma if we set $`K=\sqrt{Cb}`$. ∎
Let $`\{E(\lambda ):\lambda [0,\mathrm{})\}`$ denote the right continuous family of spectral projections of the DML $`\mathrm{\Delta }_\sigma `$. Since $`\mathrm{\Delta }_\sigma `$ commutes with magnetic translations, so does $`E(\lambda )=\chi _{[0,\lambda ]}(\mathrm{\Delta }_\sigma )`$, for any $`\lambda [0,\mathrm{})`$. Let $`F:[0,\mathrm{})[0,\mathrm{})`$ denote the spectral density function, where
$$F(\lambda )=\text{Tr}_{\mathrm{\Gamma },\sigma }(E(\lambda )).$$
Observe that the kernel of $`d`$ is also given by $`E(0)`$, and we have the following.
###### Lemma 1.4.
Let $`\mathrm{\Gamma }`$ be an infinite group and $`X`$ a connected graph. Then
$$0=F(0)=dim(ker(\mathrm{\Delta }_\sigma )).$$
###### Proof.
Observe that $`fL^2(X)`$ is in the kernel of the DML $`\mathrm{\Delta }_\sigma =d_\tau ^{}d_\tau `$ if and only if it is in the kernel of the twisted coboundary operator $`d_\tau `$. That is,
$$0=(d_\tau f)(e)=\tau (e)f(𝔱(e))\overline{\tau (e)}f(𝔬(e))$$
for all edges $`e|X|_1`$. It follows that $`|f(𝔱(e))|=|f(𝔬(e))|`$ for all edges $`e|X|_1`$. Since $`X`$ is connected, it follows that $`|f|`$ is a constant function on $`X`$. Since $`\mathrm{\Gamma }`$ is infinite, it follows that $`f`$ is identically zero. ∎
Let $`\mathrm{spec}(\mathrm{\Delta }_\sigma )`$ denote the $`L^2`$-spectrum of $`\mathrm{\Delta }_\sigma `$. Then one can show that
$$\mathrm{spec}(\mathrm{\Delta }_\sigma )=\{\lambda :F(\lambda +ϵ)F(\lambda ϵ)>0ϵ>0\}.$$
One can also show that $`\lambda `$ is an eigenvalue of $`\mathrm{\Delta }_\sigma `$ if and only if
(1.6)
$$\underset{ϵ0}{lim\; inf}\left\{F(\lambda +ϵ)F(\lambda ϵ)\right\}>0.$$
In fact if $`\lambda `$ is an eigenvalue of $`\mathrm{\Delta }_\sigma `$, and $`P_\lambda `$ is the projection onto the $`\lambda `$ eigenspace, then
$$P_\lambda =\underset{ϵ0}{lim}\{E(\lambda +ϵ)E(\lambda ϵ)\}.$$
## 2. Proofs of the approximation theorems
Let $`E_m(\lambda )`$ denote the number of eigenvalues $`\mu `$ of $`\mathrm{\Delta }_\sigma ^{(m)}`$ satisfying $`\mu \lambda `$ and which are counted with multiplicity.
We next make the following definitions,
(2.1)
$$\begin{array}{ccc}F_m(\lambda )\hfill & =& \frac{E_m(\lambda )}{N_m}\hfill \\ \overline{F}(\lambda )\hfill & =& \underset{m\mathrm{}}{lim\; sup}F_m(\lambda )\hfill \\ F(\lambda )\hfill & =& \underset{m\mathrm{}}{lim\; inf}F_m(\lambda )\hfill \\ \overline{F}^+(\lambda )\hfill & =& \underset{\delta +0}{lim}\overline{F}(\lambda +\delta )\hfill \\ F^+(\lambda )\hfill & =& \underset{\delta +0}{lim}F(\lambda +\delta ).\hfill \end{array}$$
The following lemma can be regarded as a combinatorial analogue of the principle of not feeling the boundary, cf. \[DM\].
###### Lemma 2.1.
Let $`\mathrm{\Gamma }`$ be an amenable group and let $`p(\lambda )=_{r=0}^da_r\lambda ^r`$ be a polynomial. Then,
$$\text{Tr}_{\mathrm{\Gamma },\sigma }(p(\mathrm{\Delta }_\sigma ))=\underset{m\mathrm{}}{lim}\frac{1}{N_m}\text{Tr}_{}\left(p\left(\mathrm{\Delta }_\sigma ^{(m)}\right)\right).$$
###### Proof.
First observe that if $`v\mathrm{Vert}X_m`$ is such that $`d(v,X_m)>k`$, then
$$D^k(v,v)=\mathrm{\Delta }_\sigma ^k\delta _v,\delta _v=\mathrm{\Delta }_\sigma ^{(m)k}\delta _v,\delta _v=D^{(m)k}(v,v).$$
By (1.5)
$$\text{Tr}_{\mathrm{\Gamma },\sigma }(p(\mathrm{\Delta }_\sigma ))=\frac{1}{N_m}\underset{v\mathrm{Vert}X_m}{}p(\mathrm{\Delta }_\sigma )v,v.$$
Therefore we see that
$$\begin{array}{c}\left|\text{Tr}_{\mathrm{\Gamma },\sigma }(p(\mathrm{\Delta }_\sigma ))\frac{1}{N_m}\text{Tr}_{}\left(p\left(\mathrm{\Delta }_\sigma ^{(m)}\right)\right)\right|\hfill \\ \hfill \frac{1}{N_m}\underset{r=0}{\overset{d}{}}|a_r|\underset{\begin{array}{c}v\mathrm{Vert}X_m\\ d(v,X_m)d\end{array}}{}\left(|D^r(v,v)|+|D^{(m)r}(v,v)|\right).\end{array}$$
Using Lemma 1.2, we see that there is a positive constant $`C`$ such that
$$\left|\text{Tr}_{\mathrm{\Gamma },\sigma }(p(\mathrm{\Delta }_\sigma ))\frac{1}{N_m}\text{Tr}_{}\left(p\left(\mathrm{\Delta }_\sigma ^{(m)}\right)\right)\right|2\underset{r=0}{\overset{d}{}}\frac{\mathrm{\#}\mathrm{Vert}_dX_m}{N_m}|a_r|C^r.$$
The proof of the lemma is completed by taking the limit as $`m\mathrm{}`$, as the $`X_m`$ form a regular exhaustion of $`X`$. ∎
We next recall the following abstract lemma of Lück \[Lu\] which is proved using the Lebesgue dominated convergence theorem.
###### Lemma 2.2.
Let $`p_n(\mu )`$ be a sequence of polynomials such that for the characteristic function of the interval $`[0,\lambda ]`$, $`\chi _{[0,\lambda ]}(\mu )`$, and an appropriate real number $`L`$,
$$\underset{n\mathrm{}}{lim}p_n(\mu )=\chi _{[0,\lambda ]}(\mu )\text{ and }|p_n(\mu )|L$$
holds for each $`\mu [0,\mathrm{\Delta }_\sigma ^2]`$. Then
$$\underset{n\mathrm{}}{lim}\text{Tr}_{\mathrm{\Gamma },\sigma }(p_n(\mathrm{\Delta }_\sigma ))=F(\lambda ).$$
Lemma 2.1 and Lemma 2.2 give
$$F(\lambda )=\underset{n\mathrm{}}{lim}\text{Tr}_{\mathrm{\Gamma },\sigma }(p_n(\mathrm{\Delta }_\sigma ))=\underset{n\mathrm{}}{lim}\underset{m\mathrm{}}{lim}\frac{1}{N_m}\text{Tr}_{}\left(p_n\left(\mathrm{\Delta }_\sigma ^{(m)}\right)\right).$$
If one could interchange the two limits on the right-hand side above, the proof of Theorem 0.3 would be complete. The following lemma is a variant of Lemma 2.8 in \[Lu\] and will be used in this paper to justify the interchange of the limits. Such a result is called ‘Log Hölder continuity’ in the mathematical physics literature, and was proved for the case $`\mathrm{\Gamma }=^k`$ by Craig and Simon \[CS\].
###### Lemma 2.3.
Let $`G:VV`$ be a self-adjoint linear map of the finite dimensional Hilbert space $`V`$. Let $`p(t)=det(tG)`$ be the characteristic polynomial of $`G`$. Then $`p(t)`$ can be factorized as $`p(t)=t^kq(t)`$ where $`q(t)`$ is a polynomial with $`q(0)0`$. Let $`K`$ be a real number, $`K^2\mathrm{max}\{1,G\}`$ and $`C>0`$ be a positive constant with $`|q(0)|C>0`$. Let $`E(\lambda )`$ be the number of eigenvalues $`\eta `$ of $`G`$, counted with multiplicity, satisfying $`\eta \lambda `$. Then for $`0<ϵ<1`$, the following estimate is satisfied.
$$\frac{\{E(ϵ)\}\{E(0)\}}{dim_{}V}\frac{\mathrm{log}C}{(\mathrm{log}ϵ)dim_{}V}+\frac{\mathrm{log}K^2}{\mathrm{log}ϵ}.$$
Applying this lemma to the self adjoint operator $`G\lambda `$, we obtain the following estimate for $`0<ϵ<1`$,
$$\frac{\{E(\lambda +ϵ)\}\{E(\lambda )\}}{dim_{}V}\frac{\mathrm{log}C(\lambda )}{(\mathrm{log}ϵ)dim_{}V}+\frac{\mathrm{log}(K^2+\lambda )}{\mathrm{log}ϵ}.$$
where $`C(\lambda )`$ is the lower bound for $`|q_\lambda (0)|`$, where $`p_\lambda (t)`$ is the characteristic polynomial of $`G\lambda `$, and $`p_\lambda (t)=t^{k(\lambda )}q_\lambda (t)`$ is its factorization such that $`q_\lambda (0)0`$.
To use this lemma, observe that $`\mathrm{\Delta }_\sigma ^{(m)}`$ can be regarded as a matrix with entries in the ring of integers of the number field generated by the weight function, whenever the weight function is rational. Therefore we have the following
###### Lemma 2.4.
Let $`\sigma `$ be a rational weight function, that is $`\sigma ^n=1`$ for some positive integer $`n`$. If $`\lambda `$ is a non-negative algebraic number, and $`p_{m,\lambda }(t)=det(t(\mathrm{\Delta }_\sigma ^{(m)}\lambda ))`$ is the characteristic polynomial of $`\mathrm{\Delta }_\sigma ^{(m)}\lambda `$, and $`p_{m,\lambda }(t)=t^kq_{m,\lambda }(t)`$ with $`q_{m,\lambda }(0)0`$, then there exist constants $`h`$ and $`Q`$ independent of $`m`$ such that $`|q_{m,\lambda }(0)|(aN_m)^hQ^{haN_m}`$, where $`a`$ is the size of the fundamental domain $``$.
###### Proof.
See section 3. ∎
In order to prove Theorem 0.3 for all $`\lambda `$, we will also need the following proposition characterizing the point spectra of the DML.
###### Proposition 2.5.
The point spectrum of the DML $`\mathrm{\Delta }_\sigma `$ is a subset of the union of the spectra of the restricted DMLs $`\mathrm{\Delta }_\sigma ^{(m)}`$. In particular when $`\sigma `$ is a rational weight function, if there exists $`\lambda \mathrm{spec}_{\text{point}}\mathrm{\Delta }_\sigma `$, then $`\lambda `$ is algebraic.
###### Proof.
Suppose $`\lambda `$ is not an eigenvalue of any of the $`\mathrm{\Delta }_\sigma ^{(m)}`$, that is $`\lambda _m\mathrm{spec}\mathrm{\Delta }_\sigma ^{(m)}`$. Let $`P_\lambda `$ be the projection onto the eigenspace for $`\lambda `$ of $`\mathrm{\Delta }_\sigma `$. Let $`X_m^o=_1X_m(XX_m)`$ be the outer $`1`$-boundary of $`X_m`$ and let $`P_m`$ and $`P_m^o`$ be the projections onto the functions with support in $`X_m`$ and $`X_m^o`$ respectively. One has
$`\text{Tr}_{\mathrm{\Gamma },\sigma }P_\lambda `$ $`={\displaystyle \frac{1}{N_m}}\text{Tr}_{}P_mP_\lambda `$
$`{\displaystyle \frac{1}{N_m}}dim\text{Im}P_mP_\lambda \text{as }P_mP_\lambda 1\text{,}`$
giving
$`\text{Tr}_{\mathrm{\Gamma },\sigma }P_\lambda `$ $`\underset{m\mathrm{}}{lim\; inf}{\displaystyle \frac{1}{N_m}}dim\text{Im}P_mP_\lambda .`$
Note that by lemma 1.1, $`P_m\mathrm{\Delta }_\sigma `$ differs from $`\mathrm{\Delta }_\sigma ^{(m)}P_m`$ only on $`_1X_m`$. One can write
$`P_m\mathrm{\Delta }_\sigma =P_m\mathrm{\Delta }_\sigma P_m+B_mP_m^o`$
$`=\mathrm{\Delta }_\sigma ^{(m)}P_m+B_mP_m^o`$
where $`B_m:C^0(X_m^o)C^0(X_m_1X_m)`$ encodes this difference.
Consider $`f\text{Im}P_mP_\lambda `$, $`f=P_mg`$, $`g`$ an eigenfunction of $`\mathrm{\Delta }_\sigma `$ for $`\lambda `$. From $`\mathrm{\Delta }_\sigma g=\lambda g`$ we have
$$P_m\mathrm{\Delta }_\sigma g=\lambda P_mg,$$
and so
(2.2)
$$(\mathrm{\Delta }_\sigma ^{(m)}P_m+B_mP_m^o)g=\lambda P_mg.$$
Consider two solutions $`g_1`$ and $`g_2`$ of this equation, with $`g_1|_{X_m^o}=g_2|_{X_m^o}`$. Then $`B_mP_m^o(g_1g_2)=0`$ and we have
$$\mathrm{\Delta }_\sigma ^{(m)}P_m(g_1g_2)=\lambda P_m(g_1g_2).$$
As $`\lambda \mathrm{spec}\mathrm{\Delta }_\sigma ^{(m)}`$, this implies that $`P_m(g_1g_2)=0`$. So the value of $`g`$ on $`X_m^o`$ uniquely determines $`f=P_mg`$, giving
$`dim\text{Im}P_mP_\lambda `$ $`dimC^0(X_m^o)`$
$`\mathrm{\#}_1X_m`$
and thus
$`\text{Tr}_{\mathrm{\Gamma },\sigma }P_\lambda `$ $`\underset{m\mathrm{}}{lim}{\displaystyle \frac{1}{N_m}}\mathrm{\#}_1X_m`$
$`=0,`$
demonstrating that $`\lambda `$ is not in the point spectrum of $`\mathrm{\Delta }_\sigma `$. As the restricted DMLs $`\mathrm{\Delta }_\sigma ^{(m)}`$ are all finite operators with matrix elements belonging to the set of algebraic numbers, the union of their spectra must in turn be a subset of the algebraic numbers; this result therefore implies that any $`\lambda `$ in the point spectrum of the DML must be algebraic. ∎
Remark. This lemma is easily generalized to any local operator $`A`$ in the von Neumann algebra whose components $`A(u,v)=A\delta _u,\delta _v`$ are algebraic, where we regard as local an operator $`A`$ for whom $`A(u,v)=0`$ for all $`u`$ and $`v`$ with $`d(u,v)>k`$ for some fixed $`k`$.
Theorem 0.2 follows immediately from part (i) of the following theorem, and Theorem 0.3 follows immediately from part (ii) and the preceeding proposition 2.5.
###### Theorem 2.6.
Let $`\mathrm{\Gamma }`$ be an amenable discrete group.
* In the notation of (2.1), one has
$$F(\lambda )=\overline{F}^+(\lambda )=F^+(\lambda ).$$
* If in addition, one assumes that the weight function $`\sigma `$ is rational, then
$$F(\lambda )=\overline{F}(\lambda )=F(\lambda ),$$
whenever $`\lambda `$ is an algebraic number. That is, under these assumptions, one has,
$$F(\lambda )=\underset{m\mathrm{}}{lim}F_m(\lambda ).$$
###### Proof.
Fix $`\lambda 0`$ and define for $`n1`$ a continuous function $`f_n:`$ by
$$f_n(\mu )=\{\begin{array}{ccc}1+\frac{1}{n}\hfill & \text{ if }& \mu \lambda \hfill \\ 1+\frac{1}{n}n(\mu \lambda )\hfill & \text{ if }& \lambda \mu \lambda +\frac{1}{n}\hfill \\ \frac{1}{n}\hfill & \text{ if }& \lambda +\frac{1}{n}\mu \hfill \end{array}$$
Then clearly $`\chi _{[0,\lambda ]}(\mu )<f_{n+1}(\mu )<f_n(\mu )`$ and $`f_n(\mu )\chi _{[0,\lambda ]}(\mu )`$ as $`n\mathrm{}`$ for all $`\mu [0,\mathrm{})`$. For each $`n`$, choose a polynomial $`p_n`$ such that $`\chi _{[0,\lambda ]}(\mu )<p_n(\mu )<f_n(\mu )`$ holds for all $`\mu [0,K^2]`$. We can always find such a polynomial by a sufficiently close approximation of $`f_{n+1}`$. Hence
$$\chi _{[0,\lambda ]}(\mu )<p_n(\mu )<2$$
and
$$\underset{n\mathrm{}}{lim}p_n(\mu )=\chi _{[0,\lambda ]}(\mu )$$
for all $`\mu [0,K^2]`$. Recall that $`E_m(\lambda )`$ denotes the number of eigenvalues $`\mu `$ of $`\mathrm{\Delta }_\sigma ^{(m)}`$ satisfying $`\mu \lambda `$ and counted with multiplicity. Note that $`\mathrm{\Delta }_\sigma ^{(m)}K^2`$ by Lemma 1.3.
$$\begin{array}{ccc}\frac{1}{N_m}\text{Tr}_{}\left(p_n(\mathrm{\Delta }_\sigma ^{(m)})\right)\hfill & =& \frac{1}{N_m}\underset{\mu [0,K^2]}{}p_n(\mu )\hfill \\ & =& \frac{E_m(\lambda )}{N_m}+\frac{1}{N_m}\{\underset{\mu [0,\lambda ]}{}(p_n(\mu )1)+\underset{\mu (\lambda ,\lambda +1/n]}{}p_n(\mu )\hfill \\ & & +\underset{\mu (\lambda +1/n,K^2]}{}p_n(\mu )\}\hfill \end{array}$$
Hence, we see that
(2.3)
$$F_m(\lambda )=\frac{E_m(\lambda )}{N_m}\frac{1}{N_m}\text{Tr}_{}\left(p_n(\mathrm{\Delta }_\sigma ^{(m)})\right).$$
In addition,
$$\begin{array}{ccc}\frac{1}{N_m}\text{Tr}_{}\left(p_n(\mathrm{\Delta }_\sigma ^{(m)})\right)\hfill & & \frac{E_m(\lambda )}{N_m}+\frac{1}{N_m}sup\{p_n(\mu )1:\mu [0,\lambda ]\}E_m(\lambda )\hfill \\ & +& \frac{1}{N_m}sup\{p_n(\mu ):\mu [\lambda ,\lambda +1/n]\}(E_m(\lambda +1/n)E_m(\lambda ))\hfill \\ & +& \frac{1}{N_m}sup\{p_n(\mu ):\mu [\lambda +1/n,K^2]\}(E_m(K^2)E_m(\lambda +1/n))\hfill \\ & & \frac{E_m(\lambda )}{N_m}+\frac{E_m(\lambda )}{nN_m}+\frac{(1+1/n)(E_m(\lambda +1/n)E_m(\lambda ))}{N_m}\hfill \\ & & +\frac{(E_m(K^2)E_m(\lambda +1/n))}{nN_m}\hfill \\ & & \frac{E_m(\lambda +1/n)}{N_m}+\frac{1}{n}\frac{E_m(K^2)}{N_m}\hfill \\ & & F_m(\lambda +1/n)+\frac{a}{n}\hfill \end{array}$$
since $`E_m(K^2)=dimC^0(X_m)=aN_m`$ for a positive constant $`a`$ independent of $`m`$, $`a`$ being the size of the fundamental domain $``$. It follows that
(2.4)
$$\frac{1}{N_m}\text{Tr}_{}\left(p_n(\mathrm{\Delta }_\sigma ^{(m)})\right)F_m(\lambda +1/n)+\frac{a}{n}.$$
Taking the limit inferior in (2.4) and the limit superior in (2.3), as $`m\mathrm{}`$, we get that
(2.5)
$$\overline{F}(\lambda )\text{Tr}_{\mathrm{\Gamma },\sigma }\left(p_n(\mathrm{\Delta }_\sigma )\right)F(\lambda +1/n)+\frac{a}{n}.$$
Taking the limit as $`n\mathrm{}`$ in (2.5) and using Lemma 2.2, we see that
$$\overline{F}(\lambda )F(\lambda )F^+(\lambda ).$$
For all $`\epsilon >0`$ we have
$$F(\lambda )F^+(\lambda )F(\lambda +\epsilon )\overline{F}(\lambda +\epsilon )F(\lambda +\epsilon ).$$
Since $`F`$ is right continuous, we see that
$$F(\lambda )=\overline{F}^+(\lambda )=F^+(\lambda )$$
proving part (i) of Theorem 2.6.
Next we assume that $`\sigma `$ be a rational weight function, that is $`\sigma ^n=1`$ for some positive integer $`n`$. In the notation of Lemma 2.4, let $`\lambda `$ be a non-negative algebraic number, and $`p_{m,\lambda }(t)=det(t(\mathrm{\Delta }_\sigma ^{(m)}\lambda ))`$ be the characteristic polynomial of $`\mathrm{\Delta }_\sigma ^{(m)}\lambda `$, and $`p_{m,\lambda }(t)=t^kq_{m,\lambda }(t)`$ with $`q_{m,\lambda }(0)0`$. Then $`|q_{m,\lambda }(0)|(aN_m)^hQ^{haN_m}`$, $`h`$ and $`Q`$ being constants independent of $`m`$. By Lemma 2.3 and the remarks following it, for $`ϵ>0`$ one has,
$$\frac{F_m(\lambda +ϵ)F_m(\lambda )}{a}\frac{\mathrm{log}((aN_m)^hQ^{haN_m})}{(\mathrm{log}ϵ)aN_m}+\frac{\mathrm{log}(K^2+\lambda )}{\mathrm{log}ϵ}.$$
That is,
(2.6)
$$F_m(\lambda +ϵ)F_m(\lambda )\frac{h(\mathrm{log}aN_m+aN_m\mathrm{log}Q)}{(\mathrm{log}ϵ)N_m}\frac{a\mathrm{log}(K^2+\lambda )}{\mathrm{log}ϵ}.$$
Taking limit inferior in (2.6) as $`m\mathrm{}`$ yields
$$F(\lambda +ϵ)F(\lambda )\frac{a\mathrm{log}Q}{\mathrm{log}ϵ}\frac{a\mathrm{log}(K^2+\lambda )}{\mathrm{log}ϵ}.$$
Passing to the limit as $`ϵ+0`$, we obtain $`F(\lambda )=F^+(\lambda )`$. A similar argument establishes that $`\overline{F}(\lambda )=\overline{F}^+(\lambda )`$. By part (i) of Theorem 2.6, one has $`\overline{F}^+(\lambda )=F(\lambda )=F^+(\lambda )`$. Therefore we see that
$$\overline{F}(\lambda )=\overline{F}^+(\lambda )=F(\lambda )=F^+(\lambda )=F(\lambda )$$
which proves part (ii) of Theorem 2.6. ∎
## 3. Determining a lower bound for $`q_{m,\lambda }(0)`$
That we can determine the lower bound in lemma 2.4 depends upon the components of the matrix $`\mathrm{\Delta }_\sigma ^{(m)}`$ taking values in a ring of algebraic integers. Before providing a proof of the lemma, some definitions and elementary results of algebraic number theory are presented.
The weight function $`\sigma `$ takes values in $`U(1)`$. When it is rational, that is there is some integer $`n`$ for which $`\sigma ^n=1`$, these values are $`n`$th roots of unity. The matrix elements in $`\mathrm{\Delta }_\sigma ^{(m)}`$ are integer-linear combinations of the values taken by $`\sigma `$, and so lie in $`[\zeta ]`$ where $`\zeta `$ is a primitive $`n`$th root of unity.
The natural context for examining the properties of this matrix is the algebraic number field.
###### Definition 3.1.
An extension $`F`$ of $``$ of finite degree is an algebraic number field. The degree of the extension $`|F:|=dim_{}F`$ is called the degree of $`F`$. The integral closure of $``$ in $`F`$ (that is, the subring of $`F`$ consisting of elements $`\alpha `$ that satisfy $`f(\alpha )=0`$ for some monic polynomial $`f[X]`$) is termed the ring of (algebraic) integers of $`F`$, and written $`O_F`$.
For a given algebraic number $`\lambda `$, let $`\lambda =\eta /b`$ where $`\eta `$ is an algebraic integer and $`b`$. The field in which the DML will be examined will then be $`(\zeta ,\eta )`$. The matrix elements of $`\mathrm{\Delta }_\sigma ^{(m)}`$ are elements of $`[\zeta ]`$ and thus reside within the ring of integers of this field. Let $`h`$ be the degree of this extension, which will depend on the denominator $`n`$ of the rational weight function $`\sigma `$ and this algebraic integer $`\eta `$.
The extension $`(\zeta ,\eta )`$ is $`(\alpha )`$ for some algebraic $`\alpha `$. The minimum polynomial of $`\alpha `$ will be of degree $`h`$. The other roots of this polynomial - the conjugates of $`\alpha `$ \- determine the set $`\{e_i\}`$ of $``$-preserving embeddings of $`(\alpha )`$ into $``$.
The set of embeddings of an algebraic number field determine a norm on that field.
###### Definition 3.2.
Let $`F`$ be an algebraic number field of degree $`h`$, with distinct embeddings $`e_1,e_2,\mathrm{},e_h`$ into $``$. Then the norm on $`F`$ is defined as the product
$$N(\alpha )=\underset{i=1}{\overset{h}{}}e_i(\alpha ).$$
The norm has the following important property:
###### Theorem 3.3.
Let $`N`$ be the norm on an algebraic number field $`F`$. Then $`N(\alpha )`$ for all $`\alpha F`$, and further, $`N(\alpha )`$ for all $`\alpha `$ in the ring of integers $`O_F`$.
In particular, $`N(\alpha )`$ for all $`\alpha [\zeta ]`$.
From this one can arrive at the following well-known result, presented in \[Fa\] and repeated here.
###### Lemma 3.4.
Let $`F`$ be an algebraic number field of degree $`h`$, with ring of integers $`O_F`$ and distinct embeddings $`e_1,e_2,\mathrm{},e_h`$ into $``$. Let $`\alpha `$ be a non-zero element of $`O_F`$ satisfying
$$|e_i(\alpha )|Ri.$$
Then
$$|e_i(\alpha )|R^{1h}i.$$
Proof of lemma 2.4. We follow here an argument of Farber \[Fa\]. As discussed above, let $`\{e_i\}`$ be the $`h`$ $``$-preserving embeddings of $`(\zeta ,\eta )=(\alpha )`$ into $``$, and let $`e_1`$ be the identity embedding.
Now $`\zeta ^i=r_i(\alpha )`$ for some polynomials $`r_i`$ of degree less than $`h`$. There will be an upper bound on these co-efficients, as the number of distinct $`\zeta ^i`$ is finite. As there is also a bound on the absolute values of the conjugates of $`\alpha `$, there is an upper bound $`R`$:
$`|e_j(\zeta ^i)|`$ $`Ri,j`$
$`|e_j(\eta )|`$ $`Rj`$
Let the matrix $`A_j=e_j(\mathrm{\Delta }_\sigma ^{(m)})`$, applying $`e_j`$ component-wise. Then the matrix elements of $`A_j`$ satisfy the same conditions as those of the DML in Lemma 1.1 (for a different constant $`C^{}`$) and by the same procedure by which a bound $`K^2`$ was determined for $`\mathrm{\Delta }_\sigma ^{(m)}`$, one can find $`L`$ such that $`A_jL^2`$, and thus
$`|\text{Tr}(A_j^i)|`$ $`aN_mA_j`$
$`aN_mL^{2i}`$
The $`r`$th symmetric polynomial of eigenvalues of $`A_j`$ is therefore bounded,
$`|s_r(A_j)|`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{aN_m+r1}{r}}\right).L^{2r}`$
$`2^{2aN_m1}.L^{2r}\text{(as }raN_m\text{)}`$
$`<4^{2aN_m}.L^{2aN_m}=(4L^2)^{aN_m}.`$
As the matrix elements of $`\mathrm{\Delta }_\sigma ^{(m)}`$ are algebraic integers in $`(\alpha )`$, so are the coefficients $`c_i`$ of the characteristic polynomial $`p_m(t)`$. One has $`|e_j(c_i)|=|s_i(A_j)|<(4L^2)^{aN_m}`$.
Consider $`p_m(t+\lambda )=q_{m,\lambda }(t)t^k`$. The coefficient of $`t^k`$ in $`p_m(t+\lambda )`$ then is equal to $`q_{m,\lambda }(0)`$.
$`p_m(t+\lambda )`$ $`=c_{aN_m}(t+\lambda )^{aN_m}+c_{aN_m1}(t+\lambda )^{aN_m1}+\mathrm{}+c_0`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{aN_m}{0}}\right)c_{aN_m}\lambda ^0t^{aN_m}+`$
$`\left(\left({\displaystyle \genfrac{}{}{0pt}{}{aN_m}{1}}\right)c_{aN_m}\lambda ^1+\left({\displaystyle \genfrac{}{}{0pt}{}{aN_m1}{0}}\right)c_{aN_m1}\lambda ^0\right)t^{aN_m1}+`$
$`\left(\left({\displaystyle \genfrac{}{}{0pt}{}{aN_m}{2}}\right)c_{aN_m}\lambda ^2+\left({\displaystyle \genfrac{}{}{0pt}{}{aN_m1}{1}}\right)c_{aN_m1}\lambda ^1+\left({\displaystyle \genfrac{}{}{0pt}{}{aN_m2}{0}}\right)c_{aN_m2}\lambda ^0\right)t^{aN_m2}+`$
$`\mathrm{}+`$
$`\left(\left({\displaystyle \genfrac{}{}{0pt}{}{aN_m}{aN_m}}\right)c_{aN_m}\lambda ^{aN_m}+\mathrm{}+\left({\displaystyle \genfrac{}{}{0pt}{}{0}{0}}\right)c_0\lambda ^0\right)t^0`$
where $`c_{aN_m}`$=1. So
$$q_{m,\lambda }(0)=\underset{i=0}{\overset{aN_mk}{}}\left(\genfrac{}{}{0pt}{}{i+k}{k}\right)c_{i+k}\lambda ^i.$$
Recall that $`\lambda =\eta /b`$ with $`\eta `$ an algebraic integer in $`(\alpha )`$. $`b^{aN_m}q_{m,\lambda }(0)`$ is then an integer linear combination of the $`c_i`$ and powers of $`\eta `$ which makes it an algebraic integer in $`(\alpha )`$. As $`2^{aN_m}\left(\genfrac{}{}{0pt}{}{aN_m}{k}\right)\left(\genfrac{}{}{0pt}{}{aN_ml}{k}\right)`$, we get the following upper bound on $`|e_j(b^{aN_m}q_{m,\lambda }(0))|`$:
$`|e_j(b^{aN_m}q_{m,\lambda }(0))|`$ $`=\left|e_j\left({\displaystyle \underset{i=0}{\overset{aN_mk}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{i+k}{k}}\right)c_{i+k}\eta ^ib^{aN_mi}\right)\right|`$
$`=\left|{\displaystyle \underset{i=0}{\overset{aN_mk}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{i+k}{k}}\right)b^{aN_mi}e_j(\eta )^ie_j(c_{i+k})\right|`$
$`{\displaystyle \underset{i=0}{\overset{aN_mk}{}}}2^{aN_m}b^{aN_m}R^{aN_m}(4L^2)^{aN_m}`$
$`(8RbL^2)^{aN_m}aN_m.`$
By lemma 3.4 then,
$$|e_j(b^{aN_m}q_{m,\lambda }(0))|(8RbL^2)^{haN_m}(aN_m)^hj$$
and in particular
$$|b^{aN_m}q_{m,\lambda }(0)|=|e_1(b^{aN_m}q_{m,\lambda }(0))|(8RbL^2)^{haN_m}(aN_m)^h,$$
giving
$`|q_{m,\lambda }(0)|`$ $`b^{aN_m}(8RbL^2)^{haN_m}(aN_m)^h`$
$`(8Rb^2L^2)^{haN_m}(aN_m)^h.`$
Thus for suitable choices of constants $`h`$ and $`Q`$ independent of $`m`$, one can write the inequality as
$$|q_{m,\lambda }(0)|Q^{haN_m}(aN_m)^h.$$
In the case where $`\lambda =a/b0`$ is rational, one has $`|q_{m,\lambda }(0)|(8B^2K^2)^{naN_m}(aN_m)^n`$ where $`B`$ is the maximum of $`a`$ and $`b`$, and $`K^2`$ is the bound on the norm of the DML as given by lemma 1.3.
## 4. The Fuglede-Kadison determinant of DML$`\lambda `$
In this section, we show that the Fuglede-Kadison determinant of DML$`\lambda `$ is positive for most $`\lambda `$. We begin though by deriving some corollaries to the main theorems in the previous sections, using the notation there. The following corollary should be compared to Proposition 2.5.
###### Corollary 4.1.
Let $`\mathrm{spec}(\mathrm{\Delta }_\sigma )`$ denote the spectrum of the DML $`\mathrm{\Delta }_\sigma `$ and $`\mathrm{spec}(\mathrm{\Delta }_\sigma ^{(m)})`$ denote the spectrum of the DML $`\mathrm{\Delta }_\sigma ^{(m)}`$. Then one has
$$\mathrm{spec}(\mathrm{\Delta }_\sigma )\overline{\underset{m1}{}\mathrm{spec}(\mathrm{\Delta }_\sigma ^{(m)})}$$
###### Proof.
Let $`\lambda _1,\lambda _2`$ be points of continuity of the spectral density function $`F`$ of the DML $`\mathrm{\Delta }_\sigma `$ with $`\lambda _1<\lambda _2`$. Then by Theorem 0.2, one has
$$\underset{m\mathrm{}}{lim}\left(F_m(\lambda _2)F_m(\lambda _1)\right)=F(\lambda _2)F(\lambda _1)$$
We also notice that
$$F(\lambda _2)>F(\lambda _1)\mathrm{spec}(\mathrm{\Delta }_\sigma )(\lambda _1,\lambda _2)\mathrm{}$$
and
$$F_m(\lambda _2)>F_m(\lambda _1)\mathrm{spec}(\mathrm{\Delta }_\sigma ^{(m)})(\lambda _1,\lambda _2)\mathrm{}.$$
This immediately implies the corollary. ∎
###### Corollary 4.2 (Spectral gap criterion).
The interval $`(\lambda _1,\lambda _2)`$ is in a gap in the spectrum of the DML $`\mathrm{\Delta }_\sigma `$ with rational weight function $`\sigma `$ if and only if
$$\underset{m\mathrm{}}{lim}\left(F_m(\lambda _2)F_m(\lambda _1)\right)=0$$
###### Proof.
Notice that the interval $`(\lambda _1,\lambda _2)`$ is in a gap in the spectrum of the DML $`\mathrm{\Delta }_\sigma `$ if and only if $`F(\lambda _2)=F(\lambda _1)`$. Then by Theorem 0.3, one has
$$\underset{m\mathrm{}}{lim}\left(F_m(\lambda _2)F_m(\lambda _1)\right)=F(\lambda _2)F(\lambda _1)=0.$$
The converse is proved similarly to the proof above. ∎
###### Corollary 4.3 (Spectral density estimate).
Suppose that $`\sigma `$ is a rational weight function, and $`F`$ be the spectral density function of the DML $`\mathrm{\Delta }_\sigma `$. Then there are positive constants $`C`$ and $`\delta `$ such that for all $`ϵ(0,\delta )`$, one has for any algebraic $`\mu `$,
$$F(\mu +ϵ)F(\mu )\frac{C}{\mathrm{log}ϵ}.$$
###### Proof.
This corollary follows immediately from Lemma 2.3 and Lemma 2.4, together with Theorem 0.3. ∎
### 4.1. Fuglede-Kadison determinants
Here we investigate under what conditions the Fuglede-Kadison determinant of $`\mathrm{\Delta }_\sigma \mu `$ is positive, for rational weight function $`\sigma `$. Recall that the Fuglede-Kadison determinant $`det_{\mathrm{\Gamma },\sigma }`$ of an operator $`A`$ with spectral density function $`G`$ is defined by
(4.1)
$$\mathrm{log}\underset{\mathrm{\Gamma },\sigma }{det}(A)=_{0^+<\lambda L}\mathrm{log}|\lambda |dG(\lambda )$$
for some $`L>A`$, whenever the right hand side of (4.1) is not equal to $`\mathrm{}`$, and $`det_{\mathrm{\Gamma },\sigma }(A)`$ is defined to be zero if the right hand side of (4.1) equals $`\mathrm{}`$, cf. \[FK\]. The following proposition improves on the estimate in Corollary 4.3.
###### Proposition 4.4.
Let $`\sigma `$ be a rational weight function. Then for $`\mu `$ satisfying any of
1. $`\mu \mathrm{spec}\mathrm{\Delta }_\sigma `$,
2. $`\mu =0`$,
3. $`\mu `$ algebraic and $`\mu _m\mathrm{spec}\mathrm{\Delta }_\sigma ^{(m)}`$,
the Fuglede-Kadison determinant $`det_{\mathrm{\Gamma },\sigma }(\mathrm{\Delta }_\sigma \mu )`$ is positive.
###### Proof.
Condition 1. In this instance, the operator $`\mathrm{\Delta }_\sigma \mu `$ is invertible, and positivity of the Fuglede-Kadison determinant follows immediately.
Condition 2. The argument for the positivity of $`\mathrm{log}det_{\mathrm{\Gamma },\sigma }\mathrm{\Delta }_\sigma `$ exactly parallels that of section 4 of \[DM\], using the lower bound on the modified determinant $`|det^{}|\mathrm{\Delta }_\sigma ^{(m)}=q_{m,0}(0)`$ determined in lemma 2.4.
Condition 3. This case calls for a slightly modified version of the argument found in section 4 of \[DM\], \[BFK\], but using the results in sections 2 and 3 instead.
For notational convenience, let $`A`$ be the DML offset by $`\mu `$, $`A=\mathrm{\Delta }_\sigma \mu `$, with von Neumann spectral density function $`G(\lambda )=F(\lambda +\mu )`$ where $`F`$ is the spectral density function of $`\mathrm{\Delta }_\sigma `$.
Let $`A^{(m)}`$ be the restricted DML similarly displaced by $`\mu `$, $`A^{(m)}=\mathrm{\Delta }_\sigma ^{(m)}\mu `$. Recall that the *normalized* spectral density functions of $`\mathrm{\Delta }_\sigma ^{(m)}`$
$$F_m(\lambda )=\frac{1}{N_m}E_j^{(m)}(\lambda )$$
are right continuous, and note that the corresponding normalized spectral density functions of $`A^{(m)}`$ will be given simply by
$$G_m(\lambda )=F_m(\lambda +\mu ).$$
Observe that $`G_m(\lambda )`$ are step functions. Recall from Lemma 1.3 that there exists some $`K`$ such that $`\mathrm{\Delta }_\sigma <K^2`$, and $`\mathrm{\Delta }_\sigma ^{(m)}<K^2`$ for all $`m`$. In the following, let $`L`$ be an algebraic number greater than $`K^2+|\mu |`$, so that $`A<L`$ and $`A^{(m)}<L`$.
Denote by $`|det^{}|A^{(m)}`$ the modified determinant of $`A^{(m)}`$, that is the absolute value of the product of all nonzero eigenvalues of $`A^{(m)}`$. Choose positive $`a`$ and $`b`$ such that $`a`$ is less than the least absolute value of any non-zero eigenvalue of $`A^{(m)}`$, and $`b`$ is greater than the largest absolute value of any eigenvalue of $`A^{(m)}`$. Then
(4.2)
$$\frac{1}{N_m}\mathrm{log}|\stackrel{}{det}|A^{(m)}=_{a|\lambda |b}\mathrm{log}|\lambda |dG_m(\lambda ).$$
Integration by parts transforms this Stieltjes integral as follows.
(4.3)
$$\begin{array}{c}_{a|\lambda |b}\mathrm{log}|\lambda |dG_m(\lambda )=(\mathrm{log}b)\left(G_m(b)G_m(b)\right)\hfill \\ \hfill (\mathrm{log}a)\left(G_m(a)G_m(a)\right)_{a|\lambda |b}\frac{G_m(\lambda )G_m(0)}{\lambda }𝑑\lambda .\end{array}$$
Integrating 4.1 by parts, one obtains
(4.4)
$$\begin{array}{c}\mathrm{log}\underset{\mathrm{\Gamma },\sigma }{det}(A)=(\mathrm{log}L)\left(G(L)G(L)\right)\hfill \\ \hfill +\underset{ϵ0^+}{lim}\left\{(\mathrm{log}ϵ)\left(G(ϵ)G(ϵ)\right)_{ϵ|\lambda |L}\frac{G(\lambda )G(0)}{\lambda }𝑑\lambda \right\}.\end{array}$$
Using the fact that $`lim\; inf_{ϵ0^+}(\mathrm{log}ϵ)\left(G(ϵ)G(ϵ)\right)0`$ one sees that
(4.5)
$$\mathrm{log}\underset{\mathrm{\Gamma },\sigma }{det}(A)(\mathrm{log}L)\left(G(L)G(L)\right)_{0^+|\lambda |L}\frac{G(\lambda )G(0)}{\lambda }𝑑\lambda .$$
We now complete the proof by estimating a lower bound of $`\mathrm{log}det_{\mathrm{\Gamma },\sigma }(A)`$ from a lower bound on $`\frac{1}{N_m}\mathrm{log}|det^{}|A^{(m)}`$. $`|det^{}|A^{(m)}`$ is the absolute value of the product of all the non-zero eigenvalues of $`A^{(m)}`$, and hence for a rational weight function $`\sigma `$ by lemma 2.4 (and using the notation described there)
$`|\stackrel{}{det}|A^{(m)}`$ $`=|q_{m,\mu }(0)|`$
$`Q^{haN_m}(aN_m)^h`$
and thence
(4.6)
$$\frac{1}{N_m}\mathrm{log}|\stackrel{}{det}|A^{(m)}ha\mathrm{log}Q\frac{h}{N_m}\mathrm{log}(aN_m)$$
for some positive constants $`Q`$ and $`h`$ independent of $`m`$.
The following estimate is proved exactly as in Lemma 2.6 of \[DM\], so we will omit its proof here.
(4.7)
$$_{0^+|\lambda |L}\frac{G(\lambda )G(0)}{\lambda }𝑑\lambda \underset{m\mathrm{}}{lim\; inf}_{0^+|\lambda |L}\frac{G_m(\lambda )G_m(0)}{\lambda }𝑑\lambda .$$
Combining (4.2) and (4.3) with the inequality (4.6) we obtain
$$\begin{array}{c}_{0^+|\lambda |L}\frac{G_m(\lambda )G_m(0)}{\lambda }𝑑\lambda (\mathrm{log}L)\left(G_m(L)G_m(L)\right)\hfill \\ \hfill \underset{ϵ0^+}{lim}(\mathrm{log}ϵ)(G_m(ϵ)G_m(ϵ))+ha\mathrm{log}Q+\frac{h}{N_m}\mathrm{log}(aN_m).\end{array}$$
By the condition under regard, $`\mu \mathrm{spec}\mathrm{\Delta }_\sigma ^{(m)}`$, giving
$$G_m(ϵ)G_m(ϵ)=F_m(\mu +ϵ)F_m(\mu ϵ)=0\text{for small }ϵ,$$
and so one has
(4.8)
$$_{0^+|\lambda |L}\frac{G_m(\lambda )G_m(0)}{\lambda }𝑑\lambda (\mathrm{log}L)\left(G_m(L)G_m(L)\right)+ha\mathrm{log}Q+\frac{h}{N_m}\mathrm{log}(aN_m).$$
From (4.5), (4.7) and (4.8), we conclude that
(4.9)
$$\begin{array}{c}\mathrm{log}\underset{\mathrm{\Gamma },\sigma }{det}\mathrm{\Delta }_\sigma (\mathrm{log}L)\left(G(L)G(L)\right)ha\mathrm{log}Q\hfill \\ \hfill \underset{m\mathrm{}}{lim\; inf}\left\{(\mathrm{log}L)\left(G_m(L)G_m(L)\right)+\frac{h}{N_m}\mathrm{log}(aN_m)\right\}\end{array}$$
Now $`h/N_m\mathrm{log}(aN_m)0`$ and by part (ii) of Theorem 2.6
$`G(L)=F(L+\mu )`$ $`=\underset{m\mathrm{}}{lim}F_m(L+\mu )`$
$`=\underset{m\mathrm{}}{lim}G_m(L)`$
(and similarly for $`L`$), hence
$$\mathrm{log}\underset{\mathrm{\Gamma },\sigma }{det}Ah\mathrm{log}Q,$$
that is,
$$\underset{\mathrm{\Gamma },\sigma }{det}(\mathrm{\Delta }_\sigma \mu )>0.$$
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# Reconstructing PSCz with a Generalised PIZA
## 1 Introduction
Cosmological density and velocity fields may be reconstructed using the radial velocity information from a galaxy redshift survey. To do this, assumptions must be made about the way the peculiar velocity field is produced, the relations between the velocity and density fields, and the relationship between galaxies and the underlying mass distribution. The reconstructed fields may then be compared with the directly measured fields, e.g. from a peculiar velocity survey, and used to place constraints on cosmological parameters.
The structures seen in the Universe today formed by the growth via gravitational instability of small density perturbations present in the early Universe. The growth and evolution of these structures may be studied using methods of two types: Eulerian methods that solve the gravitational instability equations (see Peebles 1980); and Lagrangian methods that follow individual galaxy displacements. These methods may be used to deduce the velocity field from the density field and vice versa.
Eulerian methods include those which use linear perturbation theory applied to the gravitational instability equations. The iterative reconstruction methods of Yahil et al. , who studied the 1.2Jy redshift survey, and Kaiser et al. and Taylor and Rowan-Robinson who studied the QDOT redshift survey, use this method to reconstruct the peculiar velocity field from galaxy redshift surveys. Branchini et al. use a similar method to reconstruct the peculiar velocity of the PSCz.
Lagrangian methods include the Zel’dovich approximation (ZA; Zel’dovich 1970), which was suggested as a way to study the evolution of structure in the nonlinear regime. Within the ZA the linear displacement field of galaxies is extrapolated from initial, Lagrangian coordinates to final, Eulerian coordinates. The galaxy peculiar velocity is simply the time derivative of the linear displacement.
The use of Hamilton’s principle, that the action is minimized during the evolution of a collection of particles under gravity, was introduced as a reconstruction method to cosmology by Peebles and developed by Peebles (1994, 1995), Shaya, Peebles and Tully , Giavalisco et al. and Susperregi and Binney . Sharpe et al. recently used Peebles’ numerical action method to reconstruct the peculiar velocities of PSCz galaxies within $`cz=2000\mathrm{k}\mathrm{m}s^1`$, while Monaco and Efstathiou have developed ZTRACE, an iterative reconstruction method based on Lagrangian perturbation theory.
A method combining the ZA and the least action principle is the Path Interchange Zel’dovich Approximation (PIZA) of Croft and Gaztañaga (hereafter CG97). They showed that using the ZA the least action solution is that which minimizes the total mean square particle displacement between initial and final galaxy positions. Hence, given a set of final galaxy positions and a random distribution of initial coordinates, the least action trajectories are those that minimize the total displacement.
Previously PIZA has been used to reconstruct the initial density fields from simulated real space density fields (CG97; Narayanan and Croft 1999). Here we wish to apply PIZA to a redshift survey in order to reconstruct the peculiar velocity and real space density fields. We apply our new PIZA method to the Point Source Catalogue Redshift survey (PSCz; Saunders at al 2000), a survey of 15,000 IRAS galaxies to a flux limit of 0.6 Jy at 60$`\mu m`$, covering 84 per cent of the sky.
The outline of the paper is as follows. We review PIZA in Section 2. In Section 3 we describe the problems encountered when applying PIZA to redshift surveys, and the generalizations we have made to the method in order to overcome these problems. In Section 4 we test PIZA using PSCz-like simulations, and compare the predictions for the fields with those given by linear theory. In Section 5 we apply our new method to the PSCz survey. In Section 6 we discuss our results and present our conclusions.
## 2 Path Interchange Zel’dovich Approximation
Here we briefly review the original PIZA method. For more details, we refer the reader to Croft & Gaztañaga .
We define a galaxy position in terms of its initial position and displacement:
$$𝒓(𝒒,t)=𝒒+\xi (𝒒,t),$$
(1)
where $`𝒓`$ is the Eulerian coordinate, $`𝒒`$ is the initial Lagrangian coordinate and $`𝝃`$ is the galaxy displacement.
CG97 (§2.1) showed that the action is proportional to the mass–weighted squared particle displacement,
$$S=\underset{i}{}m_i\xi _i^2.$$
(2)
where $`m_i`$ is the mass of the $`i^{th}`$ particle. The least action principle then requires the minimization of this quantity. In CG97 it is assumed that the galaxies are of equal mass, simplifying the minimization of equation (2). Assuming the final, Eulerian galaxy positions are known, the initial positions are drawn from a homogeneous distribution and assigned at random to the Eulerian galaxy positions. We shall refer to the initial objects as PIZA particles. It is not necessary that there is only one particle per galaxy; CG97 found that increasing the number of particles per galaxy improved the accuracy of the reconstruction. Applying PIZA is very simple. Two galaxies are then picked at random and their PIZA particles interchanged. If this decreases $`S`$ the interchange is kept, and if not the particles are swapped back. Interchanges are attempted until the rate of change of $`S`$ with attempted swaps, the cooling rate, is very low.
## 3 PIZA with realistic galaxy redshift surveys
The PIZA method cannot be immediately applied to a realistic galaxy redshift survey. In the first instance the radial positions of galaxies are estimated from their redshift, which introduces distortions in galaxy positions due to the peculiar velocities of galaxies. This distortion will depend on the rest frame of the observer, which we shall assume is the Local Group frame. To apply PIZA we must first deal with this distortion. Secondly most galaxy redshift surveys are flux limited, so that the observed number of galaxies drops off with distance as only the brightest galaxies can be seen. This affects the number density of galaxies, which must now be regarded as sampling the total galaxy population. This means that the galaxies must be assigned effective masses which must be included in the PIZA minimization. Finally, the sky coverage of galaxies may not be complete and some way must be included in the reconstruction to account for regions of incompleteness. In this Section we shall discuss each of these effects in turn.
### 3.1 PIZA in redshift space
The peculiar motions of galaxies, $`𝒗`$, mean that in the CMB frame the radial coordinate in redshift space is a mixture of spatial position and projected velocity,
$$𝒔(𝒓,t)=𝒓+(\widehat{𝒓}.𝒗(𝒓,t))\widehat{𝒓},$$
(3)
where $`𝒔`$ is the redshift position. In the ZA the velocity field of galaxies is a linear extrapolation, and related to the displacement field by
$$𝒗=\dot{𝝃}=f(\mathrm{\Omega })𝝃,$$
(4)
where $`fd\mathrm{ln}\delta /d\mathrm{ln}a\mathrm{\Omega }^{0.6}`$ is the growth rate of linear perturbations . If galaxies are a locally biased representation of the density field then $`\delta _{\mathrm{gal}}=F[\delta _{\mathrm{mass}}]b\delta _{\mathrm{mass}}`$ to first order. We can then define a galaxy displacement field which is related to the velocity field by
$$𝒗=\beta 𝝃$$
(5)
where
$$\beta =\frac{f(\mathrm{\Omega })}{b}$$
(6)
is the redshift distortion parameter. The redshift space displacement field is related to the real space displacement field by
$$\xi _i^s=𝒫_{ij}\xi _j,$$
(7)
where $`𝒫_{ij}\delta _{ij}^K+\beta \widehat{r}_i\widehat{r}_j`$ is the redshift space projection tensor, and $`\delta _{ij}^K`$ is the Kronecker tensor. Inverting equation (7), gives :
$$\xi _i=𝒫_{ij}^1\xi _j^s,$$
(8)
where
$$𝒫_{ij}^1\delta _{ij}^K\frac{\beta }{1+\beta }\widehat{r}_i\widehat{r}_j,$$
(9)
is the inverse redshift space projection tensor. Squaring this gives us the expression for the linear displacement field in terms of linear redshift displacements and $`\beta `$;
$$\xi ^2=\xi ^{s^2}(1\frac{1}{(1+\beta )^2})\xi _r^{s^2},$$
(10)
where $`\xi _r^s=𝝃.\widehat{𝒓}`$ is the displacement vector in redshift space projected along the line of sight. Thus we must assume an a priori value for $`\beta `$.
### 3.2 PIZA with Local Group frame redshifts
Redshift space distortions should also take into account the motion of the observer, in this case the Local Group motion. The transformation from real to redshift space in the Local Group rest frame is
$$\xi _i=𝒫_{ij}^1\xi _j^{s,LG}+(\delta _{ij}^K𝒫_{ij}^1)\xi _j(\mathrm{𝟎})),$$
(11)
where the super- or sub-script LG denotes a variable measured in the Local Group frame, and $`𝝃(\mathrm{𝟎})`$ is the displacement of the Local Group. The real-space displacement field can now be expressed in terms of redshift displacements in the Local Group frame as
$$\xi ^2=\xi _{LG}^{s^2}\xi _{r,LG}^{s^2}+\frac{(\beta \xi _r(\mathrm{𝟎})+\xi _{r,LG}^s)^2}{(1+\beta )^2}$$
(12)
To apply equation (12) we require the displacement of the Local Group. We find this by placing a galaxy representing the Local Group at the origin and using its displacement as the Local Group displacement.
### 3.3 Incomplete Sky Coverage
The incomplete sky coverage may cause problems with PIZA, as with other reconstruction methods, since in these regions we do not know the true galaxy distribution. A number of choices for filling in these regions present themselves. One could populate the incomplete regions randomly on the sky, assuming no information or leave the regions empty of both particles and galaxies. These are equivalent except for particles adjacent to the mask. Alternatively, one could interpolate across using correlations in the density field to fill in the gaps. In reconstructing PSCz we use the Fourier Interpolation Scheme of Saunders et al. which allows optimal, nonlinear interpolation across these regions.
### 3.4 PIZA with a selection function
In a flux limited survey, the number density of survey objects decreases with increasing redshift, as galaxies at greater distances must be brighter to be seen above the flux limit. This is quantified by the survey selection function, $`\varphi (r)`$, the expected number of galaxies seen at distance $`r`$ above the flux limit in the absence of clustering.
In order to account for this change in the number density with distance, we pick initial PIZA particles with the same selection function as the galaxies. To take account of those galaxies that fall below the flux limit we assign each galaxy a mass of $`1/\varphi (r)`$ and each PIZA particle a mass of
$$m_i=1/(\nu \varphi (r_i)),$$
(13)
where $`\nu `$ is the mean number of particles assigned per galaxy. In this work we use 10 particles per galaxy. We use $`\nu =10`$ in order to reduce the shot noise in the reconstruction. Equations 12 and 13 are used in equation 2 to minimize the displacement squared.
With the inclusion of particle masses it is possible for S to decrease but light particles to acquire unrealistically large trajectories. We have therefore added an additional constraint that the rms displacement squared in redshift space should also decrease.
### 3.5 Second Order Effects
The masses we have assigned to the galaxies are based on the selection function at their redshift distance. Because of peculiar velocities, these can be in error. We therefore developed a modified version of PIZA, MASSPIZA, whereby after PIZA is run, particles are randomly reassigned from galaxies whose mass has been overestimated to those where it has been underestimated. PIZA is then rerun and the whole process iterated until the masses are acceptable and the trajectories minimised. Unfortunately, this procedure is in general unstable to the rocket effects described by Kaiser (1987). In the section below, we show that PIZA without mass reassignment gives acceptable results in comparison with simulations; the agreement is in general not improved by the second, MASSPIZA stage. Therefore, we have not incorporated it in any of the results quoted in this paper.
## 4 Tests on simulations
In order to test the new PIZA method, we have applied it to a set of PSCz-like mock catalogues. For more details of these mock catalogues, we refer the reader to Branchini et al. . The catalogues are of two CDM cosmologies, a critical model , $`\mathrm{\Omega }_m=1`$, $`\beta =1`$, with $`\mathrm{\Gamma }=0.25`$ and a flat model with $`\mathrm{\Omega }_m=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`\beta =0.50`$ and are limited to $`v<15000kms^1`$. We test the reconstruction of the radial peculiar velocity field, comparing it with that of the simulation. To compare the fields the peculiar velocities were binned in real space and smoothed with a Gaussian smoothing kernel. We have also compared the PIZA reconstructions with the linear theory reconstructions for the mock catalogues from Branchini et al. Figure 1 shows typical results comparing PIZA and linear theory reconstructions with the simulations. The dotted line indicates the the least squares fit to the data.
As stated above, we must assume an a priori value for $`\beta `$. We find that the reconstruction works well for any sensible value of $`\beta `$. We find that generally the gradient reconstructed by PIZA is less biased than that reconstructed by linear theory but that the scatter is somewhat worse. This is probably one of the limitations of PIZA as the scatter in redshift space reconstructions produced by PIZA is quite large (see also CG97).
We find that PIZA works better for low density models as they are less dynamically evolved.
The accuracy of the reconstruction varies between different sets of initial conditions (i.e. PIZA particles) due to the random positioning of particles. There is also a scatter between reconstructions using the same initial conditions due to the random swaps made by PIZA. This probably reflects that there is not a simple minimum for the action, $`S`$, but that the configuration space is complicated by the particle discreteness.
We estimate $`\beta `$ from the trajectory of the Local Group galaxy using equation 5 and assuming the magnitude of the CMB dipole. Table 1 shows the reconstructed $`\beta `$, the dipole misalignment angle, and the slope of the peculiar velocity comparison, with the scatter due to different initial conditions. We have tested the effects of having the Galactic Plane filled or empty. This is important for PSCz as we know there are large structures, such as the Great Attractor, in the Plane. With the simulations we have tested PIZA using $`4\pi `$ sky coverage, and with the PSCz mask applied. With the mask empty we find that the average misalignment angle is $`24^{}`$, and with full-sky coverage it is $`14^{}`$.
## 5 Reconstructing PSCz
We have applied our new PIZA method to the PSCz survey (Saunders et al. 2000) out to $`R=300h^1\mathrm{Mpc}`$, using the parametrized IRAS selection function of Mann, Saunders and Taylor . For the reconstructions we have filled the masked regions with interpolated galaxies.
Figure 2 shows the reconstructed PSCz velocity field in a slice $`20h^1\mathrm{Mpc}`$ thick centred on the Supergalactic Plane. Figure 3 shows the same slice out to $`50h^1\mathrm{Mpc}`$. The arrows show the galaxy trajectories, with the arrows starting at the initial positions and the arrow heads at the galaxy real space positions. Generally the velocity field at the edge of the volume is quiet, apart from flow towards Shapley at (120,90) and the void at (0,-120). The Coma cluster is seen at (0,70). The Cetus wall extends northwards from (0,-100) towards Perseus-Pisces at (50,-20). Virgo is seen at (0,10). The Great Attractor is at (-40,20), and we see little evidence of backfall on to it on the opposite side to the G.A. from the Local Group.
### 5.1 Cosmological dipole
We find an average dipole direction of (l,b)=(264.4,41.7). We find that $`\beta =0.51\pm 0.14`$. This error has 3 sources: random error due to different initial conditions; misalignment between the reconstructed and CMB dipoles; cosmic variance and shot noise on the reconstructed dipole . A similar result is found by only using the z component of the dipole.
### 5.2 Radial velocity field
The reconstructed dipole and bulk flows are shown in Figure 4. We have compared our bulk flows with those from Mark III and the linear theory PSCz reconstruction (Branchini et al. 1999), and conclude that $`\beta =0.5\pm 0.15`$. This is consistent with the results of Tadros et al (1999) based on the measurement of the degree of redshift space distortion in the PSCz.
### 5.3 Comparison with SFI
We have undertaken a comparison of the reconstructed PSCz velocity field with the SFI data set of peculiar velocities of spiral galaxies (Giovanelli et al. 1998). We have calculated the average peculiar velocity, with respect to the Local Group, of galaxies in redshift slices, to be as closely as possible comparable to Giovanelli et al.. We find reasonable agreement in amplitude, though there are significant differences in direction, for $`\beta =0.55\pm 0.1`$. Table 2 shows the average peculiar velocity in shells for PIZA and SFI.
## 6 Discussion
We have presented a generalized version of the PIZA method for reconstructing velocity and density fields from galaxy surveys. These modifications to the original method are necessary if it is to be applied to realistic galaxy redshift surveys with a selection function. We have used the inverse redshift space operator of Taylor and Valentine to write the mean square particle displacement in terms of the redshift space positions of galaxies, and to map redshifts from the Local Group rest frame to real space. The PIZA particles obey the same selection function as the galaxies. Doing this we can reconstruct the real space density field, the galaxy initial positions and the peculiar velocity field.
We have tested the method using PSCz-like CDM simulations, testing the reconstruction of the real space positions and the peculiar velocities. We have also compared our reconstructions with linear theory reconstructions of the simulations. We have applied our method to the PSCz survey, and reconstructed the local velocity field out to $`300h^1\mathrm{Mpc}`$ from the Local Group, and have calculated the bulk flow and the dipole. We find the reconstructed dipole to be $`16^{}\pm 4`$ away from the CMB dipole. From comparison between the PIZA reconstruction and the Mark III bulk flow, we find that $`\beta =0.5\pm 0.15`$, while from a dipole analysis of the PIZA reconstruction, we find that $`\beta =0.51\pm 0.14`$.
## 7 Acknowledgements
HEMV thanks PPARC for a studentship. We thank Enzo Branchini for the simulations and linear theory reconstructions.
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# 1 Charge Quantization in the Standard Model
## 1 Charge Quantization in the Standard Model
The SM assumes a repetitive structure for each generation of quarks and leptons . Let us start by looking at the first generation of quarks and leptons (u, d, e, $`\nu `$ ) and assign them to $`SU(N_C)SU(2)_LU(1)_Y`$ (where $`N_C=3`$ ) representation. To keep things as general as possible this brings in five unknown hypercharges . Let us now define the electric charge in terms of the diagonal generators of $`SU(2)_LU(1)_Y`$ as
$$Q=T_3+bY$$
(1)
In the SM $`SU(N_C)`$ $``$ $`SU(2)_L`$ $``$ $`U(1)_Y`$ is spontaneously broken through the Higgs mechanism to the group $`SU(3)_c`$ $``$ $`U(1)_{em}`$ . In SM the Higgs is assumed to be a doublet . However we do not use this restriction either and assume the Higgs $`\varphi `$ to have any isospin T and arbitrary hypercharge $`Y_\varphi `$. The isospin $`T_3^\varphi `$ component of the Higgs develops a nonzero vacuum expectation value $`<\varphi >_o`$. Since we want the $`U(1)_{em}`$ generator Q to be unbroken we require $`Q<\varphi >_o=0`$. This right away fixes b in (3) and we get
$$Q=T_3(\frac{T_3^\varphi }{Y_\varphi })Y$$
(2)
For the SM to be renormalizable we require that the triangular anomaly be canceled. This leads to certain constraints which place some restrictions on the hypercharges. Also the use of the fact that after the spontaneous breaking of $`SU(N_C)SU(2)_LU(1)_Y`$ to $`SU(N_C)U(1)_{em}`$, the L- and R-handed charges couple identically with photon helps in pinning down hypercharges in terms of the hypercharge of the Higgs. Hence one obtains quantized electric charges as
$$Q(u)=\frac{1}{2}(1+\frac{1}{N_c})$$
$$Q(d)=\frac{1}{2}(1+\frac{1}{N_c})$$
(3)
$`Q(e)=1`$ (4)
$`Q(\nu )=0`$ (5)
For $`N_C=3`$ these are the correct charges in the SM. Note that this charge quantization in the SM holds for Higgs for arbitrary T and arbitrary hypercharge. Hence as far as charge quantization is concerned, the values of T and $`Y_\varphi `$ remain unconstrained. This point, for the special case of the Higgs doublet was already noted by the author earlier .
## 2 Higgs particle - a ghost !
Let us continue with the rest of the structure of the SM and see how our general Higgs with unconstrained and unspecified isospin T and hypercharge $`Y_\varphi `$ fits into it. We can write the covariant derivative of the SM as
$$D_\mu =_\mu +ig_1\frac{T_3^\varphi }{Y_\varphi }YB_\mu ig_2\stackrel{}{T}.\stackrel{}{W_\mu }$$
(6)
The photon field $`A_\mu `$ and the orthogonal $`Z_\mu `$ are written as
$$A_\mu =\frac{g_2B_\mu +g_1(\frac{2T_3^\varphi }{Y_\varphi }Y_l)W_\mu ^0}{\sqrt{g_2^2+(g_1\frac{2T_3^\varphi }{Y_\varphi }Y_l)^2}}$$
(7)
$$Z_\mu =\frac{g_1(\frac{2T_3^\varphi }{Y_\varphi }Y_L)B_\mu +g_2W_\mu ^0}{\sqrt{g_2^2+(g_1\frac{2T_3^\varphi }{Y_\varphi }Y_l)^2}}$$
(8)
With $`D_\mu `$ we can write the lepton part of the SM Lagrangian as
$`(lepton)=\overline{q_L}i\gamma ^\mu (ig_1{\displaystyle \frac{T_3^\varphi }{Y_\varphi }}Y_lB_\mu )q_L+\overline{e_R}i\gamma ^\mu (ig_1{\displaystyle \frac{T_3^\varphi }{Y_\varphi }}Y_eB_\mu )e_R`$
$`\overline{q_L}i\gamma ^\mu [ig_2\stackrel{}{T}.\stackrel{}{W_\mu }]q_L`$ (9)
We find that with the hypercharges as specified above the complete structure of the Standard Model stands intact . The point to be emphasized is that all this is independent of Higgs isospin and hypercharge, which all throughout remain unconstrained and undetermined. One should not fix any arbitrary values for them as nothing in the theory demands it.
We find that from rho parameter also the solutions for the isospin of the Higgs are infinite in number. Again, nothing in the theory demands that one fixes the isospin to a particular value.
The point is that the full structure of the SM stands intact without constraining the quantum numbers isospin and/or the hypercharge of the Higgs to any specific value. All the particles that have been isolated in the laboratory or have been studied by any other means, besides having a specific mass, have definite quantum numbers which identify them. In the case of Higgs here, no one knows of its mass and more importantly its quantum numbers like isospin and hypercharge, as shown above, are not specified. The hypercharge of all the other particles are specified as being proportional to the Higgs hypercharge which itself remains unconstrained. That is, all the hypercharges of particles are rooted on to the Higgs hypercharge which itself remains free and unspecified. Hence Higgs is very different from all known particles. Because of the above reasons Higgs cannot be a physical particle which may be isolated and studied. It must be just the ‘vacuum’ which sets up the structure of the whole thing. So Higgs is a ghost which shall not manifest itself as a genuine particle in the laboratory .
## 3 Superstring Theories-intrinsically flawed!
Now we ask the question, with this generalized picture what happens to the electric charge when the full Standard Model symmetry is restored.
Note that the expression for Q in (2) arose due to spontaneous symmetry breaking of $`SU(N_C)SU(2)_L\times U(1)_Y`$ (for $`N_C=3`$ ) to $`SU(N_C)\times U(1)_{em}`$ through the medium of a Higgs with arbitrary isospin T and hypercharge $`Y_\varphi `$. What happens when at higher temperature, as for example found in the early universe, the $`SU(N_C)SU(2)_LU(1)_Y`$ symmetry is restored ? Then the parameter ‘b’ in the electric charge definition remains undetermined. Note that ‘b’ was fixed above due to spontaneous symmetry breaking through Higgs. Without it ‘b’ remains unknown. Hence when the electroweak symmetry is restored, irrespective of the Higgs isospin and hypercharge the electric charge disappears as a physical quantity. Therefore we find that there was no electric charge in the early universe.
Here attention is drawn to the fact that all putative extensions of the Standard Model should reduce smoothly and consistently to the Standard Model at low energies. Not only that, all these extensions should be consistent with the predictions of the Standard Model at very high temperatures. Contrary to naive expectations, the SM does make specific predictions at very high temperatures too. For example one clear-cut prediction of the Standard Model as shown here and also shown earlier, is that at high enough temperatures (as in the early universe) when the unbroken $`SU(3)SU(2)U(1)`$ symmetry was restored, there was no electric charge. GUTs and other standard extensions of the SM are incompatible with this requirement .
What about Superstring Theories ? Quite clearly, generically in Superstring Theories electric charge exists right up to the Planck Scale . Hence as per this theory the electric charge, as an inherent property of matter, has existed right from the beginning . This is not correct in the SM. As shown here and earlier, the electric charge came into existence at a later stage in the evolution of the Universe when the $`SU(2)_LU(1)_Y`$ group was spontaneously broken to $`U(1)_{em}`$. It was never there all the time. This is because electric charge is a derived quantity. Hence we find that in this regard the Superstring Theories are inconsistent with the SM and hence flawed .
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# Effective String Theory of Vortices and Regge Trajectories
## I Introduction
The goal of this paper is to derive an effective string theory of vortices beginning with a field theory containing classical vortex solutions. The Abelian Higgs model is an example of such a theory. Nielsen and Olesen showed that this model has classical magnetic vortex solutions. These vortices are tubes of magnetic flux with constant energy per unit length.
The motivation for this work came from the dual superconductor picture of confinement . In this picture, a dual Meissner effect confines electric color flux ($`Z_3`$ flux) to narrow tubes connecting quark–antiquark pairs. Calculations with explicit models of this type have been compared both with experimental data and with Monte Carlo simulations of QCD . To a good approximation, aside from a color factor, the dual Abelian Higgs model, coupling dual potentials to a scalar Higgs field carrying magnetic charge, can be used to describe the results of these calculations. However, these calculations neglect the effect of fluctuations in the shape of the flux tube on the $`q\overline{q}`$ interactions. We show in this paper that taking account of those fluctuations leads to an effective string theory of long distance QCD.
Well before the introduction of the idea of dual superconductivity, string models had been used to understand the origin of Regge trajectories, and they have continued to be used to describe other features of hadron physics, such as the spectrum of hybrid mesons. In the dual superconductor picture, a string arises because the dual potentials couple to a quark–antiquark pair via a Dirac string whose ends are a source and sink of electric color flux. The effect of the string is to create a flux tube (or Abrikosov–Nielsen–Olesen vortex ) connecting the quark–antiquark pair. As the pair moves, this flux tube sweeps out a space time surface on which the dual Higgs field must vanish. This condition determines the location of the QCD string in the dual superconductor picture.
The effort to obtain an effective string theory for Abrikosov–Nielsen–Olesen vortices has a long history, independent of any connection to QCD. Nambu attached quarks to the ends of superconducting vortices, and found an expression for the classical action of the resulting ANO vortex in the singular London limit of infinite Higgs mass. He introduced a cutoff to render this action finite, and showed that it was proportional to the area of the worldsheet (the Nambu–Goto action).
Förster took into consideration the curvature of the worldsheet. He showed that in the strong coupling limit, with the ratio of vector and scalar masses held fixed, the effects of curvature were unimportant, and the classical action for the vortex reduced to the Nambu–Goto action. This limit can be regarded as the long distance limit, since only zero mass excitations are left in the theory. Equivalently, since the flux tube radius vanishes in this limit, all physical distances, measured in units of the flux tube radius, are becoming large. All degrees of freedom except the transverse oscillations of the vortex are frozen out.
Gervais and Sakita first considered the quantum theory of the vortices of the Abelian Higgs model in the same long distance limit. They used the results of Förster to define collective coordinates for the vortices, by means of which they constructed an effective vortex action. They also obtained a formal expression for the Feynman path integral of the Abelian Higgs model as an integration over vortex sheets. However, they were not able to write this expression as an integral over the physical degrees of freedom of the vortices.
Lüscher, Symanzik, and Weisz considered the leading semiclassical corrections to the classical Nambu–Goto action due to transverse string fluctuations, and showed how to regulate the resulting divergences. They showed that for a string of length $`R`$ with fixed ends, the leading semiclassical contribution to the heavy quark potential is $`\pi /12R`$. In a second paper, Lüscher showed that this result was unaffected by the addition of other terms to the effective string action.
Polchinski and Strominger discussed the relation of the Abelian Higgs model to fundamental string theory, regarding the theory of ANO vortices as an effective string theory. They explained how existing string quantization methods were inappropriate for quantizing the vortices. To compensate for the anomalies in these quantization methods, they introduced an additional term, the “Polchinski Strominger term,” into the effective vortex action.
Akhmedov, Chernodub, Polikarpov, and Zubkov studied the quantum theory of ANO vortices in the London limit, In particular, they studied the transformation from field degrees of freedom to vortex degrees of freedom. They showed that the Jacobian of this transformation contained the “Polchinski Strominger term” as a factor. Although they, like Gervais and Sakita, did not obtain a complete expression for the path integral, this paper provided an important stimulus to our own work.
In the current paper, we simplify and extend work done in an earlier paper . We begin with the path integral representation of a field theory having vortex solutions. It is an effective field theory describing phenomena at distances greater than the flux tube radius. We end up with an effective string theory of vortices in a form suitable for explicit calculations.
We apply this theory to calculate the energy $`E`$ and angular momentum $`J`$ of the fluctuations of a string bounded by the curve generated by the worldlines of a quark–antiquark pair separated by a fixed distance and rotating with fixed angular velocity. This gives the contribution of string fluctuations to the Regge trajectory $`J(E^2)`$, which we compare with the experimental $`\rho `$ and $`\omega `$ trajectories.
## II Outline
In section III, we rewrite the path integral over field configurations of the Abelian Higgs model containing vortices as an integral over surfaces on which the Higgs field vanishes. This introduces a Jacobian due to the change from field variables to string variables (surfaces). This Jacobian is the key to determining the action of the effective string theory, and to defining the integral over all surfaces. We next use the formalism described in section III to obtain an effective theory of ANO vortices. In section IV, we show how the Jacobian divides into a field part and a string part. The two parts of the Jacobian play different roles in the effective theory. In section V, we define an expression for the action of the effective string theory. All the dependence on the Abelian Higgs model is contained in the string action. We also obtain an expression for the path integral over vortices. In section VI, we show how to express the integral over surfaces as an integral over the two physical degrees of freedom of the vortex, and obtain the final form of the effective string theory.
In the remaining sections we compute the leading semiclassical contribution to Regge trajectories due to the fluctuations of the string. We obtain an expression for the contribution of string fluctuations to the effective action in section VII, and in sections VIII and IX describe how to regularize this expression, making use of the results of Lüscher, Symanzik, and Weisz . In section X we calculate the contribution of string fluctuations to the effective action for a straight, rotating string, and in section XI obtain the resulting corrections to Regge trajectories.
## III The Transformation from Fields to Strings
In this section we consider the Abelian Higgs model coupled via a Dirac string to a moving quark–antiquark pair. We transform the path integral over field configurations containing vortices to an integral over the surfaces $`\stackrel{~}{x}^\mu `$ determining the location of the vortices.
We denote the (dual) potentials by $`C_\mu `$ and the complex (monopole) Higgs field by $`\varphi `$. The dual coupling constant is $`g=2\pi /e`$, where $`e`$ is the Yang–Mills coupling constant ($`\alpha _s=e^2/4\pi `$). The worldlines of the quark and antiquark trajectories form the closed loop $`\mathrm{\Gamma }`$ (see Fig. 1).
The moving quark–antiquark pair couples to the dual potentials $`C_\mu `$ via a Dirac string tensor $`G_{\mu \nu }^S`$, which is nonvanishing along some line $`L`$ connecting the $`q\overline{q}`$ pair. As the pair moves, the line $`L`$ sweeps out a worldsheet $`\stackrel{~}{x}^\mu (\xi )`$ parameterized by coordinates $`\xi ^a`$, $`a=1,2`$. The field $`\varphi `$ vanishes on this worldsheet,
$$\varphi (x^\mu )=0,\text{ at }x^\mu =\stackrel{~}{x}^\mu (\xi ).$$
(1)
The corresponding Dirac string tensor $`G_{\mu \nu }^S`$ is given by
$$G_{\mu \nu }^S=ed^2\xi \frac{1}{2}ϵ^{ab}ϵ_{\mu \nu \alpha \beta }\frac{\stackrel{~}{x}^\alpha }{\xi ^a}\frac{\stackrel{~}{x}^\beta }{\xi ^b}\delta ^{(4)}\left(x^\mu \stackrel{~}{x}^\mu (\xi )\right).$$
(2)
The action $`S`$ of a field configuration which has a vortex on the sheet $`\stackrel{~}{x}^\mu (\xi )`$ is
$$S=\frac{4}{3}d^4x\left[\frac{1}{4}\left(G_{\mu \nu }\right)^2\frac{1}{2}\left|(_\mu igC_\mu )\varphi \right|^2\frac{\lambda }{4}\left(|\varphi |^2\varphi _0^2\right)^2\right],$$
(3)
where the field strength $`G_{\mu \nu }`$ is given by
$$G_{\mu \nu }=_\mu C_\nu _\nu C_\mu +G_{\mu \nu }^S.$$
(4)
The Higgs mechanism gives the vector particle (dual gluon) a mass $`M_V=g\varphi _0`$ and the scalar particle a mass $`M_S=\sqrt{2\lambda }\varphi _0`$, where $`\varphi _0`$ is the vacuum expectation value of the Higgs field. We have introduced the color factor $`\frac{4}{3}`$ in (4) because we are interested in using $`S`$ as a model for long distance QCD. We consider $`S`$ to be an effective action describing distances greater than the flux tube radius $`a`$.
The long distance $`q\overline{q}`$ interaction is determined by the Wilson loop $`W[\mathrm{\Gamma }]`$,
$$W[\mathrm{\Gamma }]=𝒟\varphi ^{}𝒟\varphi 𝒟C^\mu e^{i(S[\varphi ,C]+S_{GF})},$$
(5)
where $`S_{GF}`$ is a gauge fixing term. The functional integrals are cut off at the momentum scale $`1/a`$. The action (3) describes a field theory having classical vortex solutions. The functional integral (5) goes over all field configurations containing a vortex bounded by $`\mathrm{\Gamma }`$.
Previous calculations of $`W[\mathrm{\Gamma }]`$ were carried out in the classical approximation (corresponding to a flat vortex sheet $`\stackrel{~}{x}^\mu `$), and showed that the Landau–Ginzburg parameter $`\lambda /g^2`$ is approximately equal to $`\frac{1}{2}`$. This corresponds to a superconductor on the border between type I and type II. In this situation, both particles have the same mass $`M=M_V=M_S`$, the string tension is $`\sigma =\frac{4}{3}\pi \varphi _0^2`$, and the flux tube radius is $`a=\sqrt{2}/M`$.
To take into account the fluctuations of these vortices, we must evaluate $`W[\mathrm{\Gamma }]`$ beyond the classical approximation. We carry out the functional integration (5) in two steps: (1) We fix the location of a vortex sheet $`\stackrel{~}{x}^\mu `$, and integrate only over field configurations for which $`\varphi (x^\mu )`$ vanishes on $`\stackrel{~}{x}^\mu `$. (2) We integrate over all possible vortex sheets. To implement this procedure, we introduce into the functional integral (5) the factor one, written in the form
$$1=J[\varphi ]𝒟\stackrel{~}{x}^\mu \delta \left[\mathrm{Re}\varphi (\stackrel{~}{x}^\mu (\xi ))\right]\delta \left[\mathrm{Im}\varphi (\stackrel{~}{x}^\mu (\xi ))\right].$$
(6)
The integration $`𝒟\stackrel{~}{x}^\mu `$ is over the four functions $`\stackrel{~}{x}^\mu (\xi )`$. The functions $`\stackrel{~}{x}^\mu (\xi )`$ are a particular parameterization of the worldsheet $`\stackrel{~}{x}^\mu `$.
The expression (6) implies that the string worldsheet $`\stackrel{~}{x}^\mu `$, determined by the $`\delta `$ functions, is the surface of the zeros of the field $`\varphi `$. The factor $`J[\varphi ]`$ is a Jacobian, and is defined by Eq. (6). Inserting (6) into (5) puts the Wilson loop in the form
$$W[\mathrm{\Gamma }]=𝒟\varphi ^{}𝒟\varphi 𝒟C^\mu e^{i(S[\varphi ,C]+S_{GF})}J[\varphi ]𝒟\stackrel{~}{x}^\mu \delta \left[\mathrm{Re}\varphi (\stackrel{~}{x}^\mu (\xi ))\right]\delta \left[\mathrm{Im}\varphi (\stackrel{~}{x}^\mu (\xi ))\right].$$
(7)
We then reverse the order of the field integration and the string integration over surfaces $`\stackrel{~}{x}^\mu (\xi )`$,
$$W[\mathrm{\Gamma }]=𝒟\stackrel{~}{x}^\mu 𝒟\varphi ^{}𝒟\varphi 𝒟C^\mu J[\varphi ]\delta \left[\mathrm{Re}\varphi (\stackrel{~}{x}^\mu (\xi ))\right]\delta \left[\mathrm{Im}\varphi (\stackrel{~}{x}^\mu (\xi ))\right]e^{i(S[\varphi ,C]+S_{GF})}.$$
(8)
In Eq. (7), the $`\delta `$ functions fix $`\stackrel{~}{x}^\mu `$ to lie on the surface of the zeros of a given field $`\varphi `$, while in Eq. (8), they restrict the field $`\varphi `$ to vanish on a given surface $`\stackrel{~}{x}^\mu `$. The integral over $`\varphi `$ in Eq. (8) is therefore restricted to functions $`\varphi `$ which vanish on $`\stackrel{~}{x}^\mu `$, in contrast to the integral over $`\varphi `$ in Eq. (7), in which $`\varphi `$ can be any function.
## IV Factorization of the Jacobian
To evaluate $`W[\mathrm{\Gamma }]`$ we divide $`J[\varphi ]`$ into two parts. The Jacobian $`J[\varphi ]`$ in Eq. (8) is evaluated for field configurations $`\varphi `$ which vanish on a particular surface $`\stackrel{~}{x}^\mu `$. We make this explicit by writing (6) as
$$J[\varphi ,\stackrel{~}{x}^\mu ]^1=𝒟\stackrel{~}{y}^\mu \delta \left[\mathrm{Re}\varphi (\stackrel{~}{y}^\mu (\tau ))\right]\delta \left[\mathrm{Im}\varphi (\stackrel{~}{y}^\mu (\tau ))\right],$$
(9)
where $`\stackrel{~}{y}^\mu `$ is some other string worldsheet, distinct from $`\stackrel{~}{x}^\mu `$. The evaluation of the Jacobian is the essential new ingredient in deriving $`W[\mathrm{\Gamma }]`$.
The $`\delta `$ functions in (9) select surfaces $`\stackrel{~}{y}^\mu (\tau )`$ which lie in a neighborhood of the surface $`\stackrel{~}{x}^\mu (\xi )`$ of the zeros of $`\varphi `$. We separate $`\stackrel{~}{y}^\mu (\tau )`$ into components lying on the surface $`\stackrel{~}{x}^\mu (\xi )`$ and components lying along vectors $`n_\mu ^A(\xi )`$ normal to $`\stackrel{~}{x}^\mu (\xi )`$ at the point $`\xi `$:
$$\stackrel{~}{y}^\mu (\tau )=\stackrel{~}{x}^\mu (\xi (\tau ))+y_{}^A(\xi (\tau ))n_\mu ^A(\xi (\tau )).$$
(10)
The point $`\stackrel{~}{x}^\mu (\xi (\tau ))`$ is the point on the surface $`\stackrel{~}{x}^\mu (\xi )`$ lying closest to $`\stackrel{~}{y}^\mu (\tau )`$, and the magnitude of $`y_{}^A(\xi (\tau ))`$ is the distance from $`\stackrel{~}{y}^\mu (\tau )`$ to $`\stackrel{~}{x}^\mu (\xi (\tau ))`$ (see Fig. 2).
We evaluate the Jacobian (9) by making the change of variables
$$\stackrel{~}{y}^\mu (\tau )(\xi (\tau ),y_{}^A(\xi ))$$
(11)
defined by (10). Although the $`\delta `$ functions in (9) force $`y_{}^A`$ to vanish, the integrations over $`y_{}^A`$ give a contribution to the Jacobian. Furthermore, this contribution depends on the field variable $`\varphi `$ in a neighborhood of the surface. The integration over the reparameterizations $`\xi (\tau )`$ of the surface $`\stackrel{~}{x}^\mu (\xi )`$, on the other hand, depends upon the surface, but not on the fields. The change of variables (11) leads to a factorization of the Jacobian into a field contribution, and into a contribution depending only on the intrinsic properties of the worldsheet $`\stackrel{~}{x}^\mu (\xi )`$.
We now exhibit the factorization of the Jacobian. Under the transformation (11), the integral over $`\stackrel{~}{y}^\mu `$ becomes
$`𝒟\stackrel{~}{y}^\mu `$ $`=`$ $`\mathrm{Det}_\tau \left[ϵ^{\mu \nu \alpha \beta }{\displaystyle \frac{1}{2}}ϵ^{ab}{\displaystyle \frac{\stackrel{~}{x}^\mu }{\xi ^a}}{\displaystyle \frac{\stackrel{~}{x}^\nu }{\xi ^b}}{\displaystyle \frac{1}{2}}ϵ^{AB}n_\alpha ^An_\beta ^B\right]𝒟y_{}^A𝒟\xi `$ (12)
$`=`$ $`\mathrm{Det}_\tau \left[\sqrt{{\displaystyle \frac{1}{2}}\left(ϵ^{ab}{\displaystyle \frac{\stackrel{~}{x}^\mu }{\xi ^a}}{\displaystyle \frac{\stackrel{~}{x}^\nu }{\xi ^b}}\right)^2{\displaystyle \frac{1}{2}}\left(ϵ^{AB}n_\alpha ^An_\beta ^B\right)^2}\right]𝒟y_{}^A𝒟\xi `$ (13)
$`=`$ $`\mathrm{Det}_\tau \left[\sqrt{g(\xi )}|_{\xi =\xi (\tau )}\right]𝒟y_{}^A𝒟\xi ,`$ (14)
where $`\sqrt{g}`$ is the square root of the determinant of the induced metric
$$g_{ab}=\frac{\stackrel{~}{x}^\mu }{\xi ^a}\frac{\stackrel{~}{x}^\mu }{\xi ^b}$$
(15)
evaluated on the worldsheet $`\stackrel{~}{x}^\mu `$. Appendix A gives a summary of our notation, and of the relations used to obtain (14).
The functional determinant in (14) is the product of its argument evaluated at all points $`\tau `$ on the sheet, in the same way that the integration over $`𝒟\stackrel{~}{y}^\mu `$ is a product of integrals at all points $`\tau `$. Making the change of coordinates (10), (11) in the Jacobian (9) gives
$`J[\varphi ,\stackrel{~}{x}^\mu ]^1`$ $`=`$ $`{\displaystyle 𝒟\xi 𝒟y_{}^A\mathrm{Det}_\tau [\sqrt{g}]\delta \left[\mathrm{Re}\varphi \left(\stackrel{~}{x}^\mu (\xi (\tau ))+y_{}^A(\xi (\tau ))n_A^\mu (\xi (\tau ))\right)\right]}`$ (17)
$`\times \delta \left[\mathrm{Im}\varphi \left(\stackrel{~}{x}^\mu (\xi (\tau ))+y_{}^A(\xi (\tau ))n_A^\mu (\xi (\tau ))\right)\right].`$
Eq. (17) has the form:
$$J[\varphi ,\stackrel{~}{x}]^1=𝒟\xi (\tau )\mathrm{Det}_\tau \left[\sqrt{g}\right]J_{}[\varphi ,\stackrel{~}{x}^\mu (\xi (\tau ))]^1,$$
(18)
where
$`J_{}[\varphi ,\stackrel{~}{x}^\mu (\xi (\tau ))]^1`$ $`=`$ $`{\displaystyle 𝒟y_{}^A\delta \left[\mathrm{Re}\varphi \left(\stackrel{~}{x}^\mu (\xi (\tau ))+y_{}^A(\xi (\tau ))n_A^\mu (\xi (\tau ))\right)\right]}`$ (20)
$`\times \delta \left[\mathrm{Im}\varphi \left(\stackrel{~}{x}^\mu (\xi (\tau ))+y_{}^A(\xi (\tau ))n_A^\mu (\xi (\tau ))\right)\right]`$
contains all the dependence on $`\varphi `$. Since $`J_{}`$ is independent of the parameterization $`\xi (\tau )`$, the Jacobian factors into two parts:
$$J[\varphi ,\stackrel{~}{x}]^1=J_{}[\stackrel{~}{x}]^1J_{}[\varphi ,\stackrel{~}{x}]^1,$$
(21)
where
$$J_{}[\stackrel{~}{x}]^1=𝒟\xi \mathrm{Det}_\tau \left[\sqrt{g}\right].$$
(22)
The string part $`J_{}`$ of the Jacobian arises from the parameterization degrees of freedom. In the next section, we show that $`J_{}`$ is the Faddeev–Popov determinant for the $`\delta `$ functions in (8). This allows us to define the action of the effective string theory. In the following section, we will use $`J_{}`$ to fix the reparameterization degrees of freedom.
## V The String Action
Inserting the factorized form (21) of $`J[\varphi ]`$ into the expression (8) for $`W[\mathrm{\Gamma }]`$ gives the Wilson Loop the form
$$W[\mathrm{\Gamma }]=𝒟\stackrel{~}{x}^\mu J_{}[\stackrel{~}{x}]e^{iS_{\mathrm{eff}}},$$
(23)
where the action $`S_{\mathrm{eff}}`$ of the effective string theory is given by
$$e^{iS_{\mathrm{eff}}[\stackrel{~}{x}^\mu (\xi )]}=𝒟\varphi ^{}𝒟\varphi 𝒟C^\mu J_{}[\varphi ]\delta \left[\mathrm{Re}\varphi (\stackrel{~}{x}^\mu (\xi ))\right]\delta \left[\mathrm{Im}\varphi (\stackrel{~}{x}^\mu (\xi ))\right]e^{i(S+S_{GF})}.$$
(24)
The string action (24) was obtained previously by Gervais and Sakita . The novel feature of our result is the string integration measure of the Wilson loop (23).
The string action depends upon the field part $`J_{}`$ of the Jacobian,
$$J_{}[\varphi ,\stackrel{~}{x}^\mu ]^1=𝒟y_{}^A\delta \left[\mathrm{Re}\varphi \left(\stackrel{~}{x}^\mu +y_{}^An_{\mu A}\right)\right]\delta \left[\mathrm{Im}\varphi \left(\stackrel{~}{x}^\mu +y_{}^An_{\mu A}\right)\right].$$
(25)
The $`\delta `$ functions force $`y_{}^A`$ to be zero, so we can expand their arguments in a power series in $`y_{}^A`$,
$$\varphi \left(y^\mu \right)=\varphi \left(\stackrel{~}{x}^\mu \right)+y_{}^An_A^\nu _\nu \varphi \left(\stackrel{~}{x}^\mu \right)+O\left(y_{}^2\right).$$
(26)
The zeroth order term in (26) vanishes because $`\stackrel{~}{x}^\mu `$ is the surface of the zeros of $`\varphi `$. The integration (25) over $`y_{}^A`$ gives the result
$$J_{}[\varphi ,\stackrel{~}{x}^\mu ]^1=\mathrm{Det}_\xi ^1\left[ϵ^{AB}n_A^\mu n_B^\nu \left(_\mu \mathrm{Re}\varphi \right)\left(_\nu \mathrm{Im}\varphi \right)|_{x^\mu =\stackrel{~}{x}^\mu }\right].$$
(27)
The Jacobian $`J_{}`$ is a Faddeev–Popov determinant, which we discuss in Appendix B.
Eq. (24) gives the action $`S_{\mathrm{eff}}(\stackrel{~}{x}^\mu )`$ of the effective string theory as an integral over field configurations which have a vortex fixed at $`\stackrel{~}{x}^\mu `$. Since the vortex theory (5) is an effective long distance theory, the path integral (5) for $`W[\mathrm{\Gamma }]`$, written in terms of the fields of the Abelian Higgs model, is cut off at a scale $`\mathrm{\Lambda }`$ which is on the order of the mass $`M`$ of the dual gluon. Furthermore, the integration (23) over $`\stackrel{~}{x}^\mu `$ includes all the long distance fluctuations of the theory. Therefore, the path integral (24) contains neither short distance nor long distance fluctuations, and is determined by minimizing the field action $`S[\stackrel{~}{x}^\mu ,\varphi ,C_\mu ]`$ for a fixed position of the vortex sheet:
$$S_{\mathrm{eff}}[\stackrel{~}{x}^\mu ]=S[\stackrel{~}{x}^\mu ,\varphi ^{\mathrm{class}},C_\mu ^{\mathrm{class}}],\varphi ^{\mathrm{class}}(\stackrel{~}{x}^\mu )=0.$$
(28)
The fields $`\varphi ^{\mathrm{class}}`$ and $`C_\mu ^{\mathrm{class}}`$ are the solutions of the classical equations of motion, subject to the boundary condition $`\varphi (\stackrel{~}{x}^\mu )=0`$.
The action $`S_{\mathrm{eff}}`$ depends both on the distance $`R`$ between the quarks, and the radius of curvature $`R_V`$ of the vortex sheet bounded by $`\mathrm{\Gamma }`$. In the long distance limit, when the length of the string $`R`$ and its radius $`R_V`$ are large compared to the thickness of the flux tube $`a`$, the string action (28) becomes the Nambu–Goto action $`S_{NG}`$,
$$S_{NG}=\sigma d^2\xi \sqrt{g},$$
(29)
where $`\sigma `$ is the classical string tension, determined from the solution of the Nielsen–Olesen equations for a straight, infinitely long string.
It is convenient to separate the action (28) into its perturbative and nonperturbative parts:
$$S_{\mathrm{eff}}[\stackrel{~}{x}^\mu ]=S[\stackrel{~}{x}^\mu ,\varphi ^{\mathrm{class}},C_\mu ^{\mathrm{class}}]=S^{\mathrm{Maxwell}}[\stackrel{~}{x}^\mu ]+S^{\mathrm{NP}}[\stackrel{~}{x}^\mu ],$$
(30)
where $`S^{\mathrm{Maxwell}}`$ is the action obtained by setting $`\lambda =g=0`$ in Eq. (3). The value of $`S^{\mathrm{Maxwell}}`$ depends only upon the boundary $`\mathrm{\Gamma }`$, and is the usual electromagnetic interaction between charged particles:
$$S^{\mathrm{Maxwell}}[\mathrm{\Gamma }]=\frac{4}{3}\frac{e^2}{2}𝑑x^\mu 𝑑x_{}^{}{}_{}{}^{\mu }𝒟_{\mu \nu }(x^\mu x_{}^{}{}_{}{}^{\mu }),$$
(31)
where $`𝒟_{\mu \nu }`$ is the photon propagator.
To calculate the Wilson loop $`W[\mathrm{\Gamma }]`$ from the effective string theory (23), we must also examine $`S_{\mathrm{eff}}`$ at smaller values of $`R`$ and $`R_V`$, on the order of the string thickness $`a`$. We first consider the dependence of $`S_{\mathrm{eff}}`$ on $`R`$ for a flat string, where $`R_V\mathrm{}`$. In this case, the curve $`\mathrm{\Gamma }`$ is a rectangle of length $`T`$ in the time direction, and width $`R`$ in the space direction. In the large $`T`$ limit, the action $`S_{\mathrm{eff}}`$ reduces to the product of $`T`$ and the potential $`V^{\mathrm{class}}(R)`$ previously used to fit the energy levels of heavy quark systems.
Evaluation of (30) for a flat sheet gives a corresponding decomposition of $`V^{\mathrm{class}}(R)`$,
$$V^{\mathrm{class}}(R)=V^{\mathrm{Coulomb}}(R)+V^{\mathrm{NP}}(R).$$
(32)
For small $`R`$,
$$V^{\mathrm{class}}(R)\mathrm{@}>>R0>V^{\mathrm{Coulomb}}(R)=\frac{4}{3}\frac{\alpha _s}{R},$$
(33)
while Eq. (29) gives the large $`R`$ behavior
$$V^{\mathrm{NP}}(R)\mathrm{@}>>R\mathrm{}>\sigma R.$$
(34)
Recent numerical studies of the classical equations of motion for a flat sheet have shown that for a superconductor on the I–II border, the long distance behavior (34) of $`V^{\mathrm{NP}}(R)`$ persists to small values of $`R`$, even to values less than the string thickness $`a`$. Therefore, for a superconductor on the I–II border, $`V^{\mathrm{class}}(R)`$ is, to a good approximation, equal to the Cornell potential :
$$V^{\mathrm{class}}(R)\frac{4}{3}\frac{\alpha _s}{R}+\sigma R.$$
(35)
In other words, for a flat sheet,
$$S^{\mathrm{NP}}(R)S_{NG}.$$
(36)
Thus, for short straight strings the Nambu–Goto action remains a good approximation to the nonperturbative part of the classical action for a superconductor on the type I–II border.
Next, consider the nonperturbative contribution to the classical action for a long bent string. (The Maxwell action has the value (31) independent of the shape of the vortex.) The leading correction to the Nambu–Goto action when the string is bent is the curvature term:
$$S_{\mathrm{curvature}}=\beta d^2\xi \sqrt{g}\left(𝒦_{ab}^A\right)^2,$$
(37)
where $`𝒦_{ab}^A`$ is the extrinsic curvature.
$$𝒦_{ab}^A=n_\mu ^A_a_bx^\mu .$$
(38)
The magnitude of $`𝒦_{ab}^A`$ is of the order of $`1/R_V`$, so that $`S_{\mathrm{curvature}}(a^2/R_V^2)S_{NG}`$.
The calculation of the “rigidity” $`\beta `$ determining the size of $`S_{\mathrm{curvature}}`$ has been considered by a number of authors , but the value of $`\beta `$ for a superconductor on the I–II border was never calculated. We conjecture that the value of $`\beta `$ is small, because de Vega and Schaposnik have shown that the components of the stress tensor perpendicular to the axis of a straight Nielsen–Olesen flux tube vanish for a superconductor on the border between type I and type II. In other words, there are no “bonds” perpendicular to the field lines of a straight flux tube of infinite extent. When the flux tube is bent slightly, there are no perpendicular bonds to be stretched or compressed, and the change in the energy is just the string tension multiplied by the change in length. That is, the curvature term, which in a sense represents the attraction or repulsion between neighboring parts of the string, should vanish. A more formal argument can be made by regarding the borderline superconductor as the long distance limit of a theory where the forces between vortices become weak. Polyakov has shown, using renormalization group methods, that $`\beta `$ also vanishes in this limit.
Similar heuristic arguments give a reason for the above mentioned result that the Nambu–Goto action is a good approximation for short, straight strings on the I–II border. The bending of the field lines as the quark–antiquark separation becomes smaller causes no additional changes in the energy.
We therefore take the action of the effective string theory to be equal to the sum of the Maxwell action (31) and the Nambu–Goto action (29):
$$S_{\mathrm{eff}}[\stackrel{~}{x}^\mu ]=S^{\mathrm{Maxwell}}[\mathrm{\Gamma }]\sigma d^2\xi \sqrt{g}.$$
(39)
We use Eq. (39) for the full range of string lengths $`R`$ and radii of curvature $`R_V`$ greater than the inverse of the mass $`M`$ of the dual gluon, which is the cutoff of the effective string theory (23).
Eq. (39) for $`S_{\mathrm{eff}}[\stackrel{~}{x}^\mu ]`$ is the generalization of (35) to a general sheet. The first term, $`S^{\mathrm{Maxwell}}[\mathrm{\Gamma }]`$, is just a boundary term, independent of the fluctuating string, and we take take $`S_{\mathrm{eff}}=S_{NG}`$ for the calculations carried out in the rest of this paper. In the next section we show how to carry out the integration over $`\stackrel{~}{x}^\mu `$ in (23) by separating the degrees of freedom of the worldsheet $`\stackrel{~}{x}^\mu `$ into two physical degrees of freedom and two reparameterization degrees of freedom. This treatment makes no use of (39), and is applicable to any effective string theory of vortices.
## VI Effective Theory of Transverse String Fluctuations
We next show how to evaluate the integral over $`\stackrel{~}{x}^\mu (\xi )`$ in Eq. (23),
$$W[\mathrm{\Gamma }]=𝒟\stackrel{~}{x}^\mu J_{}[\stackrel{~}{x}^\mu ]e^{iS_{\mathrm{eff}}}.$$
(40)
The integration over $`\stackrel{~}{x}^\mu (\xi )`$ is the product of an integral over string worldsheets and an integral over reparameterizations of the coordinates of the string. The Jacobian $`J_{}`$ is the inverse of the integration (22) over reparameterization degrees of freedom. In this section, we fix the parameterization of the string, and show that $`J_{}`$ cancels the integration over reparameterizations.
Any surface $`\stackrel{~}{x}^\mu `$ has only two physical degrees of freedom. The other two degrees of freedom represent the invariance of the surface under coordinate reparameterizations. We fix the coordinate reparameterization symmetry by choosing a particular “representation” $`x_p^\mu `$ of the surface, which depends on two functions $`f^1(\xi )`$, $`f^2(\xi )`$,
$$x_p^\mu (\xi )=x_p^\mu [f^1(\xi ),f^2(\xi ),\xi ].$$
(41)
A particular example of a representation $`x_p^\mu `$ is obtained by expanding in transverse fluctuations $`x_{}^A`$ about a fixed sheet $`\overline{x}_p^\mu `$,
$$x_p^\mu (\xi )=x_p^\mu [x_{}^A(\xi ),\xi ]=\overline{x}_p^\mu (\xi )+x_{}^A(\xi )\overline{n}_A^\mu (\xi ).$$
(42)
The vectors $`\overline{n}_A^\mu `$ are orthogonal to the surface $`\overline{x}_p^\mu `$. In this example, $`f^1`$ and $`f^2`$ are the transverse coordinates $`x_{}^A`$.
Any physical surface can be expressed in terms of $`x_p^\mu `$ by a suitable choice of $`f^1`$ and $`f^2`$. In particular, the worldsheet $`\stackrel{~}{x}^\mu (\xi )`$ appearing in the integral (40) can be written in terms of a reparameterization $`\stackrel{~}{\xi }(\xi )`$ of the representation $`x_p^\mu `$,
$$\stackrel{~}{x}^\mu (\xi )=x_p^\mu [f^1(\stackrel{~}{\xi }(\xi )),f^2(\stackrel{~}{\xi }(\xi )),\stackrel{~}{\xi }(\xi )].$$
(43)
The four degrees of freedom in $`\stackrel{~}{x}^\mu (\xi )`$ are replaced by two physical degrees of freedom $`f^1(\xi )`$, $`f^2(\xi )`$ and two reparameterization degrees of freedom $`\stackrel{~}{\xi }(\xi )`$.
We can write the integral over $`\stackrel{~}{x}^\mu (\xi )`$ in (40) in terms of integrals over $`f^1(\xi ),f^2(\xi )`$ and $`\stackrel{~}{\xi }(\xi )`$,
$$𝒟\stackrel{~}{x}^\mu =\mathrm{Det}\left[ϵ_{\mu \nu \alpha \beta }\frac{1}{2}ϵ^{ab}\frac{\stackrel{~}{x}^\mu }{\stackrel{~}{\xi }^a}\frac{\stackrel{~}{x}^\nu }{\stackrel{~}{\xi }^b}\frac{\stackrel{~}{x}^\alpha }{f^1}\frac{\stackrel{~}{x}^\beta }{f^2}\right]𝒟f^1𝒟f^2𝒟\stackrel{~}{\xi }.$$
(44)
Noting that the derivative of $`\stackrel{~}{x}^\mu `$ with respect to $`\stackrel{~}{\xi }(\xi )`$ is,
$$\frac{\stackrel{~}{x}^\mu }{\stackrel{~}{\xi }^a}=\left[\frac{x_p^\mu }{f^1}\frac{f^1}{\xi ^a}+\frac{x_p^\mu }{f^2}\frac{f^2}{\xi ^a}+\frac{x_p^\mu }{\xi ^a}\right]|_{\xi =\stackrel{~}{\xi }},$$
(45)
we can write
$`𝒟\stackrel{~}{x}^\mu `$ $`=`$ $`\mathrm{Det}\left[ϵ_{\mu \nu \alpha \beta }{\displaystyle \frac{1}{2}}ϵ^{ab}{\displaystyle \frac{x_p^\alpha }{\xi ^a}}{\displaystyle \frac{x_p^\beta }{\xi ^b}}{\displaystyle \frac{x_p^\mu }{f^1}}{\displaystyle \frac{x_p^\nu }{f^2}}|_{\xi =\stackrel{~}{\xi }}\right]𝒟f^1𝒟f^2𝒟\stackrel{~}{\xi }`$ (46)
$`=`$ $`\mathrm{Det}\left[\stackrel{~}{t}_{\mu \nu }\sqrt{g^p}{\displaystyle \frac{x_p^\mu }{f^1}}{\displaystyle \frac{x_p^\nu }{f^2}}|_{\xi =\stackrel{~}{\xi }}\right]𝒟f^1𝒟f^2𝒟\stackrel{~}{\xi },`$ (47)
where
$$\stackrel{~}{t}_{\mu \nu }=\frac{1}{2}ϵ_{\mu \nu \alpha \beta }\frac{ϵ^{ab}}{\sqrt{g}}\frac{\stackrel{~}{x}^\alpha }{\xi ^a}\frac{\stackrel{~}{x}^\beta }{\xi ^b},$$
(48)
is the antisymmetric tensor normal to the worldsheet and,
$$g_{ab}^p=\frac{x_p^\mu }{\xi ^a}\frac{x_p^\mu }{\xi ^b},$$
(49)
is the induced metric of $`x_p^\mu (\xi )`$. The metric $`g_{ab}^p`$ is related to the metric $`g_{ab}`$ of $`\stackrel{~}{x}^\mu (\xi )`$,
$$g_{ab}=\frac{\stackrel{~}{\xi }^c}{\xi ^a}\frac{\stackrel{~}{\xi }^d}{\xi ^b}g_{cd}^p.$$
(50)
The induced metric $`g_{ab}`$ of the original worldsheet $`\stackrel{~}{x}^\mu (\xi )`$ does not appear in Eq. (47) because the determinant is independent of $`\stackrel{~}{\xi }`$. Only the induced metric $`g_{ab}^p`$ of the worldsheet $`x_p^\mu (\xi )`$ enters into the determinant.
With the parameterization (43) of $`\stackrel{~}{x}^\mu `$, the path integral (40) takes the form
$$W[\mathrm{\Gamma }]=𝒟\stackrel{~}{\xi }𝒟f^1𝒟f^2\mathrm{Det}\left[\stackrel{~}{t}^{\mu \nu }\frac{x_p^\mu }{f^1}\frac{x_p^\nu }{f^2}\right]\mathrm{Det}[\sqrt{g^p}]J_{}e^{iS_{\mathrm{eff}}}.$$
(51)
The action $`S_{\mathrm{eff}}`$ is parameterization independent, so it is independent of $`\stackrel{~}{\xi }(\xi )`$. The same is true for $`J_{}`$. Furthermore, $`\stackrel{~}{t}_{\mu \nu }`$ is parameterization independent, so that the product
$$𝒟f^1𝒟f^2\mathrm{Det}\left[\stackrel{~}{t}^{\mu \nu }\frac{x_p^\mu }{f^1}\frac{x_p^\nu }{f^2}\right]$$
(52)
is independent of $`\stackrel{~}{\xi }(\xi )`$. Therefore, this product, along with $`J_{}`$ and $`e^{iS_{\mathrm{eff}}}`$, can be brought outside the $`\stackrel{~}{\xi }`$ integral in (51). The path integral then takes the form
$$W[\mathrm{\Gamma }]=𝒟f^1𝒟f^2\mathrm{Det}\left[\stackrel{~}{t}^{\mu \nu }\frac{x_p^\mu }{f^1}\frac{x_p^\nu }{f^2}\right]J_{}e^{iS_{\mathrm{eff}}}𝒟\stackrel{~}{\xi }\mathrm{Det}[\sqrt{g^p}].$$
(53)
The remaining integral over reparameterizations $`\stackrel{~}{\xi }`$ is equal to $`J_{}^1`$, defined by (22), and is canceled by the explicit factor of $`J_{}`$ appearing in (53). This means we do not need to evaluate $`J_{}`$, and can avoid the complications inherent in evaluating the integral over reparameterizations of the string coordinates. The anomalies produced in string theory by evaluating this integral are not present, so we do not have a Polchinski–Strominger term in the theory. Eq. (53) gives the final result for the Wilson loop
$$W[\mathrm{\Gamma }]=𝒟f^1𝒟f^2\mathrm{Det}\left[\stackrel{~}{t}_{\mu \nu }\frac{x_p^\mu }{f^1}\frac{x_p^\nu }{f^2}\right]e^{iS_{\mathrm{eff}}},$$
(54)
as an integration over two function $`f^1(\xi )`$ and $`f^2(\xi )`$, the physical degrees of freedom of the string.
The path integral (54) is invariant under reparameterizations of the string, and describes a two dimensional field theory with two degrees of freedom, the two transverse oscillations of a two dimensional sheet. The integration (54) goes over the normal fluctuations of the string worldsheet. The components of $`f^1`$ and $`f^2`$ along the sheet are nonphysical. The determinant in (54) is a normalization factor for $`f^1`$ and $`f^2`$. This can be seen by applying the identity $`\stackrel{~}{t}_{\mu \nu }=ϵ^{AB}n_{\mu A}n_{\nu B}`$ to the determinant,
$$\mathrm{Det}\left[\stackrel{~}{t}_{\mu \nu }\frac{x_p^\mu }{f^1}\frac{x_p^\nu }{f^2}\right]=\mathrm{Det}\left[ϵ^{AB}\left(n_{\mu A}\frac{x_p^\mu }{f^1}\right)\left(n_{\nu B}\frac{x_p^\nu }{f^2}\right)\right].$$
(55)
The factors of $`n_{\mu A}\frac{x_p^\mu }{f^i}`$ determine the amount of the fluctuation $`f^i`$ which is in a direction normal to the sheet.
Eq. (54) is the string representation of any field theory containing classical vortex solutions. The expression (24) for $`S_{\mathrm{eff}}`$ had been obtained previously by Gervais and Sakita ; we are unaware of any previous derivation of the string representation (54) for the path integral. We now show how it provides a method for explicit calculations.
## VII The Semiclassical Approximation
In this section we carry out the semiclassical expansion of $`W[\mathrm{\Gamma }]`$ about a classical solution of the effective string theory, and find the leading contribution of string fluctuations to the effective action $`i\mathrm{log}W[\mathrm{\Gamma }]`$. The ends of the string follow the path $`\mathrm{\Gamma }`$ fixed by the prescribed trajectory of the quarks, and the fluctuations of the string are cutoff at the momentum scale $`M`$ of the inverse string radius. As explained in section V, we take the action $`S_{\mathrm{eff}}`$ of the effective string theory to be the Nambu–Goto action
$$S_{\mathrm{eff}}=\sigma d^2\xi \sqrt{g},$$
(56)
and (54) becomes
$$W[\mathrm{\Gamma }]=𝒟f^1(\xi )𝒟f^2(\xi )\mathrm{Det}\left[\stackrel{~}{t}^{\mu \nu }\frac{x_p^\mu }{f^1}\frac{x_p^\nu }{f^2}\right]e^{i\sigma {\scriptscriptstyle d^2\xi \sqrt{g}}}.$$
(57)
We expand (57) in small fluctuations $`f^i`$ of $`x_p^\mu [f^i,\xi ]`$ around a fixed sheet $`\overline{x}_p^\mu (\xi )`$, subject to the condition that the boundary of $`\overline{x}_p^\mu `$ lies on the curve $`\mathrm{\Gamma }`$,
$$x_p^\mu (f^i,\xi )=\overline{x}_p^\mu (\xi )+f^i\frac{x_p^\mu }{f^i}|_{f^i=0}+\frac{1}{2}f^if^j\frac{^2x_p^\mu }{f^if^j}|_{f^i=0}+O(f^3),$$
(58)
where $`\overline{x}_p^\mu (\xi )x_p^\mu (f^i=0,\xi )`$ is the position of the string worldsheet when $`f^1=f^2=0`$. Expanding $`\sqrt{g}`$ to quadratic order in small $`f^1`$ and $`f^2`$, we obtain
$`W[\mathrm{\Gamma }]`$ $`=`$ $`{\displaystyle 𝒟f^i(\xi )\mathrm{Det}\left[\frac{1}{2}\stackrel{~}{t}^{\mu \nu }ϵ^{ij}\frac{x_p^\mu }{f^i}\frac{x_p^\nu }{f^j}|_{f^i=0}\right]}`$ (60)
$`\times \mathrm{exp}\left\{i\sigma {\displaystyle d^2\xi \sqrt{\overline{g}}\left[1+\overline{g}^{ab}\frac{x_p^\mu }{\xi ^a}|_{f^i=0}\frac{}{\xi ^b}\left(\frac{x_p^\mu }{f^i}|_{f^i=0}f^i\right)+\frac{1}{2}f^iG_{ij}^1f^j\right]}\right\},`$
where $`\overline{g}_{ab}`$ is
$$\overline{g}_{ab}=\frac{\overline{x}_p^\mu }{\xi ^a}\frac{\overline{x}_p^\mu }{\xi ^b},$$
(61)
the metric of the fixed worldsheet $`\overline{x}_p^\mu `$, and
$$G_{ij}^1=\frac{1}{\sqrt{g}}\frac{^2\sqrt{g}}{f^if^j}|_{f^i=0}.$$
(62)
We choose $`\overline{x}_p^\mu `$ to be the surface which minimizes the action. Then $`\overline{x}_p^\mu `$ satisfies the “classical equation of motion”
$$\frac{x_p^\mu }{f^i}|_{f^i=0}(^2)\overline{x}_p^\mu =0,$$
(63)
where the covariant Laplacian is
$$^2=\frac{1}{\sqrt{\overline{g}}}\frac{}{\xi ^a}\overline{g}^{ab}\sqrt{\overline{g}}\frac{}{\xi ^b}.$$
(64)
Using the fact that the covariant derivative of the metric is zero, we show in Appendix A that
$$(_a\overline{x}_p^\mu )(^2\overline{x}_p^\mu )=0.$$
(65)
The vectors $`\frac{x_p^\mu }{f^i}|_{f^i=0}`$ and $`_a\overline{x}_p^\mu `$ form a complete basis, so (63) and (65) imply
$$^2\overline{x}_p^\mu =0.$$
(66)
Evaluating the $`f^i`$ integral in Eq. (60) gives
$$W[\mathrm{\Gamma }]=e^{i\sigma {\scriptscriptstyle d^2\xi \sqrt{\overline{g}}}}\mathrm{Det}\left[\frac{1}{2}\stackrel{~}{t}_{\mu \nu }ϵ^{ij}\frac{x_p^\mu }{f^i}\frac{x_p^\nu }{f^j}|_{f^i=0}\right]\mathrm{Det}^{1/2}\left[G_{ij}^1\right].$$
(67)
The inverse propagator $`G_{ij}^1`$ (62) can be shown to be
$$G_{ij}^1=\frac{x_p^\mu }{f^i}|_{f^i=0}\overline{n}_{\mu A}\left[^2\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}\right]\overline{n}_{\nu B}\frac{x_p^\nu }{f^j}|_{f^i=0},$$
(68)
where $`\overline{𝒦}_{ab}^A`$ is the extrinsic curvature tensor of the sheet $`\overline{x}_p^\mu `$. A derivation of (68) is given in appendix A of . The $`\overline{n}_{\mu A}`$ are vectors normal to the worldsheet $`\overline{x}_p^\mu `$. Eq. (68) gives
$$\mathrm{Det}^{1/2}\left[G_{ij}^1\right]=\mathrm{Det}^{1/2}\left[^2\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}\right]\mathrm{Det}^1\left[\frac{1}{2}ϵ^{AB}\overline{n}_{\mu A}\overline{n}_{\nu B}ϵ^{ij}\frac{x_p^\mu }{f^i}\frac{x_p^\nu }{f^j}|_{f^i=0}\right].$$
(69)
From the identity (A15),
$$\stackrel{~}{t}_{\mu \nu }=ϵ^{AB}\overline{n}_{\mu A}\overline{n}_{\nu B},$$
(70)
we see that the first determinant in (67) and the second determinant in (69) cancel. The determinant appearing in (54) produces exactly the correct normalization for the Green’s function. The functional integral (67) becomes
$$W[\mathrm{\Gamma }]=e^{i\sigma {\scriptscriptstyle d^2\xi \sqrt{\overline{g}}}}\mathrm{Det}^{1/2}\left[^2\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}\right].$$
(71)
We note that (71) is independent of the factors of $`n_\mu ^A\frac{x_p^\mu }{f^i}|_{f^i=0}`$ which appeared in the inverse propagator (68). These factors are the projections of the fluctuations $`f^i`$ normal to the string worldsheet. For small $`f^i`$, the worldsheet $`x_p^\mu `$ is
$$x_p^\mu =\overline{x}_p^\mu +f_i\frac{x_p^\mu }{f^i}|_{f^i=0}+O(f_{}^{i}{}_{}{}^{2}).$$
(72)
The perturbation of the worldsheet in the direction $`\overline{n}^{\mu A}`$ is
$$n_\mu ^A\left(x_p^\mu \overline{x}_p^\mu \right)=n_\mu ^A\frac{\overline{x}_p^\mu }{f^i}|_{f^i=0}f_i+O(f_{}^{i}{}_{}{}^{2}).$$
(73)
The factors of $`n_\mu ^A\frac{x_p^\mu }{f^i}|_{f^i=0}`$ in (68) project out the part of the fluctuations $`f^i`$ perpendicular to the worldsheet $`\overline{x}_p^\mu `$. Only normal fluctuations contribute to $`W[\mathrm{\Gamma }]`$, since fluctuations along the worldsheet are equivalent to a reparameterization of the sheet coordinates.
The effective action obtained from (71) is
$$i\mathrm{ln}W[\mathrm{\Gamma }]=S_{cl}+S_{\mathrm{fluc}}.$$
(74)
The first term in (74) is the Nambu–Goto action evaluated at the “classical” worldsheet $`\overline{x}_p^\mu (\xi )`$,
$$S_{cl}=\sigma d^2\xi \sqrt{\overline{g}}.$$
(75)
The semiclassical correction $`S_{\mathrm{fluc}}`$ due to the transverse string fluctuations is
$$S_{\mathrm{fluc}}=\frac{i}{2}\mathrm{Tr}\mathrm{ln}\left[^2\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}\right].$$
(76)
To summarize, we have integrated out the string fluctuations, and reduced the problem to the evaluation of the determinant in (71). This is a quantum mechanical scattering problem in the background of the solution of the classical equation (66), with appropriate boundary conditions. In the next section, we describe how to evaluate this determinant.
## VIII Regularization of String Integrals and the Role of the Lüscher Term
The argument of the logarithm in (76) is the inverse propagator for fluctuations on the string. This inverse propagator can also be obtained by direct variation of the Nambu–Goto action with respect to any transverse coordinates $`x_{}^A`$,
$$\frac{^2S}{x_{}^Ax_{}^B}=^2\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab},$$
(77)
up to an overall normalization factor. In fact, the correction (76) to the effective action has already been studied by Lüscher, Symanzik, and Weisz (LSW) in the case of a straight string with fixed ends. We describe their results, which we will use in evaluating (76). LSW used Pauli–Villars regularization to obtain a regulated form $`S_{\mathrm{reg}}`$ of the trace in (76),
$`S_{\mathrm{reg}}={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}\left(1+{\displaystyle \underset{j}{}}ϵ_je^{t_j^2}\right)\mathrm{Tr}e^{t\left(^2\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}\right)}.`$ (78)
The $`_j`$ are the masses of the regulators, and the $`ϵ_j`$ are suitably chosen coefficients. The Laplacian in (78) has been Wick rotated from Minkowski to Euclidean space.
The regulated quantity $`S_{\mathrm{reg}}`$ is separated into a divergent part $`S_{\mathrm{div}}`$ and a finite part $`S_{\mathrm{PV}}`$,
$$S_{\mathrm{reg}}=S_{\mathrm{div}}(_j,ϵ_j)+S_{\mathrm{PV}}.$$
(79)
LSW evaluated the divergent part $`S_{\mathrm{div}}(_j,ϵ_j)`$, and obtained terms which are quadratically, linearly, and logarithmically divergent in the cutoffs $`_j`$. The quadratic term is a renormalization of the string tension, the linear term is a renormalization of the quark masses, and the logarithmically divergent term is proportional to the integral over all space of the scalar curvature $``$ of the string worldsheet.
LSW also obtained a formal expression for the finite part $`S_{\mathrm{PV}}`$. They evaluated this expression only for the case of a straight string of length $`R`$ with fixed ends, and calculated a correction $`V_{\text{Lüscher}}`$ to the static potential:
$$V_{\text{Lüscher}}=\underset{T\mathrm{}}{lim}\frac{1}{T}S_{\mathrm{PV}}=\frac{\pi }{12R}.$$
(80)
We are interested in calculating $`S_{\mathrm{fluc}}`$ for rotating quarks, so we must evaluate $`S_{\mathrm{reg}}`$ for a more general surface. We break (76) into two parts:
$$S_{\mathrm{reg}}=i\mathrm{Tr}\mathrm{ln}[^2]+\frac{i}{2}\mathrm{Tr}\mathrm{ln}\left[\frac{^2\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}}{^2}\right].$$
(81)
We will evaluate the first term in (81) by generalizing the calculation of LSW. We will calculate the second term directly.
The first term in (81),
$$S_1=i\mathrm{Tr}\mathrm{ln}[^2],$$
(82)
involves the Laplacian in the curved background of the classical solution $`\overline{x}_p^\mu `$. In the flat case studied by LSW, the Laplacian is equal to $`^2`$,
$$^2=\frac{^2}{t^2}\frac{^2}{r^2}.$$
(83)
The coordinate $`t`$ is the time in the lab frame, and $`r`$ is a radial coordinate which takes the values $`R_1`$ and $`R_2`$ at the two ends of the string. The length of the string is $`R=R_1+R_2`$.
To calculate $`S_1`$ we extend the calculation of LSW to more general coordinate systems. We make a coordinate transformation $`\xi \xi ^{}`$ to conformal coordinates, where the transformed metric $`g_{ab}^{}=\eta _{ab}e^\phi `$, $`\eta _{ab}=\mathrm{diag}(1,1)`$. This transformation puts the Laplacian in a form similar to the flat sheet Laplacian (83), and allows us to evaluate (82) by extending the calculation of LSW.
To see how this works, we express $`e^{iS_1}`$ as a functional integral,
$$e^{iS_1}=𝒟f_1𝒟f_2\mathrm{exp}\left\{id^2\xi \sqrt{g}g^{ab}\frac{f^i}{\xi ^a}\frac{f^i}{\xi ^b}\right\}.$$
(84)
The transformation to conformal coordinates $`\xi ^{}`$ gives
$$e^{iS_1}=𝒟f_1𝒟f_2\mathrm{exp}\left\{id^2\xi ^{}\eta ^{ab}\frac{f^i}{\xi _{}^{}{}_{}{}^{a}}\frac{f^i}{\xi _{}^{}{}_{}{}^{b}}\right\},$$
(85)
which is of the form of $`S_{\mathrm{reg}}`$ treated by LSW.
We will need $`S_1`$ in the limit of large $`T`$, and hence are only interested in strings whose metric is time independent. To determine which metrics are time independent, we must choose a coordinate system. We choose coordinates $`r`$ and $`t`$, where $`t`$ is the time in the lab frame, and $`r`$ is orthogonal to $`t`$ ($`g_{rt}=0`$). This guarantees that $`t`$ is the physical time. From now on we consider only metrics $`g_{ab}`$ which are independent of $`t`$.
In Appendix C, we show that
$$S_1=S_{1,\mathrm{div}}+S_{1,\mathrm{finite}},$$
(86)
where $`S_{1,\mathrm{div}}`$ contains quark mass and string tension renormalizations. The finite part of $`S_1`$ is
$$S_{1,\mathrm{finite}}=T\frac{\pi }{12R_p},$$
(87)
where
$$R_p=_{R_1}^{R_2}𝑑r\sqrt{\frac{\overline{g}_{rr}(r)}{\overline{g}_{tt}(r)}}.$$
(88)
The results (87) and (88) are valid for any orthogonal coordinate system with a time independent metric.
We show in Appendix C that $`R_p`$ is equal to the classical energy of the string divided by the string tension $`\sigma `$. We call $`R_p`$ the “proper length” of the string. For a flat metric, where $`\overline{g}_{rr}=\overline{g}_{tt}=1`$, the proper length $`R_p`$ of the string reduces to the distance $`R`$ between its endpoints.
## IX Correction to the Effective Action for a Curved Sheet
In the previous section, we evaluated the finite part of $`S_1`$ for a sheet with a time independent metric using the results of LSW. In this section we evaluate the second term in (81),
$$S_2i\frac{T}{2}\mathrm{Tr}\mathrm{ln}\left[\frac{^2\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}}{^2}\right].$$
(89)
The trace in (89) is over functions of $`r`$ and $`t`$.
We first make the coordinate transformation $`rx`$,
$$\frac{dx}{dr}=\sqrt{\frac{\overline{g}_{rr}(r)}{\overline{g}_{tt}(r)}},x|_{r=0}=0.$$
(90)
The coordinate $`x`$ runs from $`X_1`$ to $`X_2`$,
$`X_1`$ $`=`$ $`{\displaystyle _{R_1}^0}𝑑r\sqrt{{\displaystyle \frac{\overline{g}_{rr}(r)}{\overline{g}_{tt}(r)}}}`$ (91)
$`X_2`$ $`=`$ $`{\displaystyle _0^{R_2}}𝑑r\sqrt{{\displaystyle \frac{\overline{g}_{rr}(r)}{\overline{g}_{tt}(r)}}},`$ (92)
and $`X_1+X_2=R_p`$. In Appendix C, we show that the metric in the system $`(x,t)`$ is conformal ($`\overline{g}_{xx}=\overline{g}_{tt}`$, $`\overline{g}_{xt}=0`$). In this coordinate system, the inverse propagator for string fluctuations is
$$^2\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}=\frac{1}{\sqrt{\overline{g}}}\left(\frac{^2}{t^2}\frac{^2}{x^2}\right)\delta _{AB}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}.$$
(93)
The string has infinite extent in time, and the curvature $`\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}`$ is independent of $`t`$, so we can take the Fourier transform with respect to the time coordinate. We express the trace in (89) over functions of $`t`$ and $`r`$ as an integral over a frequency $`\nu `$ and a trace over functions of a single variable $`x`$,
$$S_2=\frac{T}{2}_{\mathrm{}}^{\mathrm{}}\frac{d\nu }{2\pi }\mathrm{Tr}_x\mathrm{ln}\left[\frac{\left(\nu ^2\frac{^2}{x^2}\right)\delta _{AB}\sqrt{\overline{g}}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}}{\nu ^2\frac{^2}{x^2}}\right].$$
(94)
In going from (89) to (94), we have also carried out the Wick rotation $`\nu i\nu `$. The integration over $`\nu `$ gives
$$S_2=\frac{T}{2}\left[\mathrm{Tr}_x\sqrt{\frac{^2}{x^2}\delta _{AB}\sqrt{\overline{g}}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}}\mathrm{Tr}\sqrt{\frac{^2}{x^2}\delta _{AB}}\right].$$
(95)
Eq. (95) expresses $`S_2`$ as the trace of the difference of two operators. The first has the form of a Hamiltonian for a relativistic particle in the local potential $`\sqrt{\overline{g}}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}`$. The second operator has the form of a free Hamiltonian. The square roots enter because we are working with relativistic degrees of freedom.
The terms $`S_1`$ and $`S_2`$ are proportional to the time $`T`$, but are otherwise time independent. We define the “effective Lagrangian” of the string to be the effective action divided by the time $`T`$. The sum of (87) and (95) gives the effective Lagrangian $`L_{\mathrm{fluc}}`$ determining the contribution of the string fluctuations to $`W[\mathrm{\Gamma }]`$,
$`L_{\mathrm{fluc}}`$ $``$ $`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{T}}\left(S_{1,\mathrm{finite}}+S_2\right)`$ (96)
$`=`$ $`{\displaystyle \frac{\pi }{12R_p}}+{\displaystyle \frac{1}{2}}\left(\mathrm{Tr}\sqrt{{\displaystyle \frac{^2}{x^2}}\delta _{AB}\sqrt{\overline{g}}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}}\mathrm{Tr}\sqrt{{\displaystyle \frac{^2}{x^2}}\delta _{AB}}\right).`$ (97)
In the next section, we will evaluate $`L_{\mathrm{fluc}}`$ for a string of length $`R`$ rotating with angular velocity $`\omega `$.
In Appendix D, we show that, for a general sheet, $`S_2`$ is logarithmically divergent. We show that its divergent part is given by
$$S_{2,\mathrm{div}}=\frac{T}{4}\underset{n=1}{\overset{MR_p/\pi }{}}\frac{1}{\pi n}_{X_1}^{X_2}𝑑x\sqrt{\overline{g}},$$
(98)
where $``$ is the scalar curvature,
$$=\left(𝒦_a^{Aa}\right)^2\left(𝒦_{ab}^A\right)^2.$$
(99)
This result agrees, in the large time limit, with the logarithmically divergent term in the cutoff dependent part of the effective string action (C15) found by LSW.
## X Effective Lagrangian for Rotating String
We now evaluate the effective Lagrangian of a string with boundary $`\mathrm{\Gamma }`$ generated by a quark–antiquark pair separated at fixed distances $`R_1`$ and $`R_2`$ from the origin, and rotating with angular velocity $`\omega `$. This Lagrangian has two parts, the classical string Lagrangian and the contribution (97) of string fluctuations.
We evaluate the classical string Lagrangian first. The solution to the classical equations of motion (66) yields the classical, straight rotating string,
$$\overline{x}^\mu (r,t)=t\widehat{𝐞}_0^\mu +r\mathrm{cos}(\omega t)\widehat{𝐞}_1^\mu +r\mathrm{sin}(\omega t)\widehat{𝐞}_2^\mu .$$
(100)
The coordinate $`r`$ is chosen so that the velocity of the string is zero when $`r=0`$. The coordinate $`r`$ runs from $`R_1`$ to $`R_2`$, and $`t`$ runs from $`\mathrm{}`$ to $`\mathrm{}`$. The vectors $`\widehat{𝐞}_1^\mu `$ and $`\widehat{𝐞}_2^\mu `$ are two orthogonal unit vectors in the plane of rotation, and $`\widehat{𝐞}_0^\mu `$ is a unit vector in the time direction. The classical Lagrangian $`L_{cl}^{\mathrm{string}}`$ obtained from (75) is
$`L_{cl}^{\mathrm{string}}`$ $``$ $`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{T}}\sigma {\displaystyle d^2\xi \sqrt{\overline{g}}}`$ (101)
$`=`$ $`\sigma {\displaystyle _{R_1}^{R_2}}𝑑r\sqrt{1r^2\omega ^2}`$ (102)
$`=`$ $`\sigma {\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{R_i}{2}}\left({\displaystyle \frac{\mathrm{arcsin}(\omega R_i)}{\omega R_i}}+\sqrt{1R_i^2\omega ^2}\right).`$ (103)
Next, we calculate the contribution $`L_{\mathrm{fluc}}`$ (97) due to string fluctuations. The metric of the sheet (100) is
$$\overline{g}_{tt}=1+r^2\overline{\omega }^2,\overline{g}_{rr}=1,\overline{g}_{rt}=0.$$
(104)
This metric is independent of $`t`$, and $`\overline{g}_{rt}=0`$. We make the transformation (90) from coordinates $`r`$ and $`t`$ to coordinates $`x`$ and $`t`$, and find
$$x=\frac{1}{\omega }\mathrm{arcsin}\omega r.$$
(105)
The coordinate $`x`$ runs from $`X_1`$ to $`X_2`$, where
$$X_i=\frac{\mathrm{arcsin}\left(\omega R_i\right)}{\omega },$$
(106)
and the proper length $`R_p`$ of the string is
$$R_p=X_1+X_2=\underset{i}{}\frac{\mathrm{arcsin}\left(\omega R_i\right)}{\omega }.$$
(107)
Using this coordinate system , we evaluate $`L_{\mathrm{fluc}}`$ in Appendix E for the case of equal quark masses, $`R_1=R_2=R/2`$:
$$L_{\mathrm{fluc}}=\frac{v}{\mathrm{arcsin}v}\frac{\pi }{12R}\frac{2v}{\pi R}\left[v\gamma \mathrm{ln}\left(\frac{MR}{2(\gamma ^21)}\right)+v\gamma \frac{\pi }{2}\right]\frac{v^2\gamma }{\pi R}f\left(v\right),$$
(108)
where
$$v\frac{R}{2}\omega ,\gamma =\frac{1}{\sqrt{1v^2}},$$
(109)
and the function $`f(v)`$ is
$$f(v)=_0^{\mathrm{}}𝑑s\mathrm{ln}\left[\frac{s^2+2s\mathrm{coth}\left(2sv\gamma \mathrm{arcsin}v\right)+1}{(s+1)^2}\right].$$
(110)
Eq. (108) becomes the Lüscher term in the zero velocity limit.
We are interested in the large $`R`$ limit, where the quark velocity is close to the speed of light. For $`v`$ close to one, Eq. (108) becomes
$$L_{\mathrm{fluc}}=\frac{2}{\pi R}\gamma \left[\mathrm{ln}\left(\frac{MR}{2\gamma ^2}\right)+1\right]+\frac{7}{6R}+O\left(\frac{\mathrm{ln}\gamma }{\gamma R}\right).$$
(111)
Furthermore, for the semiclassical expansion to be valid, the theory must be weakly coupled. That is, $`L_{\mathrm{fluc}}`$ must be less than $`L_{cl}^{\mathrm{string}}`$ (103). For large $`R`$,
$$L_{cl}^{\mathrm{string}}=\frac{\pi }{4}\sigma R\frac{\pi }{8}\sigma R\gamma ^2+\frac{1}{6}\sigma R\gamma ^3+O(\gamma ^4\sigma R).$$
(112)
The semiclassical expansion is valid, since, as we will see, $`R`$ grows like $`\gamma ^2`$ in the $`v1`$ limit. In this case, the long distance limit where the effective theory is applicable is automatically the region of weak coupling.
## XI Regge Trajectories
We calculate classical Regge trajectories for equal mass quarks by adding a quark mass term to the string Lagrangian $`L_{cl}^{\mathrm{string}}`$,
$$L_{cl}=L_{cl}^{\mathrm{string}}2m\sqrt{1v^2}=\sigma \frac{R}{2}\left(\frac{\mathrm{arcsin}v}{v}+\gamma ^1\right)2m\gamma ^1.$$
(113)
We have used Eq. (103) with $`R_1=R_2=R/2`$. The quark velocity is $`v=\omega R/2`$, and $`\gamma =1/\sqrt{1v^2}`$ is the quark boost factor. The Lagrangian (113) is a function of $`R`$ and $`\omega `$,
$$L_{cl}=L_{cl}(R,\omega ).$$
(114)
The angular momentum of the meson is obtained by varying the Lagrangian with respect to the angular velocity,
$$J=\frac{L_{cl}}{\omega }=\sigma \frac{R^2}{4v}\left(\frac{\mathrm{arcsin}v}{v}\gamma ^1\right)+mRv\gamma ^1.$$
(115)
The meson energy is given by the Hamiltonian,
$$E=\omega \frac{L_{cl}}{\omega }L_{cl}=\sigma R\frac{\mathrm{arcsin}v}{v}+2m\gamma .$$
(116)
The classical equation of motion
$$\frac{L}{R}=0$$
(117)
for the quarks determines $`R`$ as a function of $`\omega `$,
$$\sigma \frac{R}{2}=m(\gamma ^21).$$
(118)
Eq. (118) shows that $`R`$ is proportional to $`\gamma ^2`$ for large $`\gamma `$. Expanding (115) and (116) in the large $`R`$ limit, where the quark velocity $`v`$ goes to one, yields the result:
$$\frac{J}{E_{cl}^2}=\frac{1}{2\pi \sigma }\left(1\frac{8}{3\pi }\gamma ^3+O(\gamma ^5)\right).$$
(119)
The first term in (119) is the classical formula for the slope of a Regge trajectory. The second term is the leading correction for nonzero classical quark mass, where $`\gamma ^1=\sqrt{1v^2}0`$.
We now calculate the correction to the energy obtained by considering $`L_{\mathrm{fluc}}`$ a small perturbation to the classical Lagrangian $`L_{cl}`$. The Lagrangian,
$$L(\omega )=L_{cl}(\omega )+L_{\mathrm{fluc}}(\omega ),$$
(120)
depends on only one degree of freedom, the rotation angle $`\theta `$, through its time derivative $`\omega =\dot{\theta }`$, To first order in $`L_{\mathrm{fluc}}`$, the correction to the energy is minus the correction to the Lagrangian ,
$$E(J)=\left[E_{cl}(\omega )L_{\mathrm{fluc}}(\omega )\right]|_{\omega =\omega (J)},$$
(121)
where $`\omega `$ is given as a function of $`J`$ through the classical relation (115).
The correction (121) to the energy of the meson gives a correction to the slope (119) of a Regge trajectory,
$$\frac{J}{E^2}=\frac{J}{E_{cl}^2}\frac{E_{cl}^2}{E^2}\frac{J}{E_{cl}^2}\left(1+2\frac{L_{\mathrm{fluc}}}{E}\right).$$
(122)
Using (119) for $`J/E_{cl}^2`$ and (111) for $`L_{\mathrm{fluc}}`$, we obtain
$$\frac{J}{E^2}=\frac{1}{2\pi \sigma }\frac{2}{\pi ^2\sigma RE}\gamma \left[\mathrm{ln}\left(\frac{MR}{2\gamma ^2}\right)+1\right]\frac{4}{3\pi ^2\sigma }\gamma ^3+\frac{7}{6\pi \sigma RE}+O(\gamma ^5,\frac{1}{RE\gamma }).$$
(123)
We write $`R`$ and $`\gamma `$ as functions of $`E`$ using the definition (116) of $`E`$ and the classical equation of motion (118). Because $`R`$ and $`\gamma `$ only appear in the small correction terms in the result (123), we only need their leading order dependence on $`E`$,
$$R\frac{2E}{\pi \sigma },\gamma \sqrt{\frac{E}{\pi m}}.$$
(124)
Substituting (124) in (123) gives
$$J=\frac{E^2}{2\pi \sigma }\sqrt{\frac{E}{\pi ^3m}}\left[\mathrm{ln}\left(\frac{Mm}{\sigma }\right)+1\right]\frac{4}{3\sigma }\sqrt{\frac{m^3E}{\pi }}+\frac{7}{12}+O\left(E^{1/2}\right).$$
(125)
The leading term is the classical Regge formula. The next term is the leading correction due to string fluctuations. The third term is a nonzero quark mass correction. The fourth term is another correction due to string fluctuations.
Eq. (125) gives a meson Regge trajectory $`J(E^2)`$.
We used values $`\frac{4}{3}\alpha _s=0.25`$ and $`\sigma =(455\text{ MeV})^2`$ obtained from the Cornell fits of heavy quark potentials . This gives $`M=g\varphi _0=\sqrt{\frac{3\sigma }{4\alpha _s}}=910\text{ MeV}`$. The only other parameter is the quark mass $`m`$. In Figure 3, we plot $`J`$ versus the square of the energy (121) for quark masses of 30 MeV, 100 MeV, and 300 MeV. For comparison, we also plot the classical formula $`J=E^2/2\pi \sigma `$. The points plotted on the graph are the $`\rho ^1`$, $`A^2`$, $`\rho ^3`$, $`A^4`$, $`\rho ^5`$, and $`A^6`$ mesons. The ‘x’s are the $`\omega ^1`$, $`f^2`$, $`\omega ^3`$, $`f^4`$, and $`f^6`$. We have added one to the value of the angular momentum $`J`$ in Figure 3 to account for the contribution of the spin of the quarks.
We have chosen a range of values for the quark masses in Fig. 3 in order to give a qualitative picture of the dependence of the Regge trajectory on the quark mass. Since (125) does not include the contribution of quark fluctuations to the Regge trajectory, this formula is incomplete. We are now in the process of including the quark degrees of freedom in the functional integral (54). The boundary $`\mathrm{\Gamma }`$ of the sheet $`\stackrel{~}{x}^\mu `$ becomes dynamical, and couples to the string fluctuations. It is clear that a calculation of the contribution of these degrees of freedom is essential to understanding why the classical formula for Regge trajectories works so well.
## XII Summary and Conclusions
The primary results of the paper are (54) and (125). We have expressed the path integral $`W[\mathrm{\Gamma }]`$ (5) of a renormalizable quantum field theory having classical vortex solutions as the path integral formulation of an effective string theory of vortices (54). This theory describes the two transverse fluctuations of the vortex at scales larger than the inverse mass of the lightest particle in the field theory. Our method is applicable to any field theory containing vortex solutions.
Using the string representation of $`W[\mathrm{\Gamma }]`$, we carried out a semiclassical expansion of the effective action $`i\mathrm{log}W[\mathrm{\Gamma }]`$ about a classical solution of the effective string theory. We calculated the contribution of these string fluctuations explicitly for the case where the worldline $`\mathrm{\Gamma }`$ is generated by the trajectory of a quark–antiquark pair separated by a distance $`R`$. We are now calculating the contribution to the effective action $`i\mathrm{log}W[\mathrm{\Gamma }]`$ due to the quantum fluctuations of the boundary.
## Acknowledgments
We would like to thank N. Brambilla for very helpful conversations. This work was supported in part by the U. S. Department of Energy grant DE-FG03-96ER40956.
## A Notation and the Curvature of the Vortex
We describe the string worldsheet by the function $`\stackrel{~}{x}^\mu (\xi )`$ of the coordinates $`\xi `$. The physics of the vortex should be independent of the coordinate system we choose, so we require the theory to be invariant under a reparameterization of the coordinates, $`\xi \stackrel{~}{\xi }(\xi )`$. The tangent vectors to the vortex worldsheet are defined by taking derivatives of $`\stackrel{~}{x}^\mu (\xi )`$,
$$t_a^\mu (\xi )_a\stackrel{~}{x}^\mu (\xi ),$$
(A1)
where $`_a=/\xi ^a`$ is a partial derivative with respect to one of the vortex coordinates. The induced metric on the worldsheet $`\stackrel{~}{x}^\mu (\xi )`$ is
$$g_{ab}t_a^\mu t_{\mu b}.$$
(A2)
It is also convenient to define the square root of the determinant of the metric,
$$\sqrt{g}=\sqrt{\frac{1}{2}ϵ^{ab}ϵ^{cd}g_{ac}g_{bd}}.$$
(A3)
We use the $`t_a^\mu `$ to define an antisymmetric tensor which describes the orientation of the string worldsheet,
$$t^{\mu \nu }\frac{ϵ^{ab}}{\sqrt{g}}t_a^\mu t_b^\nu .$$
(A4)
This quantity was defined by Polyakov . It is the projection of the antisymmetric tensor $`ϵ^{ab}`$ into the space of four dimensional tensors. This tensor defines the orientation of the two dimensional vortex worldsheet in four space. The quantity (A4) is also independent of the coordinate parameterization of the worldsheet $`\stackrel{~}{x}^\mu (\xi )`$.
We now describe the curvature of the vortex worldsheet. We do this by taking covariant derivatives of the tangent vectors. The covariant derivative of $`t_a^\mu `$ is
$$_bt_a^\mu =_bt_\mu ^a\mathrm{\Gamma }_{ab}^ct_c^\mu ,$$
(A5)
where the $`\mathrm{\Gamma }_{ab}^c`$ are Christoffel symbols,
$$\mathrm{\Gamma }_{ab}^c=\frac{1}{2}g^{cd}\left(_ag_{bd}+_bg_{ad}_dg_{ab}\right),$$
(A6)
The covariant derivatives of the tangent vectors are orthogonal to the worldsheet,
$$t_{\mu c}_at_b^\mu =t_{\mu c}\left(_at_b^\mu \mathrm{\Gamma }_{ab}^dt_d^\mu \right)=0.$$
(A7)
This identity is derived using the definition (A6) of the Christoffel symbols, the definition (A2) of the metric, and the relationship between derivatives of different $`t_a^\mu `$,
$$_at_b^\mu =_bt_a^\mu =_a_b\stackrel{~}{x}^\mu .$$
(A8)
The covariant derivatives of the tangent vectors are normal to the string worldsheet. We therefore define a basis of normal vectors $`n_A^\mu `$, which satisfy the conditions
$$n_{\mu A}t_a^\mu =0,n_{\mu A}n_B^\mu =\delta _{AB}.$$
(A9)
The $`n_A^\mu `$ are an orthonormal basis for the vectors normal to the worldsheet. Eq. (A7) implies that
$$_at_b^\mu =n_A^\mu 𝒦_{ab}^A$$
(A10)
for some tensor $`𝒦_{ab}^A`$. The tensor $`𝒦_{ab}^A`$ is called the extrinsic curvature tensor of the string worldsheet. With the definition (A10), the curvature tensor $`𝒦_{ab}^A`$ is
$$𝒦_{ab}^An_\mu ^A_at_b^\mu =n_\mu ^A_a_bx^\mu .$$
(A11)
It is symmetric in the indices $`a`$ and $`b`$ due to the relationship (A8) between derivatives of tangent vectors.
The extrinsic curvature of the string worldsheet can also be described using derivatives of the normal vectors. The orthogonality of the $`t_{\mu a}`$ and the $`n_A^\mu `$ implies
$$t_{\mu b}_an_A^\mu =𝒦_{ab}^A.$$
(A12)
Therefore, the derivatives of the normal vectors can be written as
$$_an_A^\mu =t^{\mu b}𝒦_{ab}^A+n_B^\mu 𝒜_a^{AB}.$$
(A13)
The tensor $`𝒜_a^{AB}`$ is called the torsion, and it describes the twisting of the basis of normal vectors as we move along the worldsheet. The torsion depends on our choice of the $`n_A^\mu `$, so we will choose them so that the torsion is zero. This is done by requiring that the $`n_A^\mu `$ satisfy the differential equation
$$_an_A^\mu =n_A^\nu (_at_{\nu b})t^{\mu b}.$$
(A14)
The equation (A14) is equivalent to the statement $`𝒜_b^{AB}=0`$. It is consistent with the conditions (A9) which define the normal vectors. As long as the normal vectors have an orthonormal basis at one point, Eq. (A14) guarantees they will be orthonormal in a neighborhood of that point. Therefore, it is always possible to find a local, orthonormal, torsion free basis for the normal vectors.
There is one additional property of the normal vectors we will use. The antisymmetric combination of the normal vectors is (with proper ordering) equal to the dual of the worldsheet orientation tensor $`t^{\mu \nu }`$ (A4),
$$ϵ^{AB}n_A^\mu n_B^\nu =\stackrel{~}{t}^{\mu \nu },$$
(A15)
where
$$\stackrel{~}{t}^{\mu \nu }=\frac{1}{2}ϵ^{\mu \nu \alpha \beta }t_{\alpha \beta }.$$
(A16)
The relationship (A15) can be understood by noting that any antisymmetric tensor is of the form
$$A^{\mu \nu }(\xi )=T(\xi )t^{\mu \nu }+N(\xi )ϵ^{AB}n_A^\mu n_B^\nu +M^{Aa}\left(n_A^\mu t_a^\nu n_A^\nu t_a^\mu \right).$$
(A17)
The tensor $`\stackrel{~}{t}^{\mu \nu }`$ is orthogonal to the $`t_a^\mu `$, so it must be proportional to $`ϵ^{AB}n_A^\mu n_B^\nu `$. Squaring both of these tensors gives
$$\left(\stackrel{~}{t}^{\mu \nu }\right)^2=\left(ϵ^{AB}n_A^\mu n_B^\nu \right)^2=2.$$
(A18)
Therefore, $`\stackrel{~}{t}^{\mu \nu }`$ and $`ϵ^{AB}n_A^\mu n_B^\nu `$ are equal up to an overall sign, which is fixed by choosing an appropriate ordering for the normal vectors.
## B Discussion of $`J_{}`$
The Jacobian $`J_{}`$ (27) is a Faddeev–Popov determinant, because fixing the position of the string in the field integrals is analogous to fixing a gauge in a gauge theory. In the string action, we fix the degrees of freedom which generate the transformation
$$\stackrel{~}{x}^\mu (\xi )\stackrel{~}{x}^\mu (\xi )+\delta x_{}^A(\xi )n_{\mu A}(\xi )$$
(B1)
which displaces the vortex.
The Jacobian $`J_{}`$ is analogous to the Faddeev–Popov determinant in a gauge theory. In a gauge theory, where the $`\delta `$ function fixes the symmetry generated by the transformation
$$A_\mu UA_\mu U^1+U_\mu U^1,$$
(B2)
the Faddeev–Popov determinant appears as a normalization for the $`\delta `$ function ,
$$Z_{\mathrm{gauge}}=𝒟A^\mu \delta \left[F(A^\mu )\right]\mathrm{\Delta }_{FP}e^S,$$
(B3)
where
$$\mathrm{\Delta }_{FP}^1=𝒟U\delta \left[F(A^\mu )\right].$$
(B4)
The Wilson loop (B3) is analogous to our equation (24) for the effective action. The determinant $`\mathrm{\Delta }_{FP}`$ is analogous to $`J_{}`$.
In the gauge theory, the Faddeev–Popov method is used to remove nonphysical degrees of freedom from the problem. The $`\delta `$ function is inserted in the path integral eqnlessrefgauge part to fix the fields in some particular gauge. This creates an integral over all gauges which appears as a normalization factor, and is removed. The $`\delta `$ function in Eq. (24), on the other hand, fixes the position of the vortex sheet, which is a physical degree of freedom.
## C Evaluation of $`S_1`$
We want to evaluate the term,
$$S_1=i\mathrm{Tr}\mathrm{ln}[^2],$$
(C1)
in the effective action for a general string worldsheet. We work in coordinates $`r`$ and $`t`$, such that $`t`$ is the time in the lab frame, and $`r`$ is orthogonal to $`t`$ ($`g_{rt}=0`$). In these coordinates, the functional integral (84) for $`e^{iS_1}`$ takes the form
$$e^{iS_1}=𝒟f_1𝒟f_2\mathrm{exp}\left\{i𝑑t_{R_1}^{R_2}𝑑r\sqrt{g_{tt}g_{rr}}\left[g^{tt}\left(\frac{f^i}{t}\right)^2+g^{rr}\left(\frac{f^i}{r}\right)^2\right]\right\}.$$
(C2)
We consider the case where the metric is independent of $`t`$, and we make the coordinate transformation $`rx`$ defined by
$$\frac{dx}{dr}=\sqrt{\frac{g_{rr}}{g_{tt}}},x|_{r=0}=0.$$
(C3)
The coordinate $`x`$ runs from $`X_1`$ to $`X_2`$,
$`X_1`$ $`=`$ $`{\displaystyle _{R_1}^0}𝑑r\sqrt{{\displaystyle \frac{g_{rr}}{g_{tt}}}}`$ (C4)
$`X_2`$ $`=`$ $`{\displaystyle _0^{R_2}}𝑑r\sqrt{{\displaystyle \frac{g_{rr}}{g_{tt}}}}.`$ (C5)
In these coordinates, the length of the string is $`X_1+X_2=R_p`$, the proper length of the string,
$$R_p=_{R_1}^{R_2}𝑑r\sqrt{\frac{g_{rr}}{g_{tt}}}.$$
(C6)
In the coordinate system $`(x,t)`$, the metric is conformal,
$`g_{xx}`$ $`=`$ $`\left({\displaystyle \frac{dx}{dr}}\right)^2g_{rr}=g_{tt},`$ (C7)
$`g_{xt}`$ $`=`$ $`\left({\displaystyle \frac{dx}{dr}}\right)^1g_{rt}=0,`$ (C8)
and
$$e^{iS_1}=𝒟f_1𝒟f_2\mathrm{exp}\left\{𝑑t_{X_1}^{X_2}𝑑x\left[\left(\frac{f^i}{t}\right)^2+\left(\frac{f^i}{x}\right)^2\right]\right\}.$$
(C9)
We evaluate (C9) in a manner analogous to our treatment of $`S_2`$ in section IX. We Fourier transform in both space and time, introducing variables $`\nu `$ and $`k_n=\pi n/R_p`$. This transformation puts the action in (C9) in a diagonal form. Doing the $`f_1`$ and $`f_2`$ integrals gives
$$S_1=\underset{\nu ,n}{}\mathrm{ln}\left[\nu ^2+\left(\frac{\pi n}{R_p}\right)^2\right],$$
(C10)
where we have Wick rotated $`\nu i\nu `$ to avoid the poles at $`\nu =\pm \pi n/R_p`$.
The length $`R_p`$ is just the classical string energy $`E_{cl}`$ divided by the string tension $`\sigma `$, since $`E_{cl}`$ is
$`E_{cl}`$ $`=`$ $`\sigma {\displaystyle \frac{1}{T}}{\displaystyle d^2\xi \frac{}{\dot{x}^0}\sqrt{g}}`$ (C11)
$`=`$ $`\sigma {\displaystyle \frac{1}{T}}{\displaystyle 𝑑t_{R_1}^{R_2}𝑑r\sqrt{\frac{g_{rr}}{g_{tt}}}}`$ (C12)
$`=`$ $`\sigma R_p.`$ (C13)
The quantity $`R_p=E_{cl}/\sigma `$ is the length of the string measured in local co-moving coordinates, which are at rest with respect to the string. This is different from the string length $`R`$ in the laboratory frame.
We will regulate $`S_1`$ using the results of LSW. Their result for $`S_{\mathrm{reg}}`$ is the following,
$`S_{\mathrm{reg}}`$ $`=`$ $`{\displaystyle \frac{d2}{4\pi }}A(𝒞){\displaystyle \underset{j}{}}ϵ_j_j^2\mathrm{ln}_j^2{\displaystyle \frac{d2}{4}}L(𝒞){\displaystyle \underset{j}{}}ϵ_j_j`$ (C15)
$`+\left({\displaystyle \frac{d2}{6}}{\displaystyle \frac{1}{4\pi }}{\displaystyle d^2\xi \sqrt{g}}\right){\displaystyle \underset{j}{}}ϵ_j\mathrm{ln}_j^2+S_{\mathrm{PV}},`$
where $`d`$ is the number of dimensions, $`A(𝒞)`$ is the area of the string worldsheet, $`L(𝒞)`$ is the length of its boundary, and $``$ is the scalar curvature of the sheet. The $`_j`$ are regulator masses, and the $`ϵ_j`$ are appropriate coefficients. The final term, $`S_{\mathrm{PV}}`$, is finite in the limit where the $`_j\mathrm{}`$.
LSW evaluate the finite term $`S_{\mathrm{PV}}`$ only for a straight string of length $`R`$ with fixed ends. In this case, the area of the sheet $`A(𝒞)=RT`$, the length of the boundary $`L(𝒞)=2T`$, and the curvature of the sheet is zero. They then obtained the explicit contribution to the heavy quark potential:
$$\underset{T\mathrm{}}{lim}\frac{1}{T}S_{\mathrm{reg}}=\frac{1}{2\pi }R\underset{j}{}ϵ_j_j^2\mathrm{ln}_j^2+\underset{j}{}ϵ_j_j\frac{\pi }{12R}.$$
(C16)
The first term in (C16) renormalizes the string tension. The second renormalizes the quark mass. The third is the well known Lüscher term in the heavy quark potential.
Since the extrinsic curvature vanishes for a flat sheet, we can identify the result (C16) with our expression (C10) for $`S_1`$, with $`R_p`$ replaced by $`R`$:
$$\underset{T\mathrm{}}{lim}\frac{1}{T}\underset{\nu ,n}{}\mathrm{ln}\left[\nu ^2+\left(\frac{\pi n}{R}\right)^2\right]=\frac{1}{2\pi }R\underset{j}{}ϵ_j_j^2\mathrm{ln}_j^2+\underset{j}{}ϵ_j_j\frac{\pi }{12R}.$$
(C17)
Eq. (C17) tells us how to regulate $`S_1`$. Replacing $`R`$ by $`R_p`$ in (C17) gives the regulated form of $`S_1`$:
$$\underset{T\mathrm{}}{lim}\frac{1}{T}S_1=\frac{1}{2\pi }R_p\underset{j}{}ϵ_j_j^2\mathrm{ln}_j+\underset{j}{}ϵ_j_j\frac{\pi }{12R_p}.$$
(C18)
The first term in (C18) is still a string tension renormalization, since both the string tension contribution to the energy (C13) and the first term in Eq. (C18) are proportional to $`R_p`$, The second term in Eq. (C18) is, as before, a renormalization of the quark mass. The finite part of the contribution of $`S_1`$ to the action is
$$S_1|_{\mathrm{finitepart}}=T\frac{\pi }{12R_p}.$$
(C19)
This is the result stated in section VIII. The result (C19) is the Lüscher term, with the distance $`R`$ between the quarks replaced by the proper length $`R_p`$ of the string. Our result is (C6), the derivation of $`R_p`$.
## D Cutoff Dependence of $`S_2`$
In this appendix, we show that the divergent part of $`S_2`$ for a general sheet is proportional to the integral of the scalar curvature $``$. This agrees with the logarithmic divergence (C15) derived by LSW by other means. To obtain the divergent part of $`S_2`$, we carry out a Fourier transform with respect to the variable $`x`$. For functions defined on the interval $`X_1<x<X_2`$, the $`\delta `$ function can be expressed as a sum of sines,
$$\delta (xx^{})=\frac{2}{R_p}\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{sin}\left(k_n(x+X_1)\right)\mathrm{sin}\left(k_n(x^{}+X_1)\right),$$
(D1)
where $`k_n=\pi n/R_p`$. The Fourier transform of an operator of the form $`^2/x^2+U(x)`$ can then be written
$`k_m|{\displaystyle \frac{^2}{x^2}}+U(x)|k_n`$ $`=`$ $`{\displaystyle \frac{2}{R_p}}{\displaystyle _{X_1}^{X_2}}𝑑x\mathrm{sin}\left(k_m(x+X_1)\right)\left({\displaystyle \frac{^2}{x^2}}+U(x)\right)\mathrm{sin}\left(k_n(x+X_1)\right)`$ (D2)
$`=`$ $`k_n^2\delta _{n,m}+{\displaystyle \frac{2}{R_p}}{\displaystyle _{X_1}^{X_2}}𝑑x\mathrm{sin}\left(k_m(x+X_1)\right)\mathrm{sin}\left(k_n(x+X_1)\right)U(x).`$ (D3)
Using the formula (LABEL:fourier\_transform) to evaluate (95) gives
$`S_2`$ $`=`$ $`{\displaystyle \frac{T}{2}}\mathrm{Tr}_n[\sqrt{k_n^2\delta _{n,m}\delta _{AB}{\displaystyle \frac{2}{R_p}}{\displaystyle _{X_1}^{X_2}}𝑑x\mathrm{sin}\left(k_m(x+X_1)\right)\mathrm{sin}\left(k_n(x+X_1)\right)\sqrt{\overline{g}}\overline{𝒦}_{ab}^A\overline{𝒦}^{Bab}}`$ (D6)
$`k_n\delta _{n,m}\delta _{AB}],`$
The trace is over indices $`A,B`$ which run from 1 to 2, and indices $`n,m`$ which run from 1 to $`\mathrm{}`$. The trace is cutoff at $`k_n=M`$, the mass of the vector particle in the original field theory.
We expand (D6) for large $`k_n`$ and obtain the cutoff dependent part of $`S_2`$,
$$S_2=\frac{T}{4}\underset{n=1}{\overset{MR_p/\pi }{}}\frac{1}{R_pk_n}_{X_1}^{X_2}𝑑x\sqrt{\overline{g}}\left(\overline{𝒦}_{ab}^A\right)^2+\text{finite.}$$
(D7)
The term $`\left(\overline{𝒦}_{ab}^A\right)^2`$ is equal to minus the scalar curvature $``$,
$$=\left(\overline{𝒦}_a^{Aa}\right)^2\left(\overline{𝒦}_{ab}^A\right)^2,$$
(D8)
since the equation of motion (66) implies $`\overline{𝒦}_a^{Aa}=0`$. The cutoff dependent part of $`S_2`$ is therefore
$$S_2=\frac{T}{4}\underset{n=1}{\overset{MR_p/\pi }{}}\frac{1}{\pi n}_{X_1}^{X_2}𝑑x\sqrt{\overline{g}}+\text{finite.}$$
(D9)
Eq. (D9) agrees with the result of LSW for the leading semiclassical logarithmic divergence.
## E Evaluation of $`S_2`$
We want to evaluate $`S_2`$,
$$S_2=\frac{T}{2}\left(\mathrm{Tr}\sqrt{\frac{^2}{x^2}\delta _{AB}\sqrt{\overline{g}}𝒦_{ab}^A𝒦^{Bab}}\mathrm{Tr}\sqrt{\frac{^2}{x^2}\delta _{AB}}\right)$$
(E1)
for the fluctuations about a straight string of length $`R`$ rotating with angular velocity $`\omega `$. To evaluate this, we must determine the value of the extrinsic curvature $`𝒦_{ab}^A`$. The definition of the extrinsic curvature is
$$𝒦_{ab}^A=n_\mu ^A_a_bx^\mu .$$
(E2)
The string $`x^\mu `$ is
$$x^\mu (x,t)=t\widehat{𝐞}_0^\mu +\frac{1}{\omega }\mathrm{sin}(\omega x)\left(\mathrm{cos}(\omega t)\widehat{𝐞}_1^\mu +\mathrm{sin}(\omega t)\widehat{𝐞}_2^\mu \right).$$
(E3)
The $`\widehat{𝐞}_i^\mu `$ are a basis of orthonormal unit vectors in Minkowski space. The $`n_A^\mu `$ are a basis for the vectors normal to $`x^\mu `$. We choose the basis
$`n_1^\mu `$ $`=`$ $`\widehat{𝐞}_3^\mu `$ (E4)
$`n_2^\mu `$ $`=`$ $`\mathrm{tan}(\omega x)\widehat{𝐞}_0^\mu +\mathrm{sec}(\omega x)\left(\mathrm{sin}(\omega t)\widehat{𝐞}_1^\mu +\mathrm{cos}(\omega t)\widehat{𝐞}_2^\mu \right).`$ (E5)
With this choice for the $`n_A^\mu `$, $`𝒦_{ab}^1`$ is zero, because the $`\widehat{𝐞}_3^\mu `$ component of $`x^\mu `$ is zero. The only nonzero component of $`𝒦_{ab}^A𝒦^{Bab}`$ is
$$\sqrt{\overline{g}}\left(𝒦_{ab}^2\right)^2=2\omega ^2\mathrm{sec}^2\omega x.$$
(E6)
Now that we know what $`𝒦_{ab}^A`$ is, we can evaluate (E1). Inserting (E6) into (E1) gives
$$S_2=\frac{T}{2}\left(\mathrm{Tr}\sqrt{\frac{^2}{x^2}+2\omega ^2\mathrm{sec}^2\omega x}\mathrm{Tr}\sqrt{\frac{^2}{x^2}}\right).$$
(E7)
The traces in (E7) are defined as sums over the eigenvalues of the given operators. Replacing the traces with explicit sums gives
$$S_2=\frac{T}{2}\underset{n=1}{\overset{\frac{\mathrm{\Lambda }R_p}{\pi }}{}}\left(\sqrt{\lambda _n}\frac{\pi n}{R_p}\right),$$
(E8)
where
$$R_p=\frac{\mathrm{arcsin}v}{v}R.$$
(E9)
The eigenvalues $`\lambda _n`$ are determined by the eigenfunction equation
$$\left(\frac{^2}{x^2}+2\omega ^2\mathrm{sec}^2\omega x\right)\psi _n(x)=\lambda _n\psi _n(x),$$
(E10)
with the boundary conditions $`\psi _n(\pm R_p/2)=0`$. The difference between the traces in Eq. (E7) is logarithmically dependent on the cutoff $`\mathrm{\Lambda }`$ (the mass of the dual gluon).
Eq. (E10) has the form of the Schröedinger equation, with the potential $`2\omega ^2\mathrm{sec}^2\omega x`$. This potential is an analytic continuation of the potential $`2\omega ^2\mathrm{sech}^2\omega x`$, whose eigenfunctions can be expressed terms of hypergeometric functions . Using this result, we find the eigenfunctions
$$\psi _n(x)=\{\begin{array}{cc}\sqrt{\lambda _n}\mathrm{cos}\left(\sqrt{\lambda _n}x\right)+\omega \mathrm{tan}\omega x\mathrm{sin}\left(\sqrt{\lambda _n}x\right)\hfill & \text{for n odd,}\hfill \\ \sqrt{\lambda _n}\mathrm{sin}\left(\sqrt{\lambda _n}x\right)\omega \mathrm{tan}\omega x\mathrm{cos}\left(\sqrt{\lambda _n}x\right)\hfill & \text{for n even.}\hfill \end{array}$$
(E11)
The eigenvalues $`\lambda _n`$ are
$$\lambda _n=\left(\frac{\left(\pi n+2\alpha _n\right)v}{R\mathrm{arcsin}v}\right)^2,$$
(E12)
where $`\alpha _n`$ satisfies the transcendental equation
$$\frac{\frac{\pi }{2}n+\alpha _n}{\mathrm{arcsin}v}=\frac{v}{\sqrt{1v^2}}\mathrm{cot}\alpha _n,$$
(E13)
and $`0<\alpha _n<\pi /2`$. There is no $`n=0`$ eigenvalue, despite the fact that $`\alpha _0=\mathrm{arcsin}v`$ satisfies (E13), because the corresponding eigenvalue $`\lambda _0=\omega ^2`$ makes $`\psi _n`$ zero everywhere.
We will carry out the sum (E8),
$$S_2=\frac{T}{2}\underset{n=1}{\overset{\frac{\mathrm{\Lambda }R_p}{\pi }}{}}\left(\sqrt{\lambda _n}\frac{\pi n}{R_p}\right),$$
(E14)
by converting it to a contour integral. We will find a function $`F_\lambda (z)`$ which has zeros whenever $`z=\pm \sqrt{\lambda _n}`$. We will find another function $`F_{R_p}(z)`$ which has zeros whenever $`z=\pm \pi n/R_p`$. We will then define a function $`F_{int}(z)`$,
$$F_{int}(z)=\frac{d\mathrm{ln}F_\lambda (z)}{dz}\frac{d\mathrm{ln}F_{R_p}(z)}{dz}.$$
(E15)
The function $`F_{int}(z)`$ has poles of residue $`1`$ when $`z=\pm \sqrt{\lambda _n}`$ and poles of residue $`1`$ when $`z=\pm \pi n/R_p`$. We then rewrite the sum (E14) as a contour integral,
$$S_2=\frac{T}{4\pi i}𝑑zzF_{int}(z).$$
(E16)
The contour of the integral (E16) lies along the imaginary axis, and on a semicircle at $`|z|=\mathrm{\Lambda }`$ with the real part of $`z`$ positive.
To write $`S_2`$ as the contour integral (E16), we need to find the functions $`F_\lambda (z)`$ and $`F_{R_p}(z)`$. The function $`F_{R_p}(z)`$ is
$$F_{R_p}(z)=\mathrm{sin}(R_pz),$$
(E17)
which is zero for $`z=\pm \pi n/R_p`$. We find $`F_\lambda (z)`$ by recalling that the eigenfunctions (E11) vanish at $`x=\pm R_p/2`$. Therefore,
$$F_{\mathrm{odd}}(z)=z\mathrm{cos}\left(\frac{R_p}{2}z\right)+\omega \mathrm{tan}\left(\frac{R_p}{2}\omega \right)\mathrm{sin}\left(\frac{R_p}{2}z\right),$$
(E18)
has zeros at $`z=\sqrt{\lambda _n}`$ for $`n`$ odd, and
$$F_{\mathrm{even}}(z)=z\mathrm{sin}\left(\frac{R_p}{2}z\right)\omega \mathrm{tan}\left(\frac{R_p}{2}\omega \right)\mathrm{cos}\left(\frac{R_p}{2}z\right),$$
(E19)
has zeros at $`z=\sqrt{\lambda _n}`$ for $`n`$ even. Thus,
$`F_\lambda (z)`$ $`=`$ $`{\displaystyle \frac{F_{\mathrm{odd}}(z)F_{\mathrm{even}}(z)}{z^2\omega ^2}}`$ (E20)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{z^2\mathrm{sin}(R_pz)2\omega z\mathrm{tan}\left(\frac{R_p}{2}\omega \right)\mathrm{cos}(R_pz)\omega ^2\mathrm{tan}^2\left(\frac{R_p}{2}\omega \right)\mathrm{sin}(R_pz)}{z^2\omega ^2}}.`$ (E21)
The factor $`(z^2\omega ^2)^1`$ removes nonphysical zeros which appear because $`F_{\mathrm{even}}(\pm \omega )=0`$. These zeros correspond to the $`n=0`$ “eigenfunction” which is zero everywhere for $`\lambda _0=\omega ^2`$. The function $`F_{int}(z)`$ is
$$F_{int}(z)=\frac{d}{dz}\mathrm{ln}\left[\frac{z^22\omega z\mathrm{tan}\left(\frac{R_p}{2}\omega \right)\mathrm{cot}(R_pz)\omega ^2\mathrm{tan}^2\left(\frac{R_p}{2}\omega \right)}{z^2\omega ^2}\right].$$
(E22)
Inserting (E22) in (E16) and integrating by parts gives
$$S_2=\frac{T}{4\pi i}𝑑z\mathrm{ln}\left[\frac{z^22\omega z\mathrm{tan}\left(\frac{R_p}{2}\omega \right)\mathrm{cot}\left(R_pz\right)\omega ^2\mathrm{tan}^2\left(\frac{R_p}{2}\omega \right)}{z^2\omega ^2}\right].$$
(E23)
Now, instead of having poles at $`z=\sqrt{\lambda _n}`$ and $`z=\pi n/R_p`$, the integrand has branch points at these points. The branch cuts run from $`\sqrt{\lambda _n}`$ to $`\pi n/R_p`$ for all $`n`$ along the real axis. There is one branch cut for each value of $`n`$. Since the contour either includes both $`\sqrt{\lambda _n}`$ and $`\pi n/R_p`$ or excludes both these points, none of these branch cuts cross the contour of integration. There are no branch points at $`z=\pm \omega `$, since both the numerator and denominator vanish there. None of the branch cuts crosses the contour, so the contour is still closed, and the integration by parts does not produce a boundary term.
The contour of the integral (E23) lies on the imaginary axis and a semicircle passing through positive real infinity. We rewrite (E23) as two integrals over the different pieces of the contour. For large values of the cutoff $`\mathrm{\Lambda }`$, the action $`S_2`$ is
$`S_2`$ $`=`$ $`{\displaystyle \frac{T}{4\pi }}{\displaystyle _\mathrm{\Lambda }^\mathrm{\Lambda }}𝑑y\mathrm{ln}\left[{\displaystyle \frac{y^2+2\omega y\mathrm{tan}\left(\frac{R_p}{2}\omega \right)\mathrm{coth}\left(R_py\right)+\omega ^2\mathrm{tan}^2\left(\frac{R_p}{2}\omega \right)}{y^2+\omega ^2}}\right]`$ (E25)
$`T{\displaystyle \frac{1}{4\pi }}{\displaystyle _{\pi /2}^{\pi /2}}𝑑\theta \mathrm{\Lambda }e^{i\theta }\left[2{\displaystyle \frac{\omega }{\mathrm{\Lambda }}}e^{i\theta }\mathrm{tan}\left({\displaystyle \frac{R_p}{2}}\omega \right)\mathrm{cot}\left(R_p\mathrm{\Lambda }e^{i\theta }\right)+O(\mathrm{\Lambda }^2)\right].`$
For large $`\mathrm{\Lambda }`$, $`\mathrm{cot}\left(R_p\mathrm{\Lambda }e^{i\theta }\right)`$ is proportional to the sign of $`\theta `$, so the $`\theta `$ integral vanishes. The $`y`$ integral is symmetric under $`yy`$. Changing variables to $`s=(y/\omega )\mathrm{cot}\left(\omega R_p/2\right)`$ gives
$$S_2=T\frac{\omega }{2\pi }\mathrm{tan}\left(\frac{R_p}{2}\omega \right)_0^{\frac{\mathrm{\Lambda }}{\omega }\mathrm{cot}\left(\omega R_p/2\right)}𝑑s\mathrm{ln}\left[\frac{s^2+2s\mathrm{coth}\left(R_p\omega \mathrm{tan}\left(\frac{R_p}{2}\omega \right)s\right)+1}{s^2+\mathrm{cot}^2\left(\frac{R_p}{2}\omega \right)}\right].$$
(E26)
The numerator of in (E26) is approximately $`(s+1)^2`$ for large values of $`s`$. We use this fact to extract the divergent part of (E26), getting
$$S_2=T\frac{2v}{\pi R}\left[v\gamma \mathrm{ln}\left(\frac{\mathrm{\Lambda }R}{2v^2\gamma }\right)+v\gamma \frac{\pi }{2}\right]T\frac{v^2\gamma }{\pi R}f\left(v\right).$$
(E27)
We have replaced $`R_p`$ with its definition
$$R_p=\frac{2}{\omega }\mathrm{arcsin}v.$$
(E28)
$`\gamma =(1v^2)^{1/2}`$ is the quark boost factor. The function $`f(v)`$ contains the rest of the integral,
$$f(v)_0^{\mathrm{}}𝑑s\mathrm{ln}\left[\frac{s^2+2s\mathrm{coth}\left(2sv\gamma \mathrm{arcsin}v\right)+1}{(s+1)^2}\right].$$
(E29)
For $`v1`$, the asymptotic value of $`f(v)`$ is
$$f(v)\frac{1}{6\gamma ^2}.$$
(E30)
The cutoff $`\mathrm{\Lambda }`$ used in (E27) is the cutoff in the $`x`$ coordinate. We must express $`\mathrm{\Lambda }`$ in terms of the cutoff $`M`$ for the $`r`$ coordinate, which measures physical distance. The cutoffs $`\mathrm{\Lambda }`$ and $`M`$ are related by the equation $`\mathrm{\Lambda }\delta x=M\delta r`$, or
$$\mathrm{\Lambda }=M\frac{dr}{dx}=M\gamma ^1.$$
(E31)
Inserting (E31) into (E27) gives
$$S_2=T\frac{2v}{\pi R}\left[v\gamma \mathrm{ln}\left(\frac{MR}{2(\gamma ^21)}\right)+v\gamma \frac{\pi }{2}\right]T\frac{v^2\gamma }{\pi R}f\left(v\right).$$
(E32)
Using the classical equation of motion (118) then gives
$$S_2=T\frac{2v}{\pi R}\left[v\gamma \mathrm{ln}\left(\frac{Mm}{\sigma }\right)+v\gamma \frac{\pi }{2}\right]T\frac{v^2\gamma }{\pi R}f\left(v\right).$$
(E33)
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# An 𝑓̇(𝑓)-frequency dynamics algorithm for gravitational waves
## I Introduction
This millennium will mark the beginning of gravitational wave astronomy with the broad band observatories LIGO and VIRGO. This opens up new opportunities to probe strongly gravitating processes in the coalescence of binaries and the birth of neutron stars or black holes. The late stages of spiral infall of a neutron star may show interesting new physics, for example, in association with tidal break-up around a companion black hole, where the onset will provide constraints on the density of neutron star matter. There may also be surprises. Advanced LIGO is expected to see events up to a Gpc, and might detect correlations or anti-correlations of gravitational radiation with cosmological $`\gamma `$-ray bursts. The early stages of observations will be focused on verification of potential gravitational wave signals and identification of source by type. In view of the diversity of source parameters and serendipity in any new window of observations, it is of interest to consider model-independent detection algorithms. These algorithms are expected to compare well with matched filtering, anticipating the inevitable model imperfections in both the wave-forms and the instrumentation noise.
Early evolution of a binary of two stars of mass $`M_1`$ and $`M_2`$ with orbital period $`P`$ is well-described by quasi-Newtonian evolution in the point-particle limit, whose luminosity $`_{GW}`$ in gravitational wave emissions in the theory of general relativity is given by the Peter and Mathews’ formula
$`_{GW}{\displaystyle \frac{32}{5}}(\omega )^{10/3}F(e)`$ (1)
in geometrized units with $`G=c=1`$, where $`\omega =2\pi /P`$ denotes the orbital angular velocity, $`F(e)`$ accounts for an orbital ellipticity $`e`$ and $`=(M_1M_2)^{3/5}/(M_1+M_2)^{1/5}`$ denotes the chirp mass. This is accompanied by a gradual rise in the orbital frequency $`f_{orb}=1/P`$:
$`\dot{f}_{orb}={\displaystyle \frac{96}{5}}(2\pi )^{8/3}^{5/3}f_{orb}^{11/3}.`$ (2)
The predicted luminosity of $`3.2\times 10^{33}`$ ergs/s agrees with the observed rate of change $`2.40\times 10^{12}`$ in the orbital period of about eight hours in the Hulse-Taylor binary system with neutron star masses $`M_1=1.42M_{}`$ and $`M_2=1.4M_{}`$. This strong frequency-luminosity relationship suggests to plot the evolution of a candidate source as a trajectory in the $`\dot{f}(f)`$-plane as the output of a detection algorithm. In doing so, we focus on burst sources whose gravitational waves show appreciable dynamics in their frequencies. This, in contrast to periodic sources, e.g.: rapidly spinning neutron stars, which we regard as a separate class of sources in regards to detection algorithms. Binary coalescence of compact objects, then, shows an initial branch $`\dot{f}f^{11/3}`$ indicative of the general relativistic relationship $`_{GW}f^{10/3}`$. These trajectories terminate on the $`f`$-axis at the quasi-normal mode frequency $`f_{QNR}`$ of a final black hole. For the expected low mass black holes $`M10M_{}`$ in compact binaries, $`f_{QNR}2800(10M_{}/M)`$ is beyond the (advanced) LIGO sensitivity range of about (20-1600Hz) 40-800Hz, however. The transition which connects the chirp to the final black hole state remains highly uncertain - here the gravitational radiation emission process is most nonlinear and potentially most interesting. For black hole-black hole coalescence, the transition from chirp to a common horizon envelope state should be smooth, may be very luminous and more frequent than neutron star-black hole coalescences. In neutron star-black hole coalescence, on the other hand, an intermediate black hole-torus state is expected if the black hole spins rapidly. This state is expected to be quiescent in its gravitational wave emissions, while luminous in emissions by the black hole in contact with the magnetic field of the torus. This predicts an anti-correlation between gravitational wave–emissions and $`\gamma `$-ray bursts from this type of catastrophic events. Alternatively, gravitational radiation is expected in accretion induced collapse of a white dwarf into a neutron star or black hole due to a bar mode instability, or the collapse of a young massive star, in which the relationship between $`\dot{f}`$ and $`f`$ is even more uncertain. If the collapse stalls with the formation of a neutron star, we might witness a negative branch $`\dot{f}<0`$ due to spin-down by gravitational wave emissions. It becomes potentially useful, therefore, to plot transient gravitational wave emissions as trajectories in the $`\dot{f}(f)`$-plane when classifying candidate detections by source-type. As gravitational wave-emissions are not expected to be significantly beamed, observations such as these by advanced LIGO can obtain definite evidence of $`\gamma `$–ray burst progenitors and their population statistics which are not otherwise available.
Here we describe an algorithm which enables accurate extraction of the frequency dynamics in the gravitational wave signal. The algorithm is based on counting zero-crossings, wherefrom an instantaneous frequency $`f(t)`$ and frequency rate of change $`\dot{f}(t)`$ can be estimated as a function of observer’s time $`t`$. Here the proposed algorithm is tested using simulated binary coalescence with Gaussian noise and noise recorded on the 40m LIGO test facility at Caltech. The algorithm is studied on signals of the type
$`X(t)=A(t)\mathrm{cos}(2\pi \mathrm{\Phi }(t))+\sigma \stackrel{~}{N}(t)`$ (3)
wth amplitude $`A(t)=(1t)^{1/4}`$ and phase
$`\mathrm{\Phi }(t)=c\left[\left(1(1t)^{5/8}\right)+p\left(1(1t)^{3/8}\right)\right],`$ (4)
where $`tϵ[0,1)`$ is the time normalized to the time of coalescence $`t=1`$, $`c=1000`$ and $`pϵ[0,1]`$ is a model parameter: $`p=0`$ in the Newtonian approximation and $`p=0.6038`$ in the 1PN approximation when $`M_1=M_2`$. The noise process $`\stackrel{~}{N}(t)`$ is normalized such that $`\sigma `$ controls the signal-to-noise ratio. All simulations are for data strings of a total of $`N=10,000`$ samples at $`t_n=n\mathrm{\Delta }t`$ ($`n=1,2,\mathrm{}N`$). The instantaneous frequency $`f(t)`$ is given by $`d\mathrm{\Phi }(t)/dt`$ and $`d^2\mathrm{\Phi }(t)/dt^2`$ denotes its rate of change $`\dot{f}(t)`$ with respect to time. The dimensionless rate of change of the period, $`f^2\dot{f}`$, expresses inversely the number of periods at a frequency $`f`$, which is of the order of a few hundred in the experiments reported below (final time $`t0.97`$).
## II Description of the algorithm
Consider the instantaneous frequency $`f(t)`$ as a finite time-series for $`tϵ[a,b]`$. We begin with a partition of $`[a,b]`$ into $`N`$ adjacent subwindows (bins) $`[s_{i1},s_i]`$ with centers at $`t_i=(s_{i1}+s_i)/2`$. Linear approximations $`f(t)\alpha _i+\beta _it`$ ($`s_{i1}<t<s_i`$) are then obtained by estimates of the coefficients $`\alpha _i`$ and $`\beta _i`$ from the noisy data $`X(t)`$ in each subwindow $`i`$. The number of periods in the $`i`$-th subwindow, for example, is $`_{s_{i1}}^{s_i}f(t)𝑑t=\alpha _i\delta +\beta _it_i\delta +O(\delta ^3)`$, where $`\delta =s_is_{i1}`$.
Let $`Z_S(s_{i1},s_i)`$ denote the number of zero-crossings of the true signal $`S(t)`$ in the subwindow $`[s_{i1},s_i]`$. Then $`Z_S(s_{i1},s_i)`$ differs from $`2_{s_{i1}}^{s_i}f(t)𝑑t`$ by at most one. In view of $`f(t)=d\mathrm{\Phi }(t)/dt`$, our algorithm hinges on the approximation
$`Z_S(s_{i1},s_i)2\delta \alpha _i+2\delta \beta _it_i.`$ (5)
The number of zero-crossings in the noisy data $`(X(u_0),X(u_1),\mathrm{},X(u_k))`$ in $`[s_{i1},s_i]`$ will be denoted by $`Z_X(s_{i1},s_i)`$, where the $`u_j`$ denote the discrete sampling times provided by the AD-converter (16 bit resolution, 10kHz). The $`X(u_j)`$ have expectation values $`S(u_j)`$ and variance $`\sigma ^2`$. The 10kHz-sampling frequency is assumed to be sufficiently high such that the true signal agrees with the zero-crossings of its linear interpolant based on the $`S(u_j)`$. In terms of the indicator function, $`I`$, defined by $`I(P)=1`$ if $`P`$ is true and $`I(P)=0`$ if $`P`$ is false, we have
$`Z_S(s_{i1},s_i)=\mathrm{\Sigma }_1^kI[S(u_{j1})S(u_j)<0],Z_X(s_{i1},s_i)=\mathrm{\Sigma }_1^kI[X(u_{j1})X(u_j)<0].`$ (6)
The probability $`p_j`$ that the observed signal $`X(u_j)`$ and the its expectation (the true signal) $`S(u_j)`$ have the same sign can be easily computed. Suppose that $`S(u_j)>0`$. Then, we have
$`P[X(u_j)>0]=P[(X(u_j)S(u_j))/\sigma >S(u_j)/\sigma ]=1F[S(u_j)/\sigma ]=F[S(u_j)/\sigma ],`$ (7)
where $`F`$ is the probability distribution of the noise process $`\stackrel{~}{N}(t)`$, which is assumed to be symmetric about zero. Similarly when $`S(u_j)<0`$, we find $`P[X(u_j)<0]=F[S(u_j)/\sigma ]`$. It follows that $`p_j=F[|S(u_j)|/\sigma ]`$, whereby
$`P[I[X(u_{j1})X(u_j)<0]I[S(u_{j1})S(u_j)<0)]]=p_{j1}(1p_j)+p_j(1p_{j1}).`$ (8)
Consider, then, $`R_i=\text{min}_j|S(u_j)|/\sigma `$, as a measure for the signal-to-noise ratio in the subwindow $`[s_{i1},s_i]`$. As $`R_i\mathrm{}`$, $`\text{min}_jp_j1`$. We then have for the expectation $`E`$ and the variance $`V`$ the limits
$`\underset{R_i\mathrm{}}{lim}E[Z_X(s_{i1},s_i)]=Z_S(s_{i1},s_i),V[Z_X(s_{i1},s_i)]=0.`$ (9)
Hence, for large signal to noise ratios, we recover $`Z_X(s_{i1},s_i)2\delta \alpha _i+2\delta \beta _it_i.`$
Estimates $`\alpha _i`$ for $`f(t_i)`$ and $`\beta _i`$ for $`\dot{f}(t_i)`$ can be calculated by linear regression. Define (as is usual in such regressions, see, e.g. ) the $`L^2`$ error
$`Q(\alpha _i,\beta _i)=\mathrm{\Sigma }_1^N\left(Z_X(s_{i1},s_i)(2\delta \alpha _i+2\delta \beta _it_i)\right)^2.`$ (10)
Let $`\widehat{\alpha }_i,\widehat{\beta }_i`$ minimize $`Q(\alpha ,\beta )`$. The regression estimates of the instantaneous frequency $`f(t_i)`$ and its rate of change $`\dot{f}(t_i)`$ are $`\widehat{\alpha }_i`$ and $`\widehat{\beta }_i`$ respectively. As these are obtained simultaneously, the $`\dot{f}`$ are on equal par with $`f`$, rather than being their derivatives. Standard regression theory allows us to compute error boxes in these estimates. Note that the method amounts to an approximation of the time-series $`f(t)`$ by a piecewise linear, and generally discontinuous function.
The method may be optimized by modifying the present $`L^2`$ norm in the error function to curve fitting, e.g.: different weights on $`f`$ and $`\dot{f}`$. Initial experiments indicate that the method can be implemented with adaptive window sizes, which may improve its efficiency when the dynamic range in the frequency $`f(t)`$ is large.
Because only the pattern of zero-crossings is measured, the algorithm targets precisely the underlying gravitational wave frequency-luminosity relationship of the source, as examplified in binary coalescence. The algorithm suppresses the amplitude of the signal, and hence it is robust to size fluctuations in the signal. Likewise, reference to the underlying cosine functions is weak. In fact, any $`2\pi `$-function with two roots in each period can replace the cosine function in (1).
## III Noise dependence
The frequency dynamics algorithm has been tested in simulated binary coalescence using the Newtonian model of spiral in ($`p=0`$) with Gaussian noise and the instrumental noise from the 40m-LIGO test facility at Caltech. The chirp extends over $`t=[0,0.97)`$, wherein the number of periods at a given frequency $`f(t)`$ decreases from above 400 to below 200. Figure 1 shows the results for both cases.
The algorithm obtains higher accuracy with the 40m-LIGO noise than with Gaussian noise when the variance is high. We believe that the reason for this lies in the high-degree of dependence in the noise from the LIGO prototype. The first-order autocorrelations from 40m-LIGO noise were of the order of 0.8. There were, on average, 15 zero-crossings per 100 recordings of the pure noise process. This is much smaller than the 50 or so zero crossings obtained for white noise of the same length. Thus the LIGO noise contributes far fewer “false crossing” due to noise as opposed to signal. This will be the case for all noise with low-frequency components in the spectrum. In such cases, our algorithms performs even better than it would with white noise. To make the last point more precise, consider a Gaussian process with a first order autocorrelation of $`\rho `$. If we observe a sequence of length $`N`$ from such a process the expected number of zero-crossings is
$`{\displaystyle \frac{(N1)}{2}}\left(1{\displaystyle \frac{2}{\pi }}\mathrm{tan}^1{\displaystyle \frac{\rho }{\sqrt{1\rho ^2}}}\right)`$
For $`\rho `$ of 0.8, as in the LIGO case, we would expect approximately 20 zero crossings per 100 recordings. Of course, the LIGO noise from the Caltech prototype is not pure Gaussian, but the principle is the same.
## IV Comparison with matched filtering
Optimal application of matched filtering requires accurate models of the system to be detected and the noise of the detector. Both aspects of prior knowledge are critical and must be described in a low number of parameters, lest matched filtering becomes unstable. Detection of gravitational waves from forementioned candidate sources poses a number of uncertainties which may enter as model imperfections (e.g., size of neutron stars). These imperfections can give rise to woefully wrong results when the true signal is not contained in the model for any choice of its parameters. It becomes of interest, therefore, to compare the model-independence aspect of the frequency dynamics algorithm with the sensitivity of matched filtering to these imperfections. The comparison would be complete with an additional on noise type and intensity which, however, falls outside the scope of this work. The comparison is made using the post-Newtonian model in the 1PN approximation as the true signal ($`p=0.6038`$), while applying matched filtering assuming the Newtonian model given by $`p=0`$. This mismatch of model to true signal is used here to simulate other model imperfections, such as due to the finite size of the objects.
Thus, the matched filtering model is perfect in case of $`p=0`$, and becomes imperfect as $`p0`$. It is expected that matched filtering will outperform the frequency dynamics algorithm for $`p=0`$, and vice-versa for certain larger values of $`p>0`$. The comparison is made quantitative using a mean absolute relative error (MARE), expressed as a percentage, in the estimated frequencies by comparison with the known signal input:
$`\text{MARE}={\displaystyle \frac{100}{N}}{\displaystyle \underset{1}{\overset{N}{}}}{\displaystyle \frac{|\widehat{f}(t_i)f(t_i)|}{f(t_i)}}`$ (11)
Using a small Monte Carlo study, the comparison is made at three levels of additive Gaussian white noise. Figure 2 shows the results of this comparison study, where each reported MARE is the average of 200 simulations.
## V Conclusions
The trajectory in the $`\dot{f}(f)`$-plane of a burst of gravitational radiation is expected to provide a signature of the source. This is anticipated from the strong - possibly characteristic - coupling between the frequency and the luminosity of its gravitational wave emissions. Binary coalescence, for example, shows the characteristic chirp $`\dot{f}f^{\alpha +1/3}`$ associated with a luminosity $`_{GW}f^\alpha `$, where $`\alpha =10/3`$ in general relativity in the point-particle limit. Other burst sources of gravitational radiation may have different frequency dynamics (e.g.: different $`\alpha `$’s). We have described a nonparametric (model-free) frequency dynamics algorithm to obtain an accurate plot of the trajectory in the $`\dot{f}(f)`$plane. It is based on zero-crossings in the gravitational wave signal, to extract these trajectories in the presence of noise. The method is robust against non-Gaussian additive noise and amplitude variations. The method shows good results at appreciable chirp rates and different types of noise. Even when the signal-to-noise level is low (down to about 2 in our experiments), we can still recover the distinctive curve of the frequency dynamics when $`\dot{f}`$ is high. The estimation appears to be much more accurate for LIGO noise than for Gaussian noise when the variance is high. We believe that the reason for this lies in the high-degree of dependence in the noise from the LIGO prototype, which suppresses “false zero-crossings” generated by the noise. This will be the case for all noise with low-frequency components in the spectrum. In such cases, our algorithm performs even better than it would with white noise. The frequency dynamics algorithm has been compared with matched filtering, and shown to be superior in the presence of model imperfections. This is particularly relevant in searches for unanticipated burst sources.
The present experiments indicate that the algorithm should perform well in extracting the frequency dynamics in the gravitational waves. Future experiments are desirable at very low signal-to-noise levels, to study the detection limit for transient sources with the present algorithm.
Acknowledgements. This research received partial support from NASA grant NAG5-7012, LIGO and an MIT Reed Award. The authors are grateful to R. Weiss for constructive comments, and A. Lazzarini for kindly providing the LIGO data.
Figure Captions
Figure 1. The $`\dot{f}(f)`$-diagram of simulated binary coalescence in the Newtonian approximation, extracted by the frequency dynamics algorithm in the presence of noise at a signal-to-noise ratio of about $`4`$ ($`s=\sigma =0.25`$). The horizontal axis shows the frequency $`f`$ \[Hz\] and the vertical axis its time rate of change $`df`$ \[Hz/s\]. The top window shows the results with Gaussian noise, and the lower window shows the results with the instrumental noise from the 40m-LIGO test facility at Caltech. The crosses in the right windows denote horizontal and vertical errors obtained from the linear regression estimates. The results indicate that the method improves when the noise is correlated. At higher frequencies, the results are similar due to a larger instantaneous signal-to-noise ratio.
Figure 2. Experiments on model-independence in terms of the normalized error MARE on the frequency, obtained by the frequency dynamics algorithms (lines) and matched filtering (dashed) at three different variances $`\sigma ^2=0.01,0.15,0.25`$. The comparisons are made for three different signals, parametrized by the value of $`p`$ (horizontal axis). Matched filtering is performed with model assumption $`p=0`$, hence $`p0`$ in the signal simulates a model imperfection. The results show that the frequency dynamics algorithm outperforms matched filtering in the presence of such appreciable model imperfections, with MARE approaching a constant.
References
Figure 1.
Figure 2.
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# The evolution of disk galaxies
## 0.1 Introduction
In the last decade, our understanding of galaxy formation and evolution grew dramatically thanks to the coming out of a theoretical framework for structure formation within the cosmological context (the inflation-inspired cold dark matter scenario, CDM), and to the great observational progress, in particular the imagery and spectroscopy of high redshift galaxies. Nevertheless, the fundamental questions of galaxy formation and evolution still await answers. The symbiosis of theory and observations is crucial in this field in order to connect what we “see” at different redshifts with the evolution of a given population of galaxies.
Most of the galaxies observed in the local universe are disk galaxies. An important question is how the structural, dynamical and luminosity properties of this population of galaxies evolved and how much they contributed in the past to global quantities such as the star formation (SF) rate and luminosity per unit of volume. Were the galaxy disks smaller and their surface brightness (SB) higher than at present? Which is the SF history of disk galaxies? Were these galaxies brighter and bluer in the past? Did the luminosity-velocity (Tully-Fisher) relation change in the past? A powerful theoretical tool for studying these questions related to the evolution of local and global properties of disk galaxies is the combination of inductive (backward) galaxy evolutionary models with initial and boundary conditions calculated from the hierarchical formation scenario (Avila-Reese, Firmani, & Hernández 1998; Avila-Reese & Firmani 2000; Firmani & Avila-Reese 2000). In this paper, we present some of the evolution predictions of these (seminumerical) models and we discuss them in light of the available observational data.
## 0.2 Disk galaxy evolutionary models in the hierarchical scenario
Disk galaxies are dynamically fragile objects. This is why, at a first approximation, we obviate the disk major mergers in their evolution and considering that disks grow inside-out gently with a gas accretion rate proportional to the hierarchical mass aggregation rate. A major advantage of our seminumerical models is that we follow locally the overall evolution of individual disks in centrifugal equilibrium, including star formation (SF) and feedback in the disk ISM, and luminosity evolution. At each epoch, the growing disk is characterized by a local infall rate of fresh gas, a gas and stellar surface density profile, a local SF rate, a color profile, and a rotation curve (including the dark matter component). The disks form with the baryon mass fraction within the growing dark halos. The angular momentum of each collapsing mass shell is calculated using the spin parameter $`\lambda `$ obtained in numerical and analytical works. The local SF is induced by a gravitational instability (Toomre) criterion and it is self-regulated by an energy balance in the ISM along the vertical disk direction. The efficiency of SF in this model is almost independent of the mass (luminosity).
The main physical factors which influence most of the local and global properties of our model galaxies and their correlations are the mass, mass aggregation history (MAH), and spin parameter $`\lambda `$. These factors and their statistical distribution are related to the cosmological initial conditions. For a quick review of the main disk galaxy properties and correlations predicted at $`z=0`$ for a typical CDM model, see Avila-Reese et al. 2000; the results in detail are presented in the papers mentioned above.
## 0.3 Evolutionary predictions
### 0.3.1 Star formation rate and luminosity evolution
The SF rate in our hierarchical models is driven by both the gas accretion rate determined by the MAH and the disk surface density determined by $`\lambda `$. In Fig. 1(a) and 1(b) we show the average hierarchical mass aggregation rate of halos of different masses and the corresponding SF rates histories for $`\lambda =0.05`$, respectively; less massive galaxies assemble faster than the massive ones. The correlation between MAH and SF history is evident. In Fig. 1(b) we also show the SF histories for $`\mathrm{5\; 10}^{11}M_{}`$ models with two extreme MAHs —early active and extended MAHs (upper and lower dashed lines)— and $`\lambda =0.05`$. The influence of $`\lambda `$ on the SF history is clearly seen in Fig. 1(c) where we plotted $`\mathrm{5\; 10}^{11}M_{}`$ models with $`\lambda =0.03`$ (upper short- dashed line) and $`\lambda =0.1`$ (lower short-dashed line) and the average MAH in both cases. The upper and lower dot-dashed lines correspond to extreme models with an early active MAH and $`\lambda =0.03`$ and with an extended MAH and $`\lambda =0.1`$, respectively.
A strong prediction of our disk galaxy models, due mainly to the hierarchical MAHs, is the shape of the SF rate history with a gentle maximum at $`z1.52.5`$ for most of the cases, and a relatively quick fall towards the present by factors $`24`$ on the average. From an observational study for large disk galaxies ($`r_B4h_{50}^1`$), Lilly et al. (1998) have found the SF rate to decrease likely a factor $`2.53.5`$ between $`z0.7`$ and $`z0`$ (h=0.5). For our corresponding models, since $`z0.7`$ the SF rate decreased a factor $`2`$ on average. However, for models with small $`\lambda `$, i.e., very high SB galaxies, this factor is $`3`$. It is possible that galaxy samples at different redshifts suffer from a significant SB selection effect since disk galaxies, as our and other model results show, change the SB distribution with $`z`$ (Bouwens & Silk 2000). Therefore, the factor Lilly et al. (1998) find could be biased towards HSB galaxies in agreement with our galaxy models.
Deep field studies show that the global (cosmic) SF history per unit of volume from $`z0`$ to 0.7 increased by a factor $`6`$ and by more than a factor of 10 up to the maximum which is attained at $`z1.52.0`$ (e.g., Madau, Pozzetti, & Dickinson 1998 and the references therein). If our models actually describe the evolution of normal disk galaxies which dominate the global SF rate today, then the increment of a factor $`10`$ detected at $`1.52`$ could not have been produced only by disk galaxies. Other galaxy populations had to contribute to the global SF rate in the past (e.g., Babul & Fergurson 1996).
Since the $`B`$band luminosity $`L_B`$ is a tracer of young and intermediate stellar populations, its evolution should be similar to the SF history. We find indeed that $`L_B`$ increases towards the past and more quickly for the less massive galaxies. At $`z=1`$, $`L_B`$ is $`1.5`$ times larger than at $`z=0`$ for a $`L_{}`$ galaxy. The integral colors of the galaxy models become bluer towards the past; on average from $`z=0`$ to $`z=1`$, $`(BV)`$ decreases by $`0.250.3`$ mag. Less massive galaxies undergo more color evolution; this is because they attain the peak in SF rate before than massive galaxies.
### 0.3.2 Size and surface brightness evolution
Disk size and SB evolution are natural in the hierarchical formation scenario. In the right panel of Fig. 1 we show the evolution of the $`B`$band scale radius $`r_B`$ for average models of 3 different masses, using the flat $`\mathrm{\Lambda }`$CDM ($`\mathrm{\Omega }_\mathrm{\Lambda }=`$h=0.7) cosmology. The size evolution is slightly more pronounced for more massive galaxies. In this figure we also plot a rough observational estimate by Lilly et al. (1998). According to theses authors, for large disk galaxies $`r_B`$ has grown not more than $`25\%`$ since $`z1`$.
Size evolution implies SB evolution. Lilly et al. (1998; dashed line), after correcting for some selection effects, found that the average SB of their large disk galaxies sample has grown $`0.5`$ mag at $`z0.7`$ with respect to $`z0`$. Similar results were obtained by e.g., Forbes et al. (1996), Vogt et al. (1997), Roche et al. (1998). The last authors conclude that the deep field observations of disk galaxies can be better explained by luminosity and size evolution models. For our hierarchical models, we find slightly more pronounced SB evolution than the observations. In a more recent observational work, Simard et al. (1999) have found that galaxies from $`z=0.1`$ to $`z=0.9`$ seem to increase their average SB by $`1.3`$ mag. However, these authors concluded that if the selection effect due to comparing low luminosity galaxies in nearby redshift bins to high luminosity galaxies in distant bins is allowed for, then no discernible evolution remains in the SB of bright disk galaxies. Using the same data of Simard et al., Bouwens & Silk (2000) have derived an increase in the SB of $`1.5`$ mag. This difference is because the last authors have introduced corrections for SB selection bias that they find to be important due to the strong evolution in the SB distribution of disk galaxies. Certainly, the disk size and SB evolution are important tests for the hierarchical formation scenario.
### 0.3.3 Dynamics: evolution of the $`H`$ and $`B`$band Tully-Fisher relations
The linking of the resolved photometric and spectroscopic data to dynamical data is crucial for the understanding of the evolution and physics of a given population of galaxies. In our hierarchical models, the mass of the galaxy grows faster than the maximum rotation velocity $`V_m`$; this is because the dense inner parts of the dark halo are not significantly affected by the “secondary” mass infall. The dark halos, and, therefore, the disks, obey a tight relation between mass and maximum circular velocity (see e.g., Avila-Reese et al. 1998,1999). For our $`\mathrm{\Lambda }`$CDM cosmology, we find that the disk stellar mass-maximum rotation velocity relation, $`M_d=AV_m^n`$, at $`z=0`$ has a slope $`n3.4`$. This relation is proportional to the $`H`$band Tully-Fisher relation (TFR). The slope $`n`$ remains almost constant with $`z`$. The coefficient $`A`$ decreases with $`z`$, i.e., as we said above, while the mass significantly decreases towards the past, $`V_m`$ decreases only a little. We find that at $`z=1`$, $`A`$ is $`1.6`$ times smaller than at $`z=0`$ or, that is the same, the zero-point of the $`H`$band TFR is $`0.5`$ mag higher.
In the $`B`$band, the TFR evolves in an opposite way: although the disk mass ($`L_H`$ luminosity) decreases with $`z`$, $`L_B`$ increases as we have seen in $`\mathrm{\S }3.1`$. This is because SF is more active in the past. The slope of the $`B`$band TFR slightly decreases with $`z`$ (less massive galaxies evolve more quickly than massive ones; in particular, $`L_B`$ increases more with $`z`$ for the former than for the latter). Assuming the slope constant with $`z`$, for our $`\mathrm{\Lambda }`$CDM model, we obtain that the zero-point of the $`B`$band TFR at $`z=1`$ is $`0.7`$ mag lower than at $`z=0`$ (taking into account HSB and LSB galaxies), i.e. $`A`$ is $`2`$ times larger.
From the observational point of view, several efforts were made in order to acquire internal kinematic (Vogt et al. 1997 and more references therein) and galaxy-galaxy lensing (Hudson et al. 1998) data for disk galaxies at high redshifts. From these works it is still difficult to draw conclusions regarding the evolution of the $`B`$band TFR. Since the interpretation of the results is sensitive to the cosmological model assumed, just in order to evaluate the data from several of these works, we used the critical model ($`q=0.5`$) with h=0.5. The results of Vogt et al. (1997) show that the zero-point of the $`B`$band TFR, from $`z=0`$ to the average redshift of the sample $`z=0.54`$, has changed by $`\mathrm{\Delta }M_B0.33`$ mag; for the Hudson et al. (1998), Simard & Pritchet (1998) and Rix et al. (1997) samples, $`z=0.6,`$ 0.35, and 0.25, and $`\mathrm{\Delta }M_B1.0,1.5`$ and $`1.5`$ mag, respectively. As one sees, these works seem to be at odds with one another. The samples used in the last two papers are dominated by small, actively star-forming galaxies, while in the two first papers, normal spiral galaxies dominate. Thus, we may compare our model results with those of Vogt et al. and Hudson et al. papers. For the SCDM cosmology (with $`\sigma _8=0.67`$ and h=0.5), $`\mathrm{\Delta }M_B0.4`$ mag at $`z=0.6`$ (in fact, the theoretical evolution of the TFR is similar for both the SCDM and $`\mathrm{\Lambda }`$CDM models). This strongly disagrees with observations in a $`q=`$h=0.5 cosmology. For low density or flat with cosmological constant models with h$`>0.5`$, the observational data are in better agreement with the models. In particular, in the flat $`\mathrm{\Lambda }`$CDM cosmology ($`\mathrm{\Omega }_\mathrm{\Lambda }=`$h=0.7), the data of Vogt et al. give only a slightly larger value for $`\mathrm{\Delta }M_B`$ than our models (i.e., the observed zero-point at $`z=0.54`$ is slightly less bright than our hierarchical models predict).
In conclusion, we find that the hierarchical formation scenario (for $`\mathrm{\Lambda }`$CDM-like models) offers robust initial and boundary conditions for disk galaxy evolution. However, the evolution of disk sizes and the zero-point of the $`B`$band TFR seems to be too exagerated with respect to some observational data. We stress that analytical approaches should take into account that the zero-point of the “structural” and $`B`$band TFRs evolve in a different way.
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# Delocalized Qubits as a Computational Basis in the System of Interacting Spins
## 1 Introduction
Nuclear magnetic resonance (NMR) provides an excellent proving ground for testing different quantum information processing (QIP) ideas. Due to combination of its good developed theory and sophisticated experimental technique the realization of the simplest quantum algorithms using standard NMR spectrometers turned out to be possible. At present time the achievements of liquid state (LS) NMR QIP are far beyond the capabilities of any other experimental methods.
This success of LS NMR QIP is due also to the fact that a clear correspondence was found between the abstract notion of a quantum mechanical bit - qubit and real physical object - nuclear spin. Namely, two qubit states $`|0`$ and $`|1`$ are mapped onto two possible spin 1/2 projections. As a computational basis the eigenfunctions $`|m`$ of the Z-components of a nuclear spin are used, which commutates with the full Hamiltonian. For example, for two spin system this computational basis is written as follows:
$$\begin{array}{cc}|00=|m_1=1/2,m_2=1/2,& |01=|m_1=1/2,m_2=+1/2,\\ |10=|m_1=+1/2,m_2=1/2,& |11=|m_1=+1/2,m_2=+1/2.\end{array}$$
(1)
where $`|m_1`$ and $`|m_2`$ are the nuclear spin Z-components eigenfunctions.
However, already now one can see the LS NMR QIP limitations and the next step is going to be the solid state (SS) NMR QIP -. The SS NMR essential feature is the existence of the spin-spin (exchange and dipole-dipole) interaction terms in the total Hamiltonian, which are not averaged to zero by thermal motions, and which contain not only Z, but also X and Y spin components. Due to this fact a single spin orientation becomes bad integral of motion. This means that such system stationary states do not correspond to definite values of the spins Z-component and the functions (1) become time dependent. To use these functions as a stable computational basis it is necessary to remove spin-spin interactions using special means . For example a WAHUHA multipulse sequence, consisting of four pulses, can be used for this purpose .
In this paper it is shown how in the system of interacting spins with spin Z-components not commuting with total Hamiltonian, one can introduce a stable computational basis, which does not require the continuous application of a multipulse sequences. The trick is that instead of the correspondence “one qubit” \- “a pair of single spin states” one makes the correspondence “one qubit” - “a pair of collective spin states”. Using the virtual spin formalism language it can be written as “one qubit” - “one virtual spin”, instead of “one qubit” - “one real spin”. Such approach is being developed in the framework of our program “Ouantum computer: many levels instead of many particles” . Unlike the previous cases, in which the virtual spins have been defined on the levels of the individual multilevel particle - on the spins 3/2 or 7/2 levels -, in the present case they are defined on the collective levels of the interacting spins system. It means that a qubit is “delocalized” and there is no direct correspondence “one qubit” - “one particle”.
## 2 A simple model of two spins interaction
A universal quantum gate set consists of single qubit rotations and two-qubit CNOT gate . In order to implement such a set it will be enough to use a simple two spin interaction model. Let it will be the system of two non equivalent nuclear spins $`I=1/2`$ and $`S=1/2`$, connected by isotropic exchange interaction.
This system Hamiltonian is
$$\begin{array}{c}=\mathrm{}\omega _0(I_z+S_z)+\mathrm{}\delta /2(I_zS_z)+\mathrm{}J(𝐈𝐒),\hfill \\ \omega _0=(1/2)(\gamma _I+\gamma _S)H_0,\delta =(\gamma _I\gamma _S)H_0\hfill \end{array}$$
(2)
where $`\gamma _I`$ and $`\gamma _S`$ \- the nuclei giromagnetic ratios, $`J`$ \- the exchange integral, $`H_0`$ \- the external constant magnetic field.
This Hamiltonian eigenfunctions and corresponding eigenvalues are
$$\begin{array}{cc}|\psi _1=|++|m_I=+1/2,m_S=+1/2,\hfill & E_1\mathrm{}\epsilon _1=\mathrm{}\omega _0+(1/4)\mathrm{}J\hfill \\ |\psi _2=p|++q|+,\hfill & E_2\mathrm{}\epsilon _2=(1/4)\mathrm{}J+(1/2)\mathrm{}\theta \hfill \\ |\psi _3=p|+q|+,\hfill & E_3\mathrm{}\epsilon _3=(1/4)\mathrm{}J(1/2)\mathrm{}\theta \hfill \\ |\psi _4=|,\hfill & E_4\mathrm{}\epsilon _4=\mathrm{}\omega _0+(1/4)\mathrm{}J\hfill \end{array}$$
(3)
where $`p=cos(\varphi /2),q=sin(\varphi /2),\theta ^2=J^2+\delta ^2`$, $`J/\delta =tg(\varphi ),\pi /2\varphi \pi /2.`$
This system has four allowed transitions with the following frequencies and relative intensities:
$$\begin{array}{cc}\epsilon _{12}=\omega _0+(1/2)J(1/2)\theta ,\hfill & P_{12}|\psi _1\left|I_x+S_x\right|\psi _2|^2=1+sin(\varphi ),\hfill \\ \epsilon _{13}=\omega _0+(1/2)J+(1/2)\theta ,\hfill & P_{13}|\psi _1\left|I_x+S_x\right|\psi _3|^2=1sin(\varphi ),\hfill \\ \epsilon _{24}=\omega _0(1/2)J+(1/2)\theta ,\hfill & P_{24}|\psi _2\left|I_x+S_x\right|\psi _4|^2=1+sin(\varphi ),\hfill \\ \epsilon _{34}=\omega _0(1/2)J(1/2)\theta ,\hfill & P_{34}|\psi _3\left|I_x+S_x\right|\psi _4|^2=1sin(\varphi ),\hfill \end{array}$$
(4)
Fig. 1 depicts the energy levels of this system. This simple system allows exact solution and manifests the specific features and advantages of the proposed information coding – the so-called “virtual spin formalism”. In a case of the more general interaction (anisotropic exchange or dipole-dipole one) all necessary expressions can be obtained using perturbation theory.
## 3 Computational basis in the virtual spin representation
The four dimensional Hilbert space, spanned on eigenfunctions (3), can be considered as a direct product of the pair of virtual spins $`R=1/2`$ and $`S=1/2`$ two-dimensional Hilbert spaces . It means that the eigenfunctions (3) of Hamiltonian (2) can be taken to form the computational basis for the two qubit system, thus:
$$\begin{array}{cc}|00=|\psi _1,\hfill & |01=|\psi _2,\hfill \\ |10=|\psi _3,\hfill & |11=|\psi _4.\hfill \end{array}$$
(5)
This notation means that the qubit corresponds to the virtual spin, so:
$$\begin{array}{cc}|0_Q=|m_Q=1/2,\hfill & |1_Q=|m_Q=+1/2,\hfill \\ |0_R=|m_R=1/2,\hfill & |1_R=|m_R=+1/2.\hfill \end{array}$$
(6)
and two qubit states correspond to the computational basis:
$$\begin{array}{cc}|00=|m_Q=1/2,m_R=1/2,\hfill & |01=|m_Q=1/2,m_R=+1/2,\hfill \\ |10=|m_Q=+1/2,m_R=1/2,\hfill & |11=|m_Q=+1/2,m_R=+1/2.\hfill \end{array}$$
(7)
Now a resonance transition between any two real states of physical system admits the interpretation as virtual spin reorientations. For example, the irradiation of the transition $`\psi _1||\psi _2`$ can be interpreted in the virtual spin representation as the Q spin rotation etc. Using the previous results - it can be shown, that in this two qubit system the universal gates set can be realized by means of the following resonance electromagnetic field pulses.
| Logic operation: | Realization by transition(s): |
| --- | --- |
| Virtual spin Q rotation | $`\psi _1||\psi _2`$ and $`\psi _3||\psi _4`$ |
| Virtual spin R rotation | $`\psi _1||\psi _3`$ and $`\psi _2||\psi _4`$ |
| Controlled $`Q`$ spin inversion $`CNOT_{RQ}`$ | $`\pi `$ \- pulse applied to $`\psi _3||\psi _4`$ |
| Controlled $`R`$ spin inversion $`CNOT_{QR}`$ | $`\pi `$ \- pulse applied to $`\psi _2||\psi _4`$ |
where $`CNOT_{RQ}`$ means that qubit $`Q`$ undergoes the $`NOT`$ operation, which is controlled by the state of qubit $`R`$, and $`\pi `$ \- pulse means the virtual spin $`Q`$ rotation through the angle $`\pi `$. The corresponding resonance pulses are depicted on Fig. 2. It can be seen that all quantum gates can be implemented using just one pulse: for virtual spin rotation a double frequency pulse is necessary, whereas for $`CNOT`$ operation - a single frequency pulse.
## 4 Conclusions
It was shown that the information coding onto virtual spins allows to implement a universal gate set in a solid state two interacting spin system. The advantages of this coding are connected with the fact, that spin-spin interactions can be large and that there is no need to use continuous irradiation with a multipulse sequence to have a stable computational basis. Large spin-spin interaction produces big resonance frequencies differences, which facilitate the selective resonance excitation of the individual transitions, desired for gates implementation. In addition, the gate operation time is under full control of an experimentalist and can be done short, whereas the coding using real spins requires time, defined by exchange interaction value (which is a given molecule property) and can be rather long.
It should be noted that the suggested approach can be useful also for LS NMR QIP in a case, when an exchange interaction is not averaged to spin Z-components. It can be used also for information coding onto any cluster of interacting particles of arbitrary nature. The only requirements are the existence of the proper selection rules for resonance transitions among cluster stationary states, on which a two virtual qubit system is defined. An example of the virtual qubit formalism applied to optical states of a single atom is given in the paper .
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# Spectrum of neutrino masses and their nature in the light of present and future experiments.
## Abstract
The present experimental data on neutrino oscillations, neutrinoless double beta decay and tritium beta decay are collected together and possible mass ranges for Dirac and Majorana neutrinos are found. Four future experimental situations are investigated: both decay experiments give only upper bounds, one of them gives a positive result ($`|m_\nu |0`$ or $`m_\beta 0`$), or finally both effective neutrino masses are different from zero ($`|m_\nu |0`$ and $`m_\beta 0`$). Each scenario gives new information on neutrino masses and nature but only the last has a chance to resolve the problem and give some additional information on $`CP`$ violation in the lepton sector.
The problem of neutrino mass spectrum and its nature is the most important issue in the lepton part of the Standard Model. What new information can we obtain from the last experimental results, and what are the future perspectives? Three kinds of experiments play a fundamental role in answering this question. Two are traditional and known for years: beta decay and neutrinoless double beta decay ($`(\beta \beta )_{0\nu }`$) of nuclei. Already Fermi in 1934 and Furry in 1939 realized that both processes are important to find the neutrino mass and nature. The third type constitute the neutrino oscillation experiments . We strongly believe that neutrino oscillations are responsible for anomalies observed in solar , atmospheric and LSND experiments. There are trials of alternative explanations of the observations but they require much more sophisticated assumptions (as for example the breaking of the equivalence principle, breaking of the special theory of relativity, the neutrino decay with life-time much below expectations or huge neutrino magnetic moments) and give much poorer fits to the data .
The end of the Curie plot in tritium beta decay has been observed since the late forties giving now a bound on the effective electron neutrino mass of $`m_\beta 2.8eV`$ and $`m_\beta 2.5eV`$ both with $`95\%`$ c.l.. Although less pronounced, the problem of negative $`m_\beta ^2`$ remains.
Trials of finding the neutrinoless double beta decay of even-even nuclei have also been conducted for years. The best result for the effective electron neutrino mass $`|m_\nu |`$ bound comes from the last experiment with $`{}_{}{}^{76}Ge`$. The Heidelberg-Moscow collaboration found that the half life time of $`{}_{}{}^{76}Ge`$ is bounded as
$$T_{1/2}^{0\nu }(Ge)>5.7\times 10^{25}years,$$
(1)
which gives a bound of $`|m_\nu |<0.2eV`$ . The derivation of this bound required the calculation of complicated nuclear matrix elements. A discrepancy of the order of a factor of 3 between the independent studies has been found. As a consequence the uncertainty of $`|m_\nu |`$ is of the order of $`\sqrt{3}`$.
Finally, the oscillation experiments give results on the basis of three different observations: solar , atmospheric and LSND . The results of the last of these, if correct and explained through neutrino oscillations requires the introduction of a fourth sterile light neutrino species. The observation of the LSND collaboration has not been confirmed by KARMEN and is partially excluded by Bugey and BNL776 experiments. In such circumstances we will assume that only three massive neutrinos exist which explain the solar and atmospheric neutrino anomalies. The case of four massive neutrinos can be studied in a similar way. Even with this assumption the oscillatory data has some ambiguities related in particular to the existence of several possible solutions of the solar neutrino problem. The values of $`\delta m^2`$ and $`\mathrm{sin}^22\theta `$ for the atmospheric and the four solar solutions (SMA MSW, LMA MSW, LOW MSW and VO) of the observed anomalies are collected in Table I. As there are definitely two scales of $`\delta m^2`$, $`\delta m_{atm}^2\delta m_{sol}^2`$, two possible neutrino mass spectra must be considered (Fig. 1). The first, known as normal mass hierarchy ($`A_3`$) where $`\delta m_{sol}^2=\delta m_{21}^2\delta m_{32}^2\delta m_{atm}^2`$ and the second, inverse mass hierarchy spectrum ($`A_3^{inv}`$) with $`\delta m_{sol}^2=\delta m_{21}^2\delta m_{atm}^2\delta m_{31}^2`$. Both schemes are not distinguishable by present experiments. There is hope that next long base line experiments (e.g. MINOS, ICANOE) will do that.
The values of the allowed range of $`\delta m^2`$ and $`sin^22\theta `$ in the Table I are presented with a $`95\%c.l.`$. As we will see the value of $`\mathrm{sin}^22\theta `$ plays a decisive role in our considerations. A value of $`\mathrm{sin}^22\theta <1`$ is crucial. If we look at the data , we see that a maximal mixing is still possible at $`99\%c.l.`$ for three solutions of the solar neutrino problem LMA, LOW and VO . However, the best fit values, contrary to the atmospheric neutrino case, are smaller than 1. Still, it is possible that this is only a fluctuation and future better data are needed to solve this problem. Here we assume that the tendency observed in experimental data is real and future experiments will confirm that $`\mathrm{sin}^22\theta <1`$.
Two elements of the first row of the mixing matrix $`|U_{e1}|`$ and $`|U_{e2}|`$ can be expressed by the third element $`|U_{e3}|`$ and the $`\mathrm{sin}^22\theta `$
$$|U_{e1}|^2=(1|U_{e3}|^2)\frac{1}{2}(1+\sqrt{1\mathrm{sin}^22\theta }),$$
(2)
and
$$|U_{e2}|^2=(1|U_{e3}|^2)\frac{1}{2}(1\sqrt{1\mathrm{sin}^22\theta }).$$
(3)
The value of the third element $`|U_{e3}|`$ is not fixed yet and only different bounds exist for it . We will take the bound directly inferred from the CHOOZ experiment
$$|U_{e3}|^2<0.04.$$
(4)
In spite of various complications and uncertainties even now we will get some information about the neutrino mass spectrum and their nature. We hope however, that such considerations with better data will give a key to the solution of the problem in the future.
For the three massive neutrino scenario (without LSND) the oscillation experiments give the lower bound on the highest mass of neutrinos
$$(m_\nu )_{max}>\sqrt{\delta m_{atm}^2}0.06eV$$
(5)
and a bound on the absolute value of the difference of two masses
$$|m_im_j|<\sqrt{\delta m_{atm}^2}.$$
(6)
No upper bound on neutrino masses can be inferred from oscillation experiments alone. So, in the schemes in Fig 1, we only know that $`m_{3,2}0.06eV`$, for $`A_3`$ and $`A_3^{inv}`$ schemes respectively.
The tritium beta decay measures the effective electron neutrino mass $`m_\beta `$
$$m_\beta =\left[\underset{i=1}{\overset{3}{}}|U_{ei}|^2m_i^2\right]^{1/2},$$
(7)
and no upper bound on neutrino masses can be given. We can only find that
$$(m_\nu )_{min}m_\beta (m_\nu )_{max}$$
(8)
where $`(m_\nu )_{min}`$ denotes the lowest neutrino mass. In practice this means that $`(m_\nu )_{min}<2.5eV`$, but $`(m_\nu )_{max}`$ can be very large. Connecting both the tritium beta decay data and the oscillatory data we can find the following upper bound on the neutrino mass
$$\sqrt{\delta m_{atm}^2}(m_\nu )_{max}\sqrt{(m_\beta )^2+\delta m_{atm}^2}.$$
(9)
This means that the mass of the heaviest neutrino must be somewhere in between $`0.06eV`$ and $`2.5eV`$.
The estimations which we get up to now do not depend on the neutrino nature. It is well known that the electron energy distribution in nuclei beta decay and flavor oscillations do not distinguish Dirac from Majorana neutrinos . Therefore, the bounds (Eqs. 5, 6 and Eq. 9) are valid for both neutrino types. This is not the case of the neutrinoless double beta decay of nuclei. This decay is only possible for a massive Majorana neutrino . Neglecting all other mechanisms which can participate in the process we can derive a bound on
$$|m_\nu |\left|\underset{i=1}{\overset{3}{}}U_{ei}^2m_i\right|.$$
(10)
Alone, this bound is of little value since it only means that
$$0|m_\nu |(m_\nu )_{max}.$$
(11)
Let us now consider however all three experiments together . If neutrinos are Dirac particles then $`|m_\nu |=0`$ and we do not have any additional information from $`(\beta \beta )_{0\nu }`$. All we can say about masses of Dirac neutrinos follows from $`{}_{1}{}^{3}H`$ decay and oscillation experiments and is given by Eqs. 5, 6 and Eq. 9.
If neutrinos are Majorana particles then the bound on $`|m_\nu |`$ works and we can find new restrictions. For three neutrinos in the scheme $`A_3`$ (Fig. 1) we have
$`|m_\nu |`$ $`=`$ $`||U_{e1}|^2(m_\nu )_{min}+|U_{e2}|^2e^{2i\varphi _2}\sqrt{(m_\nu )_{min}^2+\delta m_{sol}^2}`$ (12)
$`+`$ $`|U_{e3}|^2e^{2i\varphi _3}\sqrt{(m_\nu )_{min}^2+\delta m_{sol}^2+\delta m_{atm}^2}|,`$ (13)
and similarly in the case of the scheme $`A_3^{inv}`$. Three parameters defined above are unknown, $`(m_\nu )_{min}`$ and the Majorana $`CP`$ violating phases $`\varphi _2`$ and $`\varphi _3`$. We are not able to predict the value of $`|m_\nu |`$ but lower $`|m_\nu |_{min}`$ and upper $`|m_\nu |_{max}`$ limits as function of $`(m_\nu )_{min}`$ can be obtained. In Fig. 2 (scheme $`A_3`$) and Fig. 3 (scheme $`A_3^{inv}`$) $`|m_\nu |_{min}`$ is shown as function of $`(m_\nu )_{min}`$ for the four possible solutions of the solar neutrino anomaly . The shaded regions give the uncertainties of the results caused by the allowed region of input parameters (mostly $`\mathrm{sin}^22\theta `$). For the SMA MSW solution the $`95\%c.l.`$ range of $`|m_\nu |_{min}`$ is described by one curve in the adopted logarithmic scale.
Present Experimental Data
We see that already with the present experimental data the possible values of $`|m_\nu |_{min}`$ depicted in Fig. 2 and Fig. 3 exceed the bound on $`|m_\nu |`$
$$|m_\nu |_{min}>0.2eV,$$
(14)
for some values of $`(m_\nu )_{min}`$. This means that Majorana neutrinos with masses above some $`(m_\nu )_{min}`$ are forbidden. The limiting mass depends on the solution to the solar neutrino problem and on the adopted value of $`\mathrm{sin}^22\theta `$ and $`|U_{e3}|^2`$. For the SMA MSW solution in both $`A_3`$ and $`A_3^{inv}`$ schemes $`(m_\nu )_{min}=0.22eV`$. For LMA, LOW and VO with $`(\mathrm{sin}^22\theta )_{max}=0.98`$ and $`|U_{e3}|^2=0.01`$, $`(m_\nu )_{min}=1.5eV`$ .
We see that the bound on the effective mass $`|m_\nu |`$ given by the present $`(\beta \beta )_{0\nu }`$ experiments restricts the range of possible Majorana neutrino masses. This range depends on the maximal value of $`\mathrm{sin}^22\theta `$, and for
$$\mathrm{sin}^22\theta >\frac{12|U_{e3}|^2}{(1|U_{e3}|^2)^2}$$
(15)
is the same as for Dirac neutrinos .
Future Perspectives
Future improving results of neutrino experiments will give better information on neutrino mass spectra. The improvements that are important for our purpose and are realistic are the following
* Concerning neutrino oscillations
+ the problem of the sterile neutrino should be solved
+ a single solution of the solar neutrino problem should be found and the value of $`\mathrm{sin}^22\theta `$ should be given with better precision
+ the allowed range of $`\delta m_{atm}^2`$, $`\delta m_{sol}^2`$ and $`|U_{e3}|`$ will be reduced.
* Concerning $`(\beta \beta )_{0\nu }`$ decay
+ there are plans of going down with the effective Majorana mass down to $`|m_\nu |0.006eV`$ in two stages, first $`|m_\nu |0.02eV`$ (GENIUS I) and later $`|m_\nu |0.006eV`$ (GENIUS II).
+ two possibilities should be envisaged with different impact on this study. The pessimistic option is that the decay of $`{}_{}{}^{76}Ge`$ will not be observed and only a new bound on $`|m_\nu |`$ will be found. The optimistic, the decay is discovered and a value $`|m_\nu |(0.20.006)eV`$ is inferred.
* Concerning the $`{}_{1}{}^{3}H`$ decay
+ two Collaborations , , plan to go down with the value of $`m_\beta `$ below $`1eV`$ and perhaps even down to $`0.6eV`$.
+ Once more two scenarios can happen. A distortion in the electron kinetic energy will not be found and a new bound will follow, $`m_\beta <0.6eV`$. On the optimistic side, such a distortion will be observed, the problem with $`m_\beta ^2<0`$ will be solved and the tritium $`\beta `$ decay will give a value of $`m_\beta =\kappa ^{}(0.62.5)eV`$.
We will now discuss four different possibilities (i) two bounds $`|m_\nu |<0.006eV`$ and $`m_\beta <0.6eV`$ exist, (ii) $`|m_\nu |<0.006eV`$ and $`m_\beta \kappa ^{}`$, (iii) $`|m_\nu |\kappa `$ but $`m_\beta <0.6eV`$, and finally (iv) $`|m_\nu |\kappa `$ and $`m_\beta \kappa ^{}`$.
Ad. (i) $`|m_\nu |<0.006eV`$, $`m_\beta <0.6eV`$
Nothing special happens, the accepted range of Dirac and Majorana neutrino masses will become smaller. The bound on $`m_\beta `$ gives the possible range of Dirac neutrino masses
$$0.06eV<(m_\nu )_{max}0.6eV.$$
(16)
The bounds on $`|m_\nu |`$ (depicted in Fig. 2 and Fig. 3) give a small space of Majorana neutrino masses. Once more the maximal value of $`(m_\nu )_{min}`$ depends on $`(sin^22\theta )_{max}`$ and $`|U_{e3}|^2`$, which at that time should be known much better. Taking as an example for $`|m_\nu |<0.006eV`$ with $`|U_{e3}|^2=0.01`$, only two values of $`(sin^22\theta )_{max}`$
$`(sin^22\theta )_{max}0.01(SMA)(m_\nu )_{min}<0.02eV0.06eV(m_\nu )_{max}0.062eV,`$ (17)
$`(sin^22\theta )_{max}0.98(LMA)(m_\nu )_{min}<0.05eV0.06eV(m_\nu )_{max}0.078eV.`$ (18)
We see that a much smaller range of Majorana neutrino masses will be accepted than in the Dirac case. The scheme $`A_3^{inv}`$ will be excluded for SMA (LMA) by GENIUS I (GENIUS II) (see Fig. 3). This is a very pessimistic scenario because the problem of neutrino nature will not be solved and the chance to find the spectrum of Majorana neutrinos in a close future will be very small.
Ad. (ii)$`|m_\nu |<0.006eV`$ and $`m_\beta \kappa ^{}`$
If a value of $`m_\beta \kappa ^{}`$ is found the situation changes considerably. With the oscillation parameters and $`m_\beta \kappa ^{}`$ we can calculate the spectrum of neutrinos. The only accepted scheme is $`A_3`$. It gives
$`m_1`$ $`=`$ $`(m_\nu )_{min}=[\kappa ^2(1|U_{e1}|^2)\delta m_{solar}^2|U_{e3}|^2\delta m_{atm}^2]^{1/2},`$ (19)
and
$`m_2`$ $`=`$ $`[(m_\nu )_{min}^2+\delta m_{solar}^2]^{1/2},`$ (20)
$`m_3`$ $`=`$ $`[(m_\nu )_{min}^2+\delta m_{solar}^2+\delta m_{atm}^2]^{1/2}.`$ (21)
In Fig. 4 we show the $`(m_\nu )_{min}`$ as function of $`m_\beta \kappa ^{}`$ where the oscillation parameters change within their allowed range given in Tab. I. For $`(m_\nu )_{min}0.1eV`$, practically $`(m_\nu )_{min}\kappa ^{}`$. The difference is visible only for very small $`(m_\nu )_{min}`$. This means that the experimental error on $`\kappa ^{}`$ is the only significant source of the variation range of $`(m_\nu )_{min}`$. The mass spectrum obtained from Eq. 1920 and 21 does not depend on the neutrino nature. If neutrinos are Majorana particles, the value of $`(m_\nu )_{min}`$ obtained from $`H_1^3`$ decay can be used to find the range of possible values of $`|m_\nu |`$. In Fig. 5 we depicted $`|m_\nu |_{max}`$ and $`|m_\nu |_{min}`$ for two possible solutions of the solar neutrino problem, SMA (Fig. 5A) and LMA (Fig. 5B) in double logarithmic scales. The bound of $`(m_\nu )_{min}(0.60.8)eV`$ gives the range of possible $`|m_\nu |`$ values which follows from crossing the space allowed by oscillation experiments.
$$|m_\nu |_\beta ^{min}|m_\nu ||m_\nu |_\beta ^{max},$$
(22)
If the experimental bound on $`|m_\nu |`$ from $`(\beta \beta )_{0\nu }`$ decay is below the $`|m_\nu |_\beta ^{min}`$ range, then neutrinos can not be Majorana particles. If it is lager, the problem of neutrino nature is not solved.
Ad.(iii) $`|m_\nu |\kappa `$, $`m_\beta <0.6eV`$.
It is quite probable that this scenario will happen. In this case the neutrinos are Majorana particles. In Fig. 6A a possible band of $`|m_\nu |`$, obtained from the GENIUS experiment is given. This bound crosses the region of space allowed by oscillation data giving the possible value of $`(m_\nu )_{min}`$
$$(m_\nu )_{min}^{min(\beta \beta )_{0\nu }}(m_\nu )_{min}(m_\nu )_{min}^{max(\beta \beta )_{0\nu }},$$
(23)
If the value of $`|m_\nu |`$ is found in the second stage of the GENIUS experiment $`|m_\nu |0.010.005eV`$ the $`(m_\nu )_{min}^{min(\beta \beta )_{0\nu }}=0`$. With the range of possible $`(m_\nu )_{min}`$ the result on $`m_\beta `$ from tritium $`\beta `$ decay can predicted. In practice $`m_\beta `$ should also satisfy the inequalities given by Eq.( 23). If the experimental limit on $`m_\beta `$ is larger than $`(m_\nu )_{min}^{max(\beta \beta )_{0\nu }}`$, the theory with three Majorana neutrinos is consistent. If, what is less probable, the experimental limit on $`m_\beta `$ is smaller than $`(m_\nu )_{min}^{max(\beta \beta )_{0\nu }}`$ then the theory with three Majorana neutrinos is ruled out.
Ad.(iv) $`|m_\nu |\kappa `$, $`m_\beta \kappa ^{}`$.
This is the best, but on the other side, the least probable scenario. The value of $`|m_\nu |0`$ defines the neutrino as a Majorana particle. The value $`m_\beta 0`$ gives their mass spectrum. Comparing both bands with the region of $`|m_\nu |`$ allowed by the oscillation data (Fig. 6B) is a check of internal consistency of the theory. With precise data the crossing of the three regions can be used to specify the values of the CP breaking Majorana phases $`\phi _1`$ and $`\phi _2`$ (Eq.(12)). If the two bands $`|m_\nu |`$ and $`m_\beta `$ cross the oscillation region near $`|m_\nu |_{max}`$, two phases are equal $`\phi _1=\phi _2n\pi `$. This means that all three Majorana neutrinos have the same CP parity $`\eta _{CP}=+i`$ and the symmetry is conserved. If the two bands cross the oscillation region near $`|m_\nu |_{min}`$ once more the CP symmetry is satisfied with $`\eta _{CP}(\nu _1)=\eta _{CP}(\nu _2)=\eta _{CP}(\nu _3)=i`$. Finally, if all three regions cross somewhere in between, the phases $`\phi _1`$ are nontrivial and the CP symmetry is broken.
In conclusion
Present experimental data define the region of neutrino mass. The heaviest Dirac neutrino mass must be somewhere in the range $`0.06eV(m_\nu )_{max}2.5eV`$. The analogous range for Majorana neutrinos depends on the solution of the solar neutrino problem and is $`0.06eV(m_\nu )_{max}0.26eV`$ for SMA and $`0.06eV(m_\nu )_{max}1.5eV`$ for the LMA scenario. A better future bound on the effective mass in the neutrinoless double $`\beta `$ decay and the tritium $`\beta `$ decay would give better limits on Dirac and Majorana neutrino masses but the problem of the mass spectrum and their nature will still not be solved. The situation will change significantly if at least one of the experiments gives a positive result and $`|m_\nu |`$ or $`m_\beta `$ will be different from zero. The most difficult scenario, where both experiments give positive results is the one where (i) the nature of neutrinos, (ii) their mass spectrum, (iii) the consistency of the theory and (iv) the CP breaking Majorana phases can be found.
Acknowledgments This work was supported by the Polish Committee for Scientific Research under Grant No. 2P03B05418 and 2P03B04919, by CICYT and by Junta de Andalucia.
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# Nonzero angular momentum states of the helium atom in a strong magnetic field
## 1 Introduction
Motivated by the astrophysical discovery of strong magnetic fields on the surfaces of white dwarfs ($`10^5`$ Tesla) and neutron stars ($`10^8`$ Tesla), the behaviour and the properties of matter in strong magnetic fields has increasingly attracted interest. The theoretical description of atoms in strong magnetic fields is well covered in the literature only for the case of the hydrogen atom (see refs. and refs. therein). Until recently our knowledge about atoms with more than one electron in a strong magnetic field has been relatively sparse and definitely not sufficient for a comparison of the corresponding theoretical data with the mysterious absorption edges in the spectrum of the magnetic white dwarf GD229. Those have for the first time been identified as helium lines by Jordan et al in 1998. This overwhelming evidence was based on results of highly accurate ab initio calculations performed by Becken and Schmelcher in 1998, a part of which has been published in ref.. For a more detailed overview over the various theoretical approaches to the helium atom in a strong magnetic field and the corresponding literature before 1998 we refer the reader to and in particular the references therein. In a fully correlated configuration interaction approach has been applied to the helium atom in a strong magnetic field. Total energies for spin singlet and triplet states for both positive and negative $`z`$-parity in the subspace of vanishing magnetic quantum number $`M=0`$ have been provided, thereby covering the regime of magnetic fields strengths from $`B=0`$ to $`B=100a.u.`$ ($`B=1a.u.`$ corresponds to $`2.3510^5`$Tesla). Additionally all the transition energies within the $`M=0`$ subspace have been presented and discussed there, including the stationary components with respect to the field dependence which have been the key tool for the comparison with the observed spectra .
The aim of the present paper is to provide important data for states with the magnetic quantum number $`M=1`$. Some data for states with positive $`z`$-parity and $`M=1`$ have already successfully been used for a comparison with the astronomical observation(in , we used the corresponding stationary transitions in a certain field and wavelength regime). The complete data are presented here for the first time. Data resulting from states with negative $`z`$-parity are very recent results. The latter further extend our knowledge on the helium atom in a strong magnetic field and permit to investigate additional transitions.
For the case of triplet spin symmetry, Jones et al very recently used a released-phase quantum Monte Carlo method for calculating accurate data. However, they cover only three field strengths, investigate less excited states for common symmetries and do not study the spin singlet states at all. Nevertheless for the common values of the field strength they confirm our data to several digits. In contrast to this the energies presented by Scrinzi, though in principle variational, are in the presence of a magnetic field significantly lower than ours. Carefully comparing the energies, it appears to us that they are systematically too low, i.e. most probably due to a numerical error.
We remark that performing calculations for finite magnetic quantum numbers requires within our approach drastically improved computational techniques (in comparison to the case $`M=0`$) for keeping the CPU time affordable. A summary of these techniques is provided in the appendix of the present paper.
The starting point of the present paper is the nonrelativistic Hamiltonian of the helium atom with infinite nuclear mass in a magnetic field as given in section 2. To be self-contained we briefly discuss the Hamiltonian’s symmetries and provide a description of the basis set as well as the full configuration interaction approach (for more details see ref.). We introduce a maximal set of conserved quantities, chosen to be the total spin $`S^2`$, the $`z`$-component $`S_z`$ of the total spin, the total spatial magnetic quantum number $`M`$ and the total spatial $`z`$-parity $`\mathrm{\Pi }_z`$. These symmetries serve for classifying the results for the energies for $`M=1`$ in sec.3. In each of the two subspaces for positive and negative $`z`$-parity we present the total energies and the ionization energies of the ground state and the first four excited states for singlet and triplet spin symmetry. Additionally we consider in sec.3 all the transitions within the $`M=1`$ subspace as well as all the possible transitions between the $`M=1`$ states treated in the present paper and the $`M=0`$ states given in ref.. The wavelengths of all the stationary components are provided, being the basic ingredient for the successful comparison of theoretical data with the spectrum of magnetic white dwarfs in general and specifically for GD229.
## 2 Hamiltonian, Symmetries and basis sets
### 2.1 Hamiltonian and Symmetries
Assuming the magnetic field to point in the $`+z`$-direction, the Hamiltonian reads
$`H`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}\left(\begin{array}{c}\\ \frac{1}{2}\end{array}𝐩_i^2+\begin{array}{c}\\ \frac{1}{2}\end{array}Bl_{z}^{}{}_{i}{}^{}+\begin{array}{c}\\ \frac{B^2}{8}\end{array}(x_i^2+y_i^2)\begin{array}{c}\\ \frac{2}{|𝐫_i|}\end{array}+Bs_{z}^{}{}_{i}{}^{}\right)+{\displaystyle \frac{1}{|𝐫_2𝐫_1|}}`$ (9)
The one-particle operators in eq.(9) are the Coulomb potential energies $`\begin{array}{c}\\ \frac{2}{|𝐫_i|}\end{array}`$ of the electrons in the field of the nucleus as well as their kinetic energies, here splitted into the parts $`\begin{array}{c}\\ \frac{1}{2}\end{array}𝐩_i^2`$, the Zeeman terms $`\begin{array}{c}\\ \frac{1}{2}\end{array}Bl_{z}^{}{}_{i}{}^{}`$, the diamagnetic terms $`\begin{array}{c}\\ \frac{B^2}{8}\end{array}(x_i^2+y_i^2)`$ and their spin energies $`Bs_{z}^{}{}_{i}{}^{}`$. The electron-electron repulsion energy is represented by the two-particle operator $`\frac{1}{|𝐫_2𝐫_1|}`$. We remark that we use an electron spin $`g`$-factor equal 2, and any more accurate value for it can be simply incorporated by shifting the final total energies correspondingly. For remarks on the influence of relativistic effects and on a scaling relation taking into account the finite nuclear mass we refer the reader to ref..
Analogously to ref., we exploit that there exist four independent commuting conserved quantities: the total spin $`𝐒^2`$, the $`z`$-component $`S_z`$ of the total spin, the $`z`$\- component $`L_z`$ of the total angular momentum and the total spatial $`z`$-parity $`\mathrm{\Pi }_z`$.
### 2.2 Basis sets
For constructing a two-particle basis set of eigenfunctions of the above mentioned conserved quantities, our central ingredient is an anisotropic Gaussian basis set of one-particle functions
$`\mathrm{\Phi }_i(\rho ,\phi ,z)`$ $`=`$ $`\rho ^{n_{\rho }^{}{}_{i}{}^{}}z^{n_{z}^{}{}_{i}{}^{}}e^{\alpha _i\rho ^2\beta _iz^2}e^{im_i\phi }i=1,\mathrm{},n,`$ (10)
which are themselves eigenfunctions of the corresponding one-particle operators of the mentioned conserved quantities. The parameters $`n_{\rho }^{}{}_{i}{}^{}`$ and $`n_{z}^{}{}_{i}{}^{}`$ are restricted by
$`n_{\rho }^{}{}_{i}{}^{}`$ $`=`$ $`|m_i|+2k_i;k_i=0,1,2,\mathrm{}withm_i=\mathrm{}2,1,0,1,2,\mathrm{}`$ (11)
$`n_{z}^{}{}_{i}{}^{}`$ $`=`$ $`\pi _{z}^{}{}_{i}{}^{}+2l_i;l_i=0,1,2,\mathrm{}with\pi _{z}^{}{}_{i}{}^{}=0,1`$ (12)
whereas the nonlinear variational parameters $`\alpha _i`$ and $`\beta _i`$ are positive and have to be nonlinearly optimized for each field strength as described in ref.. For each one-particle subspace of given symmetry we used an algorithm for determining the nonlinear parameters $`\alpha _i`$ and $`\beta _i`$ such that the states of the hydrogen atom or the He<sup>+</sup> ion for that symmetry were optimally described. We emphasize that this procedure gives rise to considerable effort since it has to be repeated for each field strength separately.
We construct a basis set of spatial two-particle states by
$`|\psi _q`$ $`:=`$ $`b_i^{}b_j^{}|0i=1,\mathrm{},n,j=i,\mathrm{},n,`$ (13)
where $`b_i^{}`$ is the creation operator of the $`i`$-th one-particle state $`|i=b_i^{}|0`$ whose position representation is given by eq.(10). The spin space is spanned by spin singlet or spin triplet states, and therefore the operators $`b_i^{}`$ have to be chosen bosonic or fermionic, respectively. Selecting combinations with
$`m_i+m_j`$ $`=`$ $`M,mod(\pi _{z}^{}{}_{i}{}^{}+\pi _{z}^{}{}_{j}{}^{},2)=\mathrm{\Pi }_z,`$ (14)
we achieve the two-particle states (13) to be a basis set within the subspace for given total symmetries $`M`$ and $`\mathrm{\Pi }_z`$. The number $`N`$ of two-particle basis states is thus in general smaller than $`n(n+1)/2`$.
We perform a full Configuration Interaction (full CI) approach by representing the Hamiltonian in a basis whose spatial part is given by the in general nonorthonormal states (13). Since the spin part $`Bs_{z}^{}{}_{i}{}^{}`$ of the Hamiltonian can trivially be taken into account by a shift of the energies it is sufficient to represent the spatial part of the Hamiltonian $`H`$ and the overlap $`S`$ by
$`S_{pq}`$ $`=`$ $`\psi _p|\psi _q,H_{pq}=\psi _p|H|\psi _q`$ (15)
The matrices $`S`$ and $`H`$ are Hermitian, and the overlap $`S`$ is additionally positive definite. Furthermore the matrix elements turn out to be real. The finite-dimensional generalized real-symmetric eigenvalue problem
$`(\underset{¯}{\underset{¯}{H}}E\underset{¯}{\underset{¯}{S}})\underset{¯}{c}`$ $`=`$ $`0`$ (16)
provides eigenvalues $`E`$ which are variational upper bounds to the exact eigenvalues of the Hamiltionian (9) within each subspace of given $`M`$ and $`\mathrm{\Pi }_z`$.
### 2.3 Matrix elements
For calculating the matrix elements of the spatial part of the Hamiltonian (9), we rewrite the former in second quantization, $`\widehat{H}=\widehat{H}_I+\widehat{H}_{II}`$, where $`\widehat{H}_I`$ and $`\widehat{H}_{II}`$ denote the second-quantized counterparts of the familiar one- and two-particle operators whose position representations read
$`H_I(𝐩,𝐫)`$ $`=`$ $`\begin{array}{c}\\ \frac{1}{2}\end{array}𝐩^2+\begin{array}{c}\\ \frac{1}{2}\end{array}𝐁𝐥+\begin{array}{c}\\ \frac{1}{8}\end{array}B^2(x^2+y^2)\begin{array}{c}\\ \frac{2}{|𝐫|}\end{array}H_{II}(𝐫_1,𝐫_2)=\begin{array}{c}\\ \frac{1}{|𝐫_2𝐫_1|}\end{array}`$ (27)
Now, with $`|\psi _q:=b_i^{}b_j^{}|0`$ and $`|\psi _p:=b_k^{}b_l^{}|0`$ a straightforward calculation leads to
$`\psi _p|\psi _q`$ $`=`$ $`i|kj|l\pm i|lj|k`$ (28)
$`\psi _p|\widehat{H}_I|\psi _q`$ $`=`$ $`i|H_I|kj|l\pm i|H_I|lj|k`$ (29)
$`+j|H_I|li|k\pm j|H_I|ki|l`$
$`\psi _p|\widehat{H}_{II}|\psi _q`$ $`=`$ $`ij|H_{II}|kl\pm ij|H_{II}|lk,`$ (30)
where $`|ij:=|i|j`$ and where the sign ’$`\pm `$’ stands for ’$`+`$’ in the singlet case and for ’$``$’ in the triplet case.
For the relatively simple evaluation of the $`n(n+1)/2`$ different one-particle overlaps $`i|k`$ and matrix elements $`i|H_I|k`$ we refer the reader to appendices A,B in ref.. The two-particle matrix elements $`ij|H_{II}|kl`$ are by no means trivial, in particular in view of the fact that their accurate and fast evaluation is necessary in order to build up the Hamiltonian matrix in an affordable amount of CPU time. In ref. we discussed a method using a decomposition in Cartesian coordinates which expresses the two-particle matrix elements in series of hypergeometric functions whose evaluation has been performed by highly efficient analytical continuation formulas. The latter are necessary in order to keep the CPU time acceptable since the number of different two-particle matrix elements is of the order $`N(N+1)/2`$ rather than $`n(n+1)/2`$. However, the Cartesian decomposition becomes more and more inefficient with increasing magnetic quantum number, which is already relevant for calculations of the subspace $`M=1`$. Therefore we have developped a drastically improved procedure using cylindrical coordinates which leads to an enormous gain of speed such that the computation of the whole Hamiltonian matrix becomes even faster than its diagonalization by standard library routines. The derivation of the corresponding powerful formula for the electron-electron integral is rather lengthy and complicated: we therefore present only major steps of it in appendices A,B,C of the present paper.
## 3 Results and discussion
Throughout the paper we use the notation $`\nu _{S_z}^{2S+1}M^{(1)^{\mathrm{\Pi }_z}}`$ for a state with spin multiplicity $`(2S+1)`$ and degree of excitation $`\nu =1,2,3,\mathrm{}`$ within the subspace of given magnetic quantum number $`M`$ and $`z`$-parity $`\mathrm{\Pi }_z`$. The index $`S_z`$ will be omitted in obvious cases. The present paper is concerned with the subspaces $`{}_{}{}^{1}(1)_{}^{+}`$, $`{}_{}{}^{3}(1)_{}^{+}`$, $`{}_{}{}^{1}(1)_{}^{}`$, $`{}_{}{}^{3}(1)_{}^{}`$. The correspondence between our field notation and the common spectroscopic notation $`n_{S_z}^{2S+1}L_M`$ in field-free space is discussed in ref. (see table 1 therein). For completeness we mention that the $`5^1(\pm 1)^{}`$ states correspond to the $`6^1D_{\pm 1}`$ states whose field free energy is $`2.138982274a.u.`$, and analogously the $`5^3(\pm 1)^{}`$ states correspond to the $`6^3D_{\pm 1}`$ states with a field free energy of $`2.013901415a.u.`$ (Both values are taken from ref.).
### 3.1 Aspects for the selection of basis functions
For the $`M=0`$ states treated in ref., we have been able to achieve a considerable accuracy by choosing basis sets which can describe the shape of the exact wave function, i.e. include electronic correlation effects. The latter become less important with increasing quantum number $`|M|`$, and this manifests itself already for $`M=1`$. The reason is that bound two-particle states with nonzero values of $`M`$ are approximately one-particle excitations. Consequently, the electrons are spatially more separated than in a $`0^+`$ state, and this lowers the correlation energy. Additionally, the cusp problem is also less important for excited states $`M=1`$, and therefore fewer one-particle functions with large values for the nonlinear $`\alpha `$ and $`\beta `$ parameters are needed. We have exploited these facts and achieved even more accurate results for the $`M=1`$ states than for the $`M=0`$ states.
In detail the strategy was similar to the $`M=0`$ case. In the case of the $`(1)^+`$ subspace we used 260 optimized one-particle basis functions (for each field strength) for constructing a two-particle basis set of dimension $`N=3793`$. The latter number is identical for the singlet and the triplet subspace since it cannot occur that a two-particle state contains two identical one-particle contributions which combine to the odd total magnetic quantum number $`M=1`$. This is a principal difference to the case of the $`0^+`$ subspace. In order to describe angular correlation, we have added also $`(2)^+`$ and $`(3)^+`$ functions which are paired with $`(+1)^+`$ and $`(+2)^+`$ functions, respectively, to build up $`M=1`$. The same scheme was used for the contributions of one-particle functions of negative $`z`$-parities. There it was sufficient to use the combinations $`(1)^{}/0^{}`$ and $`(2)^{}/(+1)^{}`$. In order to describe excitations we added one-particle basis functions with quantum numbers $`m^{\pi _z}=(1)^+`$ and values $`l_i=1`$ and $`k_i=1`$ (see eqs.(11,12)). The latter have exclusively been optimized for a nuclear charge number $`Z=1`$, in contrast to all the other types of basis functions which have been optimized for $`Z=1`$ (hydrogen) and for $`Z=2`$(He<sup>+</sup>). The reason is again that all bound $`M=1`$ states are one-particle excitations in which the excited electron is associated to the one-particle quantum numbers $`m^{\pi _z}=(1)^+`$, and the effect of the nucleus on that outer electron is screened by the inner one, thereby giving rise to an effective nuclear charge close to unity.
For the $`(1)^{}`$ subspace we proceeded in a similar manner, the basis set dimension was $`N=3671`$ for both singlet and triplet states, built up from a set of 228 optimized one-particle basis functions of type (10).
In the following, we present and discuss the results for the helium energy calculations. Comparing our field-free data with the literature we observe that they are for the subspace $`(1)^+/(1)^{}`$ subspaces more accurate than in the case of the $`0^+/0^{}`$ subspaces. We will as far as possible compare our helium energies for finite field strength with the best data available in the literature.
### 3.2 Energies for finite field strengths
#### 3.2.1 Results for $`M=1`$ and even $`z`$-parity
a) The singlet states $`\nu ^1(1)^+`$
For the singlet subspace $`\nu ^1(1)^+`$ we present the energies of the ground state and the first four excitations, i.e. $`1\nu 5`$. The energies for the $`1^1(1)^+`$ state are presented in table 1, together with the values given in the literature, if available. The accuracy of the field free energy is even higher than the ones of the $`1^10^+`$ and $`1^10^{}`$ states due to the lower correlation energy for the $`M=1`$ states. We remark that even though the data for the $`(1)^+`$ symmetry presented in ref. have been accurate enough for identifying features in the spectrum of the magnetic white dwarf GD229 as absorption edges of helium, we have been able to improve the accuracy slightly.
We observe that for small values of $`B`$ the total energies of the $`1^1(1)^+`$ state lie slightly lower than for $`B=0`$ which is an effect due to the Zeeman energy which is negative for $`M=1`$. However, for fields stronger than $`B0.16a.u.`$ the total energies rise drastically with increasing field strength which has its origin in the increasing kinetic energy of the electrons which can roughly be estimated by their Landau energy amounting to $`B`$ for both electrons together.
In order to reveal the internal energetics of the atom we have to subtract such pure and overall field effects. For an analysis it is advantageous to subtract even more: one measure for the accuracy of the total energies $`E(B)`$ are their corresponding one-particle ionization energies $`|E(B)T(B)|`$ corresponding to the process He $``$ He$`{}_{}{}^{+}+e^{}`$. The threshold $`T(B)`$, i.e. the lowest possible total energy for which the system He$`{}_{}{}^{+}+e^{}`$ can exist possessing the same quantum numbers as the He state in question, is given in the fourth column in table 1 (the values for the ionization energies are trivial to compute and are not additionally listed in any of the tables). The values $`T(B)`$ for $`(1)^+`$ symmetry are identical to those given in ref. for the $`M=0`$ symmetry. The reason is that for any $`z`$-parity and spin symmetry the lowest $`M=1`$ state of the system He$`{}_{}{}^{+}+e^{}`$ is realized by the ionized electron in a Landau state with magnetic quantum number $`m=1`$ and the He<sup>+</sup> ion in its $`0^+`$ ground state: the Landau energy of the electron depends on $`(m+|m|)`$ and is therefore identical for $`m=1`$ and $`m=0`$. The alternative possibility of associating the value $`m=1`$ to the He<sup>+</sup> ion possesses a higher energy.
The energies for the excited states $`\nu ^1(1)^+`$, $`2\nu 5`$ are given in table 2. We remark that for finite field strengths there are no data about these states available in the literature so far. We observe as expected that their total energies rapidly approach the threshold $`T(B)`$ with increasing excitation.
The dependence of the one-particle ionization energies on the field strength is shown in figure 1. In contrast to the total energies of the $`1^1(1)^+`$ state its ionization energy depends monotonously on $`B`$: the outer electron becomes increasingly bound with increasing field strength. This is not in general the case for the excited states within the $`(1)^+`$ subspace although the ionization energy of the $`2^1(1)^+`$ state is monotonous (note the changes in the slope compared to the $`1^1(1)^+`$ state). The ionization energies of the states $`3^1(1)^+`$, $`4^1(1)^+`$ and $`5^1(1)^+`$ exhibit a pattern of several avoided crossings. Those occur in the field strength interval $`0.02B0.2`$ which is the regime where a rearrangement takes place: for low field strengths the states $`3^1(1)^+`$ and $`4^1(1)^+`$ (field free $`4^1F_1`$ and $`4^1P_1`$) are energetically almost degenerate since they both belong to the field free principal quantum number $`n=4`$. This degeneracy is disturbed by the magnetic field and completely destroyed for strong fields (see figure 1).
b) The triplet states $`\nu ^3(1)^+`$
For the triplet subspace $`\nu ^3(1)^+`$ we present the energies for the ground state and the first four excitations, i.e. $`1\nu 5`$. The total energies of the states with $`S_z=1`$ are given in table 3 together with data existing in the literature for the states $`1^3(1)^+`$, $`2^3(1)^+`$ and $`3^3(1)^+`$. The states $`4^3(1)^+`$ and $`5^3(1)^+`$ have not been present in the literature so far. We observe that our energies are variationally lower than any values in the literature apart from only a few exceptions with respect to the energies computed by Jones et al . However, one must take into account that the released-phase quantum Monte Carlo method performed by Jones et al leads to statistical error bars (see the corresponding numbers in parentheses in table 3) such that none of those energies is significantly stronger bound than our results. This means, either we confirm these results or our accuracy is even higher, reflecting the fact that our optimized anisotropic Gaussian basis sets excellently describes the wave functions in a magnetic field. Due to the efficiency of our method of computation we were able to cover a large number of field strengths.
Due to the spin shift $`BS_z`$ the triplet states with $`S_z=1`$ are the most strongly bound ones among the three states with $`S_z=0,\pm 1`$. This spin shift causes the $`1_1^3(1)^+`$ state to cross the low-field ground state $`1^10^+`$ at $`B0.750a.u.`$ For $`B\stackrel{>}{}0.750a.u.`$ the $`1_1^3(1)^+`$ state is the global ground state of the atom. The $`1_1^3(1)^+`$ state is for any field strength lower than the $`1_1^3(1)^{}`$ state (see sec. 3.2.2), and it lies also lower than the $`1_1^30^+`$ and $`1_1^30^{}`$ states for $`B0.750a.u.`$. The latter ones cross the $`1^10^+`$ state at $`B1.112a.u.`$ and $`B0.994a.u.`$, respectively.
Analogously to the singlet case, we show the one-particle ionization energies of the triplet states also in figure 1. We observe that the singlet-triplet splitting decreases with increasing excitation. This occurs due to the fact that for excitations the spatial separation between the electrons is large, and therefore the exchange terms in eqs.(28-30) are small which causes a small effect of the different signs of the matrix elements belonging to singlet and triplet states.
#### 3.2.2 Results for $`M=1`$ and odd $`z`$-parity
a) The singlet states $`\nu ^1(1)^{}`$
For the singlet subspace $`\nu ^1(1)^{}`$ we present also the ground state and the first four excitations, i.e. $`1\nu 5`$. The total energies are given in table 4. Like in the case of the excited $`{}_{}{}^{1}(1)_{}^{+}`$ states there exist no data for finite field strengths in the literature. Our field-free data are in good agreement with the literature. In figure 2 we show the one-particle ionization energies analogously to figure 1. The threshold $`T(B)`$ is identical to the case of $`(1)^+`$, $`0^+`$ or $`0^{}`$ symmetry because the energetically lowest way to realize the ionized system He$`{}_{}{}^{+}+e^{}`$ with $`(1)^{}`$ symmetry is to leave one electron in the $`0^+`$ state of a He atom. The other electron must then be placed in a Landau orbital with $`m=1`$ and negative $`z`$-parity, which does not affect the Landau energy $`B/2`$.
b) The triplet states $`\nu ^3(1)^{}`$
For the triplet states, several investigations exist in the literature , in contrast to the singlet case. In table 5 we have listed our total energies and the corresponding values of the literature for the states $`\nu _1^3(1)^{}`$ (i.e. $`S_z=1`$), $`1\nu 5`$. Again our results are better than almost all the reference values for finite field strengths or are at least comparable to the ones of Jones et al within their statistical error bars. For the states $`3^3(1)^+`$ and $`4^3(1)^+`$ there exist no data in the literature so far.
In figure 2 we show the ionization energies. The dashed triplet curves almost coincide with the corresponding singlet curves. Such a small singlet-triplet splitting is in good agreement with the considerations mentioned in sec.3.2.1 which predict a small singlet-triplet splitting for high excitations.
## 4 Transitions
For the comparison of the energy levels of helium with the spectra of magnetic white dwarfs in general and GD229 in particular it is necessary to determine the transition energies from our total energies. Restricting ourselves to electric dipole transitions, we have the selection rules $`\mathrm{\Delta }S=0`$, $`\mathrm{\Delta }S_z=0`$ for the spin degrees of freedom in our nonrelativistic approach and $`\mathrm{\Delta }M=0,\mathrm{\Delta }\mathrm{\Pi }_z=\pm 1`$ (for linearly polarized radiation) or $`\mathrm{\Delta }M=\pm 1,\mathrm{\Delta }\mathrm{\Pi }_z=0`$ (for circularly polarized radiation) for the spatial degrees of freedom. Whereas in ref. we were already able to present the $`\mathrm{\Delta }M=0`$ transitions $`0^+0^{}`$, we are now able to investigate three times as many transitions: firstly the $`\mathrm{\Delta }M=0`$ transitions between the $`(1)^+`$ and the $`(1)^{}`$ states, and additionally two classes of $`|\mathrm{\Delta }M|=1`$ transitions between the $`M=1`$ states and $`M=0`$ states, involving positive and negative $`z`$-parities, respectively. Altogether our data yield $`75`$ singlet and $`70`$ triplet transitions.
Due to the fact that the field strengths in the atmospheres of magnetic white dwarfs is not a constant but varies by a factor of two for a dipole geometry, transitions which behave monotonically as a function of the varying field are smeared out, i.e. are not expected to provide a signature in the observed spectrum. However, the transitions whose wavelengths are stationary with respect to the field dependence manifest themselves as absorption edges in the observable spectrum if they possess a relevant intensity. We therefore give in tables 6 to 11 a complete list of all the stationary points which resulted from our calculated transitions.
We remark that for a reliable comparison with observational spectra in some cases the finite nuclear mass effects have to be taken into account. This requires corrections of our total energies according to the scaling relation given in eq.(4) in ref. and in refs.. Due the selection rules $`\mathrm{\Delta }S=0`$, $`\mathrm{\Delta }S_z=0`$, we have for the transition energies the scaling relation $`\mathrm{\Delta }E(M_0,\mu ^2B)=\mu \mathrm{\Delta }E(\mathrm{},B)\frac{1}{M_0}\mu ^2B\mathrm{\Delta }M`$ (here $`M_0=7344a.u.`$ is the nuclear mass and $`\mu =0.999864a.u.`$ is the reduced mass). For $`\mathrm{\Delta }M=0`$ transitions the effect is always such that the position (field strength) and wavelength of a mass-corrected stationary point are related to the corresponding fixed-nucleus result by $`B(M_0)=\mu ^2B(\mathrm{})`$ and $`\lambda (M_0)`$ = $`\frac{1}{\mu }\lambda (\mathrm{})`$. If for $`|\mathrm{\Delta }M|=1`$ transitions the ratio $`\frac{B}{M_0}`$ is small compared to $`\mathrm{\Delta }E`$ it is still possible to do an approximate correction for the wavelengths of the stationary points directly with the data presented: We then have $`\lambda (M_0)=\frac{1}{\mu }\lambda (\mathrm{})(1+\mu \frac{B}{M_0}\lambda (\mathrm{})\mathrm{\Delta }M)`$. The stationary points corrected exactly within the scaling relation given above can, of course, only be obtained by separately scaling all the values of $`\mathrm{\Delta }E`$ and by interpolating them over the grid of scaled field strengths. In the argumentation above we have taken into account the normal mass correction terms. The specific mass corrections are expected to be even less significant, in particular for stronger fields and excited states.
Altogether we detected 139 stationary points; several ones among them possess large uncertainties which arise mainly due to the interpolation error with respect to the crude grid of field strengths. Most of the transitions $`(1)^+0^+`$, however, are so precise that the corresponding data have already successfully been used to explain the absorption edges in the spectrum of the white dwarf GD229, and these data together with the other stationary transitions serve as a good basis for astrophysicists to investigate the spectra of unidentified magnetic objects.
## 5 Concluding remarks and Outlook
We have investigated the fixed-nucleus electronic structure of the helium atom in a magnetic field by a fully correlated approach. Scaling laws allow to include finite-mass effects. The present work is concerned with the energy levels and transitions of the helium states with magnetic quantum number $`M=1`$. A small part of the corresponding data has already been used successfully for identifying the features in the spectrum of the white dwarf GD229 with electronic transitions in atomic helium, which represented one of the goals of our work. The enlarged data presented in this paper will now, together with the energy levels provided in ref., serve as a good starting point for the analysis of observed spectra of magnetic astrophysical objects in general.
The reliability of our wavelength data results from the high accuracy of our energy values which ranges between $`10^4a.u.`$ and $`10^6a.u.`$. This accuracy has become possible due to our appraoch by means of an optimized anisotropic Gaussian basis set. Since the spherical invariance is broken by the magnetic field, it has been necessary to use Gaussians with different length scales for the longitudinal and transversal degrees of freedom. The nonlinear parameters describing these length scales have been determined by the requirement to solve optimally the one-particle problem of the H atom or the He<sup>+</sup> ion in a magnetic field of given strength. These optimized one-particle functions have been used to construct configurations in order to represent the full fixed-nucleus Hamiltonian and the overlap as matrices separately in each subspace of fixed quantum numbers corresponding to the four conserved quantities: the total spin $`𝐒^2`$ and its $`z`$-component $`S_z`$, the $`z`$-component $`L_z`$ of the electronic angular momentum and the electronic $`z`$-parity $`\mathrm{\Pi }_z`$. The corresponding generalized eigenvalue problem provided a variational estimation for the energy eigenvalues.
Atomic energies of helium have been calculated for the ground state and the first four excitations in each subspace for $`M=1`$, i.e. for positive and negative $`z`$-parity as well as for singlet and triplet spin symmetry. We considered altogether $`20`$ different field strengths $`0B100a.u.`$, i.e. $`0B2.350510^7`$ Tesla. This series production of data has become by very efficient algorithms for the computation of the matrix elements, in particular the electon-electron matrix elements for which we presented an analytical formula derived in cylindrical coordinates, thereby making the CPU effort for the calculations for nonzero angular momentum affordable: Building up a matrix of dimension about 4000 takes less than one hour CPU time on a moderate Silicon Graphics workstation.
The comparison of energy data with observed spectra of astrophysical objects is possible by searching for stationary points of the transitions with respect to the magnetic field strength. The data for $`M=1`$ in the present paper yield $`\mathrm{\Delta }M=0`$ transitions, and together with the results for $`M=0`$ in ref. we have also been able to consider $`\mathrm{\Delta }M=1`$ transitions. Complete tables of all the detected stationary points for the mentioned transitions were given. Energy data for $`M=2`$ and $`M=3`$ are planned to be investigated in a future work.
In order to complete the treatment of bound electronic transitions of helium in a magnetic field, we will in the near future also investigate in detail the oscillator strengths of the mentioned transitions as a function of the magnetic field. Based on the calculated intensities it will be possible to produce synthetic spectra, and a comparison with the observed spectra will give important hints for models of the radiation transport and the field configuration in magnetic astrophysical objects.
Acknowledgements. The Deutsche Studienstiftung (W.B.) and the Deutsche Forschungsgemeinschaft (W.B.) are gratefully acknowledged for financial support.
## Appendix A Analytical solution to the electron-electron integral
In the following it is our aim to derive an analytical expression for the electron-electron integral which allows its efficient numerical implementation. We emphasize that an efficient treatment of the two-particle integrals is essential for the calculations on helium since the number of two-particle matrix elements is $`N(N+1)/2`$, in contrast to the one-particle matrix elements whose number is only $`n(n+1)/2`$ (here $`N4000`$ is the dimension of the two-particle Hamiltonian matrix whereas $`n200`$ is the dimension of the underlying one-particle basis set, see eqs.(10,13)). Denoting the two-particle interaction with $`V_{II}(𝐫_1,𝐫_2)=\frac{1}{|𝐫_1𝐫_2|}`$, we have to solve the integral
$`ij|V_{II}|kl`$ $`=`$ $`{\displaystyle d^3r_1d^3r_2\mathrm{\Phi }_i(𝐫_1)\mathrm{\Phi }_j(𝐫_2)\frac{1}{|𝐫_1𝐫_2|}\mathrm{\Phi }_k(𝐫_1)\mathrm{\Phi }_l(𝐫_2)},`$ (A.1)
where the one-particle orbitals $`\mathrm{\Phi }_i`$ are of type (10), obeying the constraints (11,12). For the sake of brevity, we will in the following use the index notation $`\gamma _{ik}:=\gamma _i+\gamma _k`$ for the sum of two indexed quantities.
The initial step is to apply a Singer transformation in order to remove the Coulomb singularity, thereby introducing the new variable $`u`$ according to $`\frac{1}{|𝐫_1𝐫_2|}=\frac{2}{\sqrt{\pi }}_0^{\mathrm{}}𝑑ue^{u^2(𝐫_1𝐫_2)^2}`$. Then the integrand of this new integration over $`u`$ decomposes into a transversal and a longitudinal part
$`ij|V_{II}|kl`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑uI_{\rho \phi }(u)I_z(u),`$ (A.2)
where
$`I_z(u)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z_1{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z_2z_1^{n_{z}^{}{}_{ik}{}^{}}z_2^{n_{z}^{}{}_{jl}{}^{}}e^{\beta _{ik}z_1^2u^2(z_1z_2)^2\beta _{jl}z_2^2}`$ (A.3)
and
$`I_{\rho \phi }(u)`$ $`=`$ $`{\displaystyle \rho _1𝑑\rho _1𝑑\phi _1\rho _2𝑑\rho _2𝑑\phi _2\rho _1^{n_{\rho }^{}{}_{ik}{}^{}}e^{i(m_im_k)\phi _1}e^{\alpha _{ik}\rho _1^2}e^{u^2(\rho _1^2+\rho _2^22\rho _1\rho _2\mathrm{cos}(\phi _2\phi _1))}}`$
$`\times \rho _2^{n_{\rho }^{}{}_{jl}{}^{}}e^{i(m_jm_l)\phi _2}e^{\alpha _{jl}\rho _2^2}`$
The longitudinal integral $`I_z(u)`$ is the trivial part of the matrix element (A.1). For decoupling the particles $`1`$ and $`2`$ we subsitute $`\stackrel{~}{z}_1=z_1b(u)z_2`$, $`\stackrel{~}{z}_2=z_2`$ where $`b(u)=\frac{u^2}{\beta _{ik}+u^2}`$ and $`\frac{(\stackrel{~}{z}_1,\stackrel{~}{z}_2)}{(z_1,z_2)}=1`$. The exponential factorizes, and the power $`z_1^{n_{z}^{}{}_{ik}{}^{}}`$ can be multiplied out, yielding for $`I_z(u)`$ a sum over standard integrals of the Gaussian type $`_{\mathrm{}}^{\mathrm{}}z^{n_z}e^{\gamma (u)z^2}`$ which can easily be evaluated, giving after a few steps of algebra
$`I_z(u)`$ $`=`$ $`4g_{n_{z}^{}{}_{ijkl}{}^{}}{\displaystyle \underset{\zeta =0}{\overset{\zeta \frac{n_{z}^{}{}_{ik}{}^{}}{2}}{}}}\left(\begin{array}{c}n_{z}^{}{}_{ik}{}^{}\\ 2\zeta \end{array}\right)J(n_{z}^{}{}_{ik}{}^{},n_{z}^{}{}_{jl}{}^{},\beta _{ik},\beta _{jl};2\zeta ;u)`$ (A.7)
where the prefactor $`g_l:=mod(l,2)`$ reflects the fact that the total $`z`$-parity is a conserved quantity. The function $`J`$ is defined by
$`J(n_1,n_2,a_1,a_2;v;u)`$ $`=`$ $`\begin{array}{c}\\ \frac{1}{4}\end{array}a_1^{n_1\frac{1}{2}}a_2^{\frac{n_{12}+1}{2}}\mathrm{\Gamma }(\begin{array}{c}\\ \frac{v+1}{2}\end{array})\mathrm{\Gamma }(\begin{array}{c}\\ \frac{n_1+n_2v+1}{2}\end{array})(a_1a_2)^{\frac{v}{2}}`$ (A.19)
$`\times u^{2(n_1v)}(1+\begin{array}{c}\\ \frac{1}{a_1}\end{array}u^2)^{\frac{n_2n_1}{2}}(1+\begin{array}{c}\\ \frac{a_{12}}{a_1a_2}\end{array}u^2)^{\frac{n_{12}v+1}{2}}`$
The usual procedure for the treatment of the transversal part $`I_{\rho \phi }`$ would be to represent it in Cartesian coordinates since then its decomposition into a sum of products of integrals $`I_x(u)`$ and $`I_y(u)`$ analoguos to $`I_z(u)`$ would be technically simple: The exponential would factorize automatically, and the remaining part consists of factors like $`\rho ^{n_\rho }e^{\pm im}`$ which can be expressed as $`(x^2+y^2)^k(x\pm iy)^{|m|}`$ due to the contraint (11) and which can trivially be multiplied out. We emphasize that this procedure, however, would yield expressions which are very expensive with respect to the CPU time since already for small nonzero maqnetic quantum numbers $`m`$ or excitations $`k`$ the number of binomic terms gives rise to a large number of integrals over $`u`$.
In the following we exploit that the described drawback is not at all a property of the integral (A.1) but only of the Cartesian approach. In fact, the many mentioned $`u`$-integrals are not independent from each other but the information how to resummarize their results to a more compact expression is lost in the lengthy Cartesian algebra. Therefore, it is our strategy to use cylindrical coordinates from the very beginning, although the derivation of this condensed analytical formula is technically involved.
The key for solving the integral $`I_{\rho \phi }(u)`$ is the following substitution
$`\stackrel{~}{\rho }_1`$ $`=`$ $`\sqrt{\rho _1^2+a^2\rho _2^22a\rho _1\rho _2\mathrm{cos}(\phi _2\phi _1)}\stackrel{~}{\rho }_2=\rho _2`$ (A.20)
$`\mathrm{sin}\stackrel{~}{\phi }_1`$ $`=`$ $`{\displaystyle \frac{\rho _1}{\stackrel{~}{\rho }_1}}\mathrm{sin}(\phi _2\phi _1);\mathrm{cos}\stackrel{~}{\phi }_1={\displaystyle \frac{\rho _1}{\stackrel{~}{\rho }_1}}\mathrm{cos}(\phi _2\phi _1)a{\displaystyle \frac{\rho _2}{\stackrel{~}{\rho }_1}}\stackrel{~}{\phi }_2=\begin{array}{c}\\ \frac{1}{2}\end{array}(\phi _2+\phi _1)`$ (A.23)
which results from the following ideas (here $`a(u)=\frac{u^2}{\alpha _{ik}+u^2}`$). Firstly, we use the angular part of $`I_{\rho \phi }(u)`$ for exploiting the conservation of $`L_z`$, and secondly the remaining radial part has to be decoupled similarly like in the case of the $`z`$-integration mentioned above. Using in eq.(A) the angle $`\phi :=\phi _2\phi _1`$ and the cyclic angle $`\overline{\phi }:=\frac{1}{2}(\phi _2+\phi _1)`$, the integration over $`\overline{\phi }`$ yields a Kronecker Delta reflecting the conservation of $`L_z`$:
$`I_{\rho \phi }(u)`$ $`=`$ $`2\pi \delta _{m_{ij},m_{lk}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\rho _1𝑑\rho _1{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\rho _2𝑑\rho _2\rho _1^{n_{\rho }^{}{}_{ik}{}^{}}e^{\alpha _{ik}\rho _1^2}\rho _2^{n_{\rho }^{}{}_{jl}{}^{}}e^{\alpha _{jl}\rho _2^2}e^{u^2(\rho _1^2+\rho _2^2)}`$ (A.24)
$`\times {\displaystyle \underset{0}{\overset{2\pi }{}}}d\phi e^{i(m_im_k)\phi }e^{u^2(2\rho _1\rho _2\mathrm{cos}\phi )}.`$
Now, the expression $`(\rho _1\mathrm{cos}\phi )`$ plays a similar role as $`z_1`$ in eq.(A.3), and therefore we temporarily introduce $`\xi _1=(\rho _1\mathrm{cos}\phi )`$ and $`\eta _1=(\rho _1\mathrm{sin}\phi )`$ as dummy Cartesian variables. For decoupling the particles $`1`$ and $`2`$ analogously to the treatment of $`I_z(u)`$ we substitute $`\stackrel{~}{x}_1=\xi _1a\stackrel{~}{\rho }_2`$ whereas $`\stackrel{~}{y}_1=\eta _1`$ and $`\stackrel{~}{\rho }_2=\rho _2`$ remain unchanged. The variables $`\stackrel{~}{\rho }_1`$ and $`\stackrel{~}{\phi }_1`$ given in eqs.(A.20,A.23) are now just the new cylindrical coordinates belonging to $`\stackrel{~}{x}_1`$ and $`\stackrel{~}{y}_1`$. In this new representation $`I_{\rho \phi }`$ reads
$`I_{\rho \phi }(u)`$ $`=`$ $`2\pi \delta _{m_{ij},m_{lk}}{\displaystyle \underset{0}{\overset{2\pi }{}}}𝑑\stackrel{~}{\phi }_1{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\stackrel{~}{\rho }_1\stackrel{~}{\rho }_1\rho _1^{n_{\rho }^{}{}_{ik}{}^{}}e^{(\alpha _{ik}+u^2)\stackrel{~}{\rho }_1^2}e^{i(m_im_k)\phi _1}`$ (A.25)
$`\times {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}d\stackrel{~}{\rho }_2\stackrel{~}{\rho }_2^{n_{\rho }^{}{}_{jl}{}^{}+1}e^{\frac{\alpha _{ik}\alpha _{jl}+\alpha ijklu^2}{\alpha _{ik}+u^2}\stackrel{~}{\rho }_2^2},`$
where the variables $`\rho _1`$ and $`\phi _1`$ occuring in the factor $`\rho _1^{n_{\rho }^{}{}_{ik}{}^{}}e^{i(m_im_k)\phi _1}`$ have to be considered as functions of the variables labelled with the symbol ($`\stackrel{~}{}`$). We split $`\rho _1^{n_{\rho }^{}{}_{ik}{}^{}}e^{i(m_im_k)\phi _1}=\rho _1^{2k_{ik}}\rho _1^{|m_i|+|m_k|}e^{i(m_im_k)\phi _1}`$ according to eq.(11). This is a powerful step since the relations (A.20,A.23) can explicitely be solved with respect to those two factors. After some algebra we arrive at:
$`\rho _1^{2k_{ik}}`$ $`=`$ $`{\displaystyle \underset{r_1+r_2+r_3+r_4=k_{ik}}{}}\left(\begin{array}{c}k_{ik}\\ r_1r_2r_3r_4\end{array}\right)a^{2r_2+r_3+r_4}\stackrel{~}{\rho }_1^{2r_1+r_3+r_4}\stackrel{~}{\rho }_2^{2r_2+r_3+r_4}e^{i(r_3r_4)\stackrel{~}{\phi }_1}`$ (A.28)
$`\rho _1^{|m_i|+|m_k|}e^{i(m_im_k)\phi _1}`$ $`=`$ $`{\displaystyle \underset{\mu _i=0}{\overset{|m_i|}{}}}\left(\begin{array}{c}|m_i|\\ \mu _i\end{array}\right){\displaystyle \underset{\mu _k=0}{\overset{|m_k|}{}}}\left(\begin{array}{c}|m_k|\\ \mu _k\end{array}\right)a^{|m_i|+|m_k|\mu _{ik}}\stackrel{~}{\rho }_1^{\mu _{ik}}\stackrel{~}{\rho }_2^{|m_i|+|m_k|\mu _{ik}}e^{i(s_i\mu _is_k\mu _k)\stackrel{~}{\phi }_1}`$ (A.33)
where $`s_i:=sgn(m_i)`$ are the signs of the magnetic quantum numbers. Inserting the expressions (A.28,A.33) into eq.(A.25), we can perform the $`\stackrel{~}{\phi }_1`$-integration yielding a Kronecker Delta which restricts the summation indices to $`r_3r_4=s_i\mu _i+s_k\mu _k`$. The obtained sum is a decomposition of $`I_{\rho \phi }`$ into products of each two radial Gaussian integrals completely analogous to the expressions encountered in the evaluation of the integral $`I_z`$, which allows us after several steps to express also $`I_{\rho \phi }`$ in terms of the function $`J`$ given in eq.(A.19):
$`I_{\rho \phi }(u)`$ $`=`$ $`4\pi ^2\delta _{m_{kl},m_{ij}}{\displaystyle \underset{\mu _i=0}{\overset{|m_i|}{}}}\left(\begin{array}{c}|m_i|\\ \mu _i\end{array}\right){\displaystyle \underset{\mu _k=0}{\overset{|m_k|}{}}}\left(\begin{array}{c}|m_k|\\ \mu _k\end{array}\right){\displaystyle \underset{\genfrac{}{}{0pt}{}{r_1+r_2+r_3+r_4=k_{ik}}{r_3r_4=s_i\mu _i+s_k\mu _k}}{}}\left(\begin{array}{c}k_{ik}\\ r_1r_2r_3r_4\end{array}\right)`$ (A.41)
$`\times J(n_{\rho }^{}{}_{ik}{}^{}+1,n_{\rho }^{}{}_{jl}{}^{}+1,\alpha _{ik},\alpha _{jl};\mu _{ik}+1+2r_1+r_3+r_4;u)`$
The final step is to evaluate the $`u`$-integration after inserting eqs.(A.41,A.7) into eq.(A.2). Considering
$`ij|V_{II}|kl`$ $`=`$ $`32\pi ^{\frac{3}{2}}\delta _{m_{kl},m_{ij}}g_{n_{z}^{}{}_{ijkl}{}^{}}{\displaystyle \underset{\mu _i=0}{\overset{|m_i|}{}}}\left(\begin{array}{c}|m_i|\\ \mu _i\end{array}\right){\displaystyle \underset{\mu _k=0}{\overset{|m_k|}{}}}\left(\begin{array}{c}|m_k|\\ \mu _k\end{array}\right){\displaystyle \underset{\genfrac{}{}{0pt}{}{r_1+r_2+r_3+r_4=k_{ik}}{r_3r_4=s_i\mu _i+s_k\mu _k}}{}}\left(\begin{array}{c}k_{ik}\\ r_1r_2r_3r_4\end{array}\right){\displaystyle \underset{\zeta =0}{\overset{\zeta \frac{n_{z}^{}{}_{ik}{}^{}}{2}}{}}}\left(\begin{array}{c}n_{z}^{}{}_{ik}{}^{}\\ 2\zeta \end{array}\right)`$ (A.51)
$`\times {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}duJ(n_{\rho }^{}{}_{ik}{}^{}+1,n_{\rho }^{}{}_{jl}{}^{}+1,\alpha _{ik},\alpha _{jl};\mu _{ik}+1+2r_1+r_3+r_4;u)J(n_{z}^{}{}_{ik}{}^{},n_{z}^{}{}_{jl}{}^{},\beta _{ik},\beta _{jl};2\zeta ;u)`$
shows that the solution of the electron-electron integral is now reduced to the integration
$`K(n_1,n_2,a_1,a_2;v;m_1,m_2,b_1,b_2;w)`$ $`=`$ $`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑uJ(n_1,n_2,a_1,a_2;v;u)J(m_1,m_2,b_1,b_2;w;u).`$ (A.52)
Up to a constant prefactor, the integrand in eq.(A.52) is of the type $`g(u)=u^{2n_u}(1+au^2)^{r_a}(1+bu^2)^{r_b}(1+cu^2)^{r_c}(1+du^2)^{r_d}`$ where $`a=\frac{1}{a_1}`$, $`b=\frac{1}{b_1}`$, $`c=\frac{a_{12}}{a_1a_2}`$, $`d=\frac{b_{12}}{b_1b_2}`$ are fixed real numbers given by the basis functions. The exponents depend on the summation indices which enter into the functions $`J`$, and $`n_u:=n_1+m_1vw`$ is running over positive integer values. Further $`r_a:=\frac{n_2n_1}{2}`$ and $`r_b:=\frac{m_2m_1}{2}`$ are integers whereas $`r_c:=\frac{v1n_{12}}{2}`$ is integer or half-integer and $`r_d:=\frac{w1m_{12}}{2}`$ is always half-integer. The substitution $`x:=\frac{u^2}{1/d+u^2}`$ with $`du=\frac{1}{2}d^{1/2}x^{1/2}(1x)^{3/2}dx`$ leads to
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}g(u)𝑑u`$ $`=`$ $`\frac{1}{2}d^{n_u\frac{1}{2}}{\displaystyle \underset{0}{\overset{1}{}}}x^{n_u\frac{1}{2}}(1x)^{r_{abcd}n_u\frac{3}{2}}(1+q_ax)^{r_a}(1+q_bx)^{r_b}(1+q_cx)^{r_c}𝑑x`$ (A.53)
where $`r_{abcd}:=r_a+r_b+r_c+r_d`$ and $`q_a:=\frac{a}{d}1`$, $`q_b:=\frac{b}{d}1`$, $`q_c:=\frac{c}{d}1`$. In order to reduce this integral to hypergeometric functions it is now necessary to multiply out one of the three last factors. This is possible since the basis functions can always be interchanged in such a way that one of the exponents $`r_a`$, $`r_b`$ is positive, say $`r_a`$. Then eq.(3.211) in ref. can be used to obtain
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}g(u)𝑑u`$ $`=`$ $`{\displaystyle \frac{1}{2}}d^{n_u\frac{1}{2}}{\displaystyle \underset{s=0}{\overset{r_a}{}}}\left(\begin{array}{c}r_a\\ s\end{array}\right)q_a^s\times B(r_{abcd}n_u\begin{array}{c}\\ \frac{1}{2}\end{array},n_u+s+\begin{array}{c}\\ \frac{1}{2}\end{array})`$ (A.63)
$`\times F_1(n_u+s+\begin{array}{c}\\ \frac{1}{2}\end{array},r_b,r_c,sr_{abcd};q_b,q_c)`$
where $`B`$ ist the beta function and $`F_1`$ is the Appell hypergeometric function. It is defined as double series, $`F_1(a,b,b^{};c;x,y):=_{\nu =0}^{\mathrm{}}\frac{(a,\nu )(b,\nu )}{(c,\nu )(1,\nu )}{}_{2}{}^{}F_{1}^{}(a+\nu ,b^{},c+\nu ;y)x^\nu `$ where $`{}_{2}{}^{}F_{1}^{}(a,b,c;z):=_{\mu =0}^{\mathrm{}}\frac{(a,\mu )(b,\mu )}{(c,\mu )(1,\mu )}z^\mu `$ is the Gaussian hypergeometric function and $`(a,\nu ):=\frac{\mathrm{\Gamma }(a+\nu )}{\mathrm{\Gamma }(a)}`$ is the Pochhammer symbol. Including all prefactors, we obtain for the integral $`K`$:
$`K(n_1,n_2,a_1,a_2,v;m_1,m_2,b_1,b_2,w)`$ $`=`$ $`\begin{array}{c}\\ \frac{1}{16}\end{array}a_1^{n_1\frac{1}{2}}a_2^{\frac{n_{12}+1}{2}}\mathrm{\Gamma }(\begin{array}{c}\\ \frac{v+1}{2}\end{array})\mathrm{\Gamma }(\begin{array}{c}\\ \frac{n_{12}v+1}{2}\end{array})(a_1a_2)^{\frac{v}{2}}`$ (A.75)
$`\times b_1^{m_1\frac{1}{2}}b_2^{\frac{m_{12}+1}{2}}\mathrm{\Gamma }(\begin{array}{c}\\ \frac{w+1}{2}\end{array})\mathrm{\Gamma }(\begin{array}{c}\\ \frac{m_{12}w+1}{2}\end{array})(b_1b_2)^{\frac{w}{2}}\times {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}g(u)𝑑u`$
where now the integral stands for the expression (A.63).
The final result is a real value for $`ij|V_{II}|kl`$ obtained by inserting eq.(A.75) into eq.(A.51):
$`ij|V_{II}|kl`$ $`=`$ $`\pi ^{\frac{3}{2}}\delta _{m_{kl},m_{ij}}g_{n_{z}^{}{}_{ijkl}{}^{}}\alpha _{ik}^{n_1\frac{1}{2}}\alpha _{jl}^{\frac{n_{12}+1}{2}}\beta _{ik}^{m_1\frac{1}{2}}\beta _{jl}^{\frac{m_{12}+1}{2}}{\displaystyle \underset{\mu _i=0}{\overset{|m_i|}{}}}\left(\begin{array}{c}|m_i|\\ \mu _i\end{array}\right){\displaystyle \underset{\mu _k=0}{\overset{|m_k|}{}}}\left(\begin{array}{c}|m_k|\\ \mu _k\end{array}\right)`$ (A.104)
$`\times {\displaystyle \underset{\genfrac{}{}{0pt}{}{r_1+r_2+r_3+r_4=k_{ik}}{r_3r_4=s_i\mu _i+s_k\mu _k}}{}}\left(\begin{array}{c}k_{ik}\\ r_1r_2r_3r_4\end{array}\right)\mathrm{\Gamma }(\begin{array}{c}\\ \frac{v+1}{2}\end{array})\mathrm{\Gamma }(\begin{array}{c}\\ \frac{n_{12}v+1}{2}\end{array})(\alpha _{ik}\alpha _{jl})^{\frac{v}{2}}`$
$`\times {\displaystyle \underset{\genfrac{}{}{0pt}{}{\zeta =0}{wgerade}}{\overset{\zeta \frac{n_{z}^{}{}_{ik}{}^{}}{2}}{}}}\left(\begin{array}{c}n_{z}^{}{}_{ik}{}^{}\\ 2\zeta \end{array}\right)\mathrm{\Gamma }(\begin{array}{c}\\ \frac{1}{2}\end{array}+\zeta )\mathrm{\Gamma }(\begin{array}{c}\\ \frac{m_{12}+1}{2}\end{array}\zeta )(\beta _{ik}\beta _{jl})^\zeta `$
$`\times d^{n_u\frac{1}{2}}{\displaystyle \underset{s=0}{\overset{r_a}{}}}\left(\begin{array}{c}r_a\\ s\end{array}\right)q_a^sB(r_{abcd}n_u\begin{array}{c}\\ \frac{1}{2}\end{array},n_u+s+\begin{array}{c}\\ \frac{1}{2}\end{array})`$
$`\times F_1(n_u+s+\begin{array}{c}\\ \frac{1}{2}\end{array},r_b,r_c,sr_{abcd};q_b,q_c)`$
An overview over the various parameters entering into eq.(A.104) is given in appendix C where we briefly present the algorithm for the implementation of eq.(A.104). The direct implementation of eq.(A.104) yields already a factor of 5 to 10 with respect to the increase in speed compared with the common result derived in Cartesian coordinates in ref.. However, by systematic exploitation of the symmetries the gain factor of CPU can further be raised to between 20 and 40.
## Appendix B Symmetry properties of the electron-electron integral
In the present section we list the different types of symmetries for the expressions encountered in appendix A.
### B.1 External symmetries
By an external symmetry we mean a permutation of basis functions which leave $`ij|V_{II}|kl`$ invariant (up to complex conjugation). Such symmetries are firstly $`ij|V_{II}|kl=ji|V_{II}|lk`$ since the particles are indistinguishable, and secondly we have $`ij|V_{II}|kl=kl|V_{II}|ij^{}`$ due to the hermiticity of $`V_{II}`$, in the case of our real matrix elements even $`ij|V_{II}|kl=kl|V_{II}|ij`$.
### B.2 Internal symmetries
The integral $`ij|V_{II}|kl`$ is also invariant under any manipulation of the basis functions which leave the integrals $`I_z(u)`$ and $`I_{\rho \phi }(u)`$ in eq.(A.2) invariant separately. This class of internal transformations is larger than the class of permutations of basis functions since the parameters entering in $`I_z(u)`$ and $`I_{\rho \phi }(u)`$ can be varied independently. $`I_z(u)`$ is invariant under the simultaneous exchanges $`ij`$ and $`kl`$ of indices only in the parameters $`\beta `$ and $`n_z`$, and $`I_{\rho \phi }(u)`$ stays unchanged under the analogous interchange of only the parameters $`\alpha `$, $`n_\rho `$ and $`m`$. Eq.(A) shows a further internal symmetry of $`I_{\rho \phi }(u)`$: if the difference $`m_im_k`$ remains unchanged the parameters $`m_i`$ and $`m_k`$ can be varied arbitrarily as long as the constraint (11) can be fulfilled for some $`k_i`$ and $`k_k`$, respectively (analogously for $`m_j`$ and $`m_l`$). It is an interesting fact that the matrix element $`ij|V_{II}|kl`$ in a subspace with $`M=m_i+m_j=m_k+m_l`$ possesses the same numerical value as some matrix elements involved in subspaces corresponding to some other value of $`M`$!
The internal symmetries can be exploited for achieving simultaneously $`n_{\rho }^{}{}_{jl}{}^{}n_{\rho }^{}{}_{ik}{}^{}`$ and $`n_{z}^{}{}_{jl}{}^{}n_{z}^{}{}_{ik}{}^{}`$, i.e. in the function $`g(u)`$ (see eqs.(A.52,A.53)) even both parameters $`r_a`$ and $`r_b`$ can be assumed to be non-negative. A very important consequence of this is that due to the definition of $`F_1`$ (see paragraph below eq.(A.63)) its double series reduces to a sum over a finite number of expressions involving the Gaussian hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ which is much simpler to evaluate than $`F_1(a,b,b^{};c;x,y)`$.
### B.3 Further symmetries
Assuming both parameters $`r_a`$ and $`r_b`$ to be non-negative, eq.(A.63) is also valid after interchanging $`r_ar_b`$ (and simultaneously $`q_aq_b`$), representing a new invariance transformation of $`ij|V_{II}|kl`$.
## Appendix C Implementation
In the following we present the algorithm for the implementation of eq.(A.104) for our calculations on helium.
1. Due the external symmetry $`ij|V_{II}|kl=kl|V_{II}|ij`$ it is sufficient to compute only the upper triangle of the matrix corresponding to $`V_{II}`$.
2. If $`n_{\rho }^{}{}_{jl}{}^{}n_{\rho }^{}{}_{ik}{}^{}<0`$, we use the external symmetry $`ij|V_{II}|kl=ji|V_{II}|lk`$ and interchange globally $`ij`$, $`kl`$, which achieves $`r_a0`$.
3. If then $`n_{z}^{}{}_{jl}{}^{}n_{z}^{}{}_{ik}{}^{}<0`$, we use the internal symmetry of $`I_z(u)`$ and interchange $`n_{z}^{}{}_{i}{}^{}n_{z}^{}{}_{j}{}^{}`$, $`n_{z}^{}{}_{k}{}^{}n_{z}^{}{}_{l}{}^{}`$ as well as $`\beta _i\beta _j`$, $`\beta _k\beta _l`$, which achieves $`r_b0`$.
4. We introduce the following abbreviations for the fixed parameters independent on the summation indices of eq.(A.104): name definition type property $`n_1`$ $`=`$ $`n_{\rho }^{}{}_{ik}{}^{}+1`$ integer $`1n_1`$ $`n_{12}`$ $`=`$ $`n_{\rho }^{}{}_{ik}{}^{}+n_{\rho }^{}{}_{jl}{}^{}+2`$ even $`2n_{12}`$ $`m_1`$ $`=`$ $`n_{z}^{}{}_{ik}{}^{}`$ integer $`0m_1`$ $`m_{12}`$ $`=`$ $`n_{z}^{}{}_{ik}{}^{}+n_{z}^{}{}_{jl}{}^{}`$ even $`0m_{12}`$ $`r_a`$ $`=`$ $`\begin{array}{c}\\ \frac{n_{\rho }^{}{}_{jl}{}^{}n_{\rho }^{}{}_{ik}{}^{}}{2}\end{array}`$ integer $`0r_a`$ $`r_b`$ $`=`$ $`\begin{array}{c}\\ \frac{n_{z}^{}{}_{jl}{}^{}n_{z}^{}{}_{ik}{}^{}}{2}\end{array}`$ integer $`0r_b`$ $`d`$ $`=`$ $`\begin{array}{c}\\ \frac{\beta _{ik}+\beta _{jl}}{\beta _{ik}\beta _{jl}}\end{array}`$ real $`1<d`$ $`q_a`$ $`=`$ $`1\begin{array}{c}\\ \frac{\beta _{ik}\beta _{jl}}{\alpha _{ik}(\beta _{ik}+\beta _{jl})}\end{array}`$ real $`1<q_a`$ $`q_b`$ $`=`$ $`\begin{array}{c}\\ \frac{\beta _{ik}}{\beta _{ik}+\beta _{jl}}\end{array}`$ real $`1<q_b<0`$ $`q_c`$ $`=`$ $`1\begin{array}{c}\\ \frac{(\alpha _{ik}+\alpha _{jl})\beta _{ik}\beta _{jl}}{(\beta _{ik}+\beta _{jl})\alpha _{ik}\alpha _{jl}}\end{array}`$ real $`1<q_c`$
5. If then $`r_a>r_b`$, we interchange $`r_ar_b`$, $`q_aq_b`$, thereby minimizing the number of terms of the innermost summation in eq.(A.104).
6. Applying the internal symmetry transformation of suitably shifting the values of $`m_i`$ and $`m_k`$, we replace both $`m_i`$ and $`m_k`$ by $`m_i+2s`$ and $`m_k+2s`$, respectively, where $`s`$ is the highest possible integer such that $`|m_i+2s|n_{\rho }^{}{}_{i}{}^{}`$ and $`|m_k+2s|n_{\rho }^{}{}_{k}{}^{}`$. This step is very useful for decreasing the number of summations over the indices $`r_1,r_2,r_3,r_4`$ since raising $`m_i`$ and $`m_k`$ has the consequence that $`k_i`$ and $`k_k`$ must decrease for keeping $`n_{\rho }^{}{}_{i}{}^{}`$ and $`n_{\rho }^{}{}_{k}{}^{}`$ constant.
7. We apply the formula (A.104), and the parameters occuring therein depend on the summation indices according to the following definitions: name definition type property $`v`$ $`=`$ $`\mu _{ik}+1+2r_1+r_3+r_4`$ integer $`1`$ $`v`$ $`n_{\rho }^{}{}_{ik}{}^{}+1`$ $`w`$ $`=`$ $`2\zeta `$ even $`0`$ $`w`$ $`n_{z}^{}{}_{ik}{}^{}g_{n_{z}^{}{}_{ik}{}^{}}`$ $`n_u`$ $`=`$ $`n_1+m_1vw`$ integer $`g_{n_{z}^{}{}_{ik}{}^{}}`$ $`n_u`$ $`n_{\rho }^{}{}_{ik}{}^{}+n_{z}^{}{}_{ik}{}^{}`$ $`r_c`$ $`=`$ $`\begin{array}{c}\\ \frac{v1n_{12}}{2}\end{array}`$ int. or half-int. $`\begin{array}{c}\\ \frac{n_{\rho }^{}{}_{ik}{}^{}+n_{\rho }^{}{}_{jl}{}^{}}{2}\end{array}1`$ $`r_c`$ $`\begin{array}{c}\\ \frac{n_{\rho }^{}{}_{jl}{}^{}}{2}\end{array}1`$ $`r_d`$ $`=`$ $`\begin{array}{c}\\ \frac{w1m_{12}}{2}\end{array}`$ half-integer $`\begin{array}{c}\\ \frac{1+n_{z}^{}{}_{ik}{}^{}+n_{z}^{}{}_{jl}{}^{}}{2}\end{array}`$ $`r_d`$ $`\begin{array}{c}\\ \frac{1+g_{n_{z}^{}{}_{ik}{}^{}}+n_{z}^{}{}_{jl}{}^{}}{2}\end{array}<0`$ $`r_{abcd}`$ $`=`$ $`n_1m_1+\begin{array}{c}\\ \frac{v+w}{2}\end{array}1`$ int. or half-int. $`n_{\rho }^{}{}_{ik}{}^{}n_{z}^{}{}_{ik}{}^{}\begin{array}{c}\\ \frac{3}{2}\end{array}`$ $`r_{abcd}`$ $`\begin{array}{c}\\ \frac{n_{\rho }^{}{}_{ik}{}^{}n_{z}^{}{}_{ik}{}^{}3g_{n_{z}^{}{}_{ik}{}^{}}}{2}\end{array}<0`$
8. Although in principle $`F_1(a,b,b^{};c;x,y)`$ is a double power series in the arguments $`x,y`$ with convergence radii $`|x|<1`$ and $`|y|<1`$ and although in our case even both arguments may happen to lie close to $`1`$, we do not have to care about the power series in $`x`$. The reason is that the corresponding series terminates after $`r_b`$ terms because the negative integer $`r_b`$ enters as second parameter argument into $`F_1`$ in eq.(A.104). The sum terms in $`F_1`$ each involve once the Gaussian hypergeometric function, and since their arguments are related, it is possible to use a continued fraction representation for the ratio $`\frac{{}_{2}{}^{}F_{1}^{}(a+1,b,c+1,z)}{{}_{2}{}^{}F_{1}^{}(a,b,c,z)}`$ in order to establish a recursion law which is stable over the typically $`r_b<10`$ recursions. Therefore, for each occurence of $`F_1`$, it was only necessary to evaluate a single time the Gaussian hypergeometric function $`{}_{2}{}^{}F_{1}^{}`$, which is not always simple but very efficiently possible by means of fast analytical continuation formulas. Several of them are given in eqs.(15.3.3-12) in ref., a much larger set of such continuation formulas is presented in ref..
We point out that without such a systematic analysis of the electron-electron integral as well as the Appell hypergeometric function $`F_1`$ and in particular the Gaussian hypergeometric function $`{}_{2}{}^{}F_{1}^{}`$, the computation of many excited helium states for nonzero magnetic quantum number for many different field strengths would not have been possible.
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# Contents
## 1 Introduction
The dynamics of massless open string states propagating in flat space–time can be described, at weak string coupling, by an effective action
$$S(A)\frac{1}{g_s}\text{Tr }(A)$$
which is function of a $`U(N)`$ gauge potential $`A`$ and where the lagrangian density $``$ is, in general, an expansion in $`\alpha ^{}`$ written in terms of arbitrary powers of the curvature $`F`$ and of its covariant derivatives. The action $`S`$ has various equivalent interpretations. First of all, it reproduces, at tree level, the disk amplitudes computed directly in string theory. Secondly, the equations of motion derived from $`S`$ correspond to the condition of conformal invariance of the open string sigma model with Wilson line interactions Tr$`P\mathrm{exp}\left(i_\mathrm{\Sigma }A\right)`$ on the boundary $`\mathrm{\Sigma }`$ of the string world-sheet. In fact, it was shown by Tseytlin and Andreev that the effective action itself can be identified with the partition function of the $`\sigma `$–model. Finally we recall that, using $`T`$–duality, the action $`S(A)`$ also describes the weak coupling dynamics of $`D`$-branes of all dimensions.
New progress in the understanding of the general form of the effective action $`S`$ has been made in a recent work by Seiberg and Witten. In the authors show that, if the open string propagates in a space–time with a constant background NSNS two–form field $`B`$, then the dynamics of the open–string modes can be described in two equivalent ways. Firstly, one can consider the original action $`S`$, and add to the curvature $`F`$ a central term $`B\mathrm{𝟏}`$. Alternatively, one can replace the original $`U(N)`$ gauge theory by a gauge theory on a non–commutative space, with non–commutativity parameter $`\theta `$ related to the background metric $`g`$ and two–form $`B`$. The action then takes the same form as the original action $`S`$, with the gauge potential $`A`$ replaced with the gauge potential $`\widehat{A}`$ of a non–commutative gauge theory, and with products of fields replaced with Moyal products $``$ with parameter $`\theta `$. More conjecturally, it is shown in that one can in fact choose the parameter $`\theta `$ freely by properly adjusting the central term $`\mathrm{\Phi }\mathrm{𝟏}`$ added to the non–commutative curvature $`\widehat{F}`$. For each value of $`\theta `$, the action describing the open–string dynamics has then the same exact form as the original action $`S`$. The parameter $`\theta `$ is then a redundant parameter, since different values of $`\theta `$ correspond to the same underlying physics. On the other hand, the simple fact that the form of the action at various values of $`\theta `$ is invariant imposes severe restrictions on the possible structure of the original action $`S`$. This paper is devoted to the understanding of those restrictions. Previous work on this subject is contained in .
Before describing the results let us remark a basic fact. In order to attack the problem of invariance, one need not consider the full $`U(N)`$ theory. In fact, if one restricts the attention to the $`U(1)`$ case, but considers non–zero values of the non–commutativity parameter $`\theta `$, one is then considering a theory which is effectively non–abelian. From an algebraic point of view, the requirements imposed by form–invariance of the action are identical in the $`U(1)`$ and $`U(N)`$ case as soon as $`\theta 0`$. For this reason we work throughout in the $`U(1)`$ case, but the results will, at the end, be valid in the general $`U(N)`$ setting. It is really quite remarkable that the simpler abelian theory contains, in an subtle way indeed, the complete information about the general non–abelian theory.
Let us now describe the general results of this paper. As just remarked, they are valid in the general non–abelian $`U(N)`$ setting. An invariant action $`S`$ is given as a linear combination $`S=_ic_iI_i`$ of basic actions $`I_i`$, which we call invariant blocks. A single block $`I`$ has itself the following structure. First write $`I`$ as an expansion in $`\alpha ^{}`$, as
$$I\underset{LP}{}\left(\alpha ^{}\right)^LI_L$$
where we call a term proportional to $`\left(\alpha ^{}\right)^L`$ a term of level $`L`$. The lowest level term $`I_P`$ is a pure derivative term. The precise meaning of this notion will be given later in the paper, but informally we can say that pure derivative terms are those which are invariant under addition of a central term to the curvature $`F`$ (the basic example is the $`F^2`$ term at level $`P=2`$). The higher level terms $`I_L`$ are then needed in order to achieve invariance of the full action $`I`$ under a change of the parameter $`\theta `$. The basic result of this paper is to reduce the question of invariance to a set of algebraic equations relating the various terms $`I_L`$. In particular, we will show that the requirement of invariance can be rephrased in terms of four basic algebraic operators
$$\mathrm{\Delta }\overline{\mathrm{\Delta }}\delta \overline{\delta }$$
which depend on an arbitrary antisymmetric matrix $`\mathrm{\Delta }_{ab}`$ and which satisfy the basic commutation relations
$$[\overline{\mathrm{\Delta }},\mathrm{\Delta }]=[\delta ,\overline{\delta }]=2L\frac{D}{2},$$
(1)
where $`D`$ is the dimension of space–time. The various terms $`I_L`$ then must satisfy the equations
$`\mathrm{\Delta }I_L`$ $`=`$ $`\delta I_{L+1}`$ (2)
$`\overline{\mathrm{\Delta }}I_{L+1}`$ $`=`$ $`\overline{\delta }I_L.`$
A term of lowest level satisfies then $`\mathrm{\Delta }I_P=\overline{\delta }I_P=0`$.
We also show in this paper how the above equations can be explicitly solved in $`D=2`$. This is clearly a toy model, since gauge–bosons in two dimensions do not have propagating degrees of freedom. On the other hand, the algebraic equations are perfectly well defined in dimension $`2`$, and are highly non–trivial. This model is useful for a variety of reasons. First of all, one can show that, starting from the lowest level $`F^2`$ term, one can reconstruct the full BI action, plus a minimal set of derivative corrections which are required for invariance of the action (we show that derivative corrections already enter at level $`4`$). Similarly, one expects that the invariant block built from $`F^2`$ in general dimension $`D`$ will be the minimal derivative extension of the NBI action. Finally, the construction in $`D=2`$ is a first indication of how to solve the equations (2) in the general case.
We have said that the action $`S`$ is a linear combination of invariant blocks. The specific coefficients are not constrained by the methods of this paper, and should be determined by other means. However let us note that none of the arguments which follow rely on supersymmetry, and therefore we expect supersymmetry to impose constraints on the coefficients themselves, restricting even more the set of allowed forms of the action.
We shall now describe the contents of this paper. First let me note that section 5.2 contains a concise summary of the results, which includes the main equations in the text. The structure of the paper is as follows. In section 2 we review the general structure of the action $`S`$, as can be inferred from the analysis of scattering amplitudes of gauge bosons. We then use this knowledge in section 3 to rewrite the effective action as a matrix action. In doing so, we review the general form of the Seiberg–Witten map which relates commutative and non–commutative gauge potentials, and we introduce the formal algebraic machinery which is required in the sequel of the paper. We also derive the second equation in (2). Section 4 is then devoted to rewriting the results of section 3 in an invariant operator language. The purpose is two–fold. On one side, the structure of the action becomes more transparent. Moreover, in this setting, the question of invariance from $`\theta `$ is more easily understood and solved. Section 5 then applies the results of section 4 to the specific problem at hand, and completes the derivation of the equations (2), together with the basic commutator (1). Section 6 is then devoted to the explicit solution of equations (2) in the case of $`D=2`$. We conclude in section 7 with discussion and comments on open problems for future research.
The results of this paper require, together with the general discussion, a considerable number of technical lemmata. We have tried to limit the technical discussion to a minimum in the main body of the paper, leaving the precise proofs of many statements to a rather large appendix.
## 2 The General Form of the Effective Action
We consider an open string propagating in flat $`D`$–dimensional space–time $`M`$ with coordinates $`x^a`$ and with constant metric $`g_{ab}`$. Throughout the paper we use units such that
$$2\pi \alpha ^{}=1.$$
We concentrate on the physics of the massless $`U(N)`$ gauge bosons, to lowest order in the string coupling constant $`g_s`$. The amplitude $`𝒜(p^I)`$ for the scattering of $`n`$ gluons with momenta $`p^I`$ – with $`g^{ab}p_a^Ip_b^I=0`$ – is computed starting from the disk $`n`$–point function of the corresponding vertex operators, cyclically ordered on the boundary of the string world–sheet, and then by summing over the cyclically inequivalent orderings. More precisely, one has that
$$𝒜(p^I)g_s^{n2}\underset{\begin{array}{c}\sigma \text{ cyclically}\\ \text{inequivalent}\end{array}}{}_g(p^{\sigma _1},\mathrm{},p^{\sigma _n})$$
where $`_g`$ depends on the metric $`g`$ and is invariant under cyclic permutations of the arguments<sup>1</sup><sup>1</sup>1We omit explicit reference to polarizations..
One can equivalently summarize the information about disk amplitudes by introducing an effective action $`S(A)`$, function of a $`U(N)`$ connection $`A`$ on $`M`$, such that the tree level amplitudes of $`S`$ are equal to the disk amplitudes $`𝒜`$. The general form of the action $`S`$ is well known and reads
$$S(A,g,g_s)=\frac{\mathrm{Tr}}{g_s}d^Dxdet{}_{}{}^{\frac{1}{2}}g_{ab}^{}\left(1+\frac{1}{4}F_{ab}F_{cd}g^{ac}g^{bd}+\mathrm{}\right).$$
(3)
We have absorbed any numerical prefactor in the definition of $`g_s`$. Moreover, the terms hidden in $`\mathrm{}`$ contain both higher powers of the field strength and derivative terms<sup>2</sup><sup>2</sup>2Various facts are known about the terms in (3). First of all, in the $`U(1)`$ case, the terms without derivatives resum to the Born–Infeld lagrangian
$$S(A)=\frac{1}{g_s}d^Dxdet{}_{}{}^{\frac{1}{2}}(g+F).$$
In the non–abelian case even the non–derivative terms are not completely known. At order $`F^4`$ the computation can be explicitly carried out in string theory, and the result is proportional to
$$\mathrm{Tr}\left(F_{ab}F_{cb}F_{ad}F_{cd}+\frac{1}{2}F_{ab}F_{cb}F_{cd}F_{ad}\frac{1}{4}F_{ab}F_{ab}F_{cd}F_{cd}\frac{1}{8}F_{ab}F_{cd}F_{ab}F_{cd}\right).$$
At higher orders in $`F`$, the more reasonable proposal is a natural extension of the Born–Infeld action proposed by Tseytlin in terms of a symmetrized trace prescription
$$S(A)=\frac{\mathrm{STr}}{g_s}d^Dxdet{}_{}{}^{\frac{1}{2}}(g+F).$$
(4) The above prescription not only matches (up to order $`F^4`$) with the scattering computations in superstring theory, but also matches results for $`D`$–brane actions derived within matrix theory. . On the other hand, it is known that the symmetrized trace prescription is incomplete at order $`F^6`$. In , the authors study the spectra of excitations around diagonal and intersecting D–brane configurations on tori, and find discrepancies with the prescription (4). The correction terms at order $`\alpha ^3`$ have been explicitly computed in . Some derivative corrections are known, both in the $`U\left(1\right)`$ case, as well as in the non–abelian setting. For the $`U(1)`$ theory, some derivative terms have been computed . In particular, for bosonic open string theory, the authors find terms at order $`F^2FF`$. Still in , derivative correction in superstring theory at order $`F^2FF`$ are discussed. In , the author finds derivative corrections at order $`F^5`$ (and $`F^3D^2F`$), proportional to $`\zeta \left(3\right)`$, by studying $`5`$–point disk amplitudes. Finally, in the bosonic theory, there is a known derivative term at order $`F^3`$ which is proportional to
$$\mathrm{Tr}\left(F_{ab}[F_{bc},F_{ca}]\right).$$
. As a note on conventions, in all that follows actions will always be written assuming a Euclidean signature of the metric.
Let us now introduce a constant background NS–NS two–form field $`B_{ab}`$. In the presence of open strings, a constant field $`B`$ is not pure gauge and does affect the dynamics of the gauge bosons. In particular, the effects have been very clearly analyzed in and can be summarized as follows
* The momenta of the asymptotic gluon states satisfy a modified on–shell condition. More precisely, gluons are massless with respect to an effective open string metric $`G=gB\frac{1}{g}B`$ and therefore the corresponding momenta satisfy $`G^{ab}p_ap_b=0`$.
* The effective coupling constant is modified to an open string value of $`G_s=g_sdet{}_{}{}^{\frac{1}{2}}Gdet{}_{}{}^{\frac{1}{2}}(g+B)`$
* Finally, the disk scattering amplitudes are modified by momentum dependent phase factors. In particular, in terms of the antisymmetric matrix $`\theta ^1=Bg\frac{1}{B}g`$, the amplitudes are given by
$$𝒜(p^I)G_s^{n2}\underset{\begin{array}{c}\sigma \text{ cyclically}\\ \text{inequivalent}\end{array}}{}_G(p^{\sigma _1},\mathrm{},p^{\sigma _n})e^{\frac{i}{2}\theta ^{ab}_{I>J}p_a^{\sigma _I}p_b^{\sigma _J}}$$
(5)
The above dynamics can again be summarized in tree diagrams of a modified effective action which can be written, starting from (3), in various different but equivalent ways . Let me briefly review the various options:
1. On one hand, one can follow the usual prescription by starting with the action (3) and by simply adding to the field strength $`F`$ the central term $`B\mathrm{𝟏}`$.
2. On the other hand, one can follow the ideas of . In this case, one considers the action (3) with the replacements $`g_{ab}G_{ab}`$ and $`g_sG_s`$. This correctly reproduces the modified gauge coupling and the modified mass–shell condition. Phase factors in (5) are reproduced by reinterpreting the matrix $`\theta `$ as a non–commutative scale of space–time, and substituting, in the action (3), ordinary products of fields with Moyal products in terms of $`\theta `$. Correspondingly, the $`U(N)`$ gauge potential $`A`$ is now mapped into a gauge potential
$$\widehat{A}=\widehat{A}_{SW}(A,\theta )$$
of a $`U(N)`$ non–commutative gauge theory, and the gauge group is modified accordingly. We will call the map $`\widehat{A}_{SW}`$ the Seiberg–Witten map.
NOTATION. We have denoted in (3) by $`S(A,g,g_s)`$ the effective action at zero $`B`$ field, as a function of the gauge potential, the metric, and the coupling constant. At finite $`B`$, there are two new relevant parameters in the description of the action – a possible central term added to the curvature, and a possible non–commutativity parameter. In general, the action will then depend on five parameters
$$S(\text{potential, metric, coupling, central term, NC paramter}).$$
Then the equivalence of the descriptions $`1`$ and $`2`$ above is just
$$S(A,g,g_s,B,0)=S(\widehat{A},G,G_s,0,\theta ).$$
3. Finally, one can follow a naive procedure, which is usually not considered, but which will be important for our future discussion. In fact, this procedure is the most natural one if we are given only the information about the amplitudes (5), without any reference to an underlying string theory. Specifically, we may wish to reproduce directly the amplitudes (5) using a standard $`U(N)`$ theory, without using any previous information about the theory at zero $`B`$. Firstly, the kinetic term of the theory must be
$$\frac{\mathrm{Tr}}{G_s}d^Dxdet{}_{}{}^{\frac{1}{2}}G_{ab}^{}\left(\frac{1}{4}G^{ac}G^{bd}F_{ab}F_{cd}\right)$$
(6)
in order to reproduce the correct mass–shell condition. We can write the above equation in a more suggestive form by introducing the matrix $`\mathrm{\Gamma }=g+B`$ and by first noting that, since $`\mathrm{Tr}\left(FF\right)`$ is a total derivative, then
$$d^Dx\mathrm{Tr}\left(F_{ab}F_{cd}+F_{ac}F_{db}+F_{ad}F_{bc}\right)=0.$$
This identity, together with the facts that $`2G^{ab}=\mathrm{\Gamma }^{ab}+\mathrm{\Gamma }^{ba}`$ and that $`G_s^1det{}_{}{}^{\frac{1}{2}}G=g_s^1det{}_{}{}^{\frac{1}{2}}\mathrm{\Gamma }`$, can be used to show that the kinetic term (6) is equal to
$$\frac{\mathrm{Tr}}{g_s}d^Dxdet{}_{}{}^{\frac{1}{2}}\mathrm{\Gamma }_{ab}^{}\left(\frac{1}{4}\mathrm{\Gamma }^{ac}\mathrm{\Gamma }^{db}F_{ab}F_{cd}+\frac{1}{8}\mathrm{\Gamma }^{ab}F_{ab}\mathrm{\Gamma }^{cd}F_{cd}\right).$$
The above is nothing but the quadratic term coming from the expansion of the Born–Infeld action $`\sqrt{g+B+F}=\sqrt{\mathrm{\Gamma }+F}`$. The matrix $`\mathrm{\Gamma }`$, which is not symmetric, now plays the role of the metric and therefore more tensor structures are possible (in fact, the expansion of the BI action includes all powers of $`F`$, including the odd ones). In some sense, we have traded the central term $`B\mathrm{𝟏}`$ added to the curvature $`F`$ with an addition to the metric $`gg+B`$, by allowing metrics to be non–symmetric. This invariance is natural from the point of view of the open string $`\sigma `$–model
$$_\mathrm{\Sigma }g_{ab}X^a\overline{}X^b+_\mathrm{\Sigma }B+_\mathrm{\Sigma }A$$
where we see, using $`_\mathrm{\Sigma }A=_\mathrm{\Sigma }F`$, that only the combination $`g+B+F`$ has an invariant meaning. We are then led to conclude that there is an extension of (3) in the case of a non–symmetric metric $`\mathrm{\Gamma }`$ of the general form
$$\frac{\mathrm{Tr}}{g_s}d^Dxdet{}_{}{}^{\frac{1}{2}}\mathrm{\Gamma }_{ab}^{}\left[1\frac{1}{2}\mathrm{\Gamma }^{ab}F_{ab}+\mathrm{}\right]$$
(7)
which reproduces the amplitudes (5). Moreover, we claim that all contractions of indices in $`\mathrm{}`$ are done with $`\mathrm{\Gamma }^{ab}`$ (no terms containing $`\mathrm{\Gamma }_{ab}`$). This fact can be shown starting with (5). In fact, these amplitudes are just functions of $`G^1`$ and $`\theta `$, which only depend on $`\mathrm{\Gamma }^{ab}`$. This shows that the amplitudes do not depend on $`\mathrm{\Gamma }_{ab}`$. To show the same fact for the vertices of the action, we must show that the subtractions coming from poles in the various subchannels also share the same property. The only problems could come from internal propagators $`p^2G_{ab}`$. We recall though that , for amplitudes of the form (5), a general tree graph is computed by first analyzing the graph at $`\theta =0`$, and then by multiplying it by a phase factor depending only on external momenta. Moreover, the graph at $`\theta =0`$ has all propagators $`p^2G_{ab}`$ contracted with metrics $`G^{ab}`$ on the vertices, thus proving that the subtraction is again just a function of $`G^1`$ and $`\theta `$. Using the general form (7) of the action we can then give a meaning to the function $`S`$ when the metric is not symmetric. We can then summarize the equality of descriptions $`1`$ and $`3`$ by saying that
$$S(A,g,g_s,B,0)=S(A,g+B,g_s,0,0).$$
We have then seen that we can trade a central term $`B\mathrm{𝟏}`$ with either a non–commutativity parameter $`\theta `$ or with an addition $`gg+B`$ to the metric. It is natural (following ) to conjecture that, in fact, one has a continuous family of possibilities, parametrized by a central term $`\mathrm{\Phi }\mathrm{𝟏}`$ and by a free parameter $`\theta `$. The effective non–symmetric metric $`\mathrm{\Gamma }`$ then combines with the central term $`\mathrm{\Phi }`$ into the invariant combination $`\mathrm{\Gamma }+\mathrm{\Phi }`$, which, following again , is given by
$$\frac{1}{\mathrm{\Gamma }+\mathrm{\Phi }}+\theta =\frac{1}{g+B}.$$
The effective coupling again depends only on the sum $`\mathrm{\Gamma }+\mathrm{\Phi }`$, and is given by
$$\frac{1}{G_s}det{}_{}{}^{\frac{1}{2}}(\mathrm{\Gamma }+\mathrm{\Phi })=\frac{1}{g_s}det{}_{}{}^{\frac{1}{2}}(g+B).$$
Finally, the gauge potential is given by the Seiberg–Witten map $`\widehat{A}=\widehat{A}_{SW}(A,\theta )`$. The action $`S(\widehat{A},\mathrm{\Gamma },G_s,\mathrm{\Phi },\theta )`$ is then independent of $`\mathrm{\Phi }`$ and $`\theta `$.
## 3 The Effective Action as a Matrix Action
In this section we continue our general analysis of the effective action $`S`$, but we restrict our attention to the $`U(1)`$ case. As already noted in the introduction, whenever the non–commutativity scale $`\theta `$ is non–zero, the $`U(1)`$ case contains the physics of the full $`U(N)`$ theory, and we therefore lose nothing in concentrating on the effective action for $`N=1`$.
### 3.1 Choosing the central term
In the previous section we have argued that the effective action describing the dynamics of gluons in space–time is given by a function $`S(\widehat{A},\mathrm{\Gamma },G_s,\mathrm{\Phi },\theta )`$ of five arguments – i.e. the gauge potential $`\widehat{A}`$, the generalized non–symmetric metric $`\mathrm{\Gamma }`$, the coupling $`G_s`$, the central term $`\mathrm{\Phi }`$ and the non–commutativity parameter $`\theta `$.
The arguments of the action $`S`$ are not all independent, since physically different backgrounds are parametrized only by the closed string parameters $`A,g+B`$ and $`g_s`$, which we keep fixed. In fact, the same physical situation corresponds to a family of different values of the arguments of $`S`$, parameterized by $`\mathrm{\Phi }`$ and $`\theta `$, which we consider as free parameters. The remaining variables $`\widehat{A},\mathrm{\Gamma }`$ and $`G_s`$ are then determined, in terms of the fixed closed string parameters, by the equations
$`\widehat{A}`$ $`=`$ $`\widehat{A}_{SW}(A,\theta )`$
$`{\displaystyle \frac{1}{\mathrm{\Gamma }+\mathrm{\Phi }}}+\theta `$ $`=`$ $`{\displaystyle \frac{1}{g+B}}`$ (8)
$`{\displaystyle \frac{1}{G_s}}det{}_{}{}^{\frac{1}{2}}(\mathrm{\Gamma }+\mathrm{\Phi })`$ $`=`$ $`{\displaystyle \frac{1}{g_s}}det{}_{}{}^{\frac{1}{2}}(g+B).`$ (9)
The action $`S`$ is then independent of $`\mathrm{\Phi }`$ and $`\theta `$.
We use this freedom to choose the central term $`\mathrm{\Phi }`$. Throughout the paper we will denote with $`K`$ the inverse of $`\theta `$
$$K=\frac{1}{\theta }.$$
Using the independence of $`S`$ on $`\mathrm{\Phi }`$, we set
$$\mathrm{\Phi }=K.$$
The only free parameter is then the non–commutativity scale $`\theta `$.
Equation (8) can be easily rewritten in terms of the combination
$$\gamma =g+BK$$
and reads
$$\mathrm{\Gamma }=K\frac{1}{\gamma }K.$$
(10)
Finally equation (9), which determines the coupling, becomes
$$\frac{1}{G_s}det{}_{}{}^{\frac{1}{2}}\mathrm{\Gamma }=\frac{1}{g_s}det{}_{}{}^{\frac{1}{2}}K.$$
(11)
### 3.2 A matrix action
Let us now consider the expression for the field strength $`\widehat{F}`$ and its covariant derivatives. We start by introducing the coordinate functions
$$x_a=K_{ab}x^b$$
and the combinations
$`\lambda _a`$ $`=`$ $`x_a+\widehat{A}_a`$
$`\lambda ^a`$ $`=`$ $`\theta ^{ab}\lambda _b=x^a+\theta ^{ab}\widehat{A}_b.`$
Note that we raise and lower indices with the matrix $`\theta ,K`$. Using the simple fact that, for any function $`f`$,
$$_af=i[x_a,f],$$
we quickly see that the commutator $`i[\lambda _a,\lambda _b]`$ is given by
$`i[\lambda _a,\lambda _b]`$ $`=`$ $`\widehat{F}_{ab}K_{ab}`$
$`=`$ $`\widehat{F}_{ab}+\mathrm{\Phi }_{ab}`$
and therefore computes the field strength with the addition of the correct central term. Moreover, for any function $`f`$ which transforms in the adjoint representation of the non–commutative gauge group, the commutator
$`i[\lambda _a,f]`$ $`=`$ $`_afi[\widehat{A}_a,f]`$
$`=`$ $`\widehat{D}_af`$
computes the covariant derivative $`\widehat{D}_af`$. Therefore, any expression involving products of covariant derivatives of the field strength, with the addition of the central term $`\mathrm{\Phi }=K`$, can be expressed in terms of $``$ products of the functions $`\lambda _a`$ – for example, an expression like $`(\widehat{F}+\mathrm{\Phi })_{ab}\widehat{D}_c(\widehat{F}+\mathrm{\Phi })_{de}`$ can be rewritten as $`i[\lambda _a,\lambda _b][\lambda _c,[\lambda _d,\lambda _e]]`$. We conclude that the general form of the effective action $`S`$ is
$$\frac{1}{G_s}\underset{n\text{ even}}{}d^Dxdet{}_{}{}^{\frac{1}{2}}\mathrm{\Gamma }\left(\lambda _{a_1}\mathrm{}\lambda _{a_n}\right)\eta ^{a_1\mathrm{}a_n}$$
where the coefficients $`\eta ^{a_1\mathrm{}a_n}`$ are constructed from the matrix $`\mathrm{\Gamma }^{ab}`$. We may further manipulate the above equation using (10) and (11) and raising and lowering indices with the matrix $`\theta ,K`$. We can then write
$$S=\frac{1}{g_s}\underset{n\text{ even}}{}d^Dxdet{}_{}{}^{\frac{1}{2}}K\left(\lambda ^1\mathrm{}\lambda ^n\right)\eta _{1\mathrm{}n}$$
(12)
We have used a compact notation for indices, which will be used extensively in the sequel, substituting $`a_11`$, $`a_22`$, $`\mathrm{}`$. Moreover, the symbol $`\eta _{1\mathrm{}n}`$ represents the tensor which is built from the matrix $`\gamma _{ab}`$ exactly as $`\eta ^{1\mathrm{}n}`$ is constructed starting from $`\mathrm{\Gamma }^{ab}`$. For example, if $`\eta ^{1234}=\mathrm{\Gamma }^{13}\mathrm{\Gamma }^{24}`$ then $`\eta _{1234}=\gamma _{13}\gamma _{24}`$.
REMARK . Let us recall that, in the general $`U\left(N\right)`$ effective action, the distinction between derivative and non–derivative terms is ambiguous, since a commutator of covariant derivatives $`\left(D_aD_bD_bD_a\right)\mathrm{}`$ is equivalent to a commutator with the field strength $`[F_{ab},\mathrm{}]`$. In fact, when one writes the action in matrix form as in (12), all terms (derivative and non–derivative) are included on an equal footing. In particular, the ambiguity discussed above becomes naturally the Jacobi identity of the commutator $`[\lambda ^a,\lambda ^b]`$.
### 3.3 Gauge invariance and invariance under addition of a central term
In the previous subsection we have rewritten the action $`S`$ in the compact form (12), which is more suited for discussing the action in its entirety, including terms with arbitrary powers of the field strength and with arbitrary number of derivatives. On the other hand, the gauge invariance of the original action is not immediately transparent in this new notation, and one needs to restate the requirement of gauge invariance in terms of the tensors $`\eta _{1\mathrm{}n}`$. This is easily done by noting that the functions $`\lambda _a`$, when used in covariant expressions, always appear within commutators, and therefore adding a constant $`\epsilon _a`$ to the function $`\lambda _a`$ does not change the action. Gauge invariance becomes then, within the matrix formulation (12) of the action, invariance under translations $`\lambda ^a\lambda ^a+\epsilon ^a`$. This requirement quickly translates into the following algebraic relation which must be satisfied by the tensors $`\eta `$
$$\eta _{123\mathrm{}n}+\eta _{213\mathrm{}n}+\eta _{231\mathrm{}n}+\mathrm{}+\eta _{23\mathrm{}1n}+\eta _{23\mathrm{}n1}=0.$$
(13)
We will call tensors satisfying the above equation gauge invariant (GI).
In order to rewrite the effective action as a matrix action, we had fixed, in section 3.2, the central term $`\mathrm{\Phi }`$ to $`K`$. We must then require by hand that the action (12) be independent of the choice of central term. This requirement will again be written as an algebraic identity involving the tensors $`\eta _{1\mathrm{}n}`$.
Let us then add a small central term $`\kappa _{ab}`$ to $`\widehat{F}_{ab}`$, and at the same time subtract the same $`\kappa _{ab}`$ from the effective metric $`\mathrm{\Gamma }_{ab}`$. The two effects must compensate each other, yielding a vanishing total variation of the action. Adding $`\kappa `$ to $`\widehat{F}`$ means that
$$i[\lambda ^a,\lambda ^b]i[\lambda ^a,\lambda ^b]+\mathrm{\Delta }^{ab},$$
(14)
with
$$\mathrm{\Delta }=\theta \kappa \theta .$$
Therefore, a term with $`n+2`$ coordinate functions $`\lambda ^a`$ will go into a term with $`n`$ functions $`\lambda ^a`$. More precisely, the variation of a term
$$\left(\lambda ^1\mathrm{}\lambda ^{n+2}\right)\eta _{1\mathrm{}n+2}$$
will be of the form
$$\left(\lambda ^1\mathrm{}\lambda ^n\right)\left(\overline{\mathrm{\Delta }}\eta \right)_{1\mathrm{}n}$$
where $`\left(\overline{\mathrm{\Delta }}\eta \right)_{1\mathrm{}n}`$ is again gauge invariant and depends on $`\eta _{1\mathrm{}n+2}`$ and on $`\mathrm{\Delta }^{ab}`$. We will show in the appendix (Lemma 9) that
$$\left(\overline{\mathrm{\Delta }}\eta \right)_{1\mathrm{}n}=\frac{i}{2}\mathrm{\Delta }^{ab}\left(\eta _{1\mathrm{}nab}+\eta _{1\mathrm{}anb}+\mathrm{}\right),$$
(15)
where $`\mathrm{}`$ indicates all the terms with the indices $`1,\mathrm{},n`$ in increasing order, and the two contracted indices $`a,b`$ in all possible positions with $`a`$ preceding $`b`$. Let us just note that, for $`n=2`$, the above result follows from (14), since a gauge invariant $`\eta _{ab}`$ is necessarily antisymmetric, and therefore $`\lambda ^a\lambda ^b\eta _{ab}=\frac{1}{2}[\lambda ^a,\lambda ^b]\eta _{ab}`$.
As noted previously, the variation (14) must be compensated by a corresponding change in the metric $`\mathrm{\Gamma }\mathrm{\Gamma }\kappa `$. In particular, in expression (12) this will affect both the measure of integration and the tensors $`\eta _{1\mathrm{}n}`$. The measure changes by
$$det{}_{}{}^{\frac{1}{2}}\mathrm{\Gamma }det{}_{}{}^{\frac{1}{2}}\mathrm{\Gamma }\left(1+\frac{1}{2}\kappa _{ab}\mathrm{\Gamma }^{ab}\right),$$
and the tensors by
$$\eta \eta \kappa _{ab}\frac{}{\mathrm{\Gamma }_{ab}}\eta .$$
Noting that $`\kappa _{ab}\mathrm{\Gamma }^{ab}=\gamma _{ab}\mathrm{\Delta }^{ab}`$ and that
$$\kappa _{ab}\frac{}{\mathrm{\Gamma }_{ab}}=(\gamma \mathrm{\Delta }\gamma )_{ab}\frac{}{\gamma _{ab}}$$
we then conclude that the variation of a term $`\left(\lambda ^1\mathrm{}\lambda ^n\right)\eta _{1\mathrm{}n}`$ will be
$$\left(\lambda ^1\mathrm{}\lambda ^n\right)\left(\overline{\delta }\eta \right)_{1\mathrm{}n}$$
where
$$\overline{\delta }=\frac{1}{2}\left(\gamma _{ab}\mathrm{\Delta }^{ab}\right)+(\gamma \mathrm{\Delta }\gamma )_{ab}\frac{}{\gamma _{ab}}$$
Let us now combine the two variations. To this end, recall that the sum of terms in (12) runs only over even values of $`n`$. In particular we will say that a tensor $`\eta _{1,\mathrm{},2L}`$ with $`2L`$ indices is of level $`L`$, and we will denote it with $`\eta ^L`$. The operator $`\overline{\mathrm{\Delta }}`$ lowers level by one, whereas $`\overline{\delta }`$ leaves the level invariant. In order to balance the two variations $`\overline{\mathrm{\Delta }}`$ and $`\overline{\delta }`$ and to have invariance under a change of the central term we must then have that<sup>3</sup><sup>3</sup>3We have checked invariance of the action under infinitesimal changes of $`\mathrm{\Phi }`$ around $`\mathrm{\Phi }=K`$. On the other hand, this is sufficient, since invariance under variation of the central term is a property of the structure of the action, property which is independent of the specific value of $`\mathrm{\Phi }`$.
$$\overline{\mathrm{\Delta }}\eta ^{L+1}=\overline{\delta }\eta ^L.$$
### 3.4 The Seiberg–Witten map following Jurco and Schupp
In section 3.2, we have written the action $`S`$ in terms of the functions $`\lambda ^a`$, which implicitly depend on the non–commutative gauge potential $`\widehat{A}=\widehat{A}_{SW}(A,\theta )`$. In order to analyze the independence of the action from the non–commutativity parameter $`\theta `$, it is convenient, as will become clear later, to rewrite $`S`$ in terms of the $`\theta `$–independent abelian potential $`A`$. To this end, we follow the analysis of Jurco and Schupp , whose work describes the Seiberg–Witten map $`\widehat{A}_{SW}`$ in an invariant way, which is best suited for our purposes. Most of this section is then nothing but a review of the ideas of , rewritten in the notation of this paper.
First recall that, on the manifold $`M`$, one can define, in a natural way, two distinct symplectic structures, defined by the two form $`K`$ and by the combination
$$\omega =K+F,$$
where $`F=dA`$ is the usual abelian field strength. Since $`F`$ is exact, the forms $`K`$ and $`\omega `$ define the same class in cohomology, and therefore, by Darboux’s lemma, there is a diffeomorphism $`\lambda :MM`$ such that
$$\lambda ^{}\omega =K.$$
(16)
Starting from the two symplectic structures $`K`$, $`\omega `$, one can, first of all, define the corresponding Poisson brackets
$$\{f,g\}_K=\theta ^{ab}_af_bg\{f,g\}_\omega =(\omega ^1)^{ab}_af_bg.$$
It is clear that the two brackets are related by the diffeomorphism $`\lambda `$. More precisely, for any two functions $`f,g`$, one has the trivial identity
$$\lambda ^{}\{f,g\}_\alpha =\{\lambda ^{}f,\lambda ^{}g\}_\theta $$
(17)
where $`\lambda ^{}f=f\lambda `$.
From the two symplectic structures $`K`$ and $`\omega `$ one can also construct, following Kontsevich , associated star–products
$$_K_\omega $$
In particular, $`_K`$ is nothing but the usual Moyal product, since in the coordinates $`x^a`$ the symplectic structure $`K`$ is constant. The product $`_\omega `$ is, on the other hand, the full product of Kontsevich, which is expressed in terms of a complicated diagrammatic expression involving derivatives of the Poisson structure $`\omega ^1`$, and for which there is an elegant path–integral expression, by Cattaneo and Felder . In what follows, we will not need the explicit form for $`_\omega `$. On the other hand, since $`K`$ and $`\omega `$ are related by diffeomorphism, it is a general result of Kontsevich that the two products $`_K`$ and $`_\omega `$ are equivalent. More precisely, there is a map $`T`$ defined on functions such that, for and functions $`f`$ and $`g`$,
$$T\left(f_\omega g\right)=Tf_KTg.$$
(18)
The above expression is the analogue of expression (17), and in fact one can show that<sup>4</sup><sup>4</sup>4We thank A. Cattaneo for pointing out that (19) is a simple consequence of formality, as defined in .
$$T=\lambda ^{}(1+\mathrm{}),$$
(19)
where $`\mathrm{}`$ are higher order terms (more precisely, if we replace $`K,\omega \frac{1}{\mathrm{}}K,\frac{1}{\mathrm{}}\omega `$, then the terms in $`\mathrm{}`$ are higher order in $`\mathrm{}`$).
Given these facts, one can define, in terms of $`T`$, the Seiberg–Witten map as follows
$$\lambda ^a=Tx^a=x^a+\theta ^{ab}\widehat{A}_b.$$
We need to check that the map $`\widehat{A}_{SW}`$ implicitly defined above maps gauge orbits of the abelian theory to gauge orbits of the non–commutative theory. On one side, it is clear that different abelian potentials $`A`$ which are gauge–equivalent do give the same map $`\widehat{A}`$, since $`T`$ is only defined in terms of the combination $`\omega `$, which is itself gauge–invariant. Moreover, on the non–commutative side, we note that the map $`T`$ defined in (18) is only defined up to transformations
$$Tf\mathrm{\Lambda }_KTf_K\mathrm{\Lambda }^1,$$
(20)
which leave (18) invariant<sup>5</sup><sup>5</sup>5The infinitesimal version of equation (20) is given by $`TfTf+[\rho ,Tf]_K`$. This change of $`T`$ is analogous to the fact that the map $`\lambda `$ is defined up to symplectomorphisms of the manifold $`(M,K)`$. In fact, if $`\chi :MM`$ is such that $`\chi ^{}K=K`$, then the composite map $`\lambda \chi `$ still satisfies (16). Recalling that symplectomorphisms are generated by Hamiltonian flows, the change $`\lambda ^{}\left(\lambda \chi \right)^{}`$ is given infinitesimally by $`\lambda ^{}f\lambda ^{}f+\{\rho ,\lambda ^{}f\}_K`$.. This in turn generates gauge transformations
$$\widehat{A}_a\widehat{A}_a+i\mathrm{\Lambda }_K_a\mathrm{\Lambda }^1+\mathrm{\Lambda }_K\widehat{A}_a_K\mathrm{\Lambda }^1,$$
therefore showing that the transformation $`\widehat{A}_{SW}`$ does map gauge orbits into gauge orbits.
### 3.5 Integration
In the previous section we have reviewed the Jurco–Schupp construction of the Seiberg–Witten map. In this section we wish to discuss some issues about integration of functions over $`M`$ which are closely related to the discussion in the previous section, and which will be important in our subsequent discussion.
Corresponding to the two symplectic structures $`K`$, $`\omega `$, one has two volume–forms on $`M`$, respectively $`d^Dxdet{}_{}{}^{\frac{1}{2}}K`$ and $`d^Dxdet{}_{}{}^{\frac{1}{2}}\omega `$, which are related by the map $`\lambda `$. More specifically, if $`f`$ is a generic function which vanishes at infinity, using the fact that $`\lambda ^{}\omega =K`$, it is immediate to show that
$$d^Dxdet{}_{}{}^{\frac{1}{2}}K\lambda ^{}f=d^Dxdet{}_{}{}^{\frac{1}{2}}\omega f.$$
Similarly, we may consider, recalling from (19) that $`T\lambda ^{}`$, the corresponding integral of $`Tf`$. We then have, in general, that
$$d^Dxdet{}_{}{}^{\frac{1}{2}}KTf=d^DxV(\omega )f,$$
(21)
where $`V(\omega )`$ is a volume element depending on $`\omega `$ and its derivatives, of the general form
$$V(\omega )=det{}_{}{}^{\frac{1}{2}}\omega \left(1+\mathrm{}\right),$$
where $`\mathrm{}`$ denotes, as in (19), higher order derivative corrections in $`\omega ^1`$ which vanish if $`\omega `$ is constant. Let me note that, since $`d^Dxdet{}_{}{}^{\frac{1}{2}}Kf_Kg=d^Dxdet{}_{}{}^{\frac{1}{2}}Kg_Kf`$, the ambiguity (20) in the definition of $`T`$ does not affect the definition (21) of $`V(\omega )`$, which really only depends on the symplectic structure $`\omega `$. Moreover, from the definition (18) of $`T`$, we have in general that
$$d^DxV(\omega )f_\omega g=d^DxV(\omega )g_\omega f.$$
The explicit form of $`V(\omega )`$ is not know. On the other hand we will show in the rest of the paper that some properties of $`V(\omega )`$ can be proven indirectly, and this will suffice for our purposes.
### 3.6 Back to the matrix action
In this section we use the results just discussed on the Seiberg–Witten map and on integration to rewrite the action (12)
$$\frac{1}{g_s}\underset{n\text{ even}}{}d^Dxdet{}_{}{}^{\frac{1}{2}}K\left(\lambda ^1_K\mathrm{}_K\lambda ^n\right)\eta _{1\mathrm{}n}$$
in an almost final form (we are now showing in the star–products $`_K`$ the explicit dependence on the symplectic structure). Using the facts that $`\lambda ^a=Tx^a`$ and that $`T\left(f_\omega g\right)=Tf_KTg,`$ one quickly sees that
$$\left(\lambda ^1_K\mathrm{}_K\lambda ^n\right)=T\left(x^1_\omega \mathrm{}_\omega x^n\right).$$
Using then equation (21) on integration one concludes that the action (12) can be rewritten as
$$S=\frac{1}{g_s}\underset{n\text{ even}}{}d^DxV(\omega )\left(x^1_\omega \mathrm{}_\omega x^n\right)\eta _{1\mathrm{}n}.$$
The above action is written almost exclusively in terms of the closed string parameters $`A,g+B`$ and $`g_s`$. Moreover it is explicitly a gauge invariant function of $`A`$, since the dependence on the gauge potential is uniquely through the gauge invariant expression $`\omega =K+F`$. On the other hand, the action $`S`$ above still depends on the parameter $`\theta `$, through the definition of $`\omega `$ – which effects $`_\omega `$ and $`V(\omega )`$ – and through the effective metric $`\gamma =g+BK`$ – which is the building block for the tensors $`\eta `$. However the action $`S`$ must be independent of $`\theta `$, and the analysis of this requirement will be the subject of the rest of the paper.
REMARK . We are now in a position to give some very intuitive arguments for the appearance of the full Kontsevich product in the effective action. The arguments which follow are vague and not precise. On the other hand, they provide a useful intuition, which, if made rigorous, could be of importance.
Let us start by recalling that the effective action can be considered as the partition function for an open string sigma model
$$S(A)DXe^{I_SI_A}$$
(22)
where
$`I_S`$ $`=`$ $`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle _\mathrm{\Sigma }}g_{ab}X^a\overline{}X^b`$
$`I_A`$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}B+{\displaystyle _\mathrm{\Sigma }}A`$
If we consider the naive limit $`\alpha ^{}0`$, the term $`I_S`$ dominates, and we should consider $`I_A`$ as a perturbation. On the other hand, we recall that, in , Seiberg and Witten consider the limit $`\alpha ^{},g_{ab}0`$, with $`g_{ab}/\alpha ^{}0`$. Therefore, in this case, the dominating term is $`I_A`$. We also note that
$$I_A=_\mathrm{\Sigma }\left(B+F\right)$$
(23)
and that the above is nothing but the Cattaneo–Felder model
$$\eta _adX^a+\frac{1}{2}\alpha ^{ab}(X)\eta _a\eta _b$$
in the special case of invertible Poisson structure $`\alpha ^{ab}`$, with $`\alpha ^1=B+F`$. In this case, one can integrate out the one–forms $`\eta _i`$, which appear quadratically, and recover (23). We recall that the perturbation theory of the Cattaneo–Felder model generates the Kontsevich graphs, which are the basis of the product $`_{B+F}`$. One then expects (in an undoubtedly vague way) to obtain effective actions based on the full Kontsevich product. Moreover one expects to obtain, among the various products considered in , the simplest one, defined using the harmonic angle map. Along the same lines one could also expand (22) in powers of $`I_S`$ and obtain an expansion like
$$DXe^{I_A}I_SI_Sd^DxV(\omega )[x^a,x^b]_\omega _\omega [x^c,x^d]_\omega g_{ac}g_{bd}$$
where the RHS above is nothing but the $`\widehat{F}^2`$ term in the action, which dominates in the $`\alpha ^{}0`$ limit of .
### 3.7 Behavior at infinity and cyclic tensors
We have seen that the general action describing the dynamics of gauge fields is of the form
$$\frac{1}{g_s}\underset{n\text{ even}}{}d^DxV(\omega )\left(x^1_\omega \mathrm{}_\omega x^n\right)\eta _{1\mathrm{}n}$$
(24)
Let us now proceed by first concentrating on a single term in the sum (24). In particular let us focus on the expression
$$\eta (x)=x^1_\omega \mathrm{}_\omega x^n\eta _{1\mathrm{}n}$$
The above clearly defined a function $`\eta `$ of the coordinates $`x`$, which is written in terms of $`\omega ^1`$ and its derivatives. We will assume throughout the paper that
$$F(x)0\mathrm{when}\text{ }x\mathrm{}$$
This implies that, for large $`x`$, the symplectic structure $`\omega K`$ becomes constant, and that the star–product $`_\omega `$ becomes the Moyal product $`_K`$ with respect to $`\theta `$. In general then, for $`x\mathrm{}`$, the function $`\eta (x)`$ is a polynomial of degree $`n`$ in the coordinates $`x^a`$. If we assume further that $`\eta _{1\mathrm{}n}`$ is a gauge invariant tensor, then we can quickly see that, again for $`x\mathrm{}`$,
$$\eta (x+\epsilon )\eta \left(x\right)=\epsilon ^1x^2_K\mathrm{}_Kx^n(\eta _{12\mathrm{}n}+\eta _{21\mathrm{}n}+\mathrm{})=0$$
and therefore the function $`\eta (x)`$ approaches a constant $`\eta _{\mathrm{}}`$ at infinity. Similarly, the volume form $`V(\omega )`$ converges to the constant $`det{}_{}{}^{\frac{1}{2}}K`$ as $`x\mathrm{}`$. It is then clear that the integral $`d^DxV(\omega )\eta `$ in general diverges unless $`\eta _{\mathrm{}}=0`$. In order to define the action properly we should replace
$$d^DxV(\omega )\eta d^DxV(\omega )\left[\eta \eta _{\mathrm{}}\right]$$
(25)
thereby eliminating the infinities coming from integration over an infinite world–volume. Let me note that, since $`T1=1`$, one has $`d^DxV(\omega )=d^Dxdet{}_{}{}^{\frac{1}{2}}K`$, and therefore the subtraction (25) is independent of $`F`$. Replacement (25) is then nothing but a constant addition to the action.
We will now show how the subtraction (25) can be achieved in an invariant way, without explicitly considering the behavior at infinity. First let us recall that, for functions $`f`$ and $`g`$ which vanish at infinity, we have that
$$d^DxV(\omega )f_\omega g=d^DxV(\omega )g_\omega f.$$
We are therefore tempted to say that the integral $`d^DxV(\omega )\left(x^1_\omega \mathrm{}_\omega x^n\right)`$ is invariant under cyclic permutations of the indices $`1,2,\mathrm{},n`$. This is, on the other hand, not quite correct, since the coordinate functions $`x^a`$ which enter in expression (24) clearly do not vanish for $`x\mathrm{}`$. Nonetheless let us, for the moment, blindly assume cyclicity of the integral. We may then substitute, in expression (24) for the action, the tensors $`\eta _{1\mathrm{}n}`$ with the cyclically symmetrized tensors
$$\tau _{1\mathrm{}n}=\frac{1}{n}\left(\eta _{1\mathrm{}n}+\mathrm{cyc}_{1\mathrm{}n}\right),$$
(26)
where $`\mathrm{cyc}_{1\mathrm{}n}`$ denotes the sum over cyclic permutations of the indices $`1,\mathrm{},n`$. Gauge invariance of the tensors $`\eta `$ then translates into the following algebraic property satisfied by the tensors $`\tau `$
$$\tau _{123\mathrm{}n}+\tau _{213\mathrm{}n}+\tau _{231\mathrm{}n}+\mathrm{}+\tau _{23\mathrm{}1n}=0.$$
(27)
Note that the above expression is very similar to (13), with the only difference that the moving index $`1`$ runs only over the cyclically independent orderings, and therefore the last term in (13) is absent in (27). Tensors which satisfy the above relation will be called cyclic gauge invariant (CGI). We leave the proof of (27) to the appendix (Lemma 2).
We may now consider, similarly to the previous analysis, the function
$$\tau (x)=x^1_\omega \mathrm{}_\omega x^n\tau _{1\mathrm{}n}$$
and in particular its behavior at infinity. As shown in the appendix (Lemma 3), for $`n`$ even (which is the case relevant to equation (24)) one has that
$$\tau (x)0$$
for $`x\mathrm{}`$. Therefore the integral
$$d^DxV(\omega )\tau $$
is well defined. Moreover we will show in the appendix (Lemma 5) that, generically, one has that
$$d^DxV(\omega )\tau =d^DxV(\omega )\left[\eta \eta _{\mathrm{}}\right]$$
so that we have lost nothing by assuming cyclicity<sup>6</sup><sup>6</sup>6Let me note that, although the functions $`\tau `$ and $`\eta \eta _{\mathrm{}}`$ have the same integral, and therefore define the same functional of $`A`$, one has in general that $`\tau \eta \eta _{\mathrm{}}`$.. In fact, using cyclic gauge invariant tensors, the subtraction which was needed in (25) in order to properly define the action $`S(A)`$ is automatically incorporated into the formalism. We will therefore consider, from now on, the final form of the action
$$S=\frac{1}{g_s}\underset{n\text{ even}}{}d^DxV(\omega )\left(x^1_\omega \mathrm{}_\omega x^n\right)\tau _{1\mathrm{}n}$$
(28)
where the tensors $`\tau `$ are cyclic gauge invariant.
We have seen in section 3.3 that, in order for the tensors $`\eta `$ to define an action, they had to be gauge invariant and they had to satisfy $`\overline{\mathrm{\Delta }}\eta ^{L+1}=\overline{\delta }\eta ^L`$. These two properties impose restrictions on the cyclically symmetrized tensors $`\tau `$. Gauge invariance of the $`\eta `$’s implies cyclic gauge invariance of the $`\tau `$’s. The equation $`\overline{\mathrm{\Delta }}\eta ^{L+1}=\overline{\delta }\eta ^L`$ implies a similar equation for the $`\tau `$’s, which we now describe.
Consider a gauge invariant tensor $`\eta _{1\mathrm{}n+2}`$. We can construct, given $`\mathrm{\Delta }^{ab}`$ and equation (15), the tensor $`\left(\overline{\mathrm{\Delta }}\eta \right)_{1\mathrm{}n}`$, which is also gauge invariant for any choice of $`\mathrm{\Delta }^{ab}`$. We may then consider the cyclically symmetric combination $`g_{1\mathrm{}n}=\frac{1}{n}\left(\overline{\mathrm{\Delta }}\eta _{1\mathrm{}n}+\mathrm{cyc}_{1\mathrm{}n}\right)`$, which will be a cyclic gauge invariant tensor. At first sight the tensor $`g_{1\mathrm{}n}`$ is a function of the original tensor $`\eta _{1\mathrm{}n+2}`$, but, as we will prove in the appendix (Lemma 10), it is actually just a function of the cyclically symmetrized tensor $`\tau _{1\mathrm{}n+2}`$. We will then denote the tensor $`g_{1\mathrm{}n}`$ with $`\left(\overline{\mathrm{\Delta }}\tau \right)_{1\mathrm{}n}`$, where
$$\left(\overline{\mathrm{\Delta }}\tau \right)_{1\mathrm{}n}=\frac{i}{2}\left(\frac{n+2}{n^2}\right)\mathrm{\Delta }^{ab}\left[n\tau _{1\mathrm{}nab}+\left(n1\right)\tau _{1\mathrm{}anb}+\mathrm{}+0\tau _{a1\mathrm{}nb}\right]+\mathrm{cyc}_{1\mathrm{}n}$$
Nothing on the other hand needs to be altered in the definition of the operator
$$\overline{\delta }=\frac{1}{2}\left(\gamma _{ab}\mathrm{\Delta }^{ab}\right)+(\gamma \mathrm{\Delta }\gamma )_{ab}\frac{}{\gamma _{ab}}$$
which commutes with the symmetrization (26). We then have the requirement on the tensors $`\tau `$
$$\overline{\mathrm{\Delta }}\tau ^{L+1}=\overline{\delta }\tau ^L.$$
(29)
EXAMPLE. Let us compute, as an important example, the first non–vanishing tensor $`\tau ^2`$. If we consider the expansion of the Born–Infeld action (with $`\mathrm{\Omega }=\widehat{F}K`$)
$$\sqrt{det(\mathrm{\Gamma }+\mathrm{\Omega })}1\frac{1}{2}\mathrm{\Gamma }^{ab}\mathrm{\Omega }_{ab}+\frac{1}{4}\mathrm{\Gamma }^{ac}\mathrm{\Gamma }^{db}\mathrm{\Omega }_{ab}\mathrm{\Omega }_{cd}+\frac{1}{8}\mathrm{\Gamma }^{ab}\mathrm{\Omega }_{ab}\mathrm{\Gamma }^{cd}\mathrm{\Omega }_{cd}$$
we can quickly see that
$`\eta _{12}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left(\gamma _{12}\gamma _{21}\right)`$
$`\eta _{1234}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\gamma _{13}\gamma _{42}\gamma _{23}\gamma _{41}+\gamma _{31}\gamma _{24}\gamma _{32}\gamma _{14}\right)`$
$`{\displaystyle \frac{1}{8}}\left(\gamma _{12}\gamma _{34}\gamma _{21}\gamma _{34}+\gamma _{12}\gamma _{43}\gamma _{21}\gamma _{43}\right).`$
We may then compute the cyclically symmetrized tensors $`\tau `$. Clearly $`\tau _{12}=0`$. A simple computation also shows that
$$\tau _{1234}=\frac{1}{4}g_{12}g_{34}+\frac{1}{4}g_{14}g_{23}\frac{1}{2}g_{13}g_{24},$$
(30)
where, we recall,
$$g_{ab}=\frac{1}{2}\left(\gamma _{ab}+\gamma _{ba}\right)$$
is the symmetric part of the tensor $`\gamma _{ab}`$. One can also check, given (30), that
$$\overline{\mathrm{\Delta }}\tau ^2=0,$$
which is consistent with (29) and the fact that $`\tau ^1=0`$.
## 4 Operator Description
In this section we leave momentarily the analysis of the effective action $`S`$, and we develop some formal tools which will allow us both to rewrite the various equations in a more compact and natural way, and also to tackle the problem of the independence of the action $`S`$ on the parameter $`\theta `$.
First we analyze, within a general framework, the description of the action using operators. We then study a simple $`2`$–dimensional example, which is not directly relevant to our more general situation, but which is completely tractable and which will be a useful frame of reference in discussing the general case. We then move on to the situation most relevant for this paper, and discuss, in that case, the generalization of the results obtained in the $`2`$–dimensional setting.
### 4.1 Operator representation of star products
Let us consider first a flat symplectic structure $`K`$. We introduce a set of operators $`J^a`$ with commutation relations $`[J^a,J^b]=i\theta ^{ab}`$, which can be represented on the Hilbert space $`=L^2(^{D/2})`$ as linear combinations of the standard $`p,q`$ operators in quantum mechanics (recall that we are assuming $`\theta `$ invertible). To any function $`f`$ on phase space $`M=^D`$ we can associate, using Weyl ordering, an operator $`Q_K(f)`$ acting on $``$. It is then well known that, if $`f`$ and $`g`$ are two generic functions, then $`Q_K(f)Q_K(g)=Q_K(f_Kg)`$. Moreover, if $`f`$ vanishes at infinity, one also has (up to an overall constant $`\left(2\pi \right)^{D/2}`$ which can be, for example, reabsorbed in the definition of the trace) that $`\mathrm{Tr}(Q_K(f))=d^Dx`$ $`det{}_{}{}^{\frac{1}{2}}Kf`$. We will call $`Q`$ a quantization map.
Consider now a general symplectic structure $`\omega `$. Since any two symplectic structures on $`M`$ are related by a diffeomorphism, one can follow section 3.4 and find a map $`T`$ on the space of functions such that, for general $`f,g`$, one has $`T\left(f_\omega g\right)=Tf_KTg`$. One may then define a new quantization map $`Q_\omega `$, related to the symplectic structure $`\omega `$, by the following relation
$$Q_\omega (f)=Q_K(Tf).$$
(31)
It is then simple to show that
$$Q_\omega (f)Q_\omega (g)=Q_\omega (f_\omega g)$$
(32)
and that
$$\mathrm{Tr}(Q_\omega f))=d^Dx\text{ }V(\omega )f.$$
(33)
Let us note that, for any fixed symplectic form $`\omega `$, the map $`Q_\omega `$ is actually only defined up to conjugation. Recall first that the map $`T`$ is defined up to a redefinition of the form $`Tf\stackrel{~}{T}f=\mathrm{\Lambda }_KTf_K\mathrm{\Lambda }^1`$. Using $`\stackrel{~}{T}`$ in equation (31), and letting $`\chi =T^1\mathrm{\Lambda }`$, we obtain a new map $`\stackrel{~}{Q}_\omega `$ which reads, in terms of the original $`Q_\omega `$,
$`\stackrel{~}{Q}_\omega (f)`$ $`=`$ $`Q_\omega (\chi )Q_\omega (f)Q_\omega ^1(\chi )`$
$`=`$ $`Q_\omega (\chi _\omega f_\omega \chi ^1).`$
We now use this notation to rewrite the action $`S`$ in a compact and invariant way. First define the operators
$$X^a=Q_K(\lambda ^a)=Q_\omega (x^a).$$
Then the action (28) can then be compactly written as
$$S=\frac{1}{g_s}\underset{n\text{ even}}{}\mathrm{Tr}\left(X^1\mathrm{}X^n\right)\tau _{1\mathrm{}n}.$$
The ambiguity (4.1) is reflected in a possible redefinition $`X^aOX^aO^1`$, which on the other hand does not affect the action.
The action $`S`$ written above implicitly depends on a specific choice of non–commutativity parameter $`\theta `$. The dependence is two–fold. On one hand the tensors $`\tau `$ are built starting from the metric $`\gamma `$, which linearly depends on $`K`$. On the other hand, the parameter $`K`$ enters into the definition of the symplectic structure $`\omega `$, and therefore it implicitly determines the operators $`X^a=Q_\omega (x^a)`$. It is then clear that we need to understand the variation of the quantization map $`Q_\omega `$, when we add to $`\omega _{ab}`$ a constant antisymmetric matrix $`\mathrm{\Delta }_{ab}`$. This is the subject of the next two sections. In particular, in the next section, we analyze this problem within a simple two–dimensional model, related to the general framework which we developed above. In this two–dimensional model the variation of $`Q_\omega `$ for $`\omega \omega +\mathrm{\Delta }`$ can be completely analyzed. Moreover the solution will give us the correct ansatz to tackle the general problem.
### 4.2 A simple $`2`$–dimensional example
The general framework of this section follows closely . We consider the space $`𝒱`$ of complex functions on the complex plane $``$, and the subspace $`𝒱`$ of holomorphic functions. We then make $`𝒱`$ into a Hilbert space by choosing a real positive function $`C`$ on the complex plane and by letting the inner product of two function $`\psi ,\varphi 𝒱`$ be given by
$$\psi |\varphi =d^2zC\overline{\psi }\varphi .$$
Associated to $`C`$ we have a natural symplectic form $`i\omega dzd\overline{z}`$ on $``$ with
$$\omega =\overline{}\mathrm{ln}C\text{.}$$
One may also consider the orthogonal projection
$$\pi :𝒱,$$
which clearly depends on the choice of inner product on $`𝒱`$, and therefore on $`C`$. Given a generic function $`f`$, we may then consider the corresponding operator
$$Q_\omega (f):$$
defined by
$$Q_\omega (f)\eta =\pi f\eta .$$
for $`\eta `$. In words, the operator $`Q_\omega (f)`$ first multiplies pointwise by $`f`$ – which is not assumed to be holomorphic – and then extracts the holomorphic part of the resulting function using the projector $`\pi `$. It is shown in that, given two functions $`f,g`$,
$`Q_\omega (f)Q_\omega (g)`$ $`=`$ $`Q_\omega (f_\omega g)`$ (35)
$`\mathrm{Tr}_{}(Q_\omega (f))`$ $`=`$ $`{\displaystyle d^2zV(\omega )f}`$
where $`_\omega `$ is a holomorphic star product ($`f_\omega g=fg`$ if either $`\overline{f}`$ or $`g`$ are holomorphic) related to $`\omega `$ and $`V(\omega )=\omega (1+\mathrm{})`$. The operator $`Q_\omega `$ actually depends on $`C`$, and not simply on $`\omega `$. On the other hand, a change of $`C`$ which leaves $`\omega `$ invariant changes the operators $`Q_\omega (f)`$ by conjugation, exactly as in (4.1).
We may now consider the variation $`\omega \omega +\mathrm{\Delta }`$, with $`\mathrm{\Delta }`$ an infinitesimal constant. This corresponds to
$$C\stackrel{~}{C}=Ce^{z\overline{z}\mathrm{\Delta }}.$$
We need to understand the change in $`\pi `$. Let then $`|`$ denote the original inner product with $`C`$, and let $`|n`$ be an orthonormal basis for $``$. Then $`\pi =_n|nn|`$. It is easy to show that the vectors
$$|\stackrel{~}{n}=|n+\frac{\mathrm{\Delta }}{2}\underset{m}{}|mm|z\overline{z}|n$$
satisfy $`\stackrel{~}{n}|e^{z\overline{z}\mathrm{\Delta }}|\stackrel{~}{m}=\delta _{n,m}`$ (to first order in $`\mathrm{\Delta }`$) and that the new projection $`\stackrel{~}{\pi }`$ related to $`\stackrel{~}{C}`$ is given by $`\stackrel{~}{\pi }=_n|\stackrel{~}{n}\stackrel{~}{n}|e^{z\overline{z}\mathrm{\Delta }}`$. Using the explicit expression for $`|\stackrel{~}{n}`$, one can then show that
$$\stackrel{~}{\pi }=\pi +\mathrm{\Delta }\left(\pi z\overline{z}\pi \pi z\overline{z}\right).$$
Let now $`f`$ be a generic function and $`F=Q_\omega (f)=\pi f`$. Noting that $`\pi z=z`$ one can show that
$`Q_{\omega +\mathrm{\Delta }}(f)`$ $`=`$ $`\stackrel{~}{\pi }f=F+\mathrm{\Delta }(\pi \overline{z}z\pi f\pi \overline{z}fz)`$
$`=`$ $`F+\mathrm{\Delta }(\overline{Z}ZFQ_\omega (\overline{z}f)Z)`$
where $`Z=Q_\omega (z)`$ and $`\overline{Z}=Q_\omega (\overline{z})`$. Using the fact that the star product is holomorphic and that $`Q_\omega (\overline{z}f)=Q_\omega (\overline{z}_\omega f)=\overline{Z}F`$, we arrive at the result
$$Q_{\omega +\mathrm{\Delta }}(f)=F+\mathrm{\Delta }(\overline{Z}ZF\overline{Z}FZ).$$
The exact form of the above equation depends on the specific quantization model we chose to analyze. On the other hand the general lesson that should be drawn is that the variation $`Q_{\omega +\mathrm{\Delta }}(f)Q_\omega (f)`$ contains $`F`$ and two powers of the coordinate operators $`X^a=Q_\omega (x^a)`$, with some ordering. We will use this intuition in the next section to compute $`Q_{\omega +\mathrm{\Delta }}(f)`$ in the setting of section 4.1.
### 4.3 The quantization map $`Q_{\omega +\mathrm{\Delta }}`$ in the general case
We have seen, from the previous example, that, if $`f`$ is a generic function, then the variation $`Q_{\omega +\mathrm{\Delta }}(f)Q_\omega (f)`$ is proportional to $`\mathrm{\Delta }_{ab}`$ and to the product of
$$F=Q_\omega (f)$$
and of two coordinate operators $`X^a=Q_\omega (x^a)`$ in some specific ordering. More precisely, the operator $`Q_{\omega +\mathrm{\Delta }}(f)`$ must be equal to $`F+\frac{i}{4}\mathrm{\Delta }_{ab}(aX^aX^bF+bFX^aX^b+cX^aFX^b)`$, for some choice of the coefficients $`a,b,c`$. It is shown in the appendix (Lemma 6) that, in the case in which the underlying star–product is that of Kontsevich, the correct coefficients are $`a=b=1`$, and $`c=2`$.
We have then the basic relation
$$Q_{\omega +\mathrm{\Delta }}(f)=F+\frac{i}{4}\mathrm{\Delta }_{ab}\left(X^aX^bF+FX^aX^b2X^aFX^b\right).$$
(36)
The above can be alternatively rewritten as
$$Q_{\omega +\mathrm{\Delta }}(f)=Q_\omega (f+Rf).$$
(37)
where
$$Rf=\frac{i}{4}\mathrm{\Delta }_{ab}\left(x^a_\omega x^b_\omega f+f_\omega x^a_\omega x^b2x^a_\omega f_\omega x^b\right).$$
(38)
As a consequence of the above facts, we have the following two results. First of all combining (37) and (32) we obtain
$$f_{\omega +\mathrm{\Delta }}g=f_\omega g\frac{i}{2}\mathrm{\Delta }_{ab}[x^a,f]_\omega _\omega [x^b,g]_\omega .$$
(39)
Also, taking the trace of (36) and using (33) we deduce that, for functions $`f`$ vanishing at infinity,
$$d^Dx\text{ }V(\omega +\mathrm{\Delta })f=d^DxV(\omega )\left[f+Rf\right].$$
(40)
REMARK . The fact that the variation of the star–product $`_\omega `$ under the change $`\omega \omega +\mathrm{\Delta }`$ is given by an expression involving a quadratic combination (39) of the coordinate functions $`x^a`$ can also be understood intuitively using the Cattaneo–Felder model. As in the remark in section 3.6, the argument is very vague, but it would be very useful to make it rigorous.
The star product $`f_\omega g`$ is given by the disk expectation value $`f\left(X\left(0\right)\right)g\left(X\left(1\right)\right)`$, with weight $`DX\mathrm{exp}\left(_\mathrm{\Sigma }\omega \right)`$. Therefore, under the change $`\omega \omega +\mathrm{\Delta }`$, one has that $`f_{\omega +\mathrm{\Delta }}gf_\omega g`$ is given by $`_\mathrm{\Sigma }\mathrm{\Delta }f\left(X\left(0\right)\right)g\left(X\left(1\right)\right)`$. But $`\mathrm{\Delta }=\frac{1}{2}\mathrm{\Delta }_{ab}dX^adX^b`$, thus giving a quadratic expression in the coordinate functions.
## 5 Invariance of the Action
### 5.1 Invariance under a change in $`\theta `$ and the basic commutator
We now have all the tools that we need to tackle our main problem. Let me first recall were we stand. The action is given by
$$S=\frac{1}{g_s}\underset{n\text{ even}}{}\mathrm{Tr}\left(X^1\mathrm{}X^n\right)\tau _{1\mathrm{}n}.$$
(41)
where $`X^a=Q_\omega (x^a)`$ and the tensors $`\tau `$ are cyclic gauge invariant tensors built from $`\gamma `$. The tensors $`\tau `$ satisfy the consistency condition $`\overline{\delta }\tau ^L=\overline{\mathrm{\Delta }}\tau ^{L+1}`$, where $`\tau ^L`$ denotes the tensor at level $`L`$, with $`2L`$ indices.
The action $`S`$ depends on a specific choice of non–commutativity parameter $`\theta `$. The dependence is two–fold, through $`\gamma =g+BK`$ and through $`\omega =K+F`$. We have argued in previous sections that the total dependence on $`\theta `$ should vanish, and therefore the two variations of the action under a change of $`K`$ should compensate each other and sum to zero. This clearly imposes additional restrictions on the possible forms of the tensors $`\tau `$, which we now analyze.
Let us start by considering an infinitesimal variation of $`K`$ given by
$$KK+\mathrm{\Delta },$$
where $`\mathrm{\Delta }_{ab}`$ is an arbitrary antisymmetric constant matrix. The metric $`\gamma `$ then changes as follows
$$\gamma \gamma \mathrm{\Delta }$$
therefore implying a change in the tensors $`\tau `$ given by
$$\tau \tau \delta \tau ,$$
where $`\delta `$ is the differential operator defined by
$$\delta =\mathrm{\Delta }_{ab}\frac{}{\gamma _{ab}}.$$
On the other hand, the symplectic structure $`\omega `$ changes by a constant term
$$\omega \omega +\mathrm{\Delta }.$$
Then, as discussed in the previous section, the coordinate operators $`X^c`$ change as
$$X^cX^c+\frac{i}{4}\mathrm{\Delta }_{ab}\left(X^aX^bX^c+X^cX^aX^b2X^aX^cX^b\right).$$
(42)
We can now discuss the variation of the term $`\mathrm{Tr}\left(X^1\mathrm{}X^{n2}\right)\tau _{1\mathrm{}n2}`$. It will consists of two parts, coming from the variation of the tensor $`\tau `$ and from the variation of the coordinate functions $`X^a`$. The first part is simply
$$\mathrm{Tr}\left(X^1\mathrm{}X^{n2}\right)\left(\delta \tau \right)_{1\mathrm{}n2}.$$
The second will involve the trace of $`n`$ coordinates, and will be of the general form
$$\mathrm{Tr}\left(X^1\mathrm{}X^n\right)\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}n},$$
where the tensor $`\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}n}`$ is built from $`\tau _{1\mathrm{}n2}`$ and from $`\mathrm{\Delta }_{ab}`$. Applying equation (42) to the expression $`\mathrm{Tr}\left(X^1\mathrm{}X^{n2}\right)`$ and rearranging cyclically under the trace it is easy to show that
$$\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}n}=\frac{i}{2}\left(\frac{n2}{n}\right)\left(\mathrm{\Delta }_{12}\tau _{345\mathrm{}n}\mathrm{\Delta }_{13}\tau _{245\mathrm{}n}\right)+\mathrm{cyc}_{1\mathrm{}n}$$
(43)
First we note that the above tensor $`\mathrm{\Delta }\tau `$ is cyclic gauge invariant for any choice of $`\mathrm{\Delta }_{ab}`$. This fact is proved in the appendix (Lemma 7). Moreover, the fact that both $`\tau `$ and $`\mathrm{\Delta }\tau `$ are cyclic gauge invariant is crucial in a more careful derivation of (43). In fact, the use of cyclic symmetry under the trace is formally correct, and does give the correct answer. On the other hand, one needs to check that the formal manipulations can be justified, since the coordinate functions do not vanish at infinity, and therefore one might forget important boundary terms. The detailed proof of equation (43) is again given in Lemma 7.
We can then finally state the main algebraic equation which must be satisfied by the tensors $`\tau `$ in order for the action (41) to be invariant under changes of $`\theta `$. Again indicating with $`L`$ the level of a tensor $`\tau _{1,\mathrm{},2L}`$ with $`2L`$ indices, we have the basic equation
$$\mathrm{\Delta }\tau ^L=\delta \tau ^{L+1}$$
(44)
which must be considered together with the equation
$$\overline{\mathrm{\Delta }}\tau ^{L+1}=\overline{\delta }\tau ^L$$
(45)
previously analyzed.
The above two equations involve the basic operators $`\mathrm{\Delta },\delta ,\overline{\mathrm{\Delta }},\overline{\delta }`$, which in general depend on two distinct antisymmetric matrices $`\mathrm{\Delta }_{ab}`$ and $`\mathrm{\Delta }^{ab}`$. Without loss in generality we can assume that
$$\mathrm{\Delta }_{ac}\mathrm{\Delta }^{cb}=\delta _a^b.$$
In this case one has simple relations between the operators $`\delta ,\overline{\delta }`$ and $`\mathrm{\Delta },\overline{\mathrm{\Delta }}`$, which we now describe. First we introduce the operator $`N`$ which simply counts the number of indices of a tensor, and which is defined by
$$N\tau _{1\mathrm{}n}=n\tau _{1\mathrm{}n}.$$
It is then easy to see that, since $`\tau _{1\mathrm{}2L}`$ is built from $`L`$ copies of $`\gamma `$, one has that
$$2\gamma _{ab}\frac{}{\gamma _{ab}}=N.$$
Using the above relation one then discovers quickly that
$$[\delta ,\overline{\delta }]=N\frac{D}{2},$$
where, we recall, $`D`$ is the dimension of space–time.
The relation between $`\mathrm{\Delta }`$ and $`\overline{\mathrm{\Delta }}`$ requires, on the other hand, a very long a quite technical analysis, which we leave to the appendix. Fortunately, though, the answer is very simple, and completely parallel to the above results. In fact, as shown in Lemma 12, one has that
$$[\overline{\mathrm{\Delta }},\mathrm{\Delta }]=N\frac{D}{2}.$$
We see that the structure of the seemingly different pairs of operators $`\overline{\mathrm{\Delta }},\mathrm{\Delta }`$ and $`\delta ,\overline{\delta }`$ is actually very similar and compatible, and we will use the above results heavily in the next section to solve the invariance equations for the simple case $`D=2`$.
Let us conclude this section by discussing the general form of the solution of equations (44) and (45). We introduce the concept of lowest level tensor, by which we mean a CGI tensor $`\rho `$ which satisfies
$$\overline{\mathrm{\Delta }}\rho =\delta \rho =0.$$
(46)
A general solution of (44) and (45) will then consist of a lowest level tensor $`\tau ^P`$ at level $`P`$, together with tensors $`\tau ^L`$ at higher levels $`L>P`$, which are required in order to obtain an invariant action. At each level $`L`$, $`\tau ^L`$ is determined using (44) and (45) in terms of the tensors of lower level, up to an addition $`\tau ^L\tau ^L+\rho `$ of a lowest level tensor $`\rho `$. Let us suppose that we can, given a $`\tau ^P`$ satisfying (46), construct, in a canonical way, a tower of tensors $`\tau ^L`$, $`L>P`$ so that (44) and (45) hold. We will then call the full set $`\left\{\tau ^L\right\}_{LP}`$ the invariant block generated by $`\tau ^P`$. The above discussion then shows that a solution of (44) and (45) is, in general, a linear combination of invariant blocks. We will see in the next section that, in the simple case $`D=2`$, we will indeed be able to construct canonically invariant blocks starting from generic lowest level tensors.
NOTE. Recall, from the example in section 3.7, that $`\tau ^2`$ is given by equation (30) and satisfies $`\overline{\mathrm{\Delta }}\tau ^2=0`$. It is trivial to check that also $`\delta \tau ^2=0`$ holds, since $`\tau ^2`$ depends only on the symmetric part of $`\gamma _{ab}`$. Therefore $`\tau ^2`$ is a lowest level state, as is natural to expect.
To conclude, let us comment on the question of uniqueness of the solution of the recursion equations (44) and (45). As we just discussed, we do not expect the solution to be unique, since general solutions are in one to one correspondence with lowest level tensors. An appropriate way of thinking about (44) and (45) is probably by analogy with general relativity. In that case, one is free to write actions of different type, subject only to the general principle of covariance under diffeomorphisms on the underlying space–time manifold. In a similar way, equations (44) and (45) imply that the action not only must be invariant under reparametrizations of the world–volume of the brane, but must also be invariant under a more general set of transformations, parametrized by changes in $`\theta `$. It would then be of great practical importance to have an explicitly covariant notation, for which invariance under (44) and (45) is manifest.
### 5.2 Summary of the basic results
We now summarize, in a compact but self–contained way, the results of this paper. Invariant actions are described by tensors $`\tau ^L`$ for $`L2`$ such that
1. The tensor $`\tau ^L`$ has $`2L`$ indices and is built from a basic matrix $`\gamma _{ab}`$.
2. The tensors $`\tau ^L`$ are cyclic gauge invariant – i.e. are cyclic tensors which satisfy the algebraic relation
$$\tau _{123\mathrm{}n}+\tau _{213\mathrm{}n}+\tau _{231\mathrm{}n}+\mathrm{}+\tau _{23\mathrm{}1n}=0.$$
3. For any choice of antisymmetric matrix $`\mathrm{\Delta }_{ab}`$ (with $`\mathrm{\Delta }^{ab}`$ indicating the inverse of $`\mathrm{\Delta }_{ab}`$) the tensors $`\tau ^L`$ satisfy the basic relation
$`\mathrm{\Delta }\tau ^L`$ $`=`$ $`\delta \tau ^{L+1}`$
$`\overline{\mathrm{\Delta }}\tau ^{L+1}`$ $`=`$ $`\overline{\delta }\tau ^L`$
where the differential operators $`\delta ,\overline{\delta }`$ are given by
$`\delta `$ $`=`$ $`\gamma _{ab}{\displaystyle \frac{}{\gamma _{ab}}}`$
$`\overline{\delta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\gamma _{ab}\mathrm{\Delta }^{ab}\right)+\left(\gamma \mathrm{\Delta }\gamma \right)_{ab}{\displaystyle \frac{}{\gamma _{ab}}}`$
and the algebraic operators $`\mathrm{\Delta },\overline{\mathrm{\Delta }}`$ are given by
$$\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}n}=\frac{i}{2}\left(\frac{n2}{n}\right)\left(\mathrm{\Delta }_{12}\tau _{345\mathrm{}n}\mathrm{\Delta }_{13}\tau _{245\mathrm{}n}\right)+\mathrm{cyc}_{1\mathrm{}n}$$
and
$$\left(\overline{\mathrm{\Delta }}\tau \right)_{1\mathrm{}n}=\frac{i}{2}\left(\frac{n+2}{n^2}\right)\mathrm{\Delta }^{ab}\left[n\tau _{1\mathrm{}nab}+\left(n1\right)\tau _{1\mathrm{}anb}+\mathrm{}+0\tau _{a1\mathrm{}nb}\right]+\mathrm{cyc}_{1\mathrm{}n}$$
The operators $`\mathrm{\Delta },\overline{\mathrm{\Delta }},\delta ,\overline{\delta }`$ satisfy the commutation relations
$$[\overline{\mathrm{\Delta }},\mathrm{\Delta }]\tau ^L=[\delta ,\overline{\delta }]\tau ^L=\left(2L\frac{D}{2}\right)\tau ^L.$$
## 6 A Solution in $`2`$ Dimensions
In this section we describe a general constructive solution of the invariance equations (44) and (45) in dimension $`D=2`$. This is clearly a toy model, since gauge fields have no propagating degrees of freedom in $`2`$ dimensions. On the other hand the solution is still highly non–trivial, and it exhibits many of the features which are expected to be present in the general $`D`$–dimensional case.
### 6.1 The general strategy
It is convenient, in two dimensions, to use on $`M`$ complex coordinates
$$z=x^1+ix^2\overline{z}=x^1ix^2$$
with a hermitian metric
$$g_{z\overline{z}}=g$$
and $`g_{zz}=g_{\overline{z}\overline{z}}=0`$. Similarly, the antisymmetric tensors $`B_{ab}`$ and $`K_{ab}`$ have a single independent component
$$B_{z\overline{z}}=iBK_{z\overline{z}}=iK$$
The tensor $`\gamma _{ab}`$ is then given by a single complex number
$`\gamma _{z\overline{z}}`$ $`=`$ $`g+i(BK)=x+iy`$
$`\gamma _{\overline{z}z}.`$ $`=`$ $`\overline{\gamma _{z\overline{z}}}=xiy`$
Finally we will use $`Z,\overline{Z}`$ for the operators which correspond to the coordinates $`z,\overline{z}`$ under the map $`Q_\omega `$.
We adopt, in this section, a notation which is not well suited for the general $`D`$–dimensional case treated in the remainder of the paper, but which is more economical in the present setting. Consider a general term of level $`L`$ in the action
$$\tau _{1\mathrm{}2L}\mathrm{Tr}\left(X^1\mathrm{}X^{2L}\right)$$
and let the various indices $`1,2,\mathrm{},2L`$ run over their possible values $`z,\overline{z}`$. We obtain a sum consisting of traces of monomials in $`Z,\overline{Z}`$, multiplied by polynomials of degree $`L`$ in $`x,y`$. In particular, the monomials under the trace satisfy the following two properties
1. They are constructed with $`L`$ coordinates $`Z`$ and $`L`$ coordinates $`\overline{Z}`$.
2. Monomials which differ only by a cyclic permutation of the coordinate operators $`Z,\overline{Z}`$ should be considered, as we recall from section 3.7, as identical.
We call objects which satisfy ($`1`$) and ($`2`$) cyclic words of level $`L`$ – or simply words – and we will denote with $`W_L`$ the space of their linear combinations (for example, for $`L=2`$, the space $`W_2`$ is spanned by the two words $`ZZ\overline{Z}\overline{Z}`$ and $`Z\overline{Z}Z\overline{Z}`$). Also we let $`P_L`$ be the space of polynomials in $`x,y`$ of degree $`L`$.
Among the possible cyclic words in $`W_L`$, we must consider the subspace
$$G_LW_L$$
of cyclic gauge invariant combinations. Following once more the discussion of section 3.7, we define the space $`G_L`$ as follows. First introduce canonical creation and annihilation operators $`a`$ and $`a^{}`$, which satisfy $`[a,a^{}]=1`$, and let $`O`$ be the space of polynomials in $`a,a^{}`$. Consider then a map
$$r:W_LO$$
defined by taking a word $`w`$ in $`W_L`$ and by constructing the operator $`r(w)`$ by replacing $`Z,\overline{Z}`$ with $`a,a^{}`$, and then by summing over the possible cyclic permutations. For example, we associate to the word $`w=\overline{Z}\overline{Z}ZZ`$ the operator<sup>7</sup><sup>7</sup>7Note that the sum over cyclic permutations is crucial in order to have a well–defined map $`r`$, since $`W_L`$ consists of words defined only up to cyclic permutation of the letters $`Z,\overline{Z}`$. In the above example, the same word $`w`$ can be equally represented by any of the permutations $`w=\overline{Z}\overline{Z}ZZ=\overline{Z}ZZ\overline{Z}=ZZ\overline{Z}\overline{Z}=Z\overline{Z}\overline{Z}Z`$, but the operator $`r(w)`$ is independent of the choice of representative for $`w`$.
$$r(w)=a^{}a^{}aa+aa^{}a^{}a+aaa^{}a^{}+a^{}aaa^{}.$$
We then have that
$$G_L=\mathrm{ker}r.$$
This equation restates the fact (proven in Lemma 3) that, given a cyclic gauge invariant tensors $`\tau _{1\mathrm{}2L}`$, the function $`\left(x^1_\omega \mathrm{}_\omega x^{2L}\right)\tau _{1\mathrm{}2L}`$ vanishes whenever $`\omega =K`$ is constant.
The operators $`\delta ,\overline{\delta },\mathrm{\Delta },\overline{\mathrm{\Delta }}`$ act naturally on the spaces $`P_L`$ and $`G_L`$. In particular the operators $`\delta `$ and $`\overline{\delta }`$ act on the spaces $`P_L`$ of polynomials
$`\delta `$ $`:`$ $`P_LP_{L+1}`$
$`\overline{\delta }`$ $`:`$ $`P_LP_{L1}`$
Choosing, without loss in generality, $`\mathrm{\Delta }_{z\overline{z}}=\mathrm{\Delta }^{z\overline{z}}=i`$, we have
$`\delta `$ $`=`$ $`{\displaystyle \frac{}{y}}`$
$`\overline{\delta }`$ $`=`$ $`y+2xy{\displaystyle \frac{}{x}}\left(x^2y^2\right){\displaystyle \frac{}{y}}`$
One can check explicitly that $`[\delta ,\overline{\delta }]=2x_x+2y_y1=2L1`$.
The operators $`\mathrm{\Delta }`$, $`\overline{\mathrm{\Delta }}`$ on the other hand act on the spaces $`G^L`$
$`\mathrm{\Delta }`$ $`:`$ $`G_LG_{L+1}`$
$`\overline{\mathrm{\Delta }}`$ $`:`$ $`G_LG_{L1}`$
We will be more explicit on the precise form of $`\mathrm{\Delta }`$ and $`\overline{\mathrm{\Delta }}`$ in the next subsection, but we know, from the general arguments of section 5.1, that$`[\overline{\mathrm{\Delta }},\mathrm{\Delta }]=2L1`$.
With this notation in place we can now easily construct invariant actions. In particular, we will first show how to canonically construct, starting from a lowest level term, a complete set of terms which combine into an invariant block. A general invariant action is then given, following the discussion at the end of section 5.1, by linear combinations of invariant blocks.
Let us then first describe the form of a lowest level term. In general, given the above discussion, a generic term in the action at level $`L`$ will be of the form
$$S_L=\underset{i}{}p_L^ig_L^i.$$
where the $`g_L^i`$ are a basis for $`G_L`$ and the $`p_L^i`$ are polynomials in $`P_L`$. A term $`S_P=_ip_P^ig_P^i`$ (we will reserve $`L`$ for a general level index, and $`P`$ for lowest level states) will be of lowest level if
$$\delta S_P=0\overline{\mathrm{\Delta }}S_P=0$$
(47)
The first equation implies that the polynomials $`p_P^i`$ depend uniquely on $`x`$, and therefore that $`p_P^i=c_ix^P`$. Then
$$S_P=x^Pg_P$$
where $`g_P=_ic_ig_P^i`$ satisfies, using the second equation in (47),
$$\overline{\mathrm{\Delta }}g_P=0.$$
Let us then start with a lowest level term $`S_P`$ and construct a full invariant block $`S`$ of the form
$`S`$ $`=`$ $`{\displaystyle \underset{LP}{}}S_L`$
$`S_L`$ $`=`$ $`p_Lg_L.`$
The above is invariant if
$`\mathrm{\Delta }S_L`$ $`=`$ $`\delta S_{L+1}`$ (48)
$`\overline{\delta }S_L`$ $`=`$ $`\overline{\mathrm{\Delta }}S_{L+1}`$
To solve the above constraints we construct the higher level $`g_L`$’s using $`\mathrm{\Delta }`$
$$g_{L+1}=\mathrm{\Delta }g_L$$
We must then find polynomials $`p_L`$ which satisfy
$$\delta p_{L+1}=p_L$$
(49)
and such that
$`\overline{\mathrm{\Delta }}g_{L+1}`$ $`=`$ $`c_Lg_L`$
$`\overline{\delta }p_L`$ $`=`$ $`c_Lp_{L+1}`$ (50)
for some constant $`c_L`$. First we compute $`\overline{\mathrm{\Delta }}g_{L+1}`$. Using the fact that $`[\overline{\mathrm{\Delta }},\mathrm{\Delta }]=2L1`$, and that $`\overline{\mathrm{\Delta }}g_P=0`$, we obtain
$`\overline{\mathrm{\Delta }}g_{L+1}`$ $`=`$ $`\overline{\mathrm{\Delta }}\mathrm{\Delta }^{L+1P}g_P=\left(\left(2L1\right)+\mathrm{}+(2P1)\right)g_L`$
$`=`$ $`c_Lg_L`$
with
$$c_L=(L+P1)(LP+1).$$
Using equation (50)
$$p_{L+1}=\frac{1}{(L+P1)(LP+1)}\overline{\delta }p_L$$
(51)
to define higher level polynomials, we can check, using $`[\delta ,\overline{\delta }]=2L1`$ and $`\delta p_P=0`$, that the remaining equation (49) is satisfied, and that we have indeed a solution to the invariance equations.
### 6.2 An important example
In this subsection we use the general construction described above and apply it to a specific important example. In particular we show again that the basic $`F^2`$ term is lowest level, and we construct part of the invariant block constructed from it. We recover in particular the Born–Infeld action, and we compute the first non–trivial derivative corrections at level $`4`$ which must be present in order to make the full action invariant.
We start by analyzing the explicit form of the operators $`\mathrm{\Delta },\overline{\mathrm{\Delta }}`$, recalling that $`\mathrm{\Delta }_{z\overline{z}}=\mathrm{\Delta }^{z\overline{z}}=i`$. We define for convenience the field strength $`F`$ as
$$F=[Z,\overline{Z}].$$
Then equation (42) reads, in the present case,
$`\mathrm{\Delta }Z`$ $`=`$ $`{\displaystyle \frac{1}{4}}(ZF+FZ)`$
$`\mathrm{\Delta }\overline{Z}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(\overline{Z}F+F\overline{Z}).`$
The above then defines the action of $`\mathrm{\Delta }`$ on $`G_L`$, since $`\mathrm{\Delta }`$ acts as a derivation on each coordinate forming the words in $`G_L`$. Similarly $`\overline{\mathrm{\Delta }}`$ is defined by
$$\overline{\mathrm{\Delta }}F=1.$$
More precisely, given a word $`wW_L`$, we cyclically rearrange the coordinates in each word so as to obtain a gauge invariant (not cyclic gauge invariant) form, containing only commutators. We then apply $`\overline{\mathrm{\Delta }}`$ on each fundamental commutator $`F`$ as a derivation.
Let us the consider the term
$$S_2=p_2g_2$$
with
$$p_2=x^2g_2=\frac{1}{2}F^2.$$
Clearly $`\delta p_2=0`$. Moreover $`\overline{\mathrm{\Delta }}g_2=F=0`$, since we recall that, as a word in $`W_1`$, the commutator $`F`$ is zero. Therefore $`S_2`$ is a lowest level state, and we can construct the corresponding invariant block.
First we do some computations explicitly. The polynomials $`p_3`$ and $`p_4`$ are given by
$`p_3`$ $`=`$ $`{\displaystyle \frac{1}{3}}\overline{\delta }p_2=yx^2`$
$`p_4`$ $`=`$ $`{\displaystyle \frac{1}{8}}\overline{\delta }p_3={\displaystyle \frac{1}{8}}x^4+{\displaystyle \frac{1}{2}}y^2x^2.`$
We then compute $`g_3`$ and $`g_4`$. First, applying the operation $`\mathrm{\Delta }`$ to $`F`$ we obtain
$$\mathrm{\Delta }F=\frac{1}{2}\left(F^2+ZF\overline{Z}\overline{Z}FZ\right).$$
This means that (recall that the RHS below is a word in $`W_3`$, and that cyclic rearrangements are allowed)
$`g_3`$ $`=`$ $`\mathrm{\Delta }g_2={\displaystyle \frac{1}{2}}F\left(\mathrm{\Delta }F\right)+{\displaystyle \frac{1}{2}}\left(\mathrm{\Delta }F\right)F=F\left(\mathrm{\Delta }F\right)`$
$`=`$ $`{\displaystyle \frac{1}{2}}F^3+{\displaystyle \frac{1}{2}}[FZ,F\overline{Z}]={\displaystyle \frac{1}{2}}F^3`$
The computation of $`g_4`$ is just slightly more complex, and we leave it for the appendix (Lemma 13). The result is
$$g_4=\mathrm{\Delta }g_3=F^4+\frac{1}{4}F[DF,\overline{D}F]$$
where
$$D\mathrm{}=[Z,\mathrm{}]\overline{D}\mathrm{}=[\overline{Z},\mathrm{}].$$
We can now combine the polynomials and the gauge invariant words. To make contact with standard notation we write the action for $`y=0`$ ($`\gamma _{ab}`$ symmetric) and revert to more standard notation
$`F`$ $``$ $`iF_{z\overline{z}}=iF`$ (52)
$`x`$ $``$ $`g^{z\overline{z}}={\displaystyle \frac{1}{g}}`$
$`D`$ $``$ $`iD`$
which give the following $`U(N)`$ lagrangian
$$\frac{1}{2g^2}\mathrm{Tr}\left(F^2\right)\frac{1}{8g^4}\mathrm{Tr}\left(F^4+\frac{i}{2}F[DF,\overline{D}F]\right)+\mathrm{}.$$
(53)
We see that the above action, written up to level $`4`$, contains the first part of the Born–Infeld action (in $`2`$–dimensions there is no ambiguity about ordering of the $`F^{2n}`$ terms), but already at level $`4`$ we have derivative corrections, which are required for the total invariance of the action<sup>8</sup><sup>8</sup>8It has been shown in that, at level $`4`$, the effective action can be written only as the $`F^4`$ term, with no derivative corrections. This result is not in contradiction with equation (53), since one can always allow for field redefinitions. In particular, consider, in general $`D`$ dimensions, the field redefinition $`A_aA_a+cF_{ab}D_cF_{bc}`$. This induces a change in the action at level $`4`$, coming from the $`F^2`$ term, of the form $`F_{da}D_d\left(F_{ab}D_cF_{bc}\right)=\frac{1}{2}F_{ab}[D_cF_{bc},D_dF_{ad}]`$, which in two dimensions is proportional to $`F[DF,\overline{D}F]`$. Therefore, for an appropriate choice of $`c`$, one can remove the derivative term in (53), thus resolving the apparent contradiction with ..
Let us now schematically consider the higher terms $`g_L`$. We wish to sketch how one can recover, within the $`F^2`$ invariant block, the complete Born–Infeld action (a much more detailed discussion on this point and related issues will appear in ). To this end, we first note that, in general,
$$g_L=c_LF^L+\mathrm{derivative}\mathrm{terms}.$$
We have that $`c_2=c_3=1/2`$, $`c_4=1`$,$`\mathrm{}`$. We can use the basic commutator $`[\overline{\mathrm{\Delta }},\mathrm{\Delta }]`$ and the fact that $`\overline{\mathrm{\Delta }}F^L=LF^{L1}`$ to compute all the $`c_L`$’s. In fact, applying the basic equation $`[\overline{\mathrm{\Delta }},\mathrm{\Delta }]=2L1`$ to $`g_L`$ (recalling that $`\mathrm{\Delta }g_L=g_{L+1}`$ and that $`\overline{\mathrm{\Delta }}`$ does not decrease the number of derivatives) one obtains the recursion relation
$$\left(L+1\right)\left(\frac{c_{L+1}}{c_L}\right)L\left(\frac{c_L}{c_{L1}}\right)=2L1$$
which is solved by $`c_L=(L2)c_{L1}`$, or by
$$c_L=\frac{1}{2}(L2)!$$
The $`F^2`$ invariant block then contains the sum $`_Lc_Lp_L(x,y)F^L`$. Let us consider more closely the polynomials $`p_L`$. First of all, from the general relation (51) one can easily show that the polynomials $`p_L`$ vanish for $`L`$ odd if $`y=0`$. On the other hand, for even levels, one has that
$$p_{2L}(x,0)=d_{2L}x^{2L}.$$
Therefore, the relevant part of the action is given by (again substituting $`x1/g`$ and $`FiF`$)
$$\underset{L}{}c_{2L}d_{2L}\frac{()^L}{g^{2L}}\mathrm{Tr}\left(F^{2L}\right).$$
(54)
To compute the coefficients $`d_{2L}`$, let us first note that, since $`\delta p_{2L1}=p_{2L2}`$, one must have that $`p_{2L1}=d_{2L2}x^{2L2}y+o\left(y^3\right)`$. This implies, using (51), that$`d_{2L}=\frac{1}{4L\left(L1\right)}d_{2L2}`$, which is solved by (recall that $`d_2=1`$)
$$d_{2L}=\left(\frac{1}{4}\right)^{L1}\frac{1}{L!\left(L1\right)!}.$$
Therefore, equation (54) gives the complete Born–Infeld action
$$\underset{L}{}\left(\frac{1}{2}\right)^{2L1}\frac{(2L2)!}{L!\left(L1\right)!}\frac{1}{g^{2L}}\mathrm{Tr}\left(F^{2L}\right)=\mathrm{Tr}\sqrt{1\frac{1}{g^2}F^2}.$$
## 7 Discussion
In this paper we have analyzed in detail the general structure of the non–abelian Born–Infeld action, together with the higher $`\alpha ^{}`$ derivative corrections. We have shown how the requirement of invariance of the action under a change of non–commutativity scale $`\theta `$ imposes severe restrictions on the possible terms which can appear. More specifically, we can construct invariant actions starting from invariant blocks, which are themselves obtained from a lowest level term (in a loose sense, a pure derivative term). Terms at higher level are then constructed so as to achieve invariance under a change in $`\theta `$. A general action is then a linear combination of invariant blocks, with coefficients which must be determined from a different computation. No argument in this paper assumes supersymmetry, and the results are therefore valid in bosonic open string theory, as well as in superstring theory. In particular, supersymmetry will impose restrictions on the allowed linear combinations of invariant blocks, possibly determining in part, or even completely, the effective action.
Let us now comment on interesting directions of possible future investigation.
* It is first of all important to explicitly solve the invariance equations $`\mathrm{\Delta }\tau ^L=\delta \tau ^{L+1}`$ and $`\overline{\mathrm{\Delta }}\tau ^{L+1}=\overline{\delta }\tau ^L`$ in the general $`D`$–dimensional case. Similarly to the $`2`$–dimensional case discussed in the text, we should study the algebra of operators $`\mathrm{\Delta },\delta ,\overline{\mathrm{\Delta }},\overline{\delta }`$ given by the relations $`[\overline{\mathrm{\Delta }},\mathrm{\Delta }]=[\delta ,\overline{\delta }]=2LD/2`$. The algebra now depends on more parameters, since the underlying matrix $`\mathrm{\Delta }_{ab}`$ now has $`D\left(D1\right)/2`$ components. It is important, in particular, to have a canonical construction of higher level terms, starting from the lowest level. This would in turn give a canonical definition of invariant block.
* It is important to understand how invariant blocks appear in the underlying boundary conformal field theory. In particular, the relation between the analysis of this paper, which is at the level of the effective action, and conformal field theory is of importance both conceptually and from a practical point of view.
* The results of this paper do not depend on supersymmetry. Understanding the additional constraints imposed by SUSY is an important task for the future. A first step in this direction is the following. Given an invariant action, we may use T–duality to describe the weak–coupling physics of D–branes. In particular, we expect, in a supersymmetric theory, to have minima of the effective action corresponding to holomorphic curves, surfaces, $`\mathrm{}`$. Very possibly, a careful restatement of this fact in terms of invariant blocks will impose constraints which must be satisfied in a supersymmetric theory.
* Given the invariant description of the action (41) as an operator trace, it is very tempting to resum the full series in one specific invariant block. In fact, although the action is usually written by artificially choosing a parameter $`\theta `$ and then by writing the expression in terms of coordinate operators $`X^a`$, it is nonetheless true that the operator $`O(\theta )=_LX^1\mathrm{}X^{2L}\tau _{1\mathrm{}2L}`$ has a trace $`\mathrm{Tr}\left(O\left(\theta \right)\right)`$ which is $`\theta `$–independent. It is then tempting to conjecture<sup>9</sup><sup>9</sup>9The spectral nature of actions has been very much stressed by A.Connes. that the various operators $`O(\theta )`$ not only have the same trace, but are related by a unitary transformation, and are then isospectral.
## 8 Acknowledgments
I would like to thank M. Douglas, A. Connes, A. Cattaneo and B. Shoikhet for useful discussions. Also I want to thank the Institute for Advanced Study and MIT, and in particular N. Seiberg and W. Taylor, for hospitality during the completion of this work. This research is supported by a European Post–doctoral Institute Fellowship.
## 9 Appendix
###### Definition 1
A tensor $`\eta _{1\mathrm{}n}`$ is called gauge invariant (GI) if
$$\eta _{123\mathrm{}n}+\eta _{213\mathrm{}n}+\eta _{231\mathrm{}n}+\mathrm{}+\eta _{23\mathrm{}1n}+\eta _{23\mathrm{}n1}=0.$$
A tensor $`\tau _{1\mathrm{}n}`$ is called cyclic gauge invariant (CGI) if
$$\tau _{123\mathrm{}n}+\tau _{213\mathrm{}n}+\tau _{231\mathrm{}n}+\mathrm{}+\tau _{23\mathrm{}1n}=0.$$
(55)
###### Lemma 2
(Section 3.7) Let $`\eta _{1\mathrm{}n}`$ be gauge invariant. Then
$$\tau _{1\mathrm{}n}=\frac{1}{n}\left(\eta _{1\mathrm{}n}+\mathrm{cyc}_{1\mathrm{}n}\right)$$
is cyclic gauge invariant.
PROOF. Let us write the left hand side of equation (55) in terms of $`\eta `$. Neglecting the multiplicative factor of $`1/n`$, we have the expression
$`\eta _{123\mathrm{}n}+\eta _{213\mathrm{}n}+\eta _{231\mathrm{}n}+\mathrm{}+\eta _{23\mathrm{}1n}+`$
$`\eta _{23\mathrm{}n1}+\eta _{13\mathrm{}n2}+\eta _{31\mathrm{}n2}+\mathrm{}+\eta _{3\mathrm{}1n2}+`$
$`\eta _{3\mathrm{}n12}+\eta _{3\mathrm{}n21}+\mathrm{}`$
$`\mathrm{}`$
$`\mathrm{}+\eta _{1n23\mathrm{}n1}+`$
$`\eta _{n123\mathrm{}n1}+\eta _{n213\mathrm{}n1}+\eta _{n231\mathrm{}n1}+\mathrm{}+\eta _{n23\mathrm{}n1,1}`$
The above expression contains $`\left(n1\right)n`$ terms ($`n`$ lines with $`n1`$ terms each). Consider the sequence of terms in the order written, and assemble them now into groups of $`n`$ terms. It is then easy to see that each individual group vanishes since $`\eta `$ is gauge invariant. $`\mathrm{}`$
###### Lemma 3
(Section 3.7) Let $`\tau _{1\mathrm{}n}`$ be cyclic gauge invariant and let
$$\tau (x)=x^1_\omega \mathrm{}_\omega x^n\tau _{1\mathrm{}n}.$$
Then, for $`x\mathrm{}`$, $`\tau (x)`$ grows linearly with $`x`$. Moreover, if $`n`$ is even, then
$$\tau (x)0$$
for $`x\mathrm{}.`$
PROOF. We use the fact that, at infinity, $`\omega K`$ approaches a constant, and that we can therefore compute the variation
$`\tau (x+\epsilon )\tau (x)`$ $`=`$ $`\epsilon ^1x^2_K\mathrm{}_Kx^n(\tau _{12\mathrm{}n}+\mathrm{}+\tau _{2\mathrm{}n1})`$
$`=`$ $`h(x),`$
where
$$h(x)=x^2_K\mathrm{}_Kx^n\epsilon ^1\tau _{12\mathrm{}n}.$$
It is immediate to see that, for any $`\epsilon `$, the tensor $`\epsilon ^1\tau _{12\mathrm{}n}`$ is gauge invariant in the indices $`2,\mathrm{},n`$ and that, using the results of section 3.7, the function $`h(x)`$ is constant at infinity. This implies that $`\tau (x)`$ is at most a linear function of the coordinates $`x`$, when $`x\mathrm{}`$. We now note that $`x^1_K\mathrm{}_Kx^n`$ is a polynomial of degree $`n`$, with monomials of degrees $`n,n2,n4,\mathrm{}`$. In particular, if $`n`$ is even, we do not have a linear term and the function $`\tau `$ approaches a constant at infinity. We now wish to show that the constant is $`0`$. Consider the polynomial $`x^1_K\mathrm{}_Kx^n`$. It will be of the form
$$x^1_K\mathrm{}_Kx^n=C^{1\mathrm{}n}+o(x).$$
We have just seen that we do not need to consider the $`o(x)`$ part, and we therefore just need to prove that $`C^{1\mathrm{}n}\tau _{1\mathrm{}n}=0`$. It is not difficult to show, using the Moyal product $`_K`$, that
$$C^{1\mathrm{}n}\underset{\sigma }{}()^{J(\sigma )}\theta ^{\sigma _1\sigma _2}\mathrm{}\theta ^{\sigma _{n1}\sigma _n}$$
(56)
where $`J(\sigma )`$ counts the number of pairs $`\sigma _{2i1},\sigma _{2i}`$ for which $`\sigma _{2i1}>\sigma _{2i}`$. Then, in order to finish the proof, one has to show that the quantity
$$\underset{\sigma }{}()^{J(\sigma ^1)}\tau _{\sigma _1\mathrm{}\sigma _n}$$
(57)
vanishes for cyclic $`\tau `$’s. Let us then fix a given permutation $`\sigma `$, and let me denote with $`\pi `$ the basic cyclic permutation $`(1,\mathrm{},n)(2,\mathrm{},n,1)`$. Consider then the permutations $`\rho _k=\pi ^k\sigma `$. Let us first show that $`()^{J(\rho _k^1)}`$ is alternating with $`k`$. In fact, for $`kk+1`$, almost all the pairs $`\sigma _{2i1},\sigma _{2i}`$ go into the pairs $`\sigma _{2i1}+1,\sigma _{2i}+1`$. This is, on the other hand, not true for the single pair with either $`\sigma _{2i1}`$ or $`\sigma _{2i}`$ equal to $`n`$, since $`n1`$. Only this one pair changes the ordering of its components, and therefore the sign $`()^{J(\rho _k^1)}`$ changes if $`kk+1`$. Consider then the set of permutations $`\rho _k`$, cyclic permutations of $`\sigma `$, for $`0k<n`$. This gives, in the sum (57),
$$\pm \left(\tau _{\sigma _1\sigma _2\sigma _3\mathrm{}\sigma _n}\tau _{\sigma _2\sigma _3\mathrm{}\sigma _n\sigma _1}+\tau _{\sigma _3\mathrm{}\sigma _n\sigma _1\sigma _2}\mathrm{}\tau _{\sigma _n\sigma _1\sigma _2\mathrm{}\sigma _{n1}}\right)$$
The terms come with alternating signs, and since $`n`$ is even the number of $`+`$ signs is equal to that of $``$ signs. Moreover, all the terms are actually the same, since $`\tau `$ is cyclically symmetric. The above sum then vanishes. Partitioning the set of all permutations $`\sigma `$ in sets of cyclically related permutations, we can then show that the full sum (57) vanishes. $`\mathrm{}`$
###### Remark 4
(Lemmata 5 and 7) We make a general comment on integration of commutators, which will be useful in the rest of the appendix. Consider two functions $`f`$ and $`g`$, and look at the integral
$$d^DxV(\omega )\left(f_\omega gg_\omega f\right).$$
It $`f`$ and $`g`$ vanish at infinity, then the above integral vanishes, as was discussed in the main text. If, on the other hand, $`f`$ and $`g`$ do not go to zero for $`x\mathrm{}`$, we can proceed as follows. Assume that $`\omega =K`$ outside of a compact domain $`DM`$ and consider the integral of $`[f,g]_\omega `$ over $`D`$
$$I=_Dd^Dx\text{ }V(\omega )\text{ }[f,g]_\omega $$
(58)
If $`f,g=0`$ on $`D`$ and outside of $`D`$, the above expression vanishes, and therefore, in general, the integral (58) must reduce to a boundary integral over $`D`$. We can then continuously deform $`\omega K`$ in the interior of $`D`$ without changing the integral. This means, in particular, that, for any functions $`f`$ and $`g`$,
$$I=_Dd^Dx\text{ }det{}_{}{}^{\frac{1}{2}}K\text{ }[f,g]_K.$$
From the above arguments it is also clear that the above equality holds even if $`f`$ and $`g`$ depend themselves in a local way on $`\omega `$. For example, if $`f=f_1_\omega f_2`$, then we have
$$I=_Dd^Dx\text{ }det{}_{}{}^{\frac{1}{2}}K\text{ }[f_1_Kf_2,g]_K.$$
In practice, when integrating commutators, we can replace $`\omega `$ with $`K`$ in all the expressions without changing the integral.
###### Lemma 5
(Section 3.7) Let $`\eta _{1\mathrm{}n}`$ be gauge invariant and let $`\tau _{1\mathrm{}n}=\frac{1}{n}\left(\eta _{1\mathrm{}n}+\mathrm{cyc}_{1\mathrm{}n}\right)`$. Define
$`\eta (x)`$ $`=`$ $`x^1_\omega \mathrm{}_\omega x^n\eta _{1\mathrm{}n}`$
$`\tau (x)`$ $`=`$ $`x^1_\omega \mathrm{}_\omega x^n\tau _{1\mathrm{}n}`$
with $`\eta (x)\eta _{\mathrm{}}`$ for $`x\mathrm{}`$. If $`n`$ is even, then
$$d^DxV(\omega )\tau =d^DxV(\omega )\left[\eta \eta _{\mathrm{}}\right]$$
PROOF. First it is clear that
$`\tau (x)`$ $`=`$ $`{\displaystyle \frac{1}{n}}\eta _{1\mathrm{}n}(x^1_\omega \mathrm{}_\omega x^n+`$
$`+x^2_\omega \mathrm{}_\omega x^n_\omega x^1+\mathrm{}`$
$`+x^n_\omega x^1_\omega \mathrm{}_\omega x^{n1}).`$
Therefore, the difference $`\eta \tau `$ is given by
$`\eta \tau `$ $`=`$ $`{\displaystyle \frac{1}{n}}\eta _{1\mathrm{}n}([x^1,x^2_\omega \mathrm{}_\omega x^n]_\omega `$
$`+[x^1_\omega x^2,x^3_\omega \mathrm{}_\omega x^n]_\omega +\mathrm{}`$
$`+[x^1_\omega \mathrm{}_\omega x^{n1},x^n]_\omega ).`$
Recalling remark 4, we will be done once we show that the RHS above is equal to
$$\eta _{\mathrm{}}=\eta _{1\mathrm{}n}x^1_K\mathrm{}_Kx^n$$
when we replace $`\omega K`$. We must therefore prove that
$$\eta _{\mathrm{}}=\frac{1}{n}\eta _{1\mathrm{}n}([x^1,x^2_K\mathrm{}_Kx^n]_K+\mathrm{}+[x^1_K\mathrm{}_Kx^{n1},x^n]_K).$$
(59)
Both the LHS and the RHS above are constant, since $`\eta `$ is GI. Let us introduce a compact notation
$`\eta _{1\mathrm{}n}x^1_K\mathrm{}_Kx^n`$ $``$ $`[1\mathrm{}n]`$
$`\eta _{1\mathrm{}n}x^2_K\mathrm{}_Kx^n_Kx^1`$ $``$ $`[2\mathrm{}n1]`$
$`\mathrm{}`$
Using formula (56) and the arguments which follow it, we can show that
$$[k\mathrm{}n1\mathrm{}k1]=\left(1\right)^{k1}[1\mathrm{}n]=\left(1\right)^{k1}\eta _{\mathrm{}}.$$
Then the RHS of (59) is given by
$`{\displaystyle \frac{n1}{n}}[12\mathrm{}n]{\displaystyle \frac{1}{n}}[23\mathrm{}n1]\mathrm{}{\displaystyle \frac{1}{n}}[n1\mathrm{}n1]`$
$`=`$ $`[12\mathrm{}n]\left({\displaystyle \frac{n1}{n}}+{\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{n}}+\mathrm{}+{\displaystyle \frac{1}{n}}\right)=`$
$`=`$ $`[12\mathrm{}n]`$
as was to be shown. $`\mathrm{}`$
###### Lemma 6
(Section 4.3) Let $`f`$ be a generic function, and let $`F=Q_\omega (f)`$, where $`\omega `$ is an arbitrary symplectic structure. Let also $`\mathrm{\Delta }_{ab}`$ be a constant antisymmetric matrix. Then, to first order in $`\mathrm{\Delta }`$,
$$Q_{\omega +\mathrm{\Delta }}=F+\frac{i}{4}\mathrm{\Delta }_{ab}\left(X^aX^bF+FX^aX^b2X^aFX^b\right).$$
PROOF. We start by noting that, if $`\omega =K`$ is constant, a simple computation using the Moyal product $`_K`$ shows that
$`f_{K+\mathrm{\Delta }}gf_Kg`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left(\theta \mathrm{\Delta }\theta \right)^{ab}_af_K_bg=`$
$`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{\Delta }_{ab}[x^a,f]_K_K[x^b,g]_K.`$
We must then consider the product $`f_{\omega +\mathrm{\Delta }}g`$ for general $`\omega `$. As always, we look for a map $`T`$ such that $`T\left(f_{\omega +\mathrm{\Delta }}g\right)=Tf_\omega Tg`$. If we work to first order in $`\mathrm{\Delta }`$, and accordingly let $`T=1+R`$ (with $`R`$ or order $`\mathrm{\Delta }`$), one has that
$$f_{\omega +\mathrm{\Delta }}gf_\omega g=Rf_\omega g+f_\omega RgR\left(f_\omega g\right).$$
(61)
The map $`T=1+R`$ also relates $`Q_{\omega +\mathrm{\Delta }}`$ and $`Q_\omega `$ as follows
$$Q_{\omega +\mathrm{\Delta }}(f)=Q_\omega (f+Rf).$$
We have seen from the example in section 4.2 that the we should consider general variations $`Q_{\omega +\mathrm{\Delta }}(f)Q_\omega (f)`$ of the form $`\frac{i}{4}\mathrm{\Delta }_{ab}\left(aX^aX^bF+bFX^aX^b+cX^aFX^b\right)`$ where $`a,b,c`$ are constants which we must determine. It is then clear that
$$Rf=\frac{i}{4}\mathrm{\Delta }_{ab}(ax^ax^bf+bfx^ax^b+cx^afx^b).$$
Using the above fact in the RHS of equation (61), and comparing, for $`\omega =K`$, with the RHS of equation (9), we obtain that $`a=b=1,b=2`$, as was required. $`\mathrm{}`$
###### Lemma 7
(Section 5.1) Let $`\tau _{1\mathrm{}n2}`$ be a cyclic gauge invariant tensor, and let $`\mathrm{\Delta }_{ab}`$ be a constant antisymmetric matrix. Then the tensor
$$\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}n}=\frac{i}{2}\left(\frac{n2}{n}\right)\left(\mathrm{\Delta }_{12}\tau _{345\mathrm{}n}\mathrm{\Delta }_{13}\tau _{245\mathrm{}n}\right)+\mathrm{cyc}_{1\mathrm{}n}$$
is itself cyclic gauge invariant. Moreover, if $`\omega `$ is a generic symplectic structure and $`X^a=Q_\omega (x^a)`$, then, under the variation $`\omega \omega +\mathrm{\Delta }`$, the operator $`\mathrm{Tr}\left(X^1\mathrm{}X^{n2}\right)\tau _{1\mathrm{}n2}`$ varies by $`\mathrm{Tr}\left(X^1\mathrm{}X^n\right)\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}n}`$ whenever $`n`$ is even.
PROOF. First we show that the tensor $`\mathrm{\Delta }\tau `$ is indeed cyclic gauge invariant. Written in full, we want to show that the following sum
$`\left(\mathrm{\Delta }_{12}\tau _{345\mathrm{}n}\mathrm{\Delta }_{13}\tau _{245\mathrm{}n}\right)+\mathrm{cyc}_{123\mathrm{}n}`$ (62)
$`\left(\mathrm{\Delta }_{21}\tau _{345\mathrm{}n}\mathrm{\Delta }_{23}\tau _{145\mathrm{}n}\right)+\mathrm{cyc}_{213\mathrm{}n}`$
$`\mathrm{}`$
$`\left(\mathrm{\Delta }_{23}\tau _{45\mathrm{}1n}\mathrm{\Delta }_{24}\tau _{35\mathrm{}1n}\right)+\mathrm{cyc}_{23\mathrm{}1n}`$
vanishes. To this end, we consider three significant cases, with the hope that the reader can understand from them the general line of the argument.
Let us first consider, within the above sum, terms which are proportional to $`\mathrm{\Delta }_{12}`$. They come only from the first and the last line and are
$$\tau _{345\mathrm{}n}\tau _{n34\mathrm{}n1}=0.$$
Similarly, terms proportional to $`\mathrm{\Delta }_{21}`$ come from the second and third line and exactly cancel each other. Consider now terms proportional, say, to $`\mathrm{\Delta }_{23}`$. These terms are present in every line of the above sum, and they are
$$\left(\tau _{145\mathrm{}n}\tau _{145\mathrm{}n}\right)+(\tau _{145\mathrm{}n}+\tau _{415\mathrm{}n}+\mathrm{}+\tau _{45\mathrm{}1n})$$
The above vanishes since $`\tau `$ is cyclic gauge invariant. Finally we consider terms which are proportional to $`\mathrm{\Delta }_{2n}`$. In this final case, we can check that no terms in the sum (62) contain $`\mathrm{\Delta }_{2n}`$. The reader can convince him or herself that all other combination of indices fall in one of these three cases.
###### Remark 8
Let us note that, in the above proof, we have not used the fact that $`\mathrm{\Delta }_{ab}`$ is antisymmetric. In fact, we have shown more generally that the tensor
$$\left(A_{12}\tau _{345\mathrm{}n}A_{13}\tau _{245\mathrm{}n}\right)+\mathrm{cyc}_{1\mathrm{}n}$$
is cyclic gauge invariant for any choice of $`A`$, whenever $`\tau `$ itself is CGI.
We now move to the second part of the lemma. Introduce the following two functions on $`M`$
$`A(x)`$ $`=`$ $`x^1_\omega \mathrm{}_\omega x^{n2}\tau _{1\mathrm{}n2}`$
$`B(x)`$ $`=`$ $`x^1_\omega \mathrm{}_\omega x^n\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}n}.`$
Since both $`\tau `$ and $`\mathrm{\Delta }\tau `$ are CGI, the two functions $`A`$ and $`B`$ tend to $`0`$ as $`x\mathrm{}`$ (we are assuming $`n`$ even). Recall that, under a variation $`\omega \omega +\mathrm{\Delta }`$, the star product of two functions $`f_\omega g`$ changes by $`Rfg+fRgR(fg)`$ (we do not show the explicit $`\omega `$ dependence in $``$), where $`R`$ is given in equation (38). Therefore, given three functions, the variation of $`fgh`$ is $`Rfgh+fRgh+fgRhR(fgh)`$, and similarly for products of more functions. In particular, the variation of the function $`A`$ is given by
$$AA+CRA$$
(63)
where
$$C(x)=Rx^1_\omega \mathrm{}_\omega x^{n2}\tau _{1\mathrm{}n2}+\mathrm{}+x^1_\omega \mathrm{}_\omega Rx^{n2}\tau _{1\mathrm{}n2}.$$
Since we are interested in traces of operators, we must consider also the variation coming from the change of integration measure. We use equations (40) and (63), together with the fact that $`A`$ vanishes at infinity, to show that the variation of
$$\mathrm{Tr}\left(X^1\mathrm{}X^{n2}\right)\tau _{1\mathrm{}n2}=d^DxV(\omega )A$$
is simply given by
$$d^DxV(\omega )C.$$
We must then prove that
$$d^DxV(\omega )\left(BC\right)=0.$$
It is simple to show that the function $`B`$ is obtained by cyclically rearranging the coordinate functions which build $`C`$. Following remark 4, the integral above reduces to a boundary term, and to show that it vanishes we just need to prove that $`B=C`$ whenever $`\omega =K`$ is constant. On one side, we know that $`B=0`$ for $`\omega =K`$, since $`\mathrm{\Delta }\tau `$ is cyclic gauge invariant. We then need to prove that, for constant symplectic structures, $`C=0`$. This is shown in two steps. First look at the operation $`R`$ on coordinate functions in the case of flat symplectic structure
$`Rx^a`$ $`=`$ $`{\displaystyle \frac{i}{4}}\mathrm{\Delta }_{bc}\left(x^b_Kx^c_Kx^a+x^a_Kx^b_Kx^c2x^b_Kx^a_Kx^c\right)`$
$`=`$ $`M_b^ax^b`$
with
$$M_b^a=\frac{1}{2}\theta ^{ac}\mathrm{\Delta }_{cb}.$$
It is then clear that
$`C(x)`$ $`=`$ $`x^1_\omega \mathrm{}_\omega x^{n2}\pi _{1\mathrm{}n2}`$
$`\pi _{1\mathrm{}n2}`$ $`=`$ $`M_1^a\tau _{a2\mathrm{}n2}+\mathrm{}+M_{n2}^a\tau _{1\mathrm{}n3,a}`$
It is now quite easy to show that $`\pi `$ is cyclic gauge invariant, therefore implying that $`C=0`$. $`\mathrm{}`$
###### Lemma 9
(Section 3.3) Let $`\eta _{1\mathrm{}n+2}`$ be gauge invariant, and let us consider the combination $`C=\left(\lambda ^1\mathrm{}\lambda ^{n+2}\right)\eta _{1\mathrm{}n+2}`$. If we add a central term $`\mathrm{\Delta }^{ab}`$ to the commutator $`i[\lambda ^a,\lambda ^b]`$, then the expression $`C`$ varies by $`\left(\lambda ^1\mathrm{}\lambda ^n\right)\left(\overline{\mathrm{\Delta }}\eta \right)_{1\mathrm{}n}`$, where
$$\left(\overline{\mathrm{\Delta }}\eta \right)^{1\mathrm{}n}=\frac{i}{2}\mathrm{\Delta }^{ab}\left(\eta _{1\mathrm{}nab}+\eta _{1\mathrm{}anb}+\mathrm{}\right).$$
Moreover, for any $`\mathrm{\Delta }^{ab}`$, the tensor $`\overline{\mathrm{\Delta }}\eta `$ is gauge invariant.
PROOF. First it is clear that, if $`\eta _{1\mathrm{}n}`$ and $`\nu _{1\mathrm{}m}`$ are gauge invariant, then so is $`(\eta \nu )_{1\mathrm{}n+m}=`$ $`\eta _{1\mathrm{}n}\nu _{n+1\mathrm{}n+m}`$. Moreover, if $`d_a`$ is any one–indexed tensor, then $`(d\eta )_{1\mathrm{}n+1}=d_1\eta _{2\mathrm{}n+1}\eta _{1\mathrm{}n}d_{n+1}`$ is again gauge invariant. These two facts can either be checked algebraically, or one can simply note that they correspond respectively to the product and covariant derivative of gauge invariant operators. In fact, any gauge invariant tensor is built using the two operations just described.
In order to prove the lemma we first show that it holds for $`n=0`$. Then $`\eta _{ab}`$ is an antisymmetric tensor and $`\lambda ^a\lambda ^b\eta _{ab}=\frac{1}{2}[\lambda ^a,\lambda ^b]\eta _{ab}\frac{1}{2}[\lambda ^a,\lambda ^b]\eta _{ab}+\frac{i}{2}\mathrm{\Delta }^{ab}\eta _{ab}`$. Now suppose that we have proved the result for $`\eta _{1\mathrm{}n}`$ and $`\nu _{1\mathrm{}m}`$. Then we must show that $`\overline{\mathrm{\Delta }}(\eta \nu )=(\overline{\mathrm{\Delta }}\eta )\nu +\eta (\overline{\mathrm{\Delta }}\nu )`$. This is easily done, since
$`\overline{\mathrm{\Delta }}(\eta \nu )(\overline{\mathrm{\Delta }}\eta )\nu \eta (\overline{\mathrm{\Delta }}\nu )`$ $``$ $`\mathrm{\Delta }^{ab}(\eta _{a,1,\mathrm{},n1}+\eta _{1,a,\mathrm{},n1}+\mathrm{}+\eta _{1,\mathrm{},n1,a})`$
$`(\nu _{b,n\mathrm{},n+m2}+\nu _{n,b,\mathrm{},n+m2}+\mathrm{}+\nu _{n,\mathrm{},n+m2,b})`$
which vanishes since $`\eta `$ and $`\nu `$ are gauge invariant. Finally me must show that, given a generic $`d_a`$, one has $`\overline{\mathrm{\Delta }}(d\eta )=d(\overline{\mathrm{\Delta }}\eta )`$. Again this is easy to show using the gauge invariance of $`\eta `$, since
$$\overline{\mathrm{\Delta }}(d\eta )d(\overline{\mathrm{\Delta }}\eta )\mathrm{\Delta }^{ab}d_a(\eta _{b,1,\mathrm{},n}+\eta _{1,b,\mathrm{},n}+\mathrm{}+\eta _{1,\mathrm{},n,b})=0.$$
This concludes the proof, since any gauge invariant operator is a product of covariant derivatives of the field strength. $`\mathrm{}`$
###### Lemma 10
(Section 3.7) Let $`\eta _{1\mathrm{}n+2}`$ be gauge invariant and let $`\tau _{1\mathrm{}n+2}`$ be the associated cyclic gauge invariant tensor. Let $`\mathrm{\Delta }^{ab}`$ be antisymmetric, and define
$$g_{1\mathrm{}n}=\frac{1}{n}\left(\overline{\mathrm{\Delta }}\eta _{1\mathrm{}n}+\mathrm{cyc}_{1\mathrm{}n}\right)$$
(64)
where $`\overline{\mathrm{\Delta }}\eta `$ is given by expression (15). The tensor $`g`$ is then uniquely a function of $`\tau `$, and is explicitly given by the expression
$$\frac{i}{2}\left(\frac{n+2}{n^2}\right)\mathrm{\Delta }^{ab}\left[n\tau _{1\mathrm{}nab}+\left(n1\right)\tau _{1\mathrm{}anb}+\mathrm{}+0\tau _{a1\mathrm{}nb}\right]+\mathrm{cyc}_{1\mathrm{}n}$$
(65)
which will be denoted by $`\overline{\mathrm{\Delta }}\tau `$.
PROOF. Define the following tensors
$`k_{1\mathrm{}nab}`$ $`=`$ $`\eta _{1\mathrm{}nab}+\mathrm{cyc}_{1\mathrm{}n}`$
$`k_{1\mathrm{}anb}`$ $`=`$ $`\eta _{1\mathrm{}anb}+\mathrm{cyc}_{1\mathrm{}n}`$
$`k_{1\mathrm{}n1,b,a,n}`$ $`=`$ $`\eta _{1\mathrm{}n1,b,a,n}+\mathrm{cyc}_{1\mathrm{}n}`$
$`\mathrm{}`$
which have two selected indices $`a,b`$ and are cyclically symmetric in the other indices $`1,\mathrm{},n`$. It is clear that
$`\tau _{1\mathrm{}nab}+\mathrm{cyc}_{1\mathrm{}n}`$ $`=`$ $`{\displaystyle \frac{1}{n+2}}\left(k_{1\mathrm{}nab}+\mathrm{cyc}_{1\mathrm{}nab}\right)`$
$`\tau _{1\mathrm{}anb}+\mathrm{cyc}_{1\mathrm{}n}`$ $`=`$ $`{\displaystyle \frac{1}{n+2}}\left(k_{1\mathrm{}anb}+\mathrm{cyc}_{1\mathrm{}anb}\right)`$
$`\mathrm{}`$
We can also express equation (64) for $`g`$ in terms of the tensors $`k`$ as
$`g_{1\mathrm{}n}`$ $`=`$ $`{\displaystyle \frac{i}{2n^2}}\mathrm{\Delta }^{ab}J_{1\mathrm{}nab}`$
$`J_{1\mathrm{}nab}`$ $`=`$ $`nk_{1\mathrm{}nab}+nk_{1\mathrm{}anb}+\mathrm{},`$
where, as in $`\overline{\mathrm{\Delta }}\eta `$, the indices $`a,b`$ are in all possible positions with $`a`$ preceding $`b`$. We then want to prove that the above expression is equal to (65), which we can also write in terms of the tensors $`k`$ as follows
$$\frac{i}{2n^2}\mathrm{\Delta }^{ab}\stackrel{~}{J}_{1\mathrm{}nab}$$
with
$`\stackrel{~}{J}_{1\mathrm{}nab}`$ $`=`$ $`n\left(k_{1\mathrm{}nab}+\mathrm{cyc}_{1\mathrm{}nab}\right)`$
$`+\left(n1\right)\left(k_{1\mathrm{}anb}+\mathrm{cyc}_{1\mathrm{}anb}\right)`$
$`+\mathrm{}`$
$`+0\left(k_{a1\mathrm{}nb}+\mathrm{cyc}_{a1\mathrm{}nb}\right)`$
We will actually prove that $`J_{1\mathrm{}nab}=\stackrel{~}{J}_{1\mathrm{}nab}`$. In order to do this, we use the gauge invariance of the tensor $`\eta `$, which implies various linear relations among the tensors $`k`$
$`0`$ $`=`$ $`k_{ab1\mathrm{}n}+k_{ba1\mathrm{}n}+k_{b1a\mathrm{}n}+\mathrm{}+k_{b1\mathrm{}na}`$ (66)
$`0`$ $`=`$ $`k_{a1b\mathrm{}n}+k_{1ab\mathrm{}n}+k_{1ba\mathrm{}n}+\mathrm{}+k_{1b2\mathrm{}na}`$
$`\mathrm{}`$
We then need to prove that the difference $`J\stackrel{~}{J}`$ can be written as a linear combination of the above equations.
In order to write an efficient and clear proof, we will concentrate, from now on, on the case $`n=4`$. The proof in the general case is absolutely identical, but the added notation would obscure the result without adding new ideas to the ones already contained in the special case discussed below. We introduce the following compact notation
$`k_{1234ab}`$ $``$ $`\begin{array}{cccccc}& & & & a& b\end{array}`$
$`k_{123a4b}`$ $``$ $`\begin{array}{cccccc}& & & a& & b\end{array}`$
$`\mathrm{}`$
where the dots $`\begin{array}{cccc}& & & \end{array}`$ indicate the indices $`1,2,3,4`$. We then arrange all the possible tensors $`k`$ in the following tableau
| $`\begin{array}{cccccc}& & & & a& b\\ & & & a& b& \\ & & a& b& & \\ & a& b& & & \\ a& b& & & & \end{array}`$ | $`\begin{array}{cccccc}& & & & b& a\\ & & & b& a& \\ & & b& a& & \\ & b& a& & & \\ b& a& & & & \end{array}`$ |
| --- | --- |
| $`\begin{array}{cccccc}& & & a& & b\\ & & a& & b& \\ & a& & b& & \\ a& & b& & & \end{array}`$ | $`\begin{array}{cccccc}& & & b& & a\\ & & b& & a& \\ & b& & a& & \\ b& & a& & & \end{array}`$ |
| $`\begin{array}{cccccc}& & a& & & b\\ & a& & & b& \\ a& & & b& & \end{array}`$ | $`\begin{array}{cccccc}& & b& & & a\\ & b& & & a& \\ b& & & a& & \end{array}`$ |
| $`\begin{array}{cccccc}& a& & & & b\\ a& & & & b& \end{array}`$ | $`\begin{array}{cccccc}& b& & & & a\\ b& & & & a& \end{array}`$ |
| $`\begin{array}{cccccc}a& & & & & b\end{array}`$ | $`\begin{array}{cccccc}b& & & & & a\end{array}`$ |
which has on the left all terms with $`a`$ preceding $`b`$, and on the right all terms with $`b`$ before $`a`$. From top to bottom, the terms are, on the other hand, arranged in groups with a fixed number of indices between $`a`$ and $`b`$. We then denote any linear combination of the tensors $`k`$ with a horizontal box of coefficients
| $`a_1a_2a_3a_4a_5`$ | $`b_1b_2b_3b_4`$ | $`c_1c_2c_3`$ | $`d_1d_2`$ | $`e_1`$ |
| --- | --- | --- | --- | --- |
| $`f_1f_2f_3f_4f_5`$ | $`g_1g_2g_3g_4`$ | $`h_1h_2h_3`$ | $`i_1i_2`$ | $`j_1`$ |
where the top line corresponds to the coefficients of the left column in the tableau, and the bottom line to the right column (the above box is then a compact notation for the sum $`a_1k_{1234ab}+a_2k_{123a4b}+\mathrm{}+f_1k_{1234ba}+\mathrm{}`$). The main statement which we want to prove can now be compactly written as
$$\begin{array}{ccccc}44444& 4444& 444& 44& 4\\ 00000& 0000& 000& 00& 0\end{array}=\begin{array}{ccccc}44444& 3333& 222& 11& 0\\ 00000& 1111& 222& 33& 4\end{array}$$
(69)
In fact the LHS above is nothing but $`J`$. To show that the RHS is $`\stackrel{~}{J}`$ we just note that
$`k_{1\mathrm{}nab}+\mathrm{cyc}_{1\mathrm{}nab}`$ $`=`$
| $`11111`$ | $`0000`$ | $`000`$ | $`00`$ | $`0`$ |
| --- | --- | --- | --- | --- |
| $`00000`$ | $`0000`$ | $`000`$ | $`00`$ | $`1`$ |
$`k_{1\mathrm{}anb}+\mathrm{cyc}_{1\mathrm{}anb}`$ $`=`$
| $`00000`$ | $`1111`$ | $`000`$ | $`00`$ | $`0`$ |
| --- | --- | --- | --- | --- |
| $`00000`$ | $`0000`$ | $`000`$ | $`11`$ | $`0`$ |
$`\mathrm{}`$
To prove the equality (69) we use the linear relations (66), which can also be compactly rewritten as
$$a_i=b_j=0,$$
where
$`a_1`$ $`=`$ $`\begin{array}{ccccc}10000& 1000& 100& 10& 1\\ 10000& 0000& 000& 00& 0\end{array}b_1=\begin{array}{ccccc}00001& 0001& 001& 01& 1\\ 00001& 0000& 000& 00& 0\end{array}`$
$`a_2`$ $`=`$ $`\begin{array}{ccccc}01000& 0100& 010& 01& 0\\ 01000& 1000& 000& 00& 0\end{array}b_2=\begin{array}{ccccc}00010& 0010& 010& 10& 0\\ 00010& 0001& 000& 00& 0\end{array}`$
$`a_3`$ $`=`$ $`\begin{array}{ccccc}00100& 0010& 001& 00& 0\\ 00100& 0100& 100& 00& 0\end{array}b_3=\begin{array}{ccccc}00100& 0100& 100& 00& 0\\ 00100& 0010& 001& 00& 0\end{array}`$
$`a_4`$ $`=`$ $`\begin{array}{ccccc}00010& 0001& 000& 00& 0\\ 00010& 0010& 010& 10& 0\end{array}b_4=\begin{array}{ccccc}01000& 1000& 000& 00& 0\\ 01000& 0100& 010& 01& 0\end{array}`$
$`a_5`$ $`=`$ $`\begin{array}{ccccc}00001& 0000& 000& 00& 0\\ 00001& 0001& 001& 01& 1\end{array}b_5=\begin{array}{ccccc}10000& 0000& 000& 00& 0\\ 10000& 1000& 100& 10& 1\end{array}`$
Summing either all the $`a`$’s or all the $`b`$’s we first of all obtain the following interesting identity
$$\begin{array}{ccccc}11111& 1111& 111& 11& 1\\ 11111& 1111& 111& 11& 1\end{array}=0.$$
(77)
Combining the above equation with (69) we can then reduce the statement of the lemma to the following equality
$$\begin{array}{ccccc}44444& 5555& 666& 77& 8\\ 44444& 3333& 222& 11& 0\end{array}=0.$$
The LHS above is nothing but
$$4(a_1+b_1)+3(a_2+b_2)+2(a_3+b_3)+1(a_4+b_4)+0(a_5+b_5),$$
and therefore the proof is complete. $`\mathrm{}`$
###### Remark 11
Using the box (77) we can show that the expression
$$\left[n\tau _{1\mathrm{}nab}+\left(n1\right)\tau _{1\mathrm{}anb}+\mathrm{}+0\tau _{a1\mathrm{}nb}\right]+\mathrm{cyc}_{1\mathrm{}n}$$
in (65) is antisymmetric in $`a`$ and $`b`$.
###### Lemma 12
(Section 5.1) Let $`\tau _{1\mathrm{}n}`$ be a cyclic gauge invariant tensor, and let $`\mathrm{\Delta }_{ab}`$ be an arbitrary invertible antisymmetric matrix with inverse $`\mathrm{\Delta }^{ab}`$. Then
$$[\overline{\mathrm{\Delta }},\mathrm{\Delta }]\tau =\left(n\frac{1}{2}D\right)\tau .$$
PROOF. First we recall the expression for $`\mathrm{\Delta }\tau `$, which is given by
$$\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}n+2}=\frac{i}{2}\left(\frac{n}{n+2}\right)\left(\mathrm{\Delta }_{12}\tau _{345\mathrm{}n+2}\mathrm{\Delta }_{13}\tau _{245\mathrm{}n+2}\right)+\mathrm{cyc}_{1\mathrm{}n+2}.$$
We then concentrate on the expression $`\overline{\mathrm{\Delta }}\mathrm{\Delta }\tau `$ by recalling, first of all, that the operation $`\overline{\mathrm{\Delta }}`$ on $`\mathrm{\Delta }\tau `$ consists of a contraction of two indices of the tensor $`\mathrm{\Delta }\tau `$ with the antisymmetric tensor $`\mathrm{\Delta }^{ab}`$. It is then clear, since $`\mathrm{\Delta }\tau `$ itself is built from the tensors $`\tau _{1\mathrm{}n}`$ and $`\mathrm{\Delta }_{ab}`$, that
$$\overline{\mathrm{\Delta }}\mathrm{\Delta }\tau =A+B+C,$$
where $`A`$ contains the terms where $`\mathrm{\Delta }^{ab}`$ is contracted uniquely with $`\mathrm{\Delta }_{ab}`$ and $`B`$ contains terms in which the two indices in $`\mathrm{\Delta }^{ab}`$ are contracted one with $`\mathrm{\Delta }_{ab}`$ and one with $`\tau _{1\mathrm{}n}`$. Finally $`C`$ consists of the remaining terms, with contractions of $`\mathrm{\Delta }^{ab}`$ only with $`\tau _{1\mathrm{}n}`$. We will prove in the sequel that
$`A`$ $`=`$ $`{\displaystyle \frac{1}{2}}D\tau `$
$`B`$ $`=`$ $`n\tau `$
$`C`$ $`=`$ $`\mathrm{\Delta }\overline{\mathrm{\Delta }}\tau ,`$
thus proving the lemma.
Let me start by concentrating on the terms in $`A`$. First recall the expression for $`\overline{\mathrm{\Delta }}\mathrm{\Delta }\tau `$
$$\overline{\mathrm{\Delta }}\mathrm{\Delta }\tau =\frac{i}{2}\left(\frac{n+2}{n^2}\right)\mathrm{\Delta }^{ab}\left[n\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}nab}+\mathrm{}+1\left(\mathrm{\Delta }\tau \right)_{1a2\mathrm{}nb}\right]+\mathrm{cyc}_{1\mathrm{}n}$$
(78)
The only terms in the above equation which contribute to $`A`$ are the first, second and last within the square bracket. In particular the first term reads
$`{\displaystyle \frac{i}{2}}\left({\displaystyle \frac{n+2}{n}}\right)\mathrm{\Delta }^{ab}\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}nab}+\mathrm{cyc}_{1\mathrm{}n}`$
$`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{\Delta }^{ab}\left[\left(\mathrm{\Delta }_{12}\tau _{3\mathrm{}nab}\mathrm{\Delta }_{13}\tau _{24\mathrm{}nab}\right)+\mathrm{cyc}_{1\mathrm{}nab}\right]+\mathrm{cyc}_{1\mathrm{}n}`$
$`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{\Delta }^{ab}\mathrm{\Delta }_{ab}\tau _{1\mathrm{}n}+\mathrm{cyc}_{1\mathrm{}n}+\mathrm{terms}\mathrm{not}\mathrm{in}A`$
$`=`$ $`n{\displaystyle \frac{D}{4}}\tau _{1\mathrm{}n}+\mathrm{terms}\mathrm{not}\mathrm{in}A,`$
where we have used that
$$\mathrm{\Delta }_{ab}\mathrm{\Delta }^{ab}=D$$
and that $`\tau `$ is cyclic. Similarly, the second term is given by
$`{\displaystyle \frac{i}{2}}\left({\displaystyle \frac{n+2}{n^2}}\right)(n1)\mathrm{\Delta }^{ab}\left(\mathrm{\Delta }\tau \right)_{1\mathrm{}anb}+\mathrm{cyc}_{1\mathrm{}n}`$
$`=`$ $`\left({\displaystyle \frac{n1}{4n}}\right)\mathrm{\Delta }^{ab}\left[\left(\mathrm{\Delta }_{12}\tau _{3\mathrm{}anb}\mathrm{\Delta }_{13}\tau _{24\mathrm{}anb}\right)+\mathrm{cyc}_{1\mathrm{}anb}\right]+\mathrm{cyc}_{1\mathrm{}n}`$
$`=`$ $`\left(n1\right){\displaystyle \frac{D}{4}}\tau _{1\mathrm{}n}+\mathrm{terms}\mathrm{not}\mathrm{in}A.`$
Finally the last term
$`{\displaystyle \frac{i}{2}}\left({\displaystyle \frac{n+2}{n^2}}\right)\mathrm{\Delta }^{ab}\left(\mathrm{\Delta }\tau \right)_{1a2\mathrm{}nb}+\mathrm{cyc}_{1\mathrm{}n}`$
$`=`$ $`{\displaystyle \frac{1}{4n}}\mathrm{\Delta }^{ab}\left[\left(\mathrm{\Delta }_{1a}\tau _{2\mathrm{}nb}\mathrm{\Delta }_{12}\tau _{a3\mathrm{}nb}\right)+\mathrm{cyc}_{1a2\mathrm{}nb}\right]+\mathrm{cyc}_{1\mathrm{}n}`$
$`=`$ $`{\displaystyle \frac{D}{4}}\tau _{1\mathrm{}n}+\mathrm{terms}\mathrm{not}\mathrm{in}A`$
Summing the three contributions we obtain
$$A=\frac{D}{2}\tau $$
as was to be shown.
We now move to the analysis of the terms in $`B`$. Again we consider equation (78), and we focus once again on the first term in the square brackets
$$\frac{1}{4}\mathrm{\Delta }^{ab}\left[\left(\mathrm{\Delta }_{12}\tau _{3\mathrm{}nab}\mathrm{\Delta }_{13}\tau _{24\mathrm{}nab}\right)+\mathrm{cyc}_{1\mathrm{}nab}\right]+\mathrm{cyc}_{1\mathrm{}n}$$
We concentrate on the terms in $`B`$, with $`\mathrm{\Delta }^{ab}`$ contracted both with $`\mathrm{\Delta }_{ab}`$ and with $`\tau _{1\mathrm{}n}`$
$`{\displaystyle \frac{1}{4}}\mathrm{\Delta }^{ab}\left[\mathrm{\Delta }_{na}\tau _{b1\mathrm{}n1}+\mathrm{\Delta }_{b1}\tau _{2\mathrm{}na}\right]+\mathrm{cyc}_{1\mathrm{}n}`$ (79)
$`{\displaystyle \frac{1}{4}}\mathrm{\Delta }^{ab}\left[\mathrm{\Delta }_{a1}\tau _{b2\mathrm{}n}+\mathrm{\Delta }_{nb}\tau _{a1\mathrm{}n1}+\mathrm{\Delta }_{n1,a}\tau _{nb1\mathrm{}n2}+\mathrm{\Delta }_{b2}\tau _{13\mathrm{}na}\right]+\mathrm{cyc}_{1\mathrm{}n}`$
$`=`$ $`{\displaystyle \frac{1}{4n}}\left(4n\tau _{1\mathrm{}n}n\tau _{213\mathrm{}n}n\tau _{23\mathrm{}1n}\right)+\mathrm{cyc}_{1\mathrm{}n}`$
The second term in (78)
$$\left(\frac{n1}{4n}\right)\mathrm{\Delta }^{ab}\left[\left(\mathrm{\Delta }_{12}\tau _{3\mathrm{}anb}\mathrm{\Delta }_{13}\tau _{24\mathrm{}anb}\right)+\mathrm{cyc}_{1\mathrm{}anb}\right]+\mathrm{cyc}_{1\mathrm{}n}$$
gives as contribution to $`B`$
$`\left({\displaystyle \frac{n1}{4n}}\right)\mathrm{\Delta }^{ab}\left[\mathrm{\Delta }_{b1}\tau _{23\mathrm{}an}+\mathrm{\Delta }_{nb}\tau _{1\mathrm{}n1a}+\mathrm{\Delta }_{n1,a}\tau _{nb1\mathrm{}n2}+\mathrm{\Delta }_{an}\tau _{b1\mathrm{}n1}\right]+\mathrm{cyc}_{1\mathrm{}n}`$ (80)
$`\left({\displaystyle \frac{n1}{4n}}\right)\mathrm{\Delta }^{ab}\left[\mathrm{\Delta }_{b2}\tau _{13\mathrm{}an}+\mathrm{\Delta }_{n2,a}\tau _{n1,n,b,1\mathrm{},n3}\right]+\mathrm{cyc}_{1\mathrm{}n}`$
$`=`$ $`\left({\displaystyle \frac{n1}{4n}}\right)\left[\tau _{213\mathrm{}n}+\tau _{23\mathrm{}1n}\tau _{231\mathrm{}n}\tau _{23\mathrm{},1,n1,n}2\tau _{1\mathrm{}n}\right]+\mathrm{cyc}_{1\mathrm{}n}`$
In order to write equations in a compact form let us introduce some notation. We explain the notation in the case $`n=4`$, but then we continue the proof in a general setting. We wish to consider tensors $`\tau _{\mathrm{}}`$ with the indices given by $`2,3,4`$ in increasing order, and with the index $`1`$ in a given position. For example, if the index $`1`$ is in the $`3`$rd position we are considering the tensor $`\tau _{2314}`$, which we will denote with the following box
$$\tau _{2314}\begin{array}{cccc}0& 0& 1& 0\end{array}.$$
Moreover, linear combinations of the various tensors will also be denoted by a single box in the following obvious way
| $`a`$ | $`b`$ | $`c`$ | $`d`$ |
| --- | --- | --- | --- |
$`=`$ $`a\begin{array}{cccc}1& 0& 0& 0\end{array}+b\begin{array}{cccc}0& 1& 0& 0\end{array}`$
$`+c\begin{array}{cccc}0& 0& 1& 0\end{array}+d\begin{array}{cccc}0& 0& 0& 1\end{array}.`$
Let us now return to the general proof, by first showing some simple properties of the box just introduced. Cyclicity of the tensor $`\tau `$ implies that
$$\begin{array}{ccccc}1& 0& \mathrm{}& 0& 0\end{array}=\begin{array}{ccccc}0& 0& \mathrm{}& 0& 1\end{array}.$$
(83)
Moreover, cyclic gauge invariance of $`\tau `$ implies the two equivalent identities
$$\begin{array}{ccccc}1& 1& \mathrm{}& 1& 0\end{array}=\begin{array}{ccccc}0& 1& \mathrm{}& 1& 1\end{array}=0$$
which can be summed to obtain
$$\begin{array}{ccccc}1& 2& \mathrm{}& 2& 1\end{array}=0.$$
(84)
We can use this notation to compactly rewrite the first two terms in $`\overline{\mathrm{\Delta }}\mathrm{\Delta }\tau `$ given by equations (79) and (80). They now read
$$\frac{1}{4n}\begin{array}{ccccccc}2n& n& 0& \mathrm{}& 0& n& 2n\end{array}+\mathrm{cyc}_{1\mathrm{}n}$$
and
$$\frac{1}{4n}\begin{array}{ccccccc}(n1)& n1& (n1)& \mathrm{}& (n1)& n1& (n1)\end{array}+\mathrm{cyc}_{1\mathrm{}n}$$
We notice that both terms are represented by boxes which are symmetric about a vertical axis of symmetry. Moreover we can do computations similar to the ones above to convince ourselves that the $`(k+1)`$–th term in $`\overline{\mathrm{\Delta }}\mathrm{\Delta }\tau `$ is given by the following box (we just show the left side of the box, since the right side is just its mirror copy)
$$\frac{1}{4n}\begin{array}{cccccccc}0& \mathrm{}& 0& nk& (nk)& nk& (nk)& \mathrm{}\end{array}+\mathrm{cyc}_{1\mathrm{}n}$$
with $`k2`$ zeros on the left before the term $`nk`$. Writing at once all the contributions to $`B`$ we get the following tableau
$$\frac{1}{4n}\begin{array}{ccccc}2n& n& & & \\ (n1)& n1& (n1)& & \\ n2& (n2)& n2& (n2)& \\ & n3& (n3)& n3& (n3)\end{array}+\mathrm{cyc}_{1\mathrm{}n}$$
where we must sum the coefficients in each column. The result is then
$`{\displaystyle \frac{1}{4n}}\begin{array}{ccccccc}2n1& 2& 2& \mathrm{}& 2& 2& 2n1\end{array}+\mathrm{cyc}_{1\mathrm{}n}`$
$`=`$ $`{\displaystyle \frac{1}{2}}\begin{array}{ccccc}1& 0& \mathrm{}& 0& 1\end{array}+\mathrm{cyc}_{1\mathrm{}n}=\begin{array}{cccc}1& 0& \mathrm{}& 0\end{array}+\mathrm{cyc}_{1\mathrm{}n}`$
where we have used equations (84) and (83). Going back to the usual tensor notation we have then obtained
$$B=\tau _{1\mathrm{}n}+\mathrm{cyc}_{1\mathrm{}n}=n\tau .$$
We conclude the proof by showing that $`C=\mathrm{\Delta }\overline{\mathrm{\Delta }}\tau `$. Consider again the first term in $`\overline{\mathrm{\Delta }}\mathrm{\Delta }\tau `$
$$\frac{1}{4}\mathrm{\Delta }^{ab}\left[\left(\mathrm{\Delta }_{12}\tau _{3\mathrm{}nab}\mathrm{\Delta }_{13}\tau _{24\mathrm{}nab}\right)+\mathrm{cyc}_{1\mathrm{}nab}\right]+\mathrm{cyc}_{1\mathrm{}n}$$
and concentrate on the terms which contain $`\tau _{\mathrm{}ab}`$i.e. terms for which the indices in $`\tau `$ contain $`a`$ just before $`b`$
$`{\displaystyle \frac{1}{4}}\left(\mathrm{\Delta }_{12}\mathrm{\Delta }^{ab}\left(\tau _{3\mathrm{}nab}+\mathrm{cyc}_{3\mathrm{}n}\right)+\mathrm{cyc}_{1\mathrm{}n}\right)`$ (87)
$`{\displaystyle \frac{1}{4}}\left(\mathrm{\Delta }_{13}\mathrm{\Delta }^{ab}\left(\tau _{24\mathrm{}nab}+\mathrm{cyc}_{24\mathrm{}n}\right)+\mathrm{cyc}_{1\mathrm{}n}\right)`$
$`+{\displaystyle \frac{1}{4}}\left(\mathrm{\Delta }_{12}\mathrm{\Delta }^{ab}\tau _{3\mathrm{}nab}+\mathrm{cyc}_{1\mathrm{}n}\right)`$
Terms with $`\tau _{\mathrm{}ab}`$ are also contained in the second and third contribution to $`\mathrm{\Delta }\overline{\mathrm{\Delta }}\tau `$ in (78)
$`\left({\displaystyle \frac{n1}{4n}}\right)\mathrm{\Delta }^{ab}\left[\left(\mathrm{\Delta }_{12}\tau _{3\mathrm{}anb}\mathrm{\Delta }_{13}\tau _{24\mathrm{}anb}\right)+\mathrm{cyc}_{1\mathrm{}anb}\right]+\mathrm{cyc}_{1\mathrm{}n}`$
$`\left({\displaystyle \frac{n2}{4n}}\right)\mathrm{\Delta }^{ab}\left[\left(\mathrm{\Delta }_{12}\tau _{3\mathrm{}a,n1,nb}\mathrm{\Delta }_{13}\tau _{24\mathrm{}a,n1,nb}\right)+\mathrm{cyc}_{1\mathrm{}a,n1,n,b}\right]+\mathrm{cyc}_{1\mathrm{}n}`$
They are respectively
$$2\left(\frac{n1}{4n}\right)\left(\mathrm{\Delta }_{12}\mathrm{\Delta }^{ab}\tau _{3\mathrm{}nab}+\mathrm{cyc}_{1\mathrm{}n}\right)$$
(88)
and
$$\left(\frac{n2}{4n}\right)\left(\mathrm{\Delta }_{12}\mathrm{\Delta }^{ab}\tau _{3\mathrm{}nab}+\mathrm{cyc}_{1\mathrm{}n}\right).$$
(89)
Summing equations (87),(88) and (89) we then obtain (the last line in (87) is canceled by (88) and (89))
$`{\displaystyle \frac{1}{4}}\left(\mathrm{\Delta }_{12}\mathrm{\Delta }^{ab}\left(\tau _{3\mathrm{}nab}+\mathrm{cyc}_{3\mathrm{}n}\right)+\mathrm{cyc}_{1\mathrm{}n}\right)`$
$`{\displaystyle \frac{1}{4}}\left(\mathrm{\Delta }_{13}\mathrm{\Delta }^{ab}\left(\tau _{24\mathrm{}nab}+\mathrm{cyc}_{24\mathrm{}n}\right)+\mathrm{cyc}_{1\mathrm{}n}\right).`$
But this is exactly the result which is obtained by computing terms proportional to $`\tau _{\mathrm{}ab}`$ in $`\mathrm{\Delta }\overline{\mathrm{\Delta }}\tau `$, since
$$\mathrm{\Delta }\overline{\mathrm{\Delta }}\tau =\frac{i}{2}\left(\frac{n2}{n}\right)\left[\mathrm{\Delta }_{12}\left(\overline{\mathrm{\Delta }}\tau \right)_{34\mathrm{}n}\mathrm{\Delta }_{13}\left(\overline{\mathrm{\Delta }}\tau \right)_{24\mathrm{}n}+\mathrm{cyc}_{1\mathrm{}n}\right]$$
and, if we concentrate on terms of the form $`\tau _{\mathrm{}ab}`$,
$$\left(\overline{\mathrm{\Delta }}\tau \right)_{34\mathrm{}n}=\frac{i}{2}\left(\frac{n}{n2}\right)\mathrm{\Delta }^{ab}\left(\tau _{3\mathrm{}nab}+\mathrm{cyc}_{3\mathrm{}n}\right).$$
Similar arguments can be used for terms proportional to $`\tau _{\mathrm{}a\mathrm{}b}`$, with different number of indices between $`a`$ and $`b`$. We have then shown that $`C=\mathrm{\Delta }\overline{\mathrm{\Delta }}\tau `$, thus completing the proof. $`\mathrm{}`$
###### Lemma 13
(Section 6.2) If $`F=[Z,\overline{Z}]`$, then
$$\mathrm{\Delta }F^3=2F^4+\frac{1}{2}F[DF,\overline{D}F].$$
PROOF. We compute the word $`wW_4`$
$$w=\mathrm{\Delta }F^3.$$
First we have that
$$w=3F^2\left(\mathrm{\Delta }F\right)=\frac{3}{2}\left(F^4+F[FZ,F\overline{Z}]\right).$$
We then use cyclicity to show that
$`F[FZ,F\overline{Z}]`$ $`=`$ $`F^4+F^2[F,\overline{Z}]Z+F^2[Z,F]\overline{Z}`$
$`=`$ $`F^4+ZF^2[F,\overline{Z}]+\overline{Z}F^2[Z,F]`$
$`=`$ $`F^4+F^2\overline{Z}[Z,F]+\overline{Z}F^2[Z,F]`$
and therefore that
$$w=3F^4+\frac{3}{2}\left(F^2\overline{Z}[Z,F]+\overline{Z}F^2[Z,F]\right).$$
(90)
Now, since
$$F^2\overline{Z}[Z,F]=[F^2\overline{Z},Z]F=F^4\overline{Z}F^2[Z,F]F\overline{Z}F[Z,F]$$
and since
$`F\overline{Z}F[Z,F]`$ $`=`$ $`FDF\overline{D}F\overline{Z}F^2[Z,F]`$
$`=`$ $`F\overline{D}FDFF^2\overline{Z}[Z,F]`$
$`=`$ $`{\displaystyle \frac{1}{2}}\left(F[DF,\overline{D}F]\overline{Z}F^2[Z,F]F^2\overline{Z}[Z,F]\right),`$
we obtain that
$$\frac{3}{2}\left(F^2\overline{Z}[Z,F]+\overline{Z}F^2[Z,F]\right)=F^4+\frac{1}{2}F[DF,\overline{D}F].$$
The above fact, together with (90), concludes the proof of the lemma. $`\mathrm{}`$
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# The Architecture of Idiotypic Networks: Percolation and Scaling Behaviour
## I Introduction
B-lymphocytes carry on their surface highly specific receptors, so-called antibodies. If these receptors detect complementary structures, the lymphocyte is stimulated to proliferate and after several generations and differentation into plasma cells secretes antibodies of the same specificity. B-lymphocytes and antibodies of a given specificity are said to have a certain idiotype. Complementary structures to an idiotype are antigen or other, anti-idiotypic antibodies. Between B-lymphocytes of different idiotype thus emerges a functional network of mutual stimulation and inhibition, the idiotypic network . The idiotypic network is supposed to at least partially contribute to the functionality of the immune system as, e.g., immunological memory or suppression of autoreactive clones. Though quantitative data are very hard to access by experiment there are some recent observations which underline the importance of idiotypic interactions .
New idiotypes are produced in the bone marrow or due to hypermutation during the proliferation of stimulated lymphocytes which introduces a random metadynamics of the repertoire. At a given moment the random network has a certain architecture. The aim of this paper is to give a statistical description of this architecture. Knowledge of the typical architecture of the idiotypic network is crucial for describing the population dynamics of the interacting B-lymphocytes and antibodies, cf. , which is however not the subject of this paper.
Derived from hypotheses of theoretical immunology (cf. , for a recent review see ) we present a statistical analysis of bit-string based networks, which shows, that this approach is well suited to reproduce reasonable network structures. Especially, our considerations show, that realistic network topologies can be conceived as an extension of first approaches to that problem which assumed a Bethe-lattice structure. Suggestions of how loops should be added to such structures have been made previously . For this, our work provides a very natural access.
Generally, idiotypic networks are supposed to realize a tradeoff between two basic requirements: they should contain a great number of small isolated components, but on the other hand still be able to respond to arbitrary antigen, that means being complete. Small components are thought necessary to store information about previously encountered antigen . The existence of such components does obviously demand a low connectivity of the network. With completeness on the other hand it is assumed that of a great number of antibodies each is able to detect many different types of antigen. Hence the network connectivity should not be too low.
With regard to the underlying biological problem the one great–many small clusters situation is thus worthwhile to earn special attention. Theoretical immunologists suppose that the idiotypic network should consist of a large number of small clusters. On the other hand, as a consequence of relatively high connectivities, also a large component should be contained within the idiotypic network, which is denoted as its central part. This great cluster could play an important role in the control of autoreactive clones .
In the following the bit-string model of will be explained briefly. Basically, antibodies are identified with bit chains of a given length $`d`$. Thus, there are $`2^d`$ different antibody types. The set of all conceivable antibodies, i.e. the potential repertoire, is then represented by $`\{i=(i_1,\mathrm{},i_d),i_j\{0,1\}\}`$. Estimating the probable size of networks that can be complete in the above sense shows that $`d32\mathrm{}36`$ should be a good value for realistic models .
Antibodies recognize each other if they have complementary structures, i.e. if they are represented by perfectly complementary bit strings. If there are small deviations from the exact structural complementarity a matching is still possible, though with lower affinity. This is described by so-called matching rules.
For example we imagine that antibodies react if the respective structures are complementary except one small region. Using the language of our model the corresponding bit-strings should be inverse except one bit that belongs to the non-complementary region. We call this kind of rule a one-mismatch rule. Two-mismatch rules, which allow reactions between antibodies that are complementary up to two non-matching regions, are accordingly defined by connecting bit-strings that are complementary with two exceptions allowed. Naturally, the multitude of possible rules is not exhausted by mismatch rules. Rules that express matching of mutually shifted antibody parts could be conceived as well. For the sake of simplicity this work is confined to those matching rules that are associated with reactions of highest affinity, i.e. inversion, one- and two-mismatch rules.
Mathematically, the set of all antibody types (represented by bit-strings or vertices of a hypercube) together with all possible reactions between them defines a graph. Thus, we want to call
$`G_d^1=`$ $`(\{v=(v_1,\mathrm{},v_d),v_j\{0,1\}\},\{v\text{ connected with }w`$ (1)
$`\text{if }v_i=\overline{w_j}\text{ for all }i,j=1,\mathrm{},d\text{ except one position}`$ (2)
$`\text{at the most}\})`$ (3)
the one-mismatch base graph and
$`G_d^2=`$ $`(\{v=(v_1,\mathrm{},v_d),v_j\{0,1\}\},\{v\text{ connected with }w`$ (4)
$`\text{if }v_i=\overline{w_j}\text{ for all }i,j=1,\mathrm{},d\text{ except two positions}`$ (5)
$`\text{at the most}\})`$ (6)
the two-mismatch base graph, respectively. To obtain a better impression of one-mismatch graphs they can be compared to hypercubes of the same dimension. Looking, e.g., on hypercubes of dimension $`d=3`$ edges of one-mismatch graphs are represented by all space and side diagonals. This, together with a picture of an one-mismatch graph in dimension $`4`$ is shown in Figs. 1 and 2, respectively.
The actual repertoire, i.e. the set of all types of idiotypes that are really existing within the body at a given time, is a subset of the potential repertoire. There are profound reasons to believe that the actual repertoire should be distributed randomly within the potential repertoire . A way of realizing this is to choose antibodies with a certain probability $`p`$. Then the idiotypic network is represented by the random graph that is composed of all occupied vertices (antibodies) and their connections (possible reactions). Clearly, this is a site percolation problem .
Percolation problems on standard lattices are far away from beeing a new field of research. Percolation on one-mismatch graphs, however, bears some major difference to ordinary percolation.
Percolation is always connected to an abrupt change of some system property (the standard property is the existence of a connecting path from the upper to the lower boundary on 2-dimensional lattices) if the occupation probability passes a certain value, the percolation threshold. This transition becomes sharp in the limit of infinite systems.
On lattices there is no question how to increase the system size to infinity. However, sequences of one-mismatch graphs exhibit relations between small and large systems that are essentially different from similar characteristics of lattices which are created by the multiplication of a unit cell. Understanding of these distinctions comes from studying the structure of fully occupied one-mismatch graphs (base graphs) which will be performed in the ensuing Section.
In this paper we investigate the architecture of functional networks built by constituents of randomly generated characteristics interacting with complementary constituents. Our motivation, the formulation of the problem and the interpretation of the results are all in the context of modeling the immune system using language and methods of statistical physics. However, our results are essentially independent of the specific immunological interpretation and could be of broader interest. The situation of interacting randomly generated complementary constituents is quite general. For example think of chemical reactions (origin of life), ecosystems, social networks.
Generally, the conception of the paper is as follows: in Sec. II topological properties of base-graphs will be studied. Section III gives an analysis of the underlying percolation problem, which then leads, together with applications of random graph theory from Sec. IV and a study of intrinsic structures of great clusters in Sec. V, to conclusions about parameter regimes, in which bit-string model induced random graphs resemble idiotypic networks. In the appendix another method to calculate thresholds via the renormalization of small cells will be discussed.
## II Structure of the Base Graph for the One-Mismatch Rule
The above definition of matching rules allows an easy calculation of distances between vertices of $`G_d^1`$. Let $`i,jG_d^1`$ be vertices, $`d_H`$ their Hamming distance (i.e. the number of different bits between them) and
$$d_G(i,j)=\underset{\begin{array}{c}\text{pathes }w\text{ in }G_d^1\\ \text{ connecting }i\text{ and }j\end{array}}{inf}l(w)\text{ ,}$$
(7)
where $`l(w)`$ denotes the length of a path $`w`$, a metrics. Then it holds that if
$`d_H<d/2`$ $`d_G`$ $`=\{\begin{array}{cc}d_H\hfill & \text{for }d_H\text{ even }\hfill \\ d_H+1\hfill & \text{for }d_H\text{ odd}\hfill \end{array}`$ (8)
$`d_H=d/2`$ $`d_G`$ $`=d/2\text{ for }d\text{ even}`$ (9)
$`d_H>d/2`$ $`d_G`$ $`=\{\begin{array}{cc}dd_H+1\hfill & \text{for }dd_H\text{ even }\hfill \\ dd_H\hfill & \text{for }dd_H\text{ odd ,}\hfill \end{array}`$ (10)
which shows that the maximum distance $`D_{\text{max}}(d)`$ between vertices is roughly the half of that on ordinary hypercubes. This is essentially due to the inversion rule. Furthermore, it should be noticed that $`D_{\text{max}}(d)`$ does not change when increasing the dimension by one from odd ($`d=2n+1`$) to even ($`d=2n+2`$) values.
To obtain additional information about the toplogy we distinguish vertices corresponding to their distance $`s`$ to an arbitrarily chosen origin, e.g. $`v=0=(0,\mathrm{},0)`$, whose choice is arbitrary due to the specific construction-rule of base graphs. We denote the set of vertices with distance $`d_G(v,i)`$ from $`v`$ by
$`E_s=\{iG_d^1|d_G(v,i)=s\}\text{.}`$ (11)
Applying (9) it is possible to compute the number of vertices $`|E_s|`$ belonging to $`E_s`$ (mass distribution) by simple combinatorics to
$`|E_s|=\{\begin{array}{cc}\left(\begin{array}{c}d\\ d_G1\end{array}\right)\hfill & \text{ if }d_G=\frac{d+1}{2}\hfill \\ & \\ \left(\begin{array}{c}d+1\\ d_G\end{array}\right)\hfill & \text{ if }d_G\frac{d}{2}\hfill \end{array}\text{.}`$ (12)
A visualization of Eq. (12) in Fig. 3 shows a major difference between odd and even dimensions $`d`$. Generally, $`|E_s|`$ grows monotonically with increasing distance. For odd dimensions, however, the number of vertices with maximal distance $`D_{\text{max}}`$ from $`v`$ is substantially smaller than that with distance $`D_{\text{max}}1`$
Differences in the overall structure of graphs in odd and even dimensions $`d`$ become even more obvious if the number of links connecting vertices of the same set $`E_s`$ is considered. Using (9), a somewhat lengthy but straightforward calculation shows that connected vertices in the same distance $`s`$ from the origin $`v`$ occur in
$`d`$ even only for $`s=D_{\text{max}}=d/2`$ (13)
$`d`$ odd only for $`s=D_{\text{max}}=(d+1)/2\text{ or}`$ (14)
$`s=D_{\text{max}}1\text{.}`$ (15)
Analogously, the number of links that connect vertices of $`E_s`$ with vertices belonging to $`E_{s+1}`$ may be computed to $`d+1s`$ for $`s<D_{\text{max}}1`$.
As a consequence there must exist loops of even as well as of odd length. Moreover, all short loops with length less than $`d`$ must be even (since only loops comprising links between vertices in the same distance from any loop element can be of an odd length).
An important property of standard systems in percolation theory is that the smaller system is always contained within the larger one. For one-mismatch graphs, however, this is not possible. Let $`d_1<d_2`$ be dimensions of $`G_{d_1}^1`$ and $`G_{d_2}^1`$, respectively. If $`G_{d_1}^1`$ were contained in $`G_{d_2}^1`$ it would follow directly that $`G_{d_2}^1`$ should have uneven loops with length less then $`d_1`$. This gives a contradiction since $`G_{d_2}^1`$ contains no odd loops smaller than $`d_2`$, whereas $`d_1<d_2`$ by assumption.
Resuming these results we state that percolation on one mismatch graphs differs from standard percolation. A crucial point for this distinction is the way how global system properties change if the system size approaches infinity. Nevertheless, methods of percolation theory can be applied to this kind of problem. Even more, choosing an appropriate majority rule renormalization group procedures can be extended to one-mismatch-graphs (see appendix A).
To the best of our present knowledge percolation problems on graphs with similar properties have not been dealt with. Even percolation on regular hypercubes has found relatively little attention in the physical literature .
## III Cluster size distribution
Having collected some simple properties of the underlying base graphs we consider now vertices of that graphs occupied with a given probability $`p`$. The set of all occupied vertices $`\mathrm{\Gamma }`$ together with all bonds, which connect two vertices belonging to $`\mathrm{\Gamma }`$ forms a random graph. In terms of percolation theory a maximal set of connected vertices is called cluster.
What is the probability that an arbitrarily chosen vertex belongs to a cluster of size $`|C|`$? This question has already been addressed in where three major regimes of the system have been identified (see also the numerical results in Fig. 4):
(i) The first typical situation arises for small values of the occupation probability $`p`$. Then only small clusters (whose individual size doesn’t make up a finite fraction of the whole system) are expected to appear.
(ii) Increasing $`p`$ and approaching infinite dimensions, a sharp transition (in the sequel denoted as percolation transition) to an one great cluster–many small clusters regime was found. At the percolation transition some characteristics obey scaling laws. These and an approximation to compute percolation thresholds will be dealt with in Sec. IV.
The one great–many small clusters situation is worthwhile to earn special attention (cf. Sec. I). The great cluster could play an important role in the control of autoreactive clones. To fulfill that purpose it is believed to have a certain internal structure (cf. , for experimental data cf. ), which will be discussed in more detail in Sec. VI. Thus, as stated already in the simple one–parametric bit chain model exhibits an interesting similarity to the idiotypic network.
(iii) Finally, if the occupation probability is further increased, a state where random graphs consist of one connected component only will be reached. Relying on some general results of the theory of random graphs it can be shown that this indeed marks a second (in the limit of infinite systems) sharp transition which we call denseness transition. In Sec. V these considerations are resumed and discussed with reference to the completeness of the immune system.
Exactly the same two thresholds have been considered earlier for a different class of random graphs and more recently in the context of networks of RNA secondary structures .
## IV Scaling laws and Bethe approximation
A common method to analyze percolation problems is the introduction of perimeter polynomials $`D_{|C|}`$. Using the standard notation we define
$`p(|C|)=p^{|C|}{\displaystyle \underset{S_{\text{free}}(C)}{}}(1p)^{|S_{\text{free}}(C)|}=p^{|C|}D_{|C|}(1p)\text{,}`$ (16)
where $`p(|C|)`$ denotes the probability that a vertex on the base graph belongs to a cluster of size $`|C|`$. $`S_{\text{free}}(C)`$ means the free surface, i.e. the set of all unoccupied vertices of the base graph that are adjacent to vertices of $`C`$.
Generally, there are few problems which allow the explicit calculation of $`D_{|C|}`$ for arbitrary cluster size $`|C|`$. Notwithstanding, it is always possible to compute $`D_{|C|}`$ for small values of $`|C|`$. This provides a basis for the application of series expansion techniques. We investigate the structure of small clusters up to size $`|C|=5`$ on one-mismatch-graphs. As a result for large base graphs ($`d>5`$) we find the following perimeter polynomials for one-mismatch-graphs
$`D_1(q)`$ $`=q^{d+1}\text{ ,}`$ (17)
$`D_2(q)`$ $`=(d+1)q^{2d}\text{ ,}`$ (18)
$`D_3(q)`$ $`={\displaystyle \frac{3}{2}}(d+1)dq^{3d2}\text{ ,}`$ (19)
$`D_4(q)`$ $`=q^{4d4}(2(d1)d(d+1)+{\displaystyle \frac{d(d+1)}{2}}+`$ (20)
$`+`$ $`{\displaystyle \frac{2}{3q}}(d1)d(d+1))\text{ ,}`$ (21)
$`D_5(q)`$ $`={\displaystyle \frac{5}{2}}q^{5d6}(d1)d(d+1)(d2+{\displaystyle \frac{1}{q}}(d1)+`$ (22)
$`{\displaystyle \frac{1}{q^2}}+{\displaystyle \frac{1}{12q^3}}(d2))\text{ ,}`$ (23)
where $`q=1p`$.
In the following we show that in the limit of high dimensions random graphs on one-mismatch-graphs for small values of $`p`$ are very similar to random graphs on Bethe-lattices with the same coordination number. Arguments of this kind have already been applied to percolation on hypercubic lattices making use of the fact that loops in random graphs on sparsely occupied high dimensional lattices are infrequent. Indeed, the first terms of a high dimension expansion on a hypercubic lattice of dimension $`d`$ are given by the exact values of the percolation threshold for the Bethe-lattice with coordination number $`d`$.
A closer look at the perimeter polynomials (17-22) supports the hypothesis that the situation is quite similar for one-mismatch-graphs. For example, the contribution from loops to $`p(4)=p^4D_4(q)`$ is given by $`p^{\text{loop}}(4)=p^4q^{4d4}d(d+1)/2`$. Obviously, the relative weight of this term vanishes in the high dimension limit. Additional support comes from the observation that the one-mismatch base-graph $`G_d^1`$ contains hypercubes up to dimension $`d1`$. To elucidate this property we consider the following example in $`d=4`$. We construct a hypercube of dimension $`3`$ applying the one-mismatch rule allowing for a mismatch in one of the, e.g., last three bits. Starting for instance at the origin $`(0,0,0,0)`$ this yields $`(1,1,1,0)`$, $`(1,0,1,1)`$, $`1,1,0,1)`$. Iteration gives the new vertices $`(0,1,0,1)`$, $`(0,0,1,1)`$, $`(0,1,1,0)`$, and finally $`(1,0,0,0)`$. These $`8`$ vertices are connected by the matching rules like a 3-dimensional hypercube. Then follow the arguments of .
Consequently the percolation threshold for one-mismatch-graphs can be approximated to
$`p_c^{(1)}=1/d\text{,}`$ (24)
whereas for two-mismatch graphs holds
$`p_c^{(2)}=1/(d+d(d+1)/2)\text{.}`$ (25)
Corrections to (24) and (25) are of order $`O(d^2)`$ and $`O(d^3)`$, respectively. From considerations of Sec. I we find that the number of links connecting vertices in $`E_s`$ with vertices in $`E_{s+1}`$ decreases with growing distance $`s`$. Thus, we conclude that the corrections to $`p_c^{(1)}`$ must be positive.
Typically for percolation problems, certain system characteristics obey scaling laws at the percolation threshold. This kind of laws reflects the statistical self similarity of clusters on different length scales in that parameter regime. Fig. 6 displays simulation results for the cluster size distributions obtained for one-mismatch-graphs of dimensions $`10`$ and $`14`$ . Both data sets can be well described by the finite size scaling ansatz
$`p(|C|)=|C|^\tau F(|C|/|C|^{})\text{ ,}`$ (26)
where the function $`F(x)`$ is nearly a constant for $`x1`$ and will be more rapidly declining than any power law for arguments $`x1`$. Equation (26) should apply for clusters which are neither too small nor too large. Due to the limited size of the investigated systems, very large clusters behave in a different way than clusters of ‘standard’ size. The transition between the laws applying to these separate cases is marked by $`|C|^{}`$, which in turn depends on the extent of the system.
Analogous properties can be observed for base graphs defined by other matching rules. Figure 6 illustrates scaling behaviour for the case of graphs diced on two-mismatch-base graphs for values of $`p=p_c^{(2)}`$ cf. (25). As a natural consequence of adding supplementary edges distances between vertices will become smaller in comparison to one-mismatch-graphs. This gives an explanation for the fact, that typical sizes $`|C|^{}`$ are substantially smaller than those examined on one-mismatch-graphs with the same choice of $`d`$.
Furthermore, it can be seen in Fig. 6 that the scaling law (26) even well applies to very small clusters. It is useful to define a cluster size dependent exponent $`\tau _{|C|}`$ by
$`\tau _{|C|}={\displaystyle \frac{\mathrm{ln}\frac{p(|C|+1)}{p(|C|)}}{\mathrm{ln}\frac{|C|+1}{|C|}}}\text{ .}`$ (27)
Evaluating the perimeter polynomials (17)-(22), values of $`\tau _{|C|}`$ for small $`|C|`$ can be derived. Thus we obtain a sequence $`\{\tau _{|C|}\}`$ which should approach the true value of $`\tau `$ for large $`|C|`$.
As a matter of fact (26) changes to $`p(|C|)|C|^\tau `$ in the limit of infinite systems. Performing the limits $`d\mathrm{}`$ in Eqs. (17)-(22) and (27), we observe that $`\tau _{|C|}^{\mathrm{}}`$ obeys
$`\tau _{|C|}^{\mathrm{}}={\displaystyle \frac{1}{\mathrm{ln}(1+\frac{1}{|C|})}}|C|+1\text{ ,}`$ (28)
which by comparison of (27) with (28) and solving the recursion relation implies the law
$`p(|C|)={\displaystyle \frac{e^{|C|}|C|^{|C|2}}{|C|!}}`$ (29)
for the cluster size distribution. Equation (28) leads to $`\tau =3/2`$ for large clusters. The thus computed value $`\tau =3/2`$ is in accord with the exact result on the Bethe-lattice.
Since we investigate a model for a biological system our main interest is devoted to large, yet finite systems. In the sequence $`\{\tau _{|C|}\}`$ our best approximation for the real value of $`\tau `$ is $`\tau _4`$. A comparison between values of $`\tau `$ computed by $`\tau \tau _4`$ and exponents $`\tau `$ obtained by evaluating numerical data for small systems suggests that the involved approximation becomes rapidly more accurate if the system size is increased. Hence in $`d=32`$ we rely on $`\tau _4`$ and find $`\tau 1.5`$, which gives a result that also supports the previous assumption.
Subsequently it was our aim to find a quantity which gives an overall estimation of deviations between random graphs on one-mismatch-graphs from those on Bethe-lattices of equal coordination number. For this purpose it appears appropriate to investigate the results of an edge elimination procedure which computes the number of bonds belonging to loops. Similar procedures have very recently been applied to lattices.
The above aim is achieved by individually removing every edge of every cluster and calculating the number of the connected components of the resulting graph. An edge is called fractioning if two components are obtained cutting this edge. Clearly, only edges that belong to loops are not fractioning.
Then the ratio of fractioning bonds $`f(C)`$ to the overall number of bonds $`b(C)`$ of a cluster $`C`$, $`f(C)/b(C)`$, is some indication for the importance of loops within the structure of $`C`$. Distinguishing clusters according to their size $`|C|`$ we computed the mean fractioning ratio
$`f_{|C|}=<f(C^{})/b(C^{})>_{|C^{}|=|C|}\text{ .}`$ (30)
Except for very small clusters with trivial structure we expect a finite size scaling law
$`f_{|C|}=|C|^\lambda \widehat{F}(|C|/|\widehat{C}^{}|)`$ (31)
for the fragmentation rate $`f`$ at the percolation threshold, the validity of which is illustrated in Fig. 7.
Similar to $`\tau _{|C|}`$ a cluster size dependent exponent $`\lambda _{|C|}`$ may be defined. The only 4- and 5-clusters containing non–fragmenting bonds are the 4-loop (no fractioning bond) and the 4-loop with a tail ($`1/5`$ of all bonds fractioning) which give the contributions $`p^{\text{4-loop}}(4)`$ and $`p^{\text{tailed 4-loop}}(5)=5/2p^5q^{5d7}(d1)d(d+1)`$ to $`p(4)`$ and $`p(5)`$, respectively. Making use of (20) and (22) we then find
$`f_4={\displaystyle \frac{p(4)p^{\text{4-loop}}(4)}{p(4)}}={\displaystyle \frac{1}{1+\frac{1}{4(d1)(1+1/(3q))}}}`$ (32)
and
$`f_5={\displaystyle \frac{p(4)p^{\text{tailed 4-loop}}(5)}{p(4)}}+{\displaystyle \frac{1}{5}}{\displaystyle \frac{p^{\text{tailed 4-loop}}(5)}{p(4)}}=1{\displaystyle \frac{4}{5}}{\displaystyle \frac{1}{1+1/q+(d2)(1+q+1/(12q^2))}}\text{ .}`$ (33)
. Finally, inserting $`pp_c^{(1)}=d^1+O(d^2)`$ a rough estimate for $`\lambda `$ is $`\lambda _4=(\mathrm{ln}f_5\mathrm{ln}f_4)/\mathrm{ln}(5/4)`$ which gives
$`\lambda _4(d){\displaystyle \frac{1}{\mathrm{ln}(5/4)}}\left({\displaystyle \frac{393}{210^3}}d^1+{\displaystyle \frac{4029120}{810^6}}d^2+O(d^3)\right)\text{ .}`$ (34)
For $`d=10`$ we have $`\lambda 0.11`$ in very good agreement with the numerically obtained value $`0.12`$, cf. Fig. 7. For $`d=32`$ we find $`\lambda 0.03`$ already very close to $`\lambda =0`$ which holds for Bethe-lattices. Thus the value $`\lambda >0`$ measured on finite dimensional random graphs on one-mismatch base-graphs – which is caused by a small number of loops – should quantify the deviation from random graphs on Bethe-lattices. It appears that the occupation probability at the percolation threshold is still small enough to apply the Bethe approximation.
Yet, random graphs at $`p_c`$ are not exactly Bethe–like and contain a certain fraction of loops. Otherwise it would be impossible to distinguish subclusters according to their connectivity within any considerable connected component. This, however, is likely to be necessary, to explain the role of the central part of idiotypic networks properly .
Further investigations concerning the structure of great clusters will be made in Sec. V.
## V Density threshold and the completeness of the idiotypic network
Hitherto, our main interest was devoted to the question at which occupation probability a great cluster starts to exist. If $`p`$, however, is further increased it is also imaginable that a situation occurs, where the whole system consists of one great cluster only. This is referred to as the question about the connectivity property of random graphs.
In the theory of random graphs some general results that address related problems have been derived for so called sequences of configuration spaces and earlier for a different general class of random graphs in . It can be shown that the sequence of one-mismatch-base graphs $`\{G_d^1\}`$ fulfills all requirements of these configurations spaces , due to its rather technical nature the proof will be omitted here.
General results of can be applied here showing that for infinite systems there exists a threshold
$`p_{\text{conn}}=1\underset{d\mathrm{}}{lim}|G_d^1|^{1/\gamma _d}`$ (35)
($`\gamma _d`$ being the coordination number) with the property that almost all random graphs are connected for $`p>p_{\text{conn}}`$ while the set of all connected random graphs has measure zero for $`p<p_{\text{conn}}`$.
Moreover, for sequences of configuration spaces it holds that both the connectivity and the density treshold, i.e. the threshold for the property that there is no non-occupied site without occupied neighbours on the base graph, coincide. Using (35) we find
$`p_{\text{conn}}^{(1)}=p_{\text{dense}}^{(1)}=1/2`$ (36)
in case of one-mismatch-graphs and
$`p_{\text{conn}}^{(2)}=p_{\text{dense}}^{(2)}=0`$ (37)
for two-mismatch graphs.
Density of random graphs in our model can be directly translated into biological terms. Since antibodies and antigens are represented by the same sets of bit chains the property that there is no free site without occupied neighbour means that every antigen is sure to encounter a complementary antibody. Hence the idiotypic network is able to respond to any antigen. This is meant by the completeness axiom for the immune system.
Nevertheless it is somewhat difficult to reconcile the demands for denseness of random graphs and occurance of small clusters at the same time. However, it seems unlikely that completeness should be understood in this strict way. Rather it appears to be a better solution to consider the fact, that evolution has most likely driven the idiotypic network into such a kind of arrangement, that it is able to respond to variations of really existing antigen only. Thus we argue that completeness of the idiotypic network does not exactly match density of random graphs, but requires only the probability for the density of the corresponding random graph to be somewhat below $`1`$. So, random graphs could still comprise small clusters and be complete.
Another paradoxon arises if two-mismatch graphs are considered. For those the connectivity and percolation thresholds $`p_c^{(2)}`$ and $`p_{\text{conn}}^{(2)}`$ fall together. How could then small clusters and a seperate great component coexist if all random graphs are connected? Yet, for the case of finite systems it is clear that $`p_c^{(2)}(d)<p_{\text{conn}}^{(2)}(d)`$, i.e. a great cluster has to be formed first, before it can devour his competing small rivals.
Consequently, also for the case of two-mismatch-graphs the biologically interesting regime is well defined. For finite systems there is a range of occupation probabilities $`p_c^{(2)}(d)<p<p_{\text{conn}}^{(2)}(d)`$, where all requirements are met. Due to the fact that both probabilities $`p_c^{(2)}(d)`$ and $`p_{\text{conn}}^{(2)}(d)`$ are converging to zero in the limit $`d\mathrm{}`$ it follows that the extent of this range $`|p_c^{(2)}(d)p_{\text{conn}}^{(2)}(d)|`$ will also become very small for large systems.
In the next Section our interest will be shifted towards the intrinsic structure of great clusters. Therefrom further conclusions about biologically relevant parameter regimes can be drawn.
## VI Intrinsic structure of great clusters
Resuming results for percolation on mismatch-graphs from Sec. I-IV two phase transitions have been found to occur. The ranges below $`p_c`$ (since there is no great cluster) and above $`p_{\text{conn}}`$ (since there is only one great cluster) are of no interest concerning the biological background of the model. The one great cluster–many small clusters situation between both, however, is likely to fulfill some requirements for idiotypic networks (see Sec. II). In this Section further investigations into the network structure within this parameter regime will be made.
Insights into the structure of a great cluster can be obtained from the mass distribution of this cluster $`M_c(s)=\{vC|d(v,c)=s\}`$, i.e. the information, how many vertices have a certain distance $`s`$ from a starting point $`c`$. Due to the equality of all starting points $`c`$ we define
$`M(s)={\displaystyle \frac{1}{|C|}}{\displaystyle \underset{cC}{}}M_c(s)`$ (38)
and consider the mean value of $`M(s)`$ calculated over all clusters, whose size exceeds a minimum value $`0.5\times p\times 2^d`$. For the metrics $`d(,)`$ there are two distinguished choices, viz. metrics defined by (7) allowing paths in $`G=G_d^1`$, i.e. $`d(,)=d_G`$, and such restricting paths to $`C`$ itself, i.e. $`d(,)=d_C`$. Since $`d_C>d_G`$ we confine our investigations to $`d(,)=d_C`$ which provides a better ‘resolution’. In the following we will discuss some typical cluster compositions to obtain a survey of possible conclusions from mass distributions to cluster structures.
Basically, three different situations could be imagined (for (i) and (iii) see Fig. 8).
(i) A cluster could consist of some loosely clinged high connectivity regions (heaps), whose vertices are distinguished by a great number of connections with each other. Other vertices are bound with relatively few links. Clearly, strong connectivity within heaps means, that vertices belonging to them have nearly the same distance from all other vertices. Thus, the presence of many heaps should result in a great number of local extrema, caused by ‘looking from all heaps’ to each other. On the other hand, loosely bound vertices will smooth the mass distribution, i.e. reduce the sharpness of extrema.
(ii) A cluster could contain no distinguished parts at all. From (12) we know that till $`s=D_{\text{max}}`$ the number of vertices is increasing with the distance $`s`$. For small occupation probabilities $`p<p_{\text{conn}}`$ (for which not the whole base-graph has been ‘conquered’ yet) we expect a compromise between with $`s`$ declining probabilities that a vertex of that cluster has distance $`s`$ from an origin $`v`$ and the increasing number of vertices with $`s`$. Consequently, the mass distribution should exhibit one maximum. For large $`p>p_{\text{conn}}`$ mass distributions can simply be derived by multiplying (12) with $`p`$.
(iii) As a special case of (i) a cluster could be made of one heap and a certain fraction of loosely bound vertices. Then up to two maxima, caused by the heaped vertex concentration and the competing tendencies (see (ii)), respectively, are likely to occur. The sharpness of both maxima will depend on the fraction of the loosely bound vertices. Thus, in case of a large proportion of those vertices both extrema could be smoothened to just one broad maximum.
Fig. 9 shows simulation results for the normalized mass distribution for different values of $`p`$. Clearly, all distributions are marked by only one maximum, whose sharpness is increasing with growing values of $`p`$. In the vicinity of the density threshold clusters are already stretched over the whole extent of the base graph. Hence, for $`p=0.4`$ slightly below $`p_{\text{conn}}=0.5`$ the curve of the mass distribution looks similar to the exactly known distribution for the fully occupied base-graph given by Eq. (12), cf. also Fig. 3.
More interestingly, for small $`p`$’s in the vicinity of the percolation threshold relatively broad maxima occur, which could be an indication for cluster structures as described in (iii).
To prove the validity of this hypothesis we have applied the edge-elimination algorithm (see Sec. II) to great clusters. As a slight extension of the described procedure all fractioning edges are removed, thus splitting a cluster $`C`$ into the sequence of its doubly connected components, or heaps in the above sense, $`(C_1,\mathrm{},C_t)`$. Obviously, vertices belonging to doubly connected components are distinguished by their strong cliquishness in comparison to other vertices. Thus, the number of such components $`t`$ should allow to distinguish between the situations (i)-(iii).
Fig. 10 shows simulation data for the mean value $`t`$ of the resulting non-trivial parts after edge-elimination depending on $`p`$. The distribution is marked by one maximum, which is again a consequence of two competing tendencies. For small values of $`p`$ clusters are generally not doubly connected, but increasing $`p`$ leads to a larger proportion of vertices that belong to loops. On the other hand, there is a tendency that loops get connected by loops, i.e., seperate doubly connected cluster components are growing together. Thus, for large $`p`$ almost every cluster will consist of one doubly connected component only.
Results of Fig. 10 show also that the maximum is reached slighly above the percolation threshold $`p_c`$. Then, the number of doubly connected parts $`t`$ is rapidly declining till it asymptotically approaches $`t=1`$ for $`p1`$. We argue that this behaviour is caused by one great doubly connected component which occurs first for some $`p_{}^>p_c`$ and then gradually incorporates all other doubly connected parts. In principle, these results do also apply to two-mismatch graphs (see Fig. 11).
Comparing this with the above described scenario of cluster structures, we can thus state that there is a range of values of $`p_{}^>p_c`$ where (iii) applies, i.e. clusters are made of one great doubly connected component (including several small ones) and a set of other loosley bound vertices. Accordingly, within this range of $`p`$ we define two subsets of great clusters $`C`$, namely the great doubly connected component $`B=C^{(2)}`$ and the complementary set $`P=CB`$.
As already discussed in Sec. I, the central part of idiotypic networks (corresponding to great clusters $`C`$) should contain strongly and weakly connected distinguished subsets. From the previous analysis it becomes clear, that there is a choice of the only parameter $`p`$, where bit-string models can exactly reproduce such a situation.
## VII conclusions
We investigated statistical properties of a bit-string based model for idiotypic networks and compared typical architectures of the thus defined network with axioms and hypotheses for idiotypic networks from theoretical biology. Before introducing randomness we undertook an analysis of the underlying base graphs to show a major difference to standard percolation problems which bases on the way how the size of the system is increased.
In the following the expressed antibody repertoire was identified with graphs, created by randomly occupying a matching rule defined base structure (base graphs). Concepts of percolation theory have been applied in order to determine the percolation threshold.
The immune system is a very large, yet finite system. Consequently, finite size corrections have to be taken in account. Series expansion techniques allowed the calculation of two critical exponents (for finite systems and in the limit $`d\mathrm{}`$) that characterize the scaling behaviour. A regime of values for the parameter $`p>p_c`$ has been found where random graphs consist of many small clusters and one great connected component. For choices $`p_{}^>p_c`$ our model reproduces the peripheral/central part conception for idiotypic networks.
The translation of the notation of the completeness of idiotypic networks into the language of graph theory allowed the determination of an upper threshold $`p_{\text{dense}}`$ for the occurance of the many small clusters–one great cluster situation. Furthermore, relying on general results of the theory of random graphs (for so-called sequences of configuration spaces) we calculated the – in this case coinciding – thresholds for the density ($`p_{\text{dense}}`$) and the connectivity ($`p_{\text{conn}}`$) properties.
Furthermore, we developed techniques to obtain additional information about the structure of great clusters. Analysing random graphs for $`p_{}^>p_c`$ great clusters can be decomposed into two subsets of vertices with different binding properties, namely groups of doubly connected vertices (backbone) and loosely linked vertices (peripheral part of the great cluster).
Thus, bit-string models are suited to describe a hierarchy of connectivity levels, that really existing idiotypic networks are also expected to exhibit. Our results support to some extent idealized architectures used in a mean field type model to describe the dynamics of the central part of the immune system . Contrary to other model approaches for topologies of idiotypic networks our simple few parametric bit-string model produces a non-trivial seemingly realistic network topology, without assuming a priori distinguished vertex groups.
## A Renormalization group approach
In this appendix we want to present a method to approximate the connectivity threshold for one-mismatch-graphs. The extension of ideas of the renormalization group theory could also be applicable to more complicated cases which do not allow an exact treatment.
Our approach bases on the idea of the renormalization of small cells which employs the self similarity of the system on different length scales at the percolation threshold. A condition to apply this procedure is that an appropriate grouping of sites on the original lattice leads to a renormalized lattice with the same symmetry properties as the original one. Here, treating not a real space lattice but functional networks, we adapt the idea of the renormalization group theory in the following way. We find a transformation $``$ which, by grouping of vertices, leads from the one-mismatch base graph $`G_{d+2}^1`$ in dimension $`d+2`$ to a graph $`(G_{d+2}^1)`$ which is equivalent to the base graph $`G_d^1`$ in dimension $`d`$, $`(G_{d+2}^1)G_d^1`$, thus allowing a systematic reduction of the degrees of freedom.
If a – for finite systems appropriately defined – threshold property, namely $`p_{\text{conn}}(d)`$, converges to a certain value for $`d\mathrm{}`$, differences between $`p_\text{d}`$ and $`p_{\text{d}2}`$ must be small and will disappear in the limit of infinite systems. Consequently, if a grouping of vertices to super vertices on a graph $`G`$ yields a renormalized graph $`(G)`$ of the same type, percolation thresholds on both graphs will be the same. Thus, we replace the term symmetry (as it is applied to lattices) by equivalency of graphs.
We encode an arbitrary vertex of the base graph $`G_{d+2}^1`$ by $`(A,b_0,b_1)`$ where $`A`$ is a bit chain of length $`d`$ and $`b_0`$, $`b_1`$ are single bits. Every vertex on $`G_{d+2}^1`$ belongs to a unique 4-loop $`\{(A,b_0,b_1),(\overline{A},\overline{b_0},b_1),(A,\overline{b_0},\overline{b_1}),(\overline{A},b_0,\overline{b_1})\}`$ connected by one-mismatch links. The renormalization $``$ replaces this 4-loop by the ‘super’vertex $`(A)`$ on $`(G_{d+2}^1)`$ (of course the choice of $`(\overline{A})`$ leads to identical results). If two vertices $`A`$ and $`B`$ are connected by an inversion (one-mismatch) link on $`G_d^1`$ the vertices of the corresponding 4-loops on $`G_{d+2}^1`$ are connected by 4 inversion (one-mismatch) links, too, cf. Fig. 12.
We thus define a renormalized graph $`(G_{d+2}^1)`$ composed of the above explained vertices and edges. It follows directly from our construction that $`(G_{d+2}^1)`$ is equivalent to $`G_d^1`$.
To describe random graphs we apply the following majority rule: a super vertex on $`G_d^1`$ is considered as occupied if at least two connected vertices of the corresponding 4-loop on $`G_{d+2}^1`$ are occupied, i.e. this 4-loop is said to percolate. On the basis of this rule we obtain
$`p^{}=4p^2(1p)^2+4p^3(1p)+p^4\text{ ,}`$ (A1)
where $`p^{}`$ denotes the (renormalized) probability that vertices of $`(G)`$ are occupied if vertices on $`G`$ were diced with probability $`p`$.
Calculating fixed points of eq. (A1) we obtain $`p^{}=0`$, $`(3_{}^+\sqrt{5})/2`$, and $`1`$. Thus, as the only unstable solution in $`[0,1]`$ we have $`p^{}=(3\sqrt{5})/2`$ as a first approximation for the connectivity threshold.
At this place one could pose the question whether only 4-loops are suited as a renormalization cell. Indeed, it is also possible to summarize successive renormalization steps to a single great one, ordinary hypercubes of even dimension $`k<d`$ form suitable cells as well. On the other hand, condensing odd dimensional hypercubes to super vertices does not lead to renormalized graphs $`(G)`$ that are equivalent to other base graphs of the sequence $`\{G_d^1\}`$. Altogether, this seems to be an effect of differences between even and odd-dimensional base graphs (see Sec. II).
Of course, renormalizing small cells entails some approximation. Grouping together vertices and applying a majority rule to occupy the renormalized vertices, situations can arise where clusters on the original lattice are cut or new clusters are formed .
There are different approaches to improve the involved approximations. One possibility suggested in which leads to exact results in the limit of very large cells is to summarize some elementary cells to one large cell of size $`z`$. This large cell will then be occupied if all elementary cells are occupied and connected, i.e. are said to percolate. Since ’renormalization faults’ are essentially due to cell surface effects improvements produced by the above method emanate from the declining surface-volume ratio of large cells. For one-mismatch graphs it holds, however, that the surface size of a cell $`s(z)`$ depends logarithmically on $`z`$, viz. $`s(z)d+1\mathrm{log}_2z`$ leading to only slight improvements with increasing $`z`$.
Using larger elementary cells consisting of two coupled 4-loops we obtain
$`(p^{})^2=p^8+8p^7q+24p^6q^2+32p^5q^3+12p^4q^4\text{ ,}`$ (A2)
using four coupled 4-loops yields
$`(p^{})^4`$ $`=p^{16}+16p^{15}q+112p^{14}q^2+448p^{13}q^3+1120p^{12}q^4+`$ (A3)
$`+1792p^{11}q^5+1776p^{10}q^6+1008p^9q^7+180p^8q^8`$ (A4)
where $`q=1p`$. As fixed points $`p^{}=p`$ of (A2) and (A4) we determined numerically $`p0.41`$ and $`p0.39`$, respectively. The difference of both values to $`p_{conn}=0.5`$ may be a consequence of the above mentioned slow convergence.
More interestingly, our renormalization procedure of replacing hypercubes of dimension $`k`$ by super vertices is without modifications applicable to ordinary hypercubes as well. This gives an additional argument that the connectivity thresholds of both sequences of graphs are equal which can also be verified by evaluating eq. (35).
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# 1 xpnwm
There has been considerable progress in the study of the spherically symmetric Einstein-Yang-Mills (EYM) and Einstein-Yang-Mills-Higgs (EYMH) equations, stimulated by the discovery of globally regular solutions of the EYM eqs. , for a review see e.g. . Up to now most of the attention has been focused on asymptotically flat, smooth particle-like (self-gravitating sphalerons and monopoles) and black hole solutions.
In this paper we present a new ‘cosmological’ type of globally regular solutions of the EYMH eqs. describing static, spherically symmetric spatially compact space-times. Analogous solutions (‘static universes’) have already been found in EYM theory in the presence of a cosmological constant $`\mathrm{\Lambda }`$, (EYMC). More precisely, numerical evidence has been presented in Ref. indicating the existence of a discrete family of solutions indexed by the number $`m`$ of nodes of the YM field, for a special set of values of the cosmological constant, $`\{\mathrm{\Lambda }(m),m=1,\mathrm{}\}`$. The $`m=1`$ solution is particularly simple as its energy density is constant and it has a simple analytic form . It appears that the presence of the cosmological constant is essential for the very existence of such solutions. Therefore at first sight it might appear surprising that such ‘static universe’ solutions also exist in an EYMH theory without a cosmological constant. In fact this should not come totally unexpected as the self-interaction potential of the scalar field, $`V(\mathrm{\Phi })`$, can generate the necessary energy density to support a spatially compact space-time. Intuitively $`V(\mathrm{\Phi })`$ can act as a ‘dynamical cosmological constant’. This viewpoint has been particularly stressed by Linde and Vilenkin in their work on ‘topological inflation’.
The EYMH theory has three different mass scales, the Planck mass $`M_{\mathrm{Pl}}=1/\sqrt{G}`$ and the masses $`M_\mathrm{W}`$ and $`M_\mathrm{H}`$ of the YM resp. Higgs field, giving rise to two dimensionless ratios $`\alpha =M_\mathrm{W}/M_{\mathrm{Pl}}`$ and $`\beta =M_\mathrm{H}/M_\mathrm{W}`$. It turns out to be convenient to introduce an ‘effective cosmological constant’ $`\mathrm{\Lambda }=\alpha ^4\beta ^2/4`$. Our ($`m=1,n=0,1,2,\mathrm{}`$) EYMH solutions bifurcate with the EYMC solution of for $`\alpha =\sqrt{3},\mathrm{}`$ at $`\mathrm{\Lambda }=3/4`$ with a vanishing Higgs field. Varying $`\alpha `$ leads to 1-parameter families ($`n=0,1,\mathrm{}`$) with $`\mathrm{\Lambda }`$ determined as a (possibly multi-valued) function of $`\alpha `$. Similarly for some discrete values $`\alpha _{m,n}`$ our solutions bifurcate with the higher nodes cosmological EYMC solutions of Ref. , when the Higgs field tends to zero.
Following the notation of , we write the spherically symmetric line element as:
$$ds^2=e^{2\nu (R)}dt^2e^{2\lambda (R)}dR^2r^2(R)d\mathrm{\Omega }^2.$$
(1)
The ‘minimal’ spherically symmetric Ansatz for the YM field is
$$W_\mu ^aT_adx^\mu =W(R)(T_1d\theta +T_2\mathrm{sin}\theta d\phi )+T_3\mathrm{cos}\theta d\phi ,$$
(2)
where $`T_a`$ denote the generators of $`SU(2)`$ and for the Higgs field
$$\mathrm{\Phi }^a=H(R)n^a,$$
(3)
$`n^a`$ denoting the unit vector in the radial direction. The reduced EYMH action can be expressed as
$$S=dRe^{(\nu +\lambda )}[\frac{1}{2}(1+e^{2\lambda }((r^{})^2+\nu ^{}(r^2)^{})e^{2\lambda }r^2V_1V_2],$$
(4)
with
$$V_1=\frac{(W^{})^2}{r^2}+\frac{1}{2}(H^{})^2,$$
(5)
and
$$V_2=\frac{(1W^2)^2}{2r^2}+\frac{\beta ^2r^2}{8}(H^2\alpha ^2)^2+W^2H^2.$$
(6)
Varying the reduced action (4) one obtains the EYMH equations:
$`1e^{2\lambda }\left(r^2+\nu ^{}(r^2)^{}\right)+2e^{2\lambda }r^2V_12V_2`$ $`=`$ $`0,`$ (7)
$`1+e^{2\lambda }r^22e^\lambda \left(e^\lambda rr^{}\right)^{}2e^{2\lambda }r^2V_12V_2`$ $`=`$ $`0,`$ (8)
$`e^\lambda \left(r^{}e^\lambda \right)^{}+e^{\nu \lambda }\left(e^{\nu \lambda }r\nu ^{}\right)^{}+e^{2\lambda }{\displaystyle \frac{(r^2V_1)}{r}}+r{\displaystyle \frac{V_2}{r}}`$ $`=`$ $`0,`$ (9)
$`\left(e^{\nu \lambda }W^{}\right)^{}e^{\nu +\lambda }W\left({\displaystyle \frac{W^21}{r^2}}+H^2\right)`$ $`=`$ $`0,`$ (10)
$`\left(r^2e^{\nu \lambda }H^{}\right)^{}e^{\nu +\lambda }H\left(2W^2+{\displaystyle \frac{\beta ^2r^2}{2}}(H^2\alpha ^2)\right)`$ $`=`$ $`0.`$ (11)
Introducing the combinations
$$Ne^\lambda r^{},\kappa e^\lambda (r^{}+r\nu ^{}),$$
(13)
the field eqs. (7a-e) take the form:
$`2\kappa N`$ $`=`$ $`1+N^2+2\dot{W}^2+r^2\dot{H}^22V_2,`$ (14)
$`\dot{r}`$ $`=`$ $`N,`$ (15)
$`\dot{N}`$ $`=`$ $`(\kappa N){\displaystyle \frac{N}{r}}2{\displaystyle \frac{\dot{W}^2}{r}}r\dot{H}^2,`$ (16)
$`\dot{\kappa }`$ $`=`$ $`[1\kappa ^2+2{\displaystyle \frac{\dot{W}^2}{r^2}}{\displaystyle \frac{\beta ^2r^2}{2}}(H^2\alpha ^2)^22H^2W^2]/r,`$ (17)
$`\ddot{W}`$ $`=`$ $`W\left({\displaystyle \frac{(W^21)}{r^2}}+H^2\right)(\kappa N){\displaystyle \frac{\dot{W}}{r}},`$ (18)
$`\ddot{H}`$ $`=`$ $`2H\left({\displaystyle \frac{W^2}{r^2}}+{\displaystyle \frac{\beta ^2}{4}}(H^2\alpha ^2)\right)(\kappa +N){\displaystyle \frac{\dot{H}}{r}}.`$ (19)
where $`\dot{f}:=df/d\sigma =e^\lambda f^{}`$. Note that in Eqs. (14) the remaining gauge freedom (i.e. diffeomorphisms in the variable $`R`$) is implicit in the choice of the independent variable $`\sigma `$. A convenient choice (also for the numerical integration) is $`e^\lambda =1`$.
Next we recall that solutions with a regular origin at $`r=0`$ have the expansion
$`W(r)`$ $`=`$ $`1br^2+O(r^4),`$ (21)
$`H(r)`$ $`=`$ $`ar+O(r^3),`$ (22)
$`N(r)`$ $`=`$ $`1\left(2b^2+{\displaystyle \frac{a^2}{2}}+{\displaystyle \frac{\alpha ^4\beta ^2}{24}}\right)r^2+O(r^4),`$ (23)
$`\kappa (r)`$ $`=`$ $`1+\left(2b^2{\displaystyle \frac{a^2}{2}}{\displaystyle \frac{\alpha ^4\beta ^2}{8}}\right)r^2+O(r^4),`$ (24)
where $`a`$, $`b`$ are free parameters. Solutions satisfying the conditions (21) which also stay globally regular contain (among others) asymptotically flat self-gravitating monopoles. Integrating the field equations (14) one finds that for the generic solution, however, $`N(\sigma )`$ becomes zero at some finite $`\sigma =\sigma _0`$ with finite values of all the other dependent variables. As it is immediately seen from Eq. (14b) $`N(\sigma _0)=0`$ implies stationarity of $`r`$ at $`r(\sigma _0)=r_0`$. Then in general $`r(\sigma )`$ decreases on some interval for $`\sigma >\sigma _0`$. We refer to such a point as an ‘equator’. In the case of the EYM eqs. it has been proven in Ref. that for all solutions with an equator $`r(\sigma )`$ decreases from $`r_0`$ all the way to $`r=0`$ where a curvature singularity develops (‘bag of gold’). Adding a cosmological constant radically changes this conclusion and there exists an infinite family of regular solutions with 3-sphere topology .
Since for $`H(\rho )0`$ the EYMH system reduces to an EYMC theory where the cosmological constant is given as $`\mathrm{\Lambda }=\alpha ^4\beta ^2/4`$, it is natural to search for solutions of Eqs. (7) with a nontrivial Higgs field bifurcating with the EYMC ones. To see if such a bifurcating class really exists one should first establish the existence of regular solutions of the linearized Higgs-field equation (7e) in the background of a 3-sphere type solution of the EYMC system since the bifurcation occurs for $`H0`$. To discuss the linearization of the Higgs-field equation around solutions of the EYMC equations, it is more convenient to choose a different gauge from $`e^\lambda =1`$, namely $`\nu =\lambda `$ and use the variable $`\rho `$ defined by $`dr/d\rho =r^{}=e^\lambda N`$. Then the linearization of the Higgs-field equation (7e), around an EYMC background solution ($`\mathrm{\Lambda }=\mathrm{\Lambda }_b`$, $`r=r_b(\rho )`$, $`\nu =\nu _b(\rho )`$, $`W=W_b(\rho )`$, $`H=H_b0`$) with $`\mathrm{\Lambda }_b=\alpha _b^4\beta _b^2/4`$) and $`H=h(\rho )`$, reads as:
$$\left(r_b^2h^{}\right)^{}=2e^{2\nu _b}\left(W_b^2\frac{\mathrm{\Lambda }_b}{\alpha ^2}r_b^2\right)h.$$
(26)
Let us consider first the simplest, $`m=1`$ EYMC solution :
$`W_b`$ $`=`$ $`\mathrm{cos}x,H_b=0,r_b=\sqrt{2}\mathrm{sin}x,`$ (27)
$`N_b=\kappa _b`$ $`=`$ $`\mathrm{cos}x,\nu _b=0,\mathrm{\Lambda }_b=3/4,`$ (28)
where $`x=\rho /\sqrt{2}`$. Then Eq. (26) in the background (27), reduces to the following simple equation:
$$(\mathrm{sin}^2(x)h^{}(x))^{}=\left(2\mathrm{cos}^2(x)\frac{3}{\alpha ^2}\mathrm{sin}^2(x)\right)h(x),$$
(30)
where now stands for the derivative with respect to $`x`$. We note that Eq. (30) can be transformed to a hypergeometric equation, however, the solutions of interest, i.e. regular on the background geometry and satisfying the condition (21b), can be directly found by the following trigonometric polynomial Ansatz for $`h(x)`$
$$h_n(x)=\underset{k=0}{\overset{n}{}}c_k\mathrm{sin}^{2k+1}x.$$
(31)
One then easily obtains the recursion relation for the coefficients $`c_k`$:
$$c_k=\frac{(2k1)(2k+1)23/\alpha ^2}{(2k+2)(2k+1)2}c_{k1},k=1\mathrm{}n,$$
(32)
together with the termination condition ($`c_{n+1}=0`$)
$$\alpha _{1,n}^2=\frac{3}{(2n+1)(2n+3)2},n=0,1,\mathrm{}$$
(33)
yielding $`\alpha _{1,0}^2=3`$, $`\alpha _{1,1}^2=3/13`$, $`\alpha _{1,2}^2=1/11`$, …
That is we have found an infinite family of regular solutions of the linearized Higgs field equation (bounded ‘zero modes’) in the $`m=1`$ EYMC background. These ‘zero modes’, indexed by the number of zeros of $`h(x)`$ ($`n=0,1,\mathrm{}`$), indicate the existence of a family of globally regular solutions of the EYMH equations (with $`H(x)`$ having the same number of zeros) bifurcating with the $`m=1`$ EYMC solution for $`a=0`$ and $`\alpha =\alpha _{1,n}`$.
From the $`\alpha `$-dependence of the r.h.s. of Eq. (30) it follows that the solutions $`h(x)`$ vanishing at $`x=0`$ have $`n+1`$ zeros for all $`\alpha [\alpha _{1,n+1},\alpha _{1,n})`$. According to standard theorems on Sturm-Liouville operators this number gives also the number of bound states of the differential operator in Eq. (30). Thus whenever $`\alpha `$ crosses $`\alpha _{1,n}`$ from above a new bound state appears, indicated by the existence of a zero mode for $`\alpha =\alpha _{1,n}`$.
In contrast to the case $`m=1`$, the higher node EYMC solutions ($`m2`$) are not known analytically. Nevertheless one can still conclude that for any of these EYMC backgrounds an infinite family of regular zero modes exists. To see this, it is sufficient to note that for $`\alpha `$ sufficiently small the r.h.s. of Eq. (26) becomes arbitrarily negative over a finite interval of $`\rho `$. This implies that the number of bound states tends to infinity for $`\alpha 0`$. According to the previous argument the same holds for the number of zero modes accumulating at $`\alpha =0`$. Numerical values for the first five bifurcation points with $`m=1,\mathrm{},5,10,20`$ are given in Tab. 1.
In order to verify the correctness of the above scenario we have integrated the field eqs. (14) numerically. To simplify the singular boundary value problem, we have looked only for solutions, which are (anti)symmetric about the equator. This means we impose boundary conditions at $`r=0`$ and at the equator (a regular point). In addition to the vanishing of $`\kappa `$ at the equator, we require the vanishing of the functions $`W`$ resp. $`\dot{W}`$ and $`H`$ resp. $`\dot{H}`$ depending if $`m`$ and $`n`$ are even or odd. Thus 3 functions must have a common zero with $`N`$ forcing us to tune 3 of the 4 available parameters $`\alpha `$, $`\beta `$, $`a`$, and $`b`$. Fig. 1 shows some solutions with $`m=1,2`$ zeros of $`W`$ and $`n=0,1`$ zeros of $`H`$.
In Fig. 2 we have plotted the values of the parameters $`\alpha `$ and $`\mathrm{\Lambda }`$ for the solutions with $`m,n2`$. While $`\alpha `$ runs to large values with decreasing $`\mathrm{\Lambda }`$ for the $`n=0`$ solutions, i.e. with a nodeless Higgs field, the parameters for the solutions with $`n1`$ show the opposite behaviour. This structure remains true for higher values of $`m`$, although the change occurs in general for some $`n_0>1`$ increasing with $`m`$ (e.g. $`n_0=6`$ for $`m=10`$). Similar to the case $`(2,0)`$ some of the graphs in the $`\alpha `$-$`\mathrm{\Lambda }`$ plane for $`m>2`$ show a minimum of $`\alpha `$, before they tend to large values of $`\alpha `$ (e.g. $`n=1,2,4,5`$ for $`m=10`$).
The families of numerical solutions for which $`\alpha `$ becomes large seem to have a smooth limit for $`\alpha \mathrm{}`$ with $`\mathrm{\Lambda }1/4`$. In this limit, assuming that $`H`$ stays finite, Eqs. (7) reduce to an EYMH system with a cosmological constant $`\mathrm{\Lambda }=1/4`$ instead of a Higgs potential ($`\beta =0`$). The spatial sections of the corresponding space-times are no longer compact ($`\sigma \mathrm{}`$ and $`r_0\sqrt{2}`$). We have also numerically integrated the limiting field eqs. and our results fully confirm the existence of the limit $`\alpha \mathrm{}`$.
It was argued in that the solutions of the EYMC system of are unstable. As usual the criterion for instability is the existence of imaginary modes of the linearized time-dependent field equations. Although we do not quite approve of the methods of , we nevertheless believe that their result is correct. For the case $`m=1`$ the instability has already been shown in . We expect the instability of the YM system to persist in our case with the Higgs field with the same number of unstable modes at least for small values of $`a`$. In view of the discussion of the solutions of Eq. (26) given above there are additional instabilities of the EYMC solutions viewed as solutions of the EYMH theory (with $`a=0`$) in the Higgs sector. More precisely, there are $`n+1`$ unstable modes for $`\alpha [\alpha _{m,n+1},\alpha _{m,n})`$. Turning to solutions with $`a0`$ but small the only questions is, if the number of the unstable Higgs modes of the $`(m,n)`$ solutions is equal to $`n`$ or $`n+1`$. The answer depends on the behaviour of the zero mode at the bifurcation point $`\alpha _{m,n}`$ when $`a`$ deviates from zero. If it turns into a bound state one gets $`n+1`$, if it moves into the continuum $`n`$ unstable modes.
Similar to the globally regular solutions discussed in this paper, there are families of solutions with a horizon branching off from the corresponding solutions of without a Higgs field. In fact there are even more general solutions having two horizons, one of them replacing the regular origin of the forementioned class. We plan to give a detailed account of these solutions in a forthcoming publication.
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# On the Role of Energy Conservation in High-Energy Nuclear Scattering
## 1 Introduction
The purpose of this paper is to provide the theoretical framework to treat hadron-hadron scattering and the initial stage of nucleus-nucleus collisions at ultra-relativistic energies, in particular with view to RHIC and LHC. The knowledge of these initial collisions is crucial for any theoretical treatment of parton thermalization and a possible parton-hadron phase transition, the detection of which being the ultimate aim of all the efforts of colliding heavy ions at very high energies.
So what are the currently used models? Quite popular are semi-classical treatments of either partons or hadrons, in any case completely ignoring quantum mechanical interference. This is certainly quite unrealistic, and we do not want to discuss such options any further.
There are as well considerable efforts to describe nuclear collisions via solving classical Yang-Mills equations, which allows to calculate inclusive parton distributions . This approach is to some extent orthogonal to ours being based on the assumption of the perturbative nature of the triple Pomeron coupling . The physical picture which stays behind the construction of was outlined in and corresponds to perturbative interactions between individual parton cascades in hadrons, i.e. to fusion of partons of virtualities $`Q^2>Q_0^2`$, with $`Q_0^2`$ being a reasonable cutoff for QCD being applicable. Contrary to that, we believe that parton cascades interact with each other in the non-perturbative region of parton virtualities $`Q^2<Q_0^2`$ and consider it as the interaction between soft Pomerons. In our scheme a relatively big value of the soft triple-Pomeron coupling should provide the necessary screening corrections which finally prevent the large increase of parton densities in the small $`x`$ limit and restore the unitarity, thus leaving little room for higher twist effects in the perturbative part of the interaction.
Another approach is the so-called Gribov-Regge theory (GRT) . This is an effective field theory, which allows multiple interactions to happen “in parallel”, with the phenomenological object called “Pomeron” representing an elementary interaction. Using the general rules of field theory, one may express cross sections in terms of a couple of parameters characterizing the Pomeron. Interference terms are crucial, they assure the unitarity of the theory. Here one observes an inconsistency: the fact that energy needs to be shared between many Pomerons in case of multiple scattering is well taken into account when calculating particle production (in particular in Monte Carlo applications), but energy conservation is not taken care of in cross section calculations. This is a serious problem and makes the whole approach inconsistent. Related to the above problem is the fact that different elementary interactions in case of multiple scattering are usually not treated equally, so the first interaction is usually considered to be quite different compared to the subsequent ones.
Provided factorization works for nuclear collisions, one may employ the parton model, which allows to calculate inclusive cross sections as a convolution of an elementary cross section with parton distribution functions, with these distribution functions taken from deep inelastic scattering. In order to get exclusive parton level cross sections, some additional assumptions are needed, which follow quite closely the Gribov-Regge approach, encountering the same difficulties.
As a solution of the above-mentioned problems, we present a new approach which we call “Parton-based Gribov-Regge Theory”: we have a consistent treatment for calculating cross sections and particle production considering energy conservation in both cases; in addition, we introduce hard processes in a natural way, and, compared to the parton model, we can deal with exclusive cross sections without arbitrary assumptions. A single set of parameters is sufficient to fit many basic spectra in proton-proton and lepton-nucleon scattering, as well as in electron-positron annihilation (with the exception of one parameter which needs to be changed in order to optimize electron-positron transverse momentum spectra).
The basic guideline of our approach is theoretical consistency. We cannot derive everything from first principles, but we use rigorously the language of field theory to make sure not to violate basic laws of physics, which is easily done in more phenomenological treatments (see discussion above).
There are still problems and open questions: there is clearly a problem with unitarity at very high energies, which should be cured by considering screening corrections due to so-called triple-Pomeron interactions, which we do not treat rigorously at present but which is our next project.
## 2 Problems
Before presenting new theoretical ideas, we want to discuss the open problems in the parton model approach and in Gribov-Regge theory.
### Gribov-Regge Theory
Gribov-Regge theory is by construction a multiple scattering theory. The elementary interactions are realized by complex objects called “Pomerons”, who’s precise nature is not known, and which are therefore simply parameterized: the elastic amplitude $`T`$ corresponding to a single Pomeron exchange is given as
$$T(s,t)is^{\alpha _0+\alpha ^{}t}$$
with a couple of parameters to be determined by experiment. Even in hadron-hadron scattering, several of these Pomerons are exchanged in parallel, see fig. 1.
Using general rules of field theory (cutting rules), one obtains an expression for the inelastic cross section,
$$\sigma _{\mathrm{inel}}^{h_1h_2}=d^2b\left\{1\mathrm{exp}\left(G(s,b)\right)\right\},$$
(1)
where the so-called eikonal $`G(s,b)`$ (proportional to the Fourier transform of $`T(s,t)`$) represents one elementary interaction (a thick line in fig. 1). One can generalize to nucleus-nucleus collisions, where corresponding formulas for cross sections may be derived.
In order to calculate exclusive particle production, one needs to know how to share the energy between the individual elementary interactions in case of multiple scattering. We do not want to discuss the different recipes used to do the energy sharing (in particular in Monte Carlo applications). The point is, whatever procedure is used, this is not taken into account in the calculation of cross sections discussed above. So, actually, one is using two different models for cross section calculations and for treating particle production. Taking energy conservation into account in exactly the same way will modify the cross section results considerably.
This problem has first been discussed in ,. The authors claim that following from the non-planar structure of the corresponding diagrams, conserving energy and momentum in a consistent way is crucial, and therefore the incident energy has to be shared between the different elementary interactions, both real and virtual ones.
Another very unpleasant and unsatisfactory feature of most “recipes” for particle production is the fact, that the second Pomeron and the subsequent ones are treated differently than the first one, although in the above-mentioned formula for the cross section all Pomerons are considered to be identical.
### The Parton Model
The standard parton model approach to hadron-hadron or also nucleus-nucleus scattering amounts to presenting the partons of projectile and target by momentum distribution functions, $`f_{h_1}`$ and $`f_{h_2}`$, and calculating inclusive cross sections for the production of parton jets with the squared transverse momentum $`p_{}^2`$ larger than some cutoff $`Q_0^2`$ as
$$\sigma _{\mathrm{incl}}^{h_1h_2}=\underset{ij}{}𝑑p_{}^2𝑑x^+𝑑x^{}f_{h_1}^i(x^+,p_{}^2)f_{h_2}^j(x^{},p_{}^2)\frac{d\widehat{\sigma }_{ij}}{dp_{}^2}(x^+x^{}s)\theta \left(p_{}^2Q_0^2\right),$$
where $`d\widehat{\sigma }_{ij}/dp_{}^2`$ is the elementary parton-parton cross section and $`i,j`$ represent parton flavors.
This simple factorization formula is the result of cancelations of complicated diagrams (AGK cancelations) and hides therefore the complicated multiple scattering structure of the reaction. The most obvious manifestation of such a structure is the fact that at high energies ($`\sqrt{s}10`$ GeV) the inclusive cross section in proton-(anti-)proton scattering exceeds the total one, so the average number $`\overline{N}_{\mathrm{int}}^{pp}`$ of elementary interactions must be greater than one:
$$\overline{N}_{\mathrm{int}}^{h_1h_2}=\sigma _{\mathrm{incl}}^{h_1h_2}/\sigma _{\mathrm{tot}}^{h_1h_2}>1.$$
The usual solution is the so-called eikonalization, which amounts to re-introducing multiple scattering, based on the above formula for the inclusive cross section:
$$\sigma _{\mathrm{inel}}^{h_1h_2}(s)=d^2b\left\{1\mathrm{exp}\left(A(b)\sigma _{\mathrm{incl}}^{h_1h_2}(s)\right)\right\}=\sigma _m^{h_1h_2}(s),$$
(2)
with
$$\sigma _m^{h_1h_2}(s)=d^2b\frac{\left(A(b)\sigma _{\mathrm{incl}}^{h_1h_2}(s)\right)^m}{m!}\mathrm{exp}\left(A(b)\sigma _{\mathrm{incl}}^{h_1h_2}(s)\right)$$
(3)
representing the cross section for $`n`$ scatterings; $`A(b)`$ being the proton-proton overlap function (the convolution of two proton profiles). In this way the multiple scattering is “recovered”. The disadvantage is that this method does not provide any clue how to proceed for nucleus-nucleus ($`AB`$) collisions. One usually assumes the proton-proton cross section for each individual nucleon-nucleon pair of a $`AB`$ system. We can demonstrate that this assumption is incorrect (see ).
Another problem, in fact the same one as discussed earlier for the GRT, arises in the case of exclusive calculations (event generation), since the above formulas do not provide any information on how to share the energy between many elementary interactions. The Pythia-method amounts to generating the first elementary interaction according to the inclusive differential cross section, then taking the remaining energy for the second one and so on. In this way, the event generation will reproduce the theoretical inclusive spectrum for hadron-hadron interaction (by construction), as it should be. The method is, however, very arbitrary, and - even more serious - we observe the same inconsistency as in the Gribov-Regge approach: energy conservation is not at all taken care of in the above formulas for cross section calculations.
## 3 A Solution: Parton-based Gribov-Regge Theory
In this paper, we present a new approach for hadronic interactions and for the initial stage of nuclear collisions, which is able to solve several of the above-mentioned problems. We provide a rigorous treatment of the multiple scattering aspect, such that questions as energy conservation are clearly determined by the rules of field theory, both for cross section and particle production calculations. In both (!) cases, energy is properly shared between the different interactions happening in parallel, see fig. 2
for proton-proton and fig. 3 for proton-nucleus collisions (generalization to nucleus-nucleus is obvious).
This is the most important and new aspect of our approach, which we consider to be a first necessary step to construct a consistent model for high energy nuclear scattering.
The elementary interactions, shown as the thick lines in the above figures, are in fact a sum of a soft, a hard, and a semi-hard contribution, providing a consistent treatment of soft and hard scattering. To some extend, our approach provides a link between the Gribov-Regge approach and the parton model, we call it “Parton-based Gribov-Regge Theory”.
## 4 Parton-Parton Scattering
Let us first investigate parton-parton scattering, before constructing a multiple scattering theory for hadronic and nuclear scattering.
We distinguish three types of elementary parton-parton scatterings, referred to as “soft”, “hard” and “semi-hard”, which we are going to discuss briefly in the following. The detailed derivations can be found in .
### The Soft Contribution
Let us first consider the pure non-perturbative contribution, where all virtual partons appearing in the internal structure of the diagram have restricted virtualities $`Q^2<Q_0^2`$, where $`Q_0^21`$ GeV<sup>2</sup> is a reasonable cutoff for perturbative QCD being applicable. Such soft non-perturbative dynamics is known to dominate hadron-hadron interactions at not too high energies. Lacking methods to calculate this contribution from first principles, it is simply parameterized and graphically represented as a ‘blob’, see fig. 4. It is traditionally assumed to correspond to multi-peripheral production of partons (and final hadrons) and is described by the phenomenological soft Pomeron exchange amplitude $`T_{\mathrm{soft}}(\widehat{s},t)`$ . The corresponding profile function is expressed via the amplitude $`T_{\mathrm{soft}}`$ as
$`D_{\mathrm{soft}}(\widehat{s},b)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2\widehat{s}}}{\displaystyle d^2q_{}\mathrm{exp}\left(i\stackrel{}{q}_{}\stackrel{}{b}\right)\mathrm{\hspace{0.17em}2}\mathrm{Im}T_{\mathrm{soft}}(\widehat{s},q_{}^2)}`$ (4)
$`=`$ $`{\displaystyle \frac{2\gamma _{\mathrm{part}}^2}{\lambda _{\mathrm{soft}}^{(2)}(\widehat{s}/s_0)}}\left({\displaystyle \frac{\widehat{s}}{s_0}}\right)^{\alpha _{\mathrm{soft}}(0)1}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^{(2)}(\widehat{s}/s_0)}}\right),`$
with
$$\lambda _{\mathrm{soft}}^{(n)}(z)=nR_{\mathrm{part}}^2+\alpha _{\mathrm{soft}}^{}\mathrm{ln}z,$$
where $`\widehat{s}`$ is the usual Mandelstam variable for parton-parton scattering. The parameters $`\alpha _{\mathrm{soft}}(0)`$, $`\alpha _{\mathrm{soft}}^{}`$ are the intercept and the slope of the Pomeron trajectory, $`\gamma _{\mathrm{part}}`$ and $`R_{\mathrm{part}}^2`$ are the vertex value and the slope for the Pomeron-parton coupling, and $`s_01`$ GeV<sup>2</sup> is the characteristic hadronic mass scale. The external legs of the diagram of fig. 4 are “partonic constituents”, which are assumed to be quark-anti-quark pairs.
### The Hard Contribution
Let us now consider the other extreme, when all the processes are perturbative, i.e. all internal intermediate partons are characterized by large virtualities $`Q^2>Q_0^2`$. In that case, the corresponding hard parton-parton scattering amplitude can be calculated using the perturbative QCD techniques , and the intermediate states contributing to the absorptive part of the amplitude can be defined in the parton basis. In the leading logarithmic approximation of QCD, summing up terms where each (small) running QCD coupling constant $`\alpha _s(Q^2)`$ appears together with a large logarithm $`\mathrm{ln}(Q^2/\lambda _{\mathrm{QCD}}^2)`$ (with $`\lambda _{QCD}`$ being the infrared QCD scale), and making use of the factorization hypothesis, one obtains the contribution of the corresponding cut diagram for $`t=q^2=0`$ as the cut parton ladder cross section $`\sigma _{\mathrm{hard}}^{jk}(\widehat{s},Q_0^2)`$ <sup>1</sup><sup>1</sup>1 Strictly speaking, one obtains the ladder representation for the process only using axial gauge. , as shown in fig. 5,
where all horizontal rungs are the final (on-shell) partons and the virtualities of the virtual $`t`$-channel partons increase from the ends of the ladder towards the largest momentum transfer parton-parton process (indicated symbolically by the ‘blob’ in the middle of the ladder):
$`\sigma _{\mathrm{hard}}^{jk}(\widehat{s},Q_0^2)`$ $`=`$ $`{\displaystyle \frac{1}{2\widehat{s}}}2\mathrm{I}\mathrm{m}T_{\mathrm{hard}}^{jk}(\widehat{s},t=0,Q_0^2)`$
$`=`$ $`K{\displaystyle \underset{ml}{}}{\displaystyle 𝑑x_B^+𝑑x_B^{}𝑑p_{}^2\frac{d\sigma _{\mathrm{Born}}^{ml}}{dp_{}^2}(x_B^+x_B^{}\widehat{s},p_{}^2)}`$
$`\times `$ $`E_{\mathrm{QCD}}^{jm}(Q_0^2,M_F^2,x_B^+)E_{\mathrm{QCD}}^{kl}(Q_0^2,M_F^2,x_B^{})\theta \left(M_F^2Q_0^2\right),`$
Here $`d\sigma _{\mathrm{Born}}^{ml}/dp_{}^2`$ is the differential $`22`$ parton scattering cross section, $`p_{}^2`$ is the parton transverse momentum in the hard process, $`m,l`$ and $`x_B^\pm `$ are correspondingly the types and the shares of the light cone momenta of the partons participating in the hard process, and $`M_F^2`$ is the factorization scale for the process (we use $`M_F^2=p_{}^2/4`$). The ‘evolution function’ $`E_{\mathrm{QCD}}^{jm}(Q_0^2,M_F^2,z)`$ represents the evolution of a parton cascade from the scale $`Q_0^2`$ to $`M_F^2`$, i.e. it gives the number density of partons of type $`m`$ with the momentum share $`z`$ at the virtuality scale $`M_F^2`$, resulted from the evolution of the initial parton $`j`$, taken at the virtuality scale $`Q_0^2`$. The evolution function satisfies the usual DGLAP equation with the initial condition $`E_{\mathrm{QCD}}^{jm}(Q_0^2,Q_0^2,z)=\delta _m^j\delta (1z)`$. The factor $`K1.5`$ takes effectively into account higher order QCD corrections.
In the following, we shall need to know the contribution of the uncut parton ladder $`T_{\mathrm{hard}}^{jk}(\widehat{s},t,Q_0^2)`$ with some momentum transfer $`q`$ along the ladder (with $`t=q^2`$). The behavior of the corresponding amplitudes was studied in in the leading logarithmic($`1/x`$ ) approximation of QCD. The precise form of the corresponding amplitude is not important for our application; we just use some of the results of , namely that one can neglect the real part of this amplitude and that it is nearly independent on $`t`$, i.e. that the slope of the hard interaction $`R_{\mathrm{hard}}^2`$ is negligible small, i.e. compared to the soft Pomeron slope one has $`R_{\mathrm{hard}}^20`$. So we parameterize $`T_{\mathrm{hard}}^{jk}(\widehat{s},t,Q_0^2)`$ in the region of small $`t`$ as
$$T_{\mathrm{hard}}^{jk}(\widehat{s},t,Q_0^2)=i\widehat{s}\sigma _{\mathrm{hard}}^{jk}(\widehat{s},Q_0^2)\mathrm{exp}\left(R_{\mathrm{hard}}^2t\right)$$
(5)
The corresponding profile function is obtained by calculating the Fourier transform $`\stackrel{~}{T}_{\mathrm{hard}}`$ of $`T_{\mathrm{hard}}`$ and dividing by the initial parton flux $`2\widehat{s}`$,
$$D_{\mathrm{hard}}^{jk}(\widehat{s},b)=\frac{1}{2\widehat{s}}2\mathrm{I}\mathrm{m}\stackrel{~}{T}_{\mathrm{hard}}^{jk}(\widehat{s},b),$$
which gives
$`D_{\mathrm{hard}}^{jk}(\widehat{s},b)={\displaystyle \frac{1}{8\pi ^2\widehat{s}}}{\displaystyle d^2q_{}\mathrm{exp}\left(i\stackrel{}{q}_{}\stackrel{}{b}\right)\mathrm{\hspace{0.17em}2}\mathrm{Im}T_{\mathrm{hard}}^{jk}(\widehat{s},q_{}^2,Q_0^2)}`$
$`=\sigma _{\mathrm{hard}}^{jk}(\widehat{s},Q_0^2){\displaystyle \frac{1}{4\pi R_{\mathrm{hard}}^2}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4R_{\mathrm{hard}}^2}}\right),`$ (6)
In fact, the above considerations are only correct for valence quarks, as discussed in detail in the next section. Therefore, we also talk about “valence-valence” contribution and we use $`D_{\mathrm{val}\mathrm{val}}`$ instead of $`D_{\mathrm{hard}}`$:
$$D_{\mathrm{val}\mathrm{val}}^{jk}(\widehat{s},b)D_{\mathrm{hard}}^{jk}(\widehat{s},b),$$
so these are two names for one and the same object.
### The Semi-hard Contribution
The discussion of the preceding section is not valid in case of sea quarks and gluons, since here the momentum share $`x_1`$ of the “first” parton is typically very small, leading to an object with a large mass of the order $`Q_0^2/x_1`$ between the parton and the proton . Microscopically, such ’slow’ partons with $`x_11`$ appear as a result of a long non-perturbative parton cascade, where each individual parton branching is characterized by a small momentum transfer squared $`Q^2<Q_0^2`$ . When calculating proton structure functions or high-$`p_t`$ jet production cross sections this non-perturbative contribution is usually included in parameterized initial parton momentum distributions at $`Q^2=Q_0^2`$. However, the description of inelastic hadronic interactions requires to treat it explicitly in order to account for secondary particles produced during such non-perturbative parton pre-evolution, and to describe correctly energy-momentum sharing between multiple elementary scatterings. As the underlying dynamics appears to be identical to the one of soft parton-parton scattering considered above, we treat this soft pre-evolution as the usual soft Pomeron emission, as discussed in detail in .
So for sea quarks and gluons, we consider so-called semi-hard interactions between parton constituents of initial hadrons, represented by a parton ladder with “soft ends”, see fig. 6. As in the case of soft scattering, the external legs are quark-anti-quark pairs, connected to soft Pomerons. The outer partons of the ladder are on both sides sea quarks or gluons (therefore the index “sea-sea”). The central part is exactly the hard scattering considered in the preceding section. As discussed in length in , the mathematical expression for the corresponding amplitude is given as
$`iT_{\mathrm{sea}\mathrm{sea}}(\widehat{s},t)`$ $`=`$ $`{\displaystyle \underset{jk}{}}{\displaystyle _0^1}{\displaystyle \frac{dz^+}{z^+}}{\displaystyle \frac{dz^{}}{z^{}}}\mathrm{Im}T_{\mathrm{soft}}^j({\displaystyle \frac{s_0}{z^+}},t)\mathrm{Im}T_{\mathrm{soft}}^k({\displaystyle \frac{s_0}{z^{}}},t)iT_{\mathrm{hard}}^{jk}(z^+z^{}\widehat{s},t,Q_0^2),`$
with $`z^\pm `$ being the momentum fraction of the external leg-partons of the parton ladder relative to the momenta of the initial (constituent) partons. The indices $`j`$ and $`k`$ refer to the flavor of these external ladder partons. The amplitudes $`T_{\mathrm{soft}}^j`$ are the soft Pomeron amplitudes discussed earlier, but with modified couplings, since the Pomerons are now connected to a parton ladder on one side. The arguments $`s_0/z^\pm `$ are the squared masses of the two soft Pomerons, $`z^+z^{}\widehat{s}`$ is the squared mass of the hard piece.
Performing as usual the Fourier transform to the impact parameter representation and dividing by $`2\widehat{s}`$, we obtain the profile function
$$D_{\mathrm{sea}\mathrm{sea}}(\widehat{s},b)=\frac{1}{2\widehat{s}}\mathrm{\hspace{0.17em}2}\mathrm{Im}\stackrel{~}{T}_{\mathrm{sea}\mathrm{sea}}(\widehat{s},b),$$
which may be written as
$`D_{\mathrm{sea}\mathrm{sea}}(\widehat{s},b)`$ $`=`$ $`{\displaystyle \underset{jk}{}}{\displaystyle _0^1}𝑑z^+𝑑z^{}E_{\mathrm{soft}}^j\left(z^+\right)E_{\mathrm{soft}}^k\left(z^{}\right)\sigma _{\mathrm{hard}}^{jk}(z^+z^{}\widehat{s},Q_0^2)`$ (7)
$`\times {\displaystyle \frac{1}{4\pi \lambda _{\mathrm{soft}}^{(2)}(1/(z^+z^{}))}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^{(2)}\left(1/(z^+z^{})\right)}}\right)`$
with the soft Pomeron slope $`\lambda _{\mathrm{soft}}^{(2)}`$ and the cross section $`\sigma _{\mathrm{hard}}^{jk}`$ being defined earlier. The functions $`E_{\mathrm{soft}}^j\left(z^\pm \right)`$ representing the “soft ends” are defined as
$$\mathrm{E}_{\mathrm{soft}}^j(z^\pm )=\mathrm{Im}T_{\mathrm{soft}}^j\left(\frac{s_0}{z^\pm },t=0\right).$$
We neglected the small hard scattering slope $`R_{\mathrm{hard}}^2`$ compared to the Pomeron slope $`\lambda _{\mathrm{soft}}^{(2)}`$. We call $`E_{\mathrm{soft}}`$ also the “ soft evolution”, to indicate that we consider this as simply a continuation of the QCD evolution, however, in a region where perturbative techniques do not apply any more. As discussed in , $`E_{\mathrm{soft}}^j\left(z\right)`$ has the meaning of the momentum distribution of parton $`j`$ in the soft Pomeron.
Consistency requires to also consider the mixed semi-hard contributions with a valence quark on one side and a non-valence participant (quark-anti-quark pair) on the other one, see fig. 7.
We have
$$iT_{\mathrm{val}\mathrm{sea}}^j(\widehat{s})=_0^1\frac{dz^{}}{z^{}}\underset{k}{}\mathrm{Im}T_{\mathrm{soft}}^k(\frac{s_0}{z^{}},q^2)iT_{\mathrm{hard}}^{jk}(z^{}\widehat{s},q^2,Q_0^2)$$
and
$`D_{\mathrm{val}\mathrm{sea}}^j(\widehat{s},b)`$ $`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle _0^1}𝑑z^{}E_{\mathrm{soft}}^k\left(z^{}\right)\sigma _{\mathrm{hard}}^{jk}(z^{}\widehat{s},Q_0^2)`$
$`\times {\displaystyle \frac{1}{4\pi \lambda _{\mathrm{soft}}^{(1)}(1/z^{})}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^{(1)}\left(1/z^{}\right)}}\right)`$
where $`j`$ is the flavor of the valence quark at the upper end of the ladder and $`k`$ is the type of the parton on the lower ladder end. Again, we neglected the hard scattering slope $`R_{\mathrm{hard}}^2`$ compared to the soft Pomeron slope. A contribution $`D_{\mathrm{sea}\mathrm{val}}^j(\widehat{s},b)`$, corresponding to a valence quark participant from the target hadron, is given by the same expression,
$$D_{\mathrm{sea}\mathrm{val}}^j(\widehat{s},b)=D_{\mathrm{val}\mathrm{sea}}^j(\widehat{s},b),$$
since eq. (4) stays unchanged under replacement $`z^{}z^+`$ and only depends on the total c.m. energy squared $`\widehat{s}`$ for the parton-parton system.
## 5 Hadron-Hadron Scattering
To treat hadron-hadron scattering we use parton momentum Fock state expansion of hadron eigenstates
$$|h=\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}_0^1\underset{l=1}{\overset{k}{}}dx_lf_k^h(x_1,\mathrm{}x_k)\delta \left(1\underset{j=1}{\overset{k}{}}x_j\right)a^+(x_1)\mathrm{}a^+(x_k)|0,$$
where $`f_k(x_1,\mathrm{}x_k)`$ is the probability amplitude for the hadron $`h`$ to consist of $`k`$ constituent partons with the light cone momentum fractions $`x_1,\mathrm{},x_k`$ and $`a^+\left(x\right)`$ is the creation operator for a parton with the fraction $`x`$. A general scattering process is described as a superposition of a number of pair-like scatterings between parton constituents of the projectile and target hadrons. Then hadron-hadron scattering amplitude is obtained as a convolution of individual parton-parton scattering amplitudes considered in the previous section and “inclusive” momentum distributions $`\frac{1}{n!}\stackrel{~}{F}_h^{(n)}(x_1,\mathrm{}x_n)`$ of $`n`$ “participating” parton constituents involved in the scattering process($`n1`$), with
$$\frac{1}{n!}\stackrel{~}{F}_h^{(n)}(x_1,\mathrm{}x_n)=\underset{k=n}{\overset{\mathrm{}}{}}\frac{1}{k!}\frac{k!}{n!(kn)!}_0^1\underset{l=n+1}{\overset{k}{}}dx_l\left|f_k(x_1,\mathrm{}x_k)\right|^2\delta \left(1\underset{j=1}{\overset{k}{}}x_j\right)$$
We assume that $`\stackrel{~}{F}_{h_1(h_2)}^{(n)}(x_1,\mathrm{}x_n)`$ can be represented in a factorized form as a product of the contributions $`F_{\mathrm{part}}^h(x_l)`$, depending on the momentum shares $`x_l`$ of the “participating” or “active” parton constituents, and on the function $`F_{\mathrm{remn}}^h\left(1_{j=1}^nx_j\right)`$,
representing the contribution of all “spectator” partons, sharing the remaining share $`1_jx_j`$ of the initial light cone momentum (see fig. 8):
$$\stackrel{~}{F}_h^{(n)}(x_1,\mathrm{}x_n)=\underset{l=1}{\overset{n}{}}F_{\mathrm{part}}^h(x_l)F_{\mathrm{remn}}^h\left(1\underset{j=1}{\overset{n}{}}x_j\right)$$
(9)
The participating parton constituents are assumed to be quark-anti-quark pairs (not necessarily of identical flavors), such that the baryon numbers of the projectile and of the target are conserved. Then, as discussed in detail in , the hadron-hadron amplitude may be written as
$`iT_{h_1h_2}(s,t)=8\pi ^2s{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle _0^1}{\displaystyle \underset{l=1}{\overset{n}{}}}dx_l^+dx_l^{}{\displaystyle \underset{l=1}{\overset{n}{}}}\left[{\displaystyle \frac{1}{8\pi ^2\widehat{s}_l}}{\displaystyle d^2q_l_{}iT_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_l^+,x_l^{},s,q_l_{}^2)}\right]`$
$`F_{\mathrm{remn}}^{h_1}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^+\right)F_{\mathrm{remn}}^{h_2}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^{}\right)\delta ^{(2)}\left({\displaystyle \underset{k=1}{\overset{n}{}}}\stackrel{}{q}_k_{}\stackrel{}{q}_{}\right),`$ (10)
where the partonic amplitudes are defined as
$$T_{1\mathrm{I}\mathrm{P}}^{h_1h_2}=T_{\mathrm{soft}}^{h_1h_2}+T_{\mathrm{sea}\mathrm{sea}}^{h_1h_2}+T_{\mathrm{val}\mathrm{val}}^{h_1h_2}+T_{\mathrm{val}\mathrm{sea}}^{h_1h_2}+T_{\mathrm{sea}\mathrm{val}}^{h_1h_2},$$
with the individual contributions representing the “elementary partonic interactions plus external legs”. The soft or semi-hard sea-sea contributions are given as
$`T_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}^{h_1h_2}(x^+,x^{},s,q_{}^2)=T_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}(s,q_{}^2)F_{\mathrm{part}}^{h_1}(x^+)F_{\mathrm{part}}^{h_2}(x^{})`$
$`\times \mathrm{exp}\left(\left[R_{h_1}^2+R_{h_2}^2\right]q_{}^2\right),`$ (11)
the hard contribution is
$`T_{\mathrm{val}\mathrm{val}}^{h_1h_2}(x^+,x^{},s,q^2)`$ $`=`$ $`{\displaystyle _0^{x^+}}𝑑x_v^+{\displaystyle \frac{x^+}{x_{v_l}^+}}{\displaystyle _0^x^{}}𝑑x_v^{}{\displaystyle \frac{x^{}}{x_v^{}}}{\displaystyle \underset{j,k}{}}T_{\mathrm{hard}}^{jk}(x_v^+x_v^{}s,q^2,Q_0^2)`$
$`\times \overline{F}_{\mathrm{part}}^{h_1,j}(x_v^+,x^+x_v^+)\overline{F}_{\mathrm{part}}^{h_2,k}(x_v^{},x^{}x_v^{})\mathrm{exp}\left(\left[R_{h_1}^2+R_{h_2}^2\right]q_l_{}^2\right),`$
the mixed semi-hard “val-sea” contribution is given as
$`T_{\mathrm{val}\mathrm{sea}}^{h_1h_2}(x^+,x^{},s,q^2)`$ $`=`$ $`{\displaystyle _0^{x^+}}𝑑x_v^+{\displaystyle \frac{x^+}{x_v^+}}{\displaystyle \underset{j}{}}T_{\mathrm{val}\mathrm{sea}}^j(x_v^+x^{}s,q^2,Q_0^2)`$
$`\times \overline{F}_{\mathrm{part}}^{h_1,j}(x_v^+,x^+x_v^+)F_{\mathrm{part}}^{h_2}(x^{})\mathrm{exp}\left(\left[R_{h_1}^2+R_{h_2}^2\right]q_l_{}^2\right),`$
and the contribution “sea-val” is finally obtained from “val-sea” by exchanging variables,
$$T_{\mathrm{sea}\mathrm{val}}^{h_1h_2}(x^+,x^{},s,q^2)=T_{\mathrm{val}\mathrm{sea}}^{h_2h_1}(x^{},x^+,s,q^2).$$
Here, we allow formally any number of valence type interactions (based on the fact that multiple valence type processes give negligible contribution). In the valence contributions, we have convolutions of hard parton scattering amplitudes $`T_{\mathrm{hard}}^{jk}`$ and valence quark distributions $`\overline{F}_{\mathrm{part}}^j`$ over the valence quark momentum share $`x_v^\pm `$ rather than a simple product, since only the valence quarks are involved in the interactions, with the anti-quarks staying idle (the external legs carrying momenta $`x^+`$ and $`x^{}`$ are always quark-anti-quark pairs).
The profile function $`\gamma `$ is as usual defined as
$$\gamma _{h_1h_2}(s,b)=\frac{1}{2s}2\mathrm{I}\mathrm{m}\stackrel{~}{\mathrm{T}}_{h_1h_2}(s,b),$$
which may be evaluated using the AGK cutting rules with the result (assuming imaginary amplitudes)
$`\gamma _{h_1h_2}(s,b)`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}{\displaystyle _0^1}{\displaystyle \underset{\mu =1}{\overset{m}{}}}dx_\mu ^+dx_\mu ^{}{\displaystyle \underset{\mu =1}{\overset{m}{}}}G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},s,b)`$ (12)
$`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{l!}}{\displaystyle _0^1}{\displaystyle \underset{\lambda =1}{\overset{l}{}}}d\stackrel{~}{x}_\lambda ^+d\stackrel{~}{x}_\lambda ^{}{\displaystyle \underset{\lambda =1}{\overset{l}{}}}G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(\stackrel{~}{x}_\lambda ^+,\stackrel{~}{x}_\lambda ^{},s,b)`$
$`F_{\mathrm{remn}}\left(x^{\mathrm{proj}}{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^+\right)F_{\mathrm{remn}}\left(x^{\mathrm{targ}}{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^{}\right),`$
with $`x^{\mathrm{proj}/\mathrm{targ}}=1x_\mu ^\pm `$ being the momentum fraction of the projectile/target remnant, and with a partonic profile function $`G`$ given as
$`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\lambda ^+,x_\lambda ^{},s,b)`$ $`=`$ $`{\displaystyle \frac{1}{2s}}2\mathrm{I}\mathrm{m}\stackrel{~}{T}_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\lambda ^+,x_\lambda ^{},s,b),`$ (13)
see fig. 9.
This is a very important result, allowing to express the total profile function $`\gamma _{h_1h_2}`$ via the elementary profile functions $`G_{1\mathrm{I}\mathrm{P}}^{hg_1h_2}`$.
## 6 Nucleus-Nucleus Scattering
We generalize the discussion of the last section in order to treat nucleus-nucleus scattering. In the Glauber-Gribov approach , the nucleus-nucleus scattering amplitude is defined by the sum of contributions of diagrams, corresponding to multiple scattering processes between parton constituents of projectile and target nucleons. Nuclear form factors are supposed to be defined by the nuclear ground state wave functions. Assuming the nucleons to be uncorrelated, one can make the Fourier transform to obtain the amplitude in the impact parameter representation. Then, for given impact parameter $`\stackrel{}{b}_0`$ between the nuclei, the only formal difference from the hadron-hadron case will be the averaging over nuclear ground states, which amounts to an integration over transverse nucleon coordinates $`\stackrel{}{b}_i^A`$ and $`\stackrel{}{b}_j^B`$ in the projectile and in the target respectively. We write this integration symbolically as
$$𝑑T_{AB}:=\underset{i=1}{\overset{A}{}}d^2b_i^AT_A(b_i^A)\underset{j=1}{\overset{B}{}}d^2b_j^BT_B(b_j^B),$$
(14)
with $`A,B`$ being the nuclear mass numbers and with the so-called nuclear thickness function $`T_A(b)`$ being defined as the integral over the nuclear density $`\rho _{A(B)}`$:
$$T_A(b):=𝑑z\rho _A(b_x,b_y,z).$$
(15)
It is convenient to use the transverse distance $`b_k`$ between the two nucleons from the $`k`$-th nucleon-nucleon pair, i.e.
$$b_k=\left|\stackrel{}{b}_0+\stackrel{}{b}_{\pi (k)}^A\stackrel{}{b}_{\tau (k)}^B\right|,$$
where the functions $`\pi (k)`$ and $`\tau (k)`$ refer to the projectile and the target nucleons participating in the $`k^{\mathrm{th}}`$ interaction (pair $`k`$). In order to simplify the notation, we define a vector $`b`$ whose components are the overall impact parameter $`b_0`$ as well as the transverse distances $`b_1,\mathrm{},b_{AB}`$ of the nucleon pairs,
$$b=\{b_0,b_1,\mathrm{},b_{AB}\}.$$
Then the nucleus-nucleus interaction cross section can be obtained applying the cutting procedure to elastic scattering diagram and written in the form
$$\sigma _{\mathrm{inel}}^{AB}(s)=d^2b_0𝑑T_{AB}\gamma _{AB}(s,b),$$
(16)
where the so-called nuclear profile function $`\gamma _{AB}`$ represents an interaction for given transverse coordinates of the nucleons.
The calculation of the profile function $`\gamma _{AB}`$ as the sum over all cut diagrams of the type shown in fig. 10 does not differ from the hadron-hadron case and follows the rules formulated
in the preceding section:
* For a remnant carrying the light cone momentum fraction $`x`$ ($`x^+`$ in case of projectile, or $`x^{}`$ in case of target), one has a factor $`F_{\mathrm{remn}}(x)`$.
* For each cut elementary diagram (real elementary interaction = dashed vertical line) attached to two participants with light cone momentum fractions $`x^+`$ and $`x^{}`$, one has a factor $`G(x^+,x^{},s,b)`$. Apart from $`x^+`$ and $`x^{}`$, $`G`$ is also a function of the total squared energy $`s`$ and of the relative transverse distance $`b`$ between the two corresponding nucleons (we use $`G`$ as an abbreviation for $`G_{1\mathrm{I}\mathrm{P}}^{NN}`$ for nucleon-nucleon scattering).
* For each uncut elementary diagram (virtual emissions = full vertical line) attached to two participants with light cone momentum fractions $`x^+`$ and $`x^{}`$, one has a factor $`G(x^+,x^{},s,b),`$ with the same $`G`$ as used for the cut diagrams.
* Finally one sums over all possible numbers of cut and uncut Pomerons and integrates over the light cone momentum fractions.
So we find
$`\gamma _{AB}(s,b)`$ $`=`$ $`{\displaystyle \underset{m_1l_1}{}}\mathrm{}{\displaystyle \underset{m_{AB}l_{AB}}{}}(1\delta _{0\mathrm{\Sigma }m_k}){\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\mu =1}{\overset{m_k}{}}dx_{k,\mu }^+dx_{k,\mu }^{}\underset{\lambda =1}{\overset{l_k}{}}d\stackrel{~}{x}_{k,\lambda }^+d\stackrel{~}{x}_{k,\lambda }^{}\right\}}`$ (17)
$`\times `$ $`{\displaystyle \underset{k=1}{\overset{AB}{}}}\left\{{\displaystyle \frac{1}{m_k!}}{\displaystyle \frac{1}{l_k!}}{\displaystyle \underset{\mu =1}{\overset{m_k}{}}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b_k){\displaystyle \underset{\lambda =1}{\overset{l_k}{}}}G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b_k)\right\}`$
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F_{\mathrm{remn}}\left(x_i^+{\displaystyle \underset{\pi (k)=i}{}}\stackrel{~}{x}_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F_{\mathrm{remn}}\left(x_j^{}{\displaystyle \underset{\tau (k)=j}{}}\stackrel{~}{x}_{k,\lambda }^{}\right)`$
with
$`x_i^{\mathrm{proj}}`$ $`=`$ $`1{\displaystyle \underset{\pi (k)=i}{}}x_{k,\mu ,}^+`$
$`x_j^{\mathrm{targ}}`$ $`=`$ $`1{\displaystyle \underset{\tau (k)=j}{}}x_{k,\mu }^{}.`$
The summation indices $`m_k`$ refer to the number of cut elementary diagrams and $`l_k`$ to number of uncut elementary diagrams, related to nucleon pair $`k`$. For each possible pair $`k`$ (we have altogether $`AB`$ pairs), we allow any number of cut and uncut diagrams. The integration variables $`x_{k,\mu }^\pm `$ refer to the $`\mu ^{\mathrm{th}}`$ elementary interaction of the $`k^{\mathrm{th}}`$ pair for the cut elementary diagrams, the variables $`\stackrel{~}{x}_{k,\lambda }^\pm `$ refer to the corresponding uncut elementary diagrams. The arguments of the remnant functions $`F_{\mathrm{remn}}`$ are the remnant light cone momentum fractions, i.e. unity minus the momentum fractions of all the corresponding elementary contributions (cut and uncut ones). We also introduce the variables $`x_i^{\mathrm{proj}}`$and $`x_j^{\mathrm{targ}}`$, defined as unity minus the momentum fractions of all the corresponding cut contributions, in order to integrate over the uncut ones (see below).
The expression for $`\gamma _{AB}`$ sums up all possible numbers of cut Pomerons $`m_k`$ with one exception due to the factor $`(1\delta _{0\mathrm{\Sigma }m_k})`$: one does not consider the case of all $`m_k`$’s being zero, which corresponds to “no interaction” and therefore does not contribute to the inelastic cross section. We may therefore define a quantity $`\overline{\gamma }_{AB}`$, representing “no interaction”, by taking the expression for $`\gamma _{AB}`$ with $`(1\delta _{0\mathrm{\Sigma }m_k})`$ replaced by $`(\delta _{0\mathrm{\Sigma }m_k})`$:
$`\overline{\gamma }_{AB}(s,b)`$ $`=`$ $`{\displaystyle \underset{l_1}{}}\mathrm{}{\displaystyle \underset{l_{AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\lambda =1}{\overset{l_k}{}}d\stackrel{~}{x}_{k,\lambda }^+d\stackrel{~}{x}_{k,\lambda }^{}\right\}\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{l_k!}\underset{\lambda =1}{\overset{l_k}{}}G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b_k)\right\}}`$ (18)
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F^+\left(1{\displaystyle \underset{\pi (k)=i}{}}\stackrel{~}{x}_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F^{}\left(1{\displaystyle \underset{\tau (k)=j}{}}\stackrel{~}{x}_{k,\lambda }^{}\right).`$
One now may consider the sum of “interaction” and “no interaction”, and one obtains easily
$$\gamma _{AB}(s,b)+\overline{\gamma }_{AB}(s,b)=1.$$
(19)
Based on this important result, we consider $`\gamma _{AB}`$ to be the probability to have an interaction and correspondingly $`\overline{\gamma }_{AB}`$ to be the probability of no interaction, for fixed energy, impact parameter and nuclear configuration, specified by the transverse distances $`b_k`$ between nucleons, and we refer to eq. (19) as “unitarity relation”. But we want to go even further and use an expansion of $`\gamma _{AB}`$ in order to obtain probability distributions for individual processes, which then serves as a basis for the calculations of exclusive quantities.
The expansion of $`\gamma _{AB}`$ in terms of cut and uncut Pomerons as given above represents a sum of a large number of positive and negative terms, including all kinds of interferences, which excludes any probabilistic interpretation. We have therefore to perform summations of interference contributions – sum over any number of virtual elementary scatterings (uncut Pomerons) – for given non-interfering classes of diagrams with given numbers of real scatterings (cut Pomerons). Let us write the formulas explicitly. We have
$`\gamma _{AB}(s,b)`$ $`=`$ $`{\displaystyle \underset{m_1}{}}\mathrm{}{\displaystyle \underset{m_{AB}}{}}(1\delta _{0{\scriptscriptstyle m_k}}){\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\mu =1}{\overset{m_k}{}}dx_{k,\mu }^+dx_{k,\mu }^{}\right\}}`$ (20)
$`\times `$ $`{\displaystyle \underset{k=1}{\overset{AB}{}}}\left\{{\displaystyle \frac{1}{m_k!}}{\displaystyle \underset{\mu =1}{\overset{m_k}{}}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b_k)\right\}\mathrm{\Phi }_{AB}(x^{\mathrm{proj}},x^{\mathrm{targ}},s,b),`$
where the function $`\mathrm{\Phi }`$ representing the sum over virtual emissions (uncut Pomerons) is given by the following expression
$`\mathrm{\Phi }_{AB}(x^{\mathrm{proj}},x^{\mathrm{targ}},s,b)`$ $`=`$ $`{\displaystyle \underset{l_1}{}}\mathrm{}{\displaystyle \underset{l_{AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\lambda =1}{\overset{l_k}{}}d\stackrel{~}{x}_{k,\lambda }^+d\stackrel{~}{x}_{k,\lambda }^{}\right\}\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{l_k!}\underset{\lambda =1}{\overset{l_k}{}}G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b_k)\right\}}`$ (21)
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F_{\mathrm{remn}}\left(x_i^{\mathrm{proj}}{\displaystyle \underset{\pi (k)=i}{}}\stackrel{~}{x}_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F_{\mathrm{remn}}\left(x_j^{\mathrm{targ}}{\displaystyle \underset{\tau (k)=j}{}}\stackrel{~}{x}_{k,\lambda }^{}\right).`$
This summation has to be carried out, before we may use the expansion of $`\gamma _{AB}`$ to obtain probability distributions. This is far from trivial, the necessary methods are described in . To make the notation more compact, we define matrices $`X^+`$ and $`X^{}`$, as well as a vector $`m`$, via
$`X^+`$ $`=`$ $`\left\{x_{k,\mu }^+\right\},`$
$`X^{}`$ $`=`$ $`\left\{x_{k,\mu }^{}\right\},`$
$`m`$ $`=`$ $`\{m_k\},`$
which leads to
$`\gamma _{AB}(s,b)`$ $`=`$ $`{\displaystyle \underset{m}{}}(1\delta _{0m}){\displaystyle 𝑑X^+𝑑X^{}\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})},`$
$`\overline{\gamma }_{AB}(s,b)`$ $`=`$ $`\mathrm{\Omega }_{AB}^{(s,b)}(0,0,0),`$
with
$$\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})=\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{m_k!}\underset{\mu =1}{\overset{m_k}{}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b_k)\right\}\mathrm{\Phi }_{AB}(x^{\mathrm{proj}},x^{\mathrm{targ}},s,b).$$
This allows to rewrite the unitarity relation eq. (19) in the following form,
$$\underset{m}{}𝑑X^+𝑑X^{}\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})=1.$$
This equation is of fundamental importance, because it allows us to interpret $`\mathrm{\Omega }^{(s,b)}(m,X^+,X^{})`$ as probability density of having an interaction configuration characterized by $`m`$, with the light cone momentum fractions of the Pomerons being given by $`X^+`$ and $`X^{}`$.
## 7 Virtual Emissions and Markov Chain Techniques
What did we achieve so far? We have formulated a well defined model, introduced by using the language of field theory, solving in this way the severe consistency problems of the most popular current approaches. To proceed further, one needs to solve two fundamental problems:
* the sum $`\mathrm{\Phi }_{AB}`$ over virtual emissions has to be performed,
* tools have to be developed to deal with the multidimensional probability distribution $`\mathrm{\Omega }_{AB}^{(s,b)}`$,
both being very difficult tasks. Introducing new numerical techniques, we were able to solve both problems, as discussed in very detail in .
Calculating the sum over virtual emissions ($`\mathrm{\Phi }_{AB}`$) is done by parameterizing the functions $`G`$ as analytical functions and performing analytical calculations. By studying the properties of $`\mathrm{\Phi }_{AB}`$, we find that at very high energies the theory is no longer unitary without taking into account screening corrections due to triple Pomeron interactions. In this sense, we consider our work as a first step to construct a consistent model for high energy nuclear scattering, but there is still work to be done.
Concerning the multidimensional probability distribution $`\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})`$, we employ methods well known in statistical physics (Markov chain techniques). So finally, we are able to calculate the probability distribution $`\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})`$, and are able to generate (in a Monte Carlo fashion) configurations $`(m,X^+,X^{})`$ according to this probability distribution.
## 8 Summary
What are finally the principal features of our basic results, summarized in eqs. (16, 20, 21)? Contrary to the traditional treatment (Gribov-Regge approach or parton model), all individual elementary contributions $`G`$ depend explicitly on the light-cone momenta of the elementary interactions, with the total energy-momentum being precisely conserved. Another very important feature is the explicit dependence of the screening contribution $`\mathrm{\Phi }_{AB}`$ (the contribution of virtual emissions) on the remnant momenta. The direct consequence of properly taking into account energy-momentum conservation in the multiple scattering process is the validity of the so-called AGK-cancelations in hadron-hadron and nucleus-nucleus collisions in the entire kinematical region.
The formulas (16, 20, 21) allow to develop a consistent scheme to simulate high energy nucleus-nucleus interactions. The corresponding Monte Carlo procedure is exactly based on the cross section formulas so that the entire model is fully self-consistent.
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# Contents
## 1 Introduction
A recent paper gives a systematic account of the invariant symmetric and skewsymmetric primitive tensors that may be constructed on a compact simple Lie algebra $`𝔤`$ of rank $`l`$. The new family of symmetric tensors introduced in allows the direct construction of the $`l`$ primitive Racah-Casimir operators for $`𝔤`$; the antisymmetric tensors determine the $`l`$ primitive Lie algebra cohomology cocycles of $`𝔤`$. We refer e.g, to and references therein for the definitions and explanation of the significance of the invariant skewsymmetric tensors
$$\mathrm{\Omega }_{}^{(2m1)}{}_{i_1i_2\mathrm{}i_{2m1}}{}^{}\mathrm{\Omega }_{i_1i_2\mathrm{}i_{2m1}}(i_1,i_2\mathrm{}\{1,2,\mathrm{},\mathrm{dim}𝔤\}),$$
(1)
that are associated with the primitive cocycles of $`𝔤`$, of $`l`$ different orders $`q=2m1`$, where $`m`$ is the order of the associated Racah-Casimir operators. The $`l`$ allowed values of $`q`$ for the different Lie algebra cohomology groups of each $`𝔤`$ are well-known (see e.g. for tables and further references). For $`su(n)=A_l,l=n1`$, $`m=2,3,\mathrm{},n`$, the cocycles $`\mathrm{\Omega }^{(2m1)}`$ have orders $`q=3,5,\mathrm{},2m1`$ and indices $`i_1,\mathrm{},i_{(2m1)}=1,\mathrm{},n^21`$. Each tensor (1) is used in to define exactly one member of a family of maximally traceless (Sec.2) fully symmetric tensors $`t`$. These have a very favourable status within the set of totally symmetric tensors, because each of them can be used to define one primitive Racah-Casimir operator for $`𝔤`$, $`l`$ in all, and no more. For example, if one specialises to $`su(3)`$ the general $`su(n)`$ definition of an $`m`$-th order tensor $`t^{(m)}`$ of our favoured family of symmetric tensors, one finds that the definition collapses to zero for any $`m>3`$, in keeping with the fact that $`su(3)`$ has no primitive Racah-Casimir operators of any order greater than three.
These matters are fully discussed and illustrated in , which pay especial attention to the case of the $`SU(n)`$ group. However, if one wants to make application of tensors like (1), for example, in the construction of higher supercharges and in the quantum mechanics of particles with $`SU(n)`$ colour, or in the definition of BRST-like operators and higher order Hodge analysis , one finds the need for more identities involving them than are to be found in or elsewhere. There are also other areas in which good control over the properties of the Omega-tensors/cocycles (1) is valuable. One of these is in the discussion of multi-bracket generalisations of Lie algebras and higher order linear Poisson structures of the type introduced in . Another is in the construction of Wess-Zumino terms for effective actions in space-times of various dimensions (see and references therein) and, in general, in the group theory factors that may appear in particle physics. A recent study of this last subject is .
The aim in this paper is to collect all that we currently know regarding the properties of and the identities involving the $`su(n)`$ algebra skewsymmetric Omega tensors (see (9)–(12) below). Our approach divides itself into three stages. The first presents a discussion of the Omega tensors of $`su(n)`$ that sets out from their definitions, and utilises only the properties of the $`f`$\- and $`d`$-tensors to deliver its output, which is then employed in the second stage. This stage is based on the fact that the Omega tensors play a central role in the discussion of the algebra of totally antisymmetrised products of an even or odd number of lambda-matrices of $`su(n)`$
$$\lambda _{[ijk\mathrm{}s]}=\lambda _{[i}\lambda _j\lambda _k\mathrm{}\lambda _{s]},i,j,\mathrm{},s=1,\mathrm{},n^21,$$
(2)
and accordingly we provide an extensive discussion of results within this algebra. The completeness properties and trace identities for such products thus obtained then give rise to the powerful approach contained in the third stage of our programme, which allows us to derive, amongst other results, one of special importance. Defining the fully contracted scalar
$$(\mathrm{\Omega }^{(2s+1)}){}_{}{}^{2}=\mathrm{\Omega }_{i_1j_1\mathrm{}i_sj_sk}\mathrm{\Omega }_{i_1j_1\mathrm{}i_sj_sk},$$
(3)
we derive the recursion relation (195)
$$(\mathrm{\Omega }^{(2s+1)}){}_{}{}^{2}=\frac{4}{2s(2s1)}(n^2s^2)(\mathrm{\Omega }^{(2s1)}){}_{}{}^{2},$$
(4)
and its explicit solution (196).
The content of this paper is organised as follows. Sec. 2 gives the basic definitions of Omega tensors. Then Sec. 3 gives results, within the first stage of our study, classified roughly according to types, e.g. Jacobi or $`Ad`$-invariance results, contractions including those with the structure constants $`f_{ijk}`$ of $`su(n)`$ (that define the three-cocycle $`\mathrm{\Omega }^{(3)}`$), recursive relations, duality relations, product relations, uniqueness questions for antisymmetric tensors, etc. Some derivations (of formulas not derived elsewhere) are given in Sec. 4 in order to illustrate methods employed within stage one of our work. In Sec. 5, we turn to the development of the algebra of the quantities (2), using results for Omega tensors deduced independently of it. This then enables the attack in Sec. 6 on the identities (3) and (4), using completeness and trace properties of the products (2) found in Sec. 5. Within this approach some auxiliary results are merely quoted in Sec. 6, and proved in Sec. 7. Some critical questions raised within the discussion in Sec. 3.7 of the uniqueness of Omega tensors are answered in Sec. 8, using lambda-matrix methods.
There have been in the literature very many discussions of the invaritensors, symmetric and skewsymmetric, associated with a simple Lie algebra $`𝔤`$ of rank $`l`$. A significant recent one is ; lists of references are given in . However such studies often do not consider the full set of invariant tensors for $`𝔤`$, neglecting the $`(l1)`$ higher order Lie algebra cocycles i.e., the invariant higher order antisymmetric Omega tensors. Our paper emphasises the crucial role these Omega tensors play not only in the method we advocate to define Racah-Casimir operators, but also in our discussion of their eigenvalues and the corresponding generalised Dynkin indices. One additional motivation for the present paper is to make readily available a comprehensive listing of results involving Omega tensors that are needed for that programme.
## 2 Definitions of $`su(n)`$ $`d,\mathrm{\Omega }`$ and $`t`$ tensors
We start with a family of symmetric invariant tensors, the $`d`$-family. It is easily defined recursively starting from the standard Gell-Mann totally symmetric tensor $`d_{ijk}`$ (see eq. (110)). First, one constructs
$$d^{(r+1)}{}_{i_1\mathrm{}i_{r+1}}{}^{}=d^{(r)}{}_{i_1\mathrm{}i_{r1}j}{}^{}d_{}^{(3)}{}_{ji_ri_{r+1}}{}^{},r=3,4,\mathrm{}.$$
(5)
For $`r3`$, eq. (5) does not define totally symmetric tensors. The $`d`$-family of symmetric tensors is obtained by symmetrising over all free indices in (5) and hence is defined by
$$d^{(r)}{}_{(i_1\mathrm{}i_r)}{}^{},$$
(6)
where the round brackets indicate symmetrisation with unit weight over the set of indices enclosed. This should be done as economically as possible, e.g.
$$d^{(4)}{}_{(ijkl)}{}^{}=\frac{1}{3}(d_{ijt}d_{klt}+d_{jkt}d_{ilt}+d_{kit}d_{jlt}).$$
(7)
The lowest order symmetric tensor, the Cartan-Killing metric (since $`𝔤`$ is compact and the generators hermitian it will be taken as the unity) may be viewed as the order two member of the $`d`$-family (6), $`d_{ij}\delta _{ij}`$. Since the iteration process (5),(6) can go on idefinitely, it is clear that not all tensors of the $`d`$-family are primitive, since for a simple algebra $`𝔤`$ of rank $`l`$ there are only $`l`$ invariant primitive symmetric tensors (or, equivalently, $`l`$ primitive Racah-Casimir operators).
We now turn to the totally antisymmetric Omega tensors (1), referring to for an explanation of their cohomological origin. Thus we define
$`\mathrm{\Omega }^{(3)}{}_{ijk}{}^{}f_{ijk}`$ $`=`$ $`f^a{}_{ij}{}^{}d_{ak}^{},`$ (8)
$`\mathrm{\Omega }^{(5)}{}_{ijklm}{}^{}\mathrm{\Omega }_{ijklm}`$ $`=`$ $`f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl]}{}^{}d_{abm}^{},`$ (9)
$`\mathrm{\Omega }^{(7)}{}_{ijklmpq}{}^{}\mathrm{\Omega }_{ijklmpq}`$ $`=`$ $`f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{mp]}{}^{}d_{}^{(4)}{}_{(abcq)}{}^{},`$ (10)
$`\mathrm{\Omega }^{(9)}{}_{ijklmpqrs}{}^{}\mathrm{\Omega }_{ijklmpqrs}`$ $`=`$ $`f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{mp}{}^{}f_{}^{d}{}_{qr]}{}^{}d_{}^{(5)}{}_{(abcds)}{}^{},`$ (11)
$`\mathrm{\Omega }^{(11)}{}_{ijklmpqrsuv}{}^{}\mathrm{\Omega }_{ijklmpqrsuv}`$ $`=`$ $`f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{mp}{}^{}f_{}^{d}{}_{qr}{}^{}f_{}^{e}{}_{su]}{}^{}d_{}^{(6)}{}_{(abcdev)}{}^{},`$ (12)
and so on. Here, the square brackets imply total unit weight antisymmetrisation over the set of indices which they enclose. The raising of indices is trivial from a metric point of view, and is usually used in this paper in order to exempt certain indices from the antisymmetrisation (or symmetrisation) effect of the square (round) brackets. We note that the $`\mathrm{\Omega }^{(2m1)}`$ tensor is fully skewsymmetric in its $`(2m1)`$ indices in spite of the fact that only $`(2m2)`$ indices are explicitly antisymmetrised in the r.h.s. above . Some further discussion of the properties of the $`d`$-tensors and their role in the definitions of the Omega tensors is given below as Sec. 2.1.
Next we use the Omega tensors to define a second family of invariant fully symmetric tensors, the $`t`$-tensors, as follows
$`\mathrm{\Omega }_{ijm}f_{ija}`$ $`=`$ $`t^{(2)}{}_{am}{}^{},`$ (13)
$`\mathrm{\Omega }_{ijklm}f_{ija}f_{klb}`$ $`=`$ $`t^{(3)}{}_{abm}{}^{},`$ (14)
$`\mathrm{\Omega }_{ijklmpq}f_{ija}f_{klb}f_{mpc}`$ $`=`$ $`t^{(4)}{}_{abcq}{}^{},`$ (15)
$`\mathrm{\Omega }_{ijklmpqrs}f_{ija}f_{klb}f_{mpc}f_{qrd}`$ $`=`$ $`t^{(5)}{}_{abcds}{}^{},`$ (16)
etc.; they are fully symmetric on account of the skewsymmetry of the $`\mathrm{\Omega }`$’s. The tensors on the right hand side of (14)–(16) have been evaluated before <sup>1</sup><sup>1</sup>1Eqs. (18) and (19) above correct the overall factors of (6.13) and (6.14) in ; the $`i`$ difference in is due to the fact that here we take the generators of $`g`$ hermitian, see (128).. We give below the expression of the lower order $`su(n)`$ $`t`$ tensors in terms of members of the $`d`$-family,
$`t^{(2)}_{ij}`$ $`=`$ $`n\delta _{ij},`$ (17)
$`t^{(3)}_{ijk}`$ $`=`$ $`\frac{1}{3}n^2d_{ijk},`$ (18)
$`t^{(4)}_{ijkl}`$ $`=`$ $`\frac{1}{15}(n(n^2+1)d^{(4)}{}_{(ijkl)}{}^{}2(n^24)\delta _{(ij}\delta _{kl)}),`$ (19)
$`t^{(5)}_{ijklm}`$ $`=`$ $`\lambda (n)\left(n(n^2+5)d^{(5)}{}_{(ijklm)}{}^{}2(3n^220)d_{(ijk}\delta _{lm)}\right),`$ (20)
where the function $`\lambda (n)`$, not determined in , turns out from the work of Sec. 6 to be given by
$$\lambda (n)=\frac{n}{105}.$$
(21)
The tensors (18) and higher collapse to zero when their order $`m`$ is larger than $`n`$. While Eq. (20) can be indeed be used as it stands (and will be to avoid circularity of argumentation below), much of the information we need will be seen to follow from the definitions (14)–(16) and the properties of Omega tensors.
The $`t`$-tensors are totally symmetric and, unlike the higher ($`m>3`$) order $`d`$-tensors, they are orthogonal to all other $`t`$-tensors of different order (Lemma 3.3 of ). For instance, for $`t^{(4)}`$ this means
$$t^{(4)}{}_{ijkl}{}^{}\delta _{ij}^{}=0,t^{(4)}{}_{ijkl}{}^{}t_{}^{(3)}{}_{ijk}{}^{}=0.$$
(22)
In contrast, since trace formulas for $`d`$-tensors easily give
$$d^{(4)}{}_{(ijkl)}{}^{}d_{ijm}^{}=\frac{2}{3}\frac{(n^28)}{n}d_{klm},$$
(23)
the contraction of only two indices gives
$$t^{(4)}{}_{ijkl}{}^{}t_{}^{(3)}{}_{ijm}{}^{}=\frac{1}{3}n^2t^{(4)}{}_{ijkl}{}^{}d_{ijm}^{}=\frac{2}{45}n^2(n^29)t^{(3)}{}_{klm}{}^{}.$$
(24)
For $`t^{(5)}`$ we have
$$t^{(5)}{}_{ijklm}{}^{}\delta _{ij}^{}=0,t^{(5)}{}_{ijklm}{}^{}t_{}^{(3)}{}_{ijk}{}^{}=0,t^{(5)}{}_{ijklm}{}^{}t_{}^{(4)}{}_{ijkl}{}^{}=0$$
(25)
or, in other words, the maximal contraction of the indices of two $`t`$-tensors of different order is zero.
Another issue concerns the claims made above that $`t^{(4)}`$ vanishes identically for $`n=3`$, and that $`t^{(5)}`$ vanishes identically for $`n=3`$ and $`n=4`$. This is a point one will see clearly illustrated in many places, and is the result of factors becoming zero or of relations expressing the $`d^{(m)}`$ tensors in terms of primitive ones when $`m>n`$. The necessary identities special to $`su(3)`$ and $`su(4)`$ are noted in . Indeed almost all we need here in the way of identities involving the $`d`$\- and $`f`$-tensors of $`su(n)`$ is presented in , especially in the appendix. See also and .
### 2.1 More about $`d`$-tensors
Below we shall need the identities that express the $`Ad`$-invariance of the $`d`$-tensors, namely,
$`d_{t(ij}f_{k)st}`$ $`=`$ $`0,`$ (26)
$`d_{}^{(4)}{}_{t(ijk}{}^{}f_{l)st}`$ $`=`$ $`0,`$ (27)
$`d_{}^{(5)}{}_{t(ijkl}{}^{}f_{m)st}`$ $`=`$ $`0,`$ (28)
and so on.
One use of (26) is as follows. Referring first to (9), we note that the symmetry properties of the $`SU(n)`$-invariant $`d`$\- and $`f`$-tensors permit the left-hand square bracket to be moved, without altering the definition, one place to the right, so as to enclose only $`j,k`$ and $`l`$. Then application of (26) allows one to show that the right side of (9) is antisymmetric under the interchange of $`i`$ and $`m`$, and hence, as mentioned, indeed defines a totally antisymmetric quantity.
Referring next to (10), we note that it may be simplified in either of two ways but not simultaneously both: one of these is the move of the left hand square bracket one place to the right (or alternatively the right hand one to the left), the other employs the result
$$d^{(4)}{}_{(abcq)}{}^{}=d^{(4)}{}_{(abc)q}{}^{},$$
(29)
and thereby allows $`d^{(4)}_{(abcq)}`$ to be replaced in (10) by one of the terms of $`d^{(4)}_{(abc)q}`$, e.g. $`d_{xab}d_{cqx}`$. It is the use of (29) that deserves close attention. It displays a simplifying feature of the $`d^{(4)}`$ situation that does not generalise systematically to $`d^{(r)}`$ for $`r>4`$. For $`r=5`$, we have
$$d^{(5)}{}_{(abcdq)}{}^{}=\frac{1}{5}d^{(5)}{}_{(ab}{}^{q}{}_{cd)}{}^{}+\frac{4}{5}d^{(5)}{}_{(abcd)q}{}^{},$$
(30)
where
$`d^{(5)}{}_{(ab}{}^{}{}_{}{}^{q}_{cd)}`$ $`=`$ $`d_{(ab}{}_{}{}^{x}d_{}^{xqy}d^y{}_{cd)}{}^{},`$ (31)
$`d^{(5)}_{(abcd)q}`$ $`=`$ $`d_{(ab}{}_{}{}^{x}d_{}^{x}{}_{c}{}^{}{}_{}{}^{y}d_{}^{y}{}_{d)q}{}^{}.`$ (32)
Inspection of the evident tree-diagram representation of the tensors occurring here makes clear the fractions seen in (30).
In (30) we meet an obstacle to extending, to (12) and beyond, the simple proof that allowed the right hand round bracket in (10) to be moved one place to the left. It is nevertheless a generally allowed step, providing a valuable simplification of the definitions of $`\mathrm{\Omega }^{2s+1)}`$ for $`s4`$. However, to obtain a convenient proof of this, we need to have recourse to lambda matrix methods, and so, will return to the matter in Sec.5.2.
Similarly
$$d^{(6)}{}_{(abcdeq)}{}^{}=\frac{1}{3}d^{(6)}{}_{(ab}{}^{q}{}_{cde)}{}^{}+\frac{2}{3}d^{(6)}{}_{(abcde)q}{}^{}.$$
(33)
The fractions in the RHS of (33) arise because the tree diagrams in use are trees with four equivalent end twigs and two equivalent non-end twigs.
An additional complication enters for symmetric tensors of order six, one that has already been observed in . The tensor $`d^{(6)}`$ that enters the definition (11) of $`\mathrm{\Omega }^{(11)}`$ is the $`r=6`$ member of the family (6). But it is not the only primitive symmetric sixth order tensor that can be defined. One also has $`d^{(6)}`$ given by
$$d^{(6)}{}_{(abcdef)}{}^{}=d_{(ab}{}_{}{}^{x}d_{cd}^{}{}_{}{}^{y}d_{ef)}^{}{}_{}{}^{z}d_{xyz}^{}.$$
(34)
The tensors $`d^{(6)}_{(abcdef)}`$ and $`d^{(6)}_{(abcdef)}`$ are for our purposes equivalent. It is shown in (below eq. (A.21) there) that they differ by non-primitive terms which are symmetrised products of lower order $`d`$-tensors. The claimed equivalence follows from the fact that such non-primitive terms cannot contribute to the definition (11) of $`\mathrm{\Omega }^{(11)}`$ because of Jacobi identities.
Inspection of the relevant tree diagram shows that
$$d^{(6)}{}_{(abcdeq)}{}^{}=d^{(6)}{}_{(abcde)q}{}^{}.$$
(35)
## 3 Identities involving the Omega tensors
These are mostly displayed for Omega tensors of lower order for obvious reasons. But one can often see patterns that would guide an attack on higher order analogues that may now seem out of reasonable reach, or perhaps just until the need of a specific application provides the necessary motivation. The trace methods for products of the hermitian $`D`$\- and $`F`$-matrices , where $`(D_i)_{jk}=d_{ijk}`$ and $`(F_i)_{jk}=if_{ijk}`$, such as are seen in use in the derivations presented in Sec. 4, become discouraging when one cannot avoid doing a trace that is more than of fourth order, unless one can harness computational skills like those of . Also, finding a viable path through increasing complication becomes progressively more taxing. The necessity for going on in later sections to develop an alternative approach – that which makes systematic use of lambda matrices – come into evidence in this way.
### 3.1 Jacobi identities
First, we consider the Jacobi identities which express the $`Ad`$-invariance of the Omega tensors. The $`Ad`$-invariance of $`\mathrm{\Omega }^{(3)}{}_{ijk}{}^{}=f_{ijk}`$ is expressed by the Jacobi identity,
$$f_{t[ij}f_{k]lt}=f_{t[ij}f_{k]ls}\delta _{st}=0.$$
(36)
For higher $`\mathrm{\Omega }`$’s $`Ad`$-invariance gives
$`\mathrm{\Omega }_{t[ijkl}f_{p]qt}`$ $`=`$ $`0,`$ (37)
$`\mathrm{\Omega }_{t[ijklpq}f_{r]st}`$ $`=`$ $`0,`$ (38)
$`\mathrm{\Omega }_{t[ijklpqrs}f_{u]vt}`$ $`=`$ $`0.`$ (39)
In analogy with the second way of writing the Jacobi identity (36), we may usefully expand (37)–(39) in terms of the higher members of the $`d`$-family (6) getting
$`f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{p]q}{}^{}d_{abc}^{}=0,`$ (40)
$`f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{pq}{}^{}f_{}^{d}{}_{r]s}{}^{}d_{}^{(4)}{}_{(abcd)}{}^{}=0.`$ (41)
In the last two results, obviously the right hand square bracket can be moved one place to the right. In the case of (41) this allows the round symmetrising brackets to be taken off $`d^{(4)}`$ since $`f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{pq}{}^{}f_{}^{d}_{rs]}`$ is fully symmetric in $`abcd`$.
Also, using (38), eq. (41) can be rearranged to read
$$\mathrm{\Omega }_{t[ijkl}\mathrm{\Omega }_{pqr]st}=0,$$
(42)
an expression that may be understood as the generalised Jacobi identity (GJI) for a higher (fourth) order multibracket algebra. In the general case, the GJI reads
$$\mathrm{\Omega }_{t[i_1\mathrm{}i_{2m2}}\mathrm{\Omega }_{i_{2m1}\mathrm{}i_{4m5}]i_{4m4}t}=0.$$
(43)
This is an identity that follows directly if one takes the coordinates of $`\mathrm{\Omega }^{(2m1)}`$ as the generalised structure constants (with $`(2m1)`$ antisymmetric indices) of a multi-bracket Lie algebra of order $`(2m2)`$. As in the standard case (for $`m=2`$, eq. (43) reduces to the JI, eq. (36)), only the associativity of the $`(2m2)`$ entries in the fully skewsymmetric multibracket is required to obtain eq. (43).
The importance of the Jacobi identities can hardly be overemphasised: they are essential to the simplification of other identities in all classes, as will be seen below. Certain other results which bear a close resemblance to (37)–(39) are also valid, namely
$`\mathrm{\Omega }_{ti[jkl}f_{pq]t}`$ $`=`$ $`0,`$ (44)
$`\mathrm{\Omega }_{ti[jklpq}f_{rs]t}`$ $`=`$ $`0,`$ (45)
$`\mathrm{\Omega }_{ti[jklpqrs}f_{uv]t}`$ $`=`$ $`0,`$ (46)
etc. One proves these results by inserting the definitions of the Omega tensors into their left sides and making simple rearrangements that allow (36) to be used to produce the answers zero. Relations of this type may be understood as mixed GJI identities. Their general form is
$$\mathrm{\Omega }_{t[i_1\mathrm{}i_r}\mathrm{\Omega }_{i_{r+1}\mathrm{}i_{r+s1}]i_{r+s}t}=0(r\mathrm{and}s\mathrm{even}),$$
(47)
and constitute consistency relations that must be satisfied by the generalised structure constants, and hence have the same origin as the GJI (see ). Here we have followed the converse path, showing that these relations follow from the definition of the $`\mathrm{\Omega }`$ tensors.
It is also worth noting that many of the results of this section can be identified as cases of the general result Lemma 3.1 of :
$$f^{p_1}{}_{[i_1j_1}{}^{}\mathrm{}f^{p_s}{}_{i_sj_s]}{}^{}k_{(p_1\mathrm{}p_s)}^{}=0,$$
(48)
where $`k_{(p_1\mathrm{}p_s)}`$ is any $`Ad`$-invariant totally symmetric tensor of order $`s`$.
### 3.2 On the definition of the $`\mathrm{\Omega }`$ tensors
Since we have introduced the Omega tensors using the recursively defined $`d`$-tensors (6), and then used the Omega tensors to obtain the preferred family of $`t`$-tensors (eqs.(13)–(16), (17)–(20) etc.), one might well ask why we did not need the latter in order to start the process off in the first place. The answer is that the non-primitive product terms that appear as the tails of the $`t`$-tensors cannot contribute to the Omega tensors at all in virtue of Jacobi identities of the type given in Sec. 3.1. In fact, non-primitive invariant symmetric tensors do not contribute to the $`\mathrm{\Omega }`$ tensors, making the definitions (9)–(12) unique (see , Cor. 3.1). For example (cf. (10)), there is no need to contemplate a contribution to $`\mathrm{\Omega }^{(7)}`$ proportional to
$$f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{mp]}{}^{}\delta _{(ab}^{}\delta _{cq)},$$
(49)
since it vanishes by eq. (36). To see that a similar state of affairs applies to a putative contribution to $`\mathrm{\Omega }^{(9)}`$ like
$$f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{mp}{}^{}f_{}^{d}{}_{qr]}{}^{}d_{(abc}^{}\delta _{ds)},$$
(50)
requires the use of both (36) and (40), depending upon where $`s`$ occurs in the five terms of the expansion of $`d_{(abc}\delta _{ds)}`$. Considerations like those described often employ steps like $`[ijklmp\mathrm{}]=[i[jkl]mp\mathrm{}]`$ ; unit weighted brackets are convenient for such use. Thus apart from overall normalisation, replacing $`d`$-tensors (6) by $`t`$-tensors (see eq. (18)) in the definitions of the Omega tensors has of no effect since, by virtue of (48), the non-primitive parts in which they differ do not contribute.
### 3.3 Recursive identities
We note the important results relating $`\mathrm{\Omega }^{(5)}`$ to $`\mathrm{\Omega }^{(7)}`$ and $`\mathrm{\Omega }^{(7)}`$ to $`\mathrm{\Omega }^{(9)}`$ respectively
$`\mathrm{\Omega }^{(7)}_{ijklmpq}`$ $`=`$ $`\mathrm{\Omega }^{(5)}{}_{t[ijkl}{}^{}f_{mp}^{s}d_{q]st},`$ (51)
$`\mathrm{\Omega }^{(9)}_{ijklmpqrs}`$ $`=`$ $`\mathrm{\Omega }^{(7)}{}_{t[ijklmp}{}^{}f_{qr}^{}{}_{}{}^{u}d_{s]ut}^{},`$ (52)
which are the two lowest versions of a general result (, eq. (7.6)). Having written these results, one sees that the definition of $`\mathrm{\Omega }^{(5)}`$ provides the first member of the series, the identification
$$\mathrm{\Omega }^{(3)}{}_{ijk}{}^{}=f_{ijk},$$
(53)
having been noted already in (8).
Each of the two results just displayed can usefully be presented in a different form
$`\mathrm{\Omega }^{(7)}_{ijklmpq}`$ $`=`$ $`f_{x[ij}\mathrm{\Omega }^{(5)}{}_{klmp}{}^{}{}_{}{}^{y}d_{q]xy}^{},`$ (54)
$`\mathrm{\Omega }^{(9)}_{ijklmpqrs}`$ $`=`$ $`\mathrm{\Omega }^{(5)}{}_{x[ijkl}{}^{}\mathrm{\Omega }_{}^{(5)}{}_{mpqr}{}^{}{}_{}{}^{y}d_{s]xy}^{};`$ (55)
evident generalisations may be expected to hold. It is easy also to use Jacobi identities to show if one replaces $`d`$-tensors by $`f`$-tensors on the right sides of (54) and (55), one gets the answer zero.
### 3.4 Contraction of higher order Omega tensors with lower order ones
In view of the antisymmetry of the Omega tensors, it is clear that these are amongst the most important contractions to be considered. We find
$`f_{ijk}f_{ijl}`$ $`=`$ $`n\delta _{kl}=t_{}^{(2)}{}_{kl}{}^{},`$ (56)
$`\mathrm{\Omega }_{tijkl}f_{ijs}`$ $`=`$ $`\frac{1}{2}nf_{u[kl}d_{t]us},`$ (57)
$`\mathrm{\Omega }_{tijkl}f_{iju}f_{klv}`$ $`=`$ $`\frac{1}{3}n^2d_{tuv}=t^{(3)}{}_{tuv}{}^{},`$ (58)
$`\mathrm{\Omega }_{tijklpq}f_{iju}f_{klv}f_{pqw}`$ $`=`$ $`\frac{1}{15}\left(n(n^2+1)d^{(4)}{}_{(uvwt)}{}^{}2(n^24)\delta _{(uv}\delta _{wt)}\right)=t^{(4)}{}_{tuvw}{}^{}.`$ (59)
The last two of course are just the definitions of $`t^{(3)}`$ and $`t^{(4)}`$ seen from a different viewpoint. We may contract (57) further, obtaining
$$\mathrm{\Omega }_{tijkl}f_{ijk}=0.$$
(60)
It is sometimes relevant to observe the absence of the display of formulas that might naively be guessed as, e.g., for $`\mathrm{\Omega }_{tijklpq}f_{ijs}`$, cf. (57). This will not reduce simply to a multiple of $`\mathrm{\Omega }_{u[klpq}d_{t]us}`$, since in this case there are other quantities with the required symmetries available to complicate matters. Although a useful reduction can be achieved, the result is not clean enough to be displayed.
Families of more complicated but still useful contractions include the following
$`\mathrm{\Omega }^{(5)}{}_{ijkls}{}^{}\mathrm{\Omega }_{}^{(7)}_{ijklpqr}`$ $`=`$ $`\frac{n}{15}(n^29)f_{u[pq}d_{r]us}`$ (61)
$`\mathrm{\Omega }^{(7)}{}_{ijklpqr}{}^{}\mathrm{\Omega }_{}^{(9)}_{ijklpqstu}`$ $`=`$ $`\frac{2n}{105}(n^29)(n^216)f_{v[st}d_{u]vr}`$ (62)
$`\mathrm{\Omega }^{(5)}{}_{ijkls}{}^{}\mathrm{\Omega }_{}^{(7)}{}_{ijklpqr}{}^{}f_{pqt}^{}`$ $`=`$ $`\frac{2}{45}n^2(n^29)d_{rst}`$ (63)
$`\mathrm{\Omega }^{(7)}{}_{ijklpqr}{}^{}\mathrm{\Omega }_{}^{(9)}{}_{ijklpqstu}{}^{}f_{stv}^{}`$ $`=`$ $`\frac{4}{315}n^2(n^29)(n^216)d_{ruv}`$ (64)
where factors $`(n^29)`$ and $`(n^216)`$ reflect the respective facts that $`\mathrm{\Omega }^{(7)}`$ is absent for $`n=3`$, and $`\mathrm{\Omega }^{(9)}`$ is absent for $`n=4`$. Eq. (61) can be rewritten also as
$$\mathrm{\Omega }^{(5)}{}_{ijkls}{}^{}\mathrm{\Omega }_{}^{(7)}{}_{ijklpqr}{}^{}=\frac{2}{15}(n^29)f_{ijs}\mathrm{\Omega }^{(5)}{}_{ijpqr}{}^{},$$
(65)
which is a recursion relation of sorts, one that be generalised along obvious lines, as indeed (61) itself has been in the production of (62). Eq. (63) is an easy consequence of (61), which also implies
$$\mathrm{\Omega }^{(5)}{}_{ijklp}{}^{}\mathrm{\Omega }_{}^{(7)}{}_{ijklpqr}{}^{}=0.$$
(66)
The same applies to (64) and (62), so that also
$$\mathrm{\Omega }^{(7)}{}_{ijklpqs}{}^{}\mathrm{\Omega }_{}^{(9)}{}_{ijklpqstu}{}^{}=0.$$
(67)
Results such as (60), (66) and (67) are evident enough, since there is no $`SU(n)`$-invariant totally antisymmetric tensor of order two. They point strongly towards an analogue to the orthogonality result for the $`t`$-tensors, eqs. (14)–(16), and suggest that the maximal contraction of two Omega tensors of different order is zero. We do not have a general proof, but as questions regarding it arise it will be shown that this is indeed the case. Thus one might ask about the claim
$$\mathrm{\Omega }^{(7)}{}_{ijklpqr}{}^{}f_{pqr}^{}=0.$$
(68)
Using the methods of this section it is indeed possible, but not easy, to verify this by direct calculation. Alternatively, we may have recourse to the assertion that there exist no $`SU(n)`$-invariant totally antisymmetric tensors of order four. Similarly, eq. (95) below indicates the absence of any $`SU(n)`$-invariant totally antisymmetric tensors of rank six. It follows that we may write
$$\mathrm{\Omega }^{(9)}{}_{ijklpqrst}{}^{}f_{rst}^{}=0,$$
(69)
and similarly
$$\mathrm{\Omega }^{(11)}{}_{ijklpqrstuv}{}^{}\mathrm{\Omega }_{}^{(5)}{}_{ijklp}{}^{}=0.$$
(70)
Such arguments fail for the proof of
$$\mathrm{\Omega }^{(11)}{}_{ijklpqrstuv}{}^{}f_{tuv}^{}=0,$$
(71)
because for $`su(n)`$ with $`n>5`$ (so that $`\mathrm{\Omega }^{(11)}`$ exists) there is a (non-primitive) totally antisymmetric $`SU(n)`$-invariant tensor of order eight. Such matters are discussed systematically in Sec 3.7, the eigth-order tensor being there displayed as (92). Eq. (71) is nevertheless true, although we need the methods of later sections, lambda-matrix methods, to obtain a convenient approach (see Sec. 8) to its proof. Thus, as is known to be true for the $`t`$-tensors, so also it seems for the Omega tensors that the only invariants that can be built out of them are their fully contracted squares. We begin the task of evaluating these in the next paragraph.
### 3.5 Product identities
We begin with the results
$`f_{ijs}f_{ijt}`$ $`=`$ $`\varphi _3(n)\delta _{st},`$ (72)
$`\mathrm{\Omega }_{ijkls}\mathrm{\Omega }_{ijklt}`$ $`=`$ $`\varphi _5(n)\delta _{st},`$ (73)
$`\mathrm{\Omega }_{ijklpqs}\mathrm{\Omega }_{ijklpqt}`$ $`=`$ $`\varphi _7(n)\delta _{st},`$ (74)
that define a family of quantities of which the first few members are
$`\varphi _3(n)`$ $`=`$ $`n,`$ (75)
$`\varphi _5(n)`$ $`=`$ $`\frac{1}{3}n(n^24),`$ (76)
$`\varphi _7(n)`$ $`=`$ $`\frac{2}{15}(n^29)\varphi _5(n).`$ (77)
These imply the consequences
$`f_{ijs}f_{ijs}`$ $`=`$ $`\psi _3(n),`$ (78)
$`\mathrm{\Omega }_{ijkls}\mathrm{\Omega }_{ijkls}`$ $`=`$ $`\psi _5(n),`$ (79)
$`\mathrm{\Omega }_{ijklpqs}\mathrm{\Omega }_{ijklpqs}`$ $`=`$ $`\psi _7(n),`$ (80)
where
$`\psi _3(n)`$ $`=`$ $`n(n^21),`$ (81)
$`\psi _5(n)`$ $`=`$ $`\frac{1}{3}n(n^21)(n^24),`$ (82)
$`\psi _7(n)`$ $`=`$ $`\frac{2}{45}n(n^21)(n^24)(n^29).`$ (83)
One can speculate with confidence on the extension of these results to higher Omega tensors. As is discussed in related contexts in , results collapse to zero equal zero for low values of $`n`$, for which the relevant primitive cocycles do not exist. For example $`su(2)`$ has no cocycle of order higher than three and $`su(3)`$ none higher than five. The right hand sides of (82) and (83) have explicit factors zero for the corresponding $`n`$-values. The proofs of (82) and (83) are given in Sec. 4.
One often needs results more general than those so far mentioned. We have
$`\mathrm{\Omega }_{ijkpq}\mathrm{\Omega }_{ijkrs}`$ $`=`$ $`\frac{1}{6}\left((n^26)f_{pqt}f_{rst}+n(\delta _{pr}\delta _{qs}\delta _{ps}\delta _{qr})\right),`$ (84)
$`\mathrm{\Omega }_{ijklpqr}\mathrm{\Omega }_{ijklpst}`$ $`=`$ $`\frac{2}{135}(n^29)\left((n^28)f_{qru}f_{stu}+2n(\delta _{qs}\delta _{rt}\delta _{qt}\delta _{rs})\right),`$ (85)
which imply (73) and (74), as they should. The right hand sides of (84) and (85) involve linear combination of the only two fourth order tensors with the correct symmetries (so that these are a basis in the vector space in question). The $`\mathrm{\Omega }^{(7)}`$ result collapses for $`n=3`$ ($`su(3)`$ has no seven-cocycle) because of the explicit factor $`n^29`$. Although there is no factor $`(n^24)`$ in (84) causing it to collapse for $`su(2)`$, the right side of (84) nevertheless vanishes because, for $`su(2)`$, we have $`f_{ijk}=ϵ_{ijk}`$, and $`ϵ_{ijk}ϵ_{ipq}=\delta _{jp}\delta _{kq}\delta _{jq}\delta _{kp}`$.
### 3.6 Duality results
We give results here mainly for $`su(3)`$ although in principle the analysis could be extended to higher $`su(n)`$.
We mention first results involving the totally antisymmetric eighth order epsilon tensor:
$`ϵ_{ijklmnpq}ϵ_{ijklmnpq}`$ $`=`$ $`8!,`$
$`ϵ_{ijklmnpq}ϵ_{ijklmnpt}`$ $`=`$ $`7!\delta _{qt},`$
$`ϵ_{ijklmnpq}ϵ_{ijklmnst}`$ $`=`$ $`6!(\delta _{ps}\delta _{qt}\delta _{pt}\delta _{qs}),`$
$`ϵ_{ijklmnpq}ϵ_{ijklmrst}`$ $`=`$ $`5!3!(\delta _n^{[r}\delta _p^s\delta _q^{t]}).`$ (86)
We note here that the factor $`3!`$ is present because the square brackets imply antisymmetrisation with unit weight. Next, from , we note
$`12\sqrt{3}\mathrm{\Omega }_{ijklm}`$ $`=`$ $`ϵ_{ijklmpqr}f_{pqr},`$ (87)
$`20\sqrt{3}f_{stu}`$ $`=`$ $`ϵ_{ijklmstu}\mathrm{\Omega }_{ijlkm}.`$ (88)
Again results more general than the above are often called for, as in . From we quote
$`ϵ_{ijklrstu}\mathrm{\Omega }_{ijlkm}`$ $`=`$ $`16\sqrt{3}\delta _{m[r}f_{stu]},`$ (89)
$`ϵ_{ijklmpqr}f_{qrs}`$ $`=`$ $`24\sqrt{3}\delta _{s[p}\mathrm{\Omega }_{ijklm]}.`$ (90)
One may check that (90) implies (87), and that (89) implies (88). However to prove (89) one must insert (87) and use an identity from the family (86). Similarly insertion of (88) allows proof of (90). Further, one may use (86) to show that (88) follows from (87).
For an $`su(4)`$ result, see Sec. 8 of . Recent work of the authors actually uses duality to obtain information about $`\mathrm{\Omega }^{(9)}`$ for $`su(5)`$, having used MAPLE programs for data about the lower Omega tensors.
### 3.7 Non-primitive antisymmetric tensors
In Sec. 2, we defined for $`su(n)`$ its Omega tensors, which are a set of $`l=(n1)`$ primitive antisymmetric tensors of orders
$$3,5,7,\mathrm{},(2n1).$$
(91)
We have described the fundamental role they play in the discussion of primitive Racah-Casimir operators of $`su(n)`$ (see also ) but they are not the only antisymmetric tensors that can be defined. One can form non-primitive, tilded tensors $`\stackrel{~}{\mathrm{\Omega }}`$, as totally antisymmetrised products of primitive tensors $`\mathrm{\Omega }`$, e.g.
$$\stackrel{~}{\mathrm{\Omega }}^{(8)}{}_{ijkpqrst}{}^{}=\mathrm{\Omega }^{(3)}{}_{[ijk}{}^{}\mathrm{\Omega }_{}^{(5)}{}_{pqrst]}{}^{},$$
(92)
which for $`su(3)`$ is a multiple of the eigth order $`ϵ`$-tensor. In terms of forms, eq. (92) determines a non-primitive de Rham cocycle on the $`SU(n)`$ group manifold (e.g., the volume form on $`SU(3)`$, ignoring factors). A more interesting example arises for $`su(8)`$, $`l=7`$, which has seven Omega tensors of orders $`3,5,7,9,11,13,15`$. In this case, one can form a second, tilded antisymmetric tensor of order $`15`$
$$\stackrel{~}{\mathrm{\Omega }}^{(15)}{}_{i_1\mathrm{}i_{15}}{}^{}=\mathrm{\Omega }^{(3)}{}_{[i_1i_2i_3}{}^{}\mathrm{\Omega }_{}^{(5)}{}_{i_4\mathrm{}i_8}{}^{}\mathrm{\Omega }_{}^{(7)}{}_{i_9\mathrm{}i_{15}]}{}^{},$$
(93)
which is non-primitive and not maximal on the $`(n^21)=63`$-dimensional space (manifold in the case of forms on $`SU(n)`$). The discussion of Sec. 5.1 suggests that it should have zero full contraction with $`\mathrm{\Omega }^{(15)}_{i_1\mathrm{}i_{15}}`$, in virtue of results like (68). This example illustrates a significant restriction on the construction of non-trivial non-primitive antisymmetric tensors: all the primitive Omega tensors used must have different orders. To see this, notice that we may write, e.g.,
$$\mathrm{\Omega }^{(3)}=\mathrm{\Omega }^{(3)}{}_{ijk}{}^{}\omega _{}^{i}\omega ^j\omega ^k,$$
(94)
where the $`\omega ^i`$ are the left-invariant (LI) Maurer-Cartan one-forms on the $`SU(n)`$ groups manoifold, dual to the LI $`su(n)`$ generators, so that $`\mathrm{\Omega }^{(3)}`$ is the invariant de Rham three-cocycle of coordinates $`\mathrm{\Omega }^{(3)}_{ijk}`$. Obviously, $`\mathrm{\Omega }^{(3)}\mathrm{\Omega }^{(3)}=0`$ and hence
$$\mathrm{\Omega }^{(3)}{}_{[ijk}{}^{}\mathrm{\Omega }_{}^{(3)}{}_{pqr]}{}^{}=f_{[ijk}f_{pqr]}=0.$$
(95)
In general, the skewsymmetrisation of two copies of the same Omega tensor is zero since this corresponds to taking the wedge product of a primitive $`SU(n)`$ de Rham cocycle by itself, which is zero because all these cocyles are represented by odd, ($`2m1`$)-forms on the $`SU(n)`$ group manifold.
### 3.8 9-cocycle results
We used the ninth order Omega tensor to define the fifth order $`t`$-tensor $`t^{(5)}`$, quoting the result for it as (20). This enables us to calculate
$$\mathrm{\Omega }^{(9)}{}_{ijklmpqrs}{}^{}\mathrm{\Omega }_{}^{(9)}{}_{ijklmpqrs}{}^{},$$
(96)
to within an overall normalisation constant. We may use (16) to derive
$`\mathrm{\Omega }_{ijklmpqrs}\mathrm{\Omega }_{ijklmpqrs}`$ $`=`$ $`\mathrm{\Omega }_{ijklmpqrs}f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{pq}{}^{}f_{}^{d}{}_{rs]}{}^{}d_{}^{(5)}_{(abcdm)}`$ (97)
$`=`$ $`t^{(5)}{}_{abcdm}{}^{}d_{}^{(5)}_{(abcdm)}`$
$`=`$ $`\lambda (n){\displaystyle \underset{r=1}{\overset{4}{}}}(n^2r^2).`$
Here the first line uses the definition (12) and the second one uses (16). The last line may be evaluated from the second one using (20) and the two following results, the first of which is the last equation of the appendix in , and the other is much more easily obtained:
$`d^{(5)}{}_{(abcdm)}{}^{}d_{}^{(5)}_{(abcdm)}`$ $`=`$ $`\frac{1}{15n^3}(n^21)(n^24)(5n^496n^2+480),`$ (98)
$`d^{(5)}{}_{(abcdm)}{}^{}\delta _{(ab}^{}d_{cdm)}`$ $`=`$ $`\frac{1}{5n^2}(n^21)(n^24)(3n^220).`$ (99)
Equation (97) already displays the essential factors anticipated in Sec. 3.5. We confirm its correctness in Sec. 6, using lambda-matrix techniques which allow the factor $`\lambda (n)`$ to be determined, the result having been given above as (21).
## 4 Selected Derivations
### 4.1 Equations (84), (73) and (79)
The first result involving $`\mathrm{\Omega }^{(5)}`$ that is not straightforward to derive from definitions is (84). We develop
$`\mathrm{\Omega }_{ijklm}\mathrm{\Omega }_{ijkpq}`$ $`=`$ $`\mathrm{\Omega }_{mijkl}\mathrm{\Omega }_{qijkp}`$ (100)
$`=`$ $`f^a{}_{m[i}{}^{}f_{}^{b}{}_{jk]}{}^{}d_{abl}^{}f^x{}_{q[i}{}^{}f_{}^{y}{}_{jk]}{}^{}d_{xyp}^{}.`$
The second set of square brackets can now be dropped, leaving behind the sum of three sixth order products of $`d`$\- and $`f`$-tensors to be reduced by trace methods using formulas from the appendix of . Two of the three terms coincide, and nothing more than the evaluation of four-fold traces needs to be done. A rough graphical representation of any term helps (here and elsewhere) to see the best way to use trace formulas. In this vertices correspond in evident manner to $`d`$\- and $`f`$-tensors, while closed loops indicate the traces. The result comes out initially in terms of products of $`\delta \delta `$ and $`dd`$ terms. But there is an identity valid for all $`su(n)`$ (, eq. (2.10)) which allows the latter to be given in terms of $`\delta \delta `$ and $`ff`$ terms as displayed.
It is not difficult to reduce (84) to confirm the correctness of (73)–(77) and hence of (79)–(83). However, a direct attack on either of the latter along the lines just indicated is a good way to get up to speed on methods useful in the current study.
### 4.2 Equations (85), (74) and (80)
Eq. (85) perhaps discourages such a direct approach as Sec. 4.1 uses, so one adopts a different approach. This requires, as a preliminary, the knowledge of (80). Hence we first develop
$`\mathrm{\Omega }_{ijklmpq}\mathrm{\Omega }_{ijklmpq}`$ $`=`$ $`\mathrm{\Omega }_{ijklmpq}f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{mp]}{}^{}d_{}^{(4)}_{(abcq)}`$ (101)
$`=`$ $`t^{(4)}{}_{abcq}{}^{}d_{}^{(4)}_{(abcq)}`$
$`=`$ $`t^{(4)}{}_{abcq}{}^{}d_{tab}^{}d_{cqt}`$
$`=`$ $`\frac{2}{45}n^2(n^29)d_{cqt}d_{cqt}.`$
The first line uses the definition (12), from which the square brackets can be dropped, so that (15) can be used in the second line. Next the symmetry properties of the $`t`$-tensors allow the replacement of $`d^{(4)}`$ by one of its terms, whereupon (24) may be used. Since
$$d_{abc}d_{abd}=\frac{(n^24)}{n}\delta _{cd},$$
(102)
eq. (80) follows.
Returning now to (85), we use that the two terms on the right side of (84) are a basis for the vector space of tensors with the the required symmetry properties. Hence we seek a result of the form
$$\mathrm{\Omega }_{ijklpqr}\mathrm{\Omega }_{ijklpst}=b(n)f_{qru}f_{stu}+a(n)(\delta _{qs}\delta _{rt}\delta _{qt}\delta _{rs}).$$
(103)
To determine the coefficients, we must perform contractions with $`\delta _{qs}`$ and with $`f_{stg}`$. As we show below, this gives equations
$`(n^22)a(n)+b(n)n`$ $`=`$ $`\varphi _7(n)`$
$`2a(n)+b(n)n`$ $`=`$ $`\frac{1}{3}\varphi _7(n),`$ (104)
which can be solved to complete the derivation of (85). It is easy to get all the terms here except the one on the right side of the second equation. A viable starting point is elusive. Consider therefore
$$\mathrm{\Omega }_{ijklpst}f_{gst}f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{p]q}{}^{}d_{}^{(4)}{}_{(abcr)}{}^{}.$$
(105)
The square brackets can again be dropped. This lead us to
$$f_{cpq}d^{(4)}{}_{(abcr)}{}^{}t_{}^{(4)}_{abgp}$$
(106)
using (15). Now, we need to check
$$d^{(4)}{}_{(abcr)}{}^{}\delta _{(ab}^{}\delta _{gp)}=\frac{2}{9}\frac{(n^24)}{n}\delta _{gp}\delta _{cr},$$
(107)
use (23), and complete the computation
$$f_{cpq}d^{(4)}{}_{(abcr)}{}^{}d_{age}^{}d_{bpe}=\frac{1}{9n^2}(n^24)(n^28)f_{qrg}.$$
(108)
This last one requires use of the four-fold $`d`$-tensor trace of (A.10) of , some simpler results also found there and some patience. Then all the pieces of the calculation have to be put together to complete the derivation of (85). One is guided through something of a morass to an answer one knows is right by the fact that both the unknowns $`a(n)`$ and $`b(n)`$ must contain a factor $`(n^29)`$ that vanishes at $`n=3`$.
## 5 Omega tensors and $`su(n)`$ lambda-matrices
### 5.1 Antisymmetrised products of $`su(n)`$ lambda-matrices
We use the lambda-matrices of ref. which are subject to
$$\mathrm{Tr}\lambda _i=0,\mathrm{Tr}\lambda _i\lambda _j=2\delta _{ij},\lambda _{i}^{}{}_{}{}^{}=\lambda _i,$$
(109)
$$\lambda _i\lambda _j=\frac{2}{n}\delta _{ij}+(d+if)_{ijk}\lambda _k,d_{ijk}\delta _{ij}=0,$$
(110)
and define totally antisymmetrised products of unit weight of lambda-matrices
$$\lambda _{[ijk\mathrm{}s]}=\lambda _{[i}\lambda _j\lambda _k\mathrm{}\lambda _{s]}.$$
(111)
Simple computations using (9) and (36) lead directly to
$`\lambda _{[ijk]}`$ $`=`$ $`\frac{2}{n}i(f_{ijk}+\mathrm{\Omega }_{ijkab}f_{abq}\lambda _q)=\frac{2}{n}if_{ijk}+if_{s[ij}d_{k]sq}\lambda _q,`$ (112)
$`\lambda _{[ijkl]}`$ $`=`$ $`\mathrm{\Omega }_{ijklt}\lambda _t.`$ (113)
These imply the trace results
$`\mathrm{Tr}\lambda _{[ijkl]}`$ $`=`$ $`0,`$ (114)
$`\mathrm{Tr}\lambda _{[ijk]l}`$ $`=`$ $`\frac{4}{n}i\mathrm{\Omega }_{ijkab}f_{abl},`$ (115)
which may be contrasted with
$`\mathrm{Tr}\lambda _{[ijk]}`$ $`=`$ $`\mathrm{Tr}\lambda _{i[jk]}=2if_{ijk},`$ (116)
$`\mathrm{Tr}\lambda _{[ijkpq]}`$ $`=`$ $`\mathrm{Tr}\lambda _{[ijkp]q}=2\mathrm{\Omega }_{ijkpq},`$ (117)
where the first two equalities use the cyclic nature of the trace. For odd traces as (116), (117) terms related by cyclicity add up, whereas for even ones they cancel pairwise and indeed completely:
$$\mathrm{Tr}\lambda _{[i_1i_2\mathrm{}i_{2s}]}=0.$$
(118)
For the five-fold case one finds using $`\lambda _{[ijkpq]}=\lambda _{[[ijkp]q]}`$ and (38) that
$$\lambda _{[ijkpq]}=\frac{2}{n}\mathrm{\Omega }_{ijkpq}\mathrm{\Omega }_{r[ijkp}d_{q]rt}\lambda _t,$$
(119)
from which (117) can be recovered. For the six-fold antisymmetrised product of $`\lambda `$’s we use $`\lambda _{[ijkpqr]}=\lambda _{[[ijkp][qr]]}`$ to deduce
$`\lambda _{[ijkpqr]}`$ $`=`$ $`i\mathrm{\Omega }_{s[ijkp}f_{qr]t}\lambda _s\lambda _t`$ (120)
$`=`$ $`i\mathrm{\Omega }_{ijkpqrs}\lambda _s,`$ (121)
and
$$\mathrm{Tr}\lambda _{[ijkpq]s}=2\mathrm{\Omega }_{t[ijkp}d_{q]st}.$$
(122)
Here, to derive (121), we used (110) for $`\lambda _s\lambda _t`$, $`Ad`$-invariance (eq. (37)) to discard the first term, and the following steps to discard another term,
$`\mathrm{\Omega }_{s[ijkp}f_{qr]t}f_{stm}`$ $`=`$ $`f^x{}_{[ij}{}^{}f_{}^{y}{}_{kp}{}^{}f_{}^{t}{}_{qr]}{}^{}d_{xys}^{}f_{stu}`$ (123)
$`=`$ $`f^x{}_{[ij}{}^{}f_{}^{y}{}_{kp}{}^{}f_{}^{t}{}_{qr]}{}^{}d_{s(xy}^{}f_{t)ms}=0,`$
upon using (26). Similar steps, using (7) and (10), show how the $`d`$-term of (110) features in the production of (121).
For the seven-fold product we find
$$\lambda _{[ijklpqr]}=\frac{2}{n}i\mathrm{\Omega }^{(7)}{}_{ijklpqr}{}^{}i\mathrm{\Omega }^{(7)}{}_{s[ijklpq}{}^{}d_{r]ts}^{}\lambda _t,$$
(124)
with the aid of (38) and
$$\mathrm{Tr}\lambda _{[ijkpqrt]}=2i\mathrm{\Omega }^{(7)}{}_{ijkpqrt}{}^{}.$$
(125)
Eq. (124) is a natural generalisation of the odd traces (112) and (119), and we may infer the result for the odd case
$$\lambda _{[i_1i_2\mathrm{}i_{2s}k]}=\frac{2}{n}i^s\mathrm{\Omega }^{(2s+1)}{}_{i_1i_2\mathrm{}i_{2s}k}{}^{}+i^s\mathrm{\Omega }^{(2s+1)}{}_{p[i_1i_2\mathrm{}i_{2s}}{}^{}d_{k]pq}^{}\lambda _q.$$
(126)
Also we may use (125) and
$$\mathrm{Tr}\lambda _{[ijkpqrt]}=\mathrm{Tr}\lambda _{[ijkpqr]t}=\mathrm{Tr}\lambda _{[[ijkp][qr]]t},$$
(127)
to check our work by reproducing the recursive identity (51).
Writing the elementary result
$$\lambda _{[ij]}=\frac{1}{2}[\lambda _i,\lambda _j]=if_{ijk}\lambda _k,$$
(128)
and comparing it also with the even case (113), and (121), one gets for the antisymmetrised product of an even number of $`\lambda `$’s the result
$$\lambda _{[i_1\mathrm{}i_{2s}]}=i^s\mathrm{\Omega }^{(2s+1)}{}_{i_1\mathrm{}i_{2s}k}{}^{}\lambda _{k}^{},$$
(129)
which implies (118) and
$$\mathrm{Tr}\lambda _{[i_1\mathrm{}i_{2s}]k}=2i^s\mathrm{\Omega }^{(2s+1)}{}_{i_1\mathrm{}i_{2s}k}{}^{}.$$
(130)
We note that (129), and in particular (113) and (121), provide an explicit realisation of the ($`2m2`$)-bracket Lie algebras for $`su(n)`$. As mentioned, the coordinates of the $`\mathrm{\Omega }^{(2m1)}`$ determine the associated higher order structure constants (above, $`s=m1`$), and satisfy the GJI (42).
### 5.2 General proofs of results for the $`\lambda _{[i_1\mathrm{}i_s]}`$
The above results are general, due to the nature of the $`\mathrm{\Omega }^{(2m1)}`$ tensors as generalised structure constants (for instance, eq. (129) may be looked at as a consequence of Th. 3.1 in ). However, the above eqs., and in particular (126), were presented on the basis of inspection of a modest number of low value cases. It is thus necessary to show their general validity, particularly since, as defined above by eqs. (10)–(11), $`\mathrm{\Omega }^{(5)}`$ and $`\mathrm{\Omega }^{(7)}`$ involve $`d`$-tensors with simple properties that do not generalise straightforwardly to the $`d`$-tensors involved in the definition of higher $`\mathrm{\Omega }`$-tensors.
Let us look first at (129). We write
$$\lambda _{[i_1j_1\mathrm{}i_sj_s]}=i^sf^{p_1}{}_{[i_1j_1}{}^{}\mathrm{}f^{p_s}{}_{i_sj_s]}{}^{}\lambda _{(p_1\mathrm{}p_s)}^{}.$$
(131)
If we apply (110) repeatedly, making full use of the symmetry properties that are implied by the round brackets, we can establish a result of the form
$$\lambda _{(p_1\mathrm{}p_s)}=\stackrel{~}{k}{}_{(p_1\mathrm{}p_s)t}{}^{}\lambda _{t}^{}+k{}_{(p_1\mathrm{}p_st)}{}^{},$$
(132)
where the $`k`$-tensors are $`Ad`$-invariant tensors with the indicated symmetries. Eq. (48) tells us that $`k`$ does not contribute to (131). Also $`\stackrel{~}{k}_{(p_1\mathrm{}p_s)t}`$ differs from $`d^{(s+1)}_{(p_1\mathrm{}p_s)t}`$ only by some linear combination of non-primitive terms, which, also by (48), do not contribute to (131). Further
$$\stackrel{~}{k}{}_{(p_1\mathrm{}p_s)t}{}^{}=\frac{1}{2}\mathrm{Tr}\lambda _{(p_1\mathrm{}p_s)t}=\frac{1}{2}\mathrm{Tr}\lambda _{(p_1\mathrm{}p_st)}=\stackrel{~}{k}{}_{(p_1\mathrm{}p_st)}{}^{},$$
(133)
all of which allows us to replace the lambda-matrix factor of (131) by $`d^{(s+1)}{}_{(p_1\mathrm{}p_st)}{}^{}\lambda _{t}^{}`$, so that also
$$\lambda _{[i_1j_1\mathrm{}i_sj_s]}=\mathrm{\Omega }^{(2s+1)}{}_{i_1j_1\mathrm{}i_sj_st}{}^{}\lambda _{t}^{}.$$
(134)
The trace of (134) now confirms (118).
Eq. (130) also follows easily. To obtain (126), we multiply (134) by $`\lambda _k`$ and use (110). Then (126) follows directly, after the use of the $`Ad`$-invariance of $`\mathrm{\Omega }^{(2s+1)}`$ to drop the contribution of the $`f`$-term of (110).
Inspection of (133) shows that it is tantamount to the statement that, in the definition (11) of $`\mathrm{\Omega }^{(9)}`$, e.g., one is, after all, allowed to move the right hand round bracket one place to the left. It is of interest to see this explicitly, because, amongst other things, a further class of identities for $`d`$-tensors emerges as a by-product. We illustrate this for $`s=4`$. One evaluation of the trace involved leads to
$$\lambda _{(abcd)}=\frac{4}{n^2}\delta _{(ab}\delta _{cd)}+\frac{2}{n}d^{(4)}{}_{(abcd)}{}^{}+\frac{2}{n}d_{(abc}\lambda _{d)}+\frac{2}{n}\delta _{(ab}d_{cd)y}\lambda _y+d^{(5)}{}_{(abcd)y}{}^{}\lambda _{y}^{}.$$
(135)
The key trace result (cf. (133))
$$\mathrm{Tr}\lambda _{(abcd)e}=\mathrm{Tr}\lambda _{(abcde)},$$
(136)
now leads to
$$d^{(5)}{}_{(abcde)}{}^{}=d^{(5)}{}_{(abcd)e}{}^{}+\frac{1}{n}\delta _{(ab}d_{cde)}\frac{1}{n}\delta _{(ab}d_{cd)e}.$$
(137)
The difference between the two $`d^{(5)}`$ tensors here, and as in the general discussion above, makes no contribution to the evaluation of $`\lambda _{[i_1j_1\mathrm{}i_4j_4]}`$. It follows then that in the definition (11), we can replace $`d^{(5)}_{(abcde)}`$ by $`d^{(5)}_{(abcd)e}`$, which of course has fewer terms.
Another question arises here: how does (136) relate to (30)? To answer, we note that a different way of evaluating the trace gives rise to
$$\lambda _{(abcd)}=\frac{4}{n^2}\delta _{(ab}\delta _{cd)}+\frac{2}{n}d^{(4)}{}_{(abcd)}{}^{}+\frac{4}{n}\delta _{(ab}d_{cd)y}\lambda _y+d^{(5)}{}_{(ab}{}^{y}{}_{cd)}{}^{}\lambda _{y}^{},$$
(138)
and hence
$$d^{(5)}{}_{(abcde)}{}^{}=d^{(5)}{}_{(ab}{}^{e}{}_{cd)}{}^{}\frac{4}{n}\delta _{(ab}d_{cde)}+\frac{4}{n}\delta _{(ab}d_{cd)e}.$$
(139)
Now (30) follows obviously form (135) and (139).
There is another instructive way to make the point that the three $`d^{(5)}`$ tensors can be used interchangeably in the definition (11) of $`\mathrm{\Omega }^{(9)}`$. It follows by comparison of
$`\lambda _{[ijklpqrs]}`$ $`=`$ $`\lambda _{[[ijklpq][rs]]}`$ (140)
$`=`$ $`\mathrm{\Omega }_{x[ijklpq}f_{rs]y}\lambda _x\lambda _y`$
$`=`$ $`f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{pq}{}^{}f_{}^{y}{}_{rs]}{}^{}d_{}^{(5)}{}_{(abcy)t}{}^{}\lambda _t,`$
and
$`\lambda _{[ijklpqrs]}`$ $`=`$ $`\lambda _{[[ijk][lpqrs]]}`$ (141)
$`=`$ $`f^a{}_{[ij}{}^{}f_{}^{b}{}_{kl}{}^{}f_{}^{c}{}_{pq}{}^{}f_{}^{y}{}_{rs]}{}^{}d_{}^{(5)}{}_{(ab}{}^{t}{}_{cy)}{}^{}\lambda _{t}^{}.`$
The discussion just given for $`\mathrm{\Omega }^{(9)}`$ generalises naturally for higher Omega tensors.
### 5.3 Use of the completeness relation for the $`su(n)`$ lambda-matrices
We set out from the result well-known for $`su(n)`$
$$\lambda _{iab}\lambda _{icd}=2\delta _{ad}\delta _{cb}\frac{2}{n}\delta _{ab}\delta _{cd},$$
(142)
and note also its consequences
$`if_{ijk}\lambda _{jab}\lambda _{kcd}`$ $`=`$ $`\lambda _{iad}\delta _{cb}\lambda _{icb}\delta _{ad}`$ (143)
$`d_{ijk}\lambda _{jab}\lambda _{kcd}`$ $`=`$ $`\lambda _{iad}\delta _{bc}+\lambda _{icb}\delta _{ad}\frac{2}{n}\left(\lambda _{iab}\delta _{cd}+\lambda _{icd}\delta _{ab}\right).`$ (144)
From (142) we may compute
$`\lambda _{[ij]ab}\lambda _{[ij]cd}`$ $`=`$ $`n\lambda _{iab}\lambda _{icd}`$ (145)
$`\lambda _{[ijk]ab}\lambda _{[ijk]cd}`$ $`=`$ $`\frac{2}{3}(n^24)\lambda _{iab}\lambda _{icd}\frac{4}{n}(n^21)\delta _{ab}\delta _{cd}`$ (146)
$`\lambda _{[ijkl]ab}\lambda _{[ijkl]cd}`$ $`=`$ $`\frac{n}{3}(n^24)\lambda _{iab}\lambda _{icd}`$ (147)
$`\lambda _{[ijklm]ab}\lambda _{[ijklm]cd}`$ $`=`$ $`\frac{2}{15}(n^24)(n^26)\lambda _{iab}\lambda _{icd}+\frac{4}{3n}(n^21)(n^24)\delta _{ab}\delta _{cd},`$ (148)
and so on. One can make checks on these results by putting $`b=c`$ to reach
$`\lambda _i\lambda _i`$ $`=`$ $`\frac{2}{n}(n^21)I`$ (149)
$`\lambda _{[ij]}\lambda _{[ij]}`$ $`=`$ $`2(n^21)I`$ (150)
$`\lambda _{[ijk]}\lambda _{[ijk]}`$ $`=`$ $`\frac{4}{3n}(n^21)^2I`$ (151)
$`\lambda _{[ijkl]}\lambda _{[ijkl]}`$ $`=`$ $`\frac{2}{3}(n^21)(n^24)I`$ (152)
$`\lambda _{[ijklm]}\lambda _{[ijklm]}`$ $`=`$ $`\frac{4}{15n}(n^21)^2(n^24)I.`$ (153)
These results are of use themselves and may be verified by other means. It is tempting to speculate on the nature of results beyond (148), but it gets increasingly hard to compute directly the $`n`$-dependences. Use of $`\mathrm{Tr}I=n`$ yields obvious trace formulas.
For the purpose, central to the aims of this paper, of computing the quantities $`(\mathrm{\Omega }^{(2m+1)})^2`$ explicitly in closed form, it is enough to analyse the traced analogues of (145)–(148), obtained by putting $`b=c`$ and $`d=a`$. For this analysis, one of the approaches available employs another set of lemmas that follow from (142). For any $`n`$-dimensional matrix $`M`$, eq. (142) gives
$$(\lambda _iM\lambda _i)_{ab}=2\delta _{ab}\mathrm{Tr}M\frac{2}{n}M_{ab}.$$
(154)
This provides us with a method for obtaining the results
$`\lambda _i\lambda _j\lambda _i`$ $`=`$ $`\frac{2}{n}\lambda _j`$ (155)
$`\lambda _i\lambda _{[jk]}\lambda _i`$ $`=`$ $`\frac{2}{n}\lambda _{[jk]}`$ (156)
$`\lambda _i\lambda _{[jkl]}\lambda _i`$ $`=`$ $`4if_{jkl}\frac{2}{n}\lambda _{[jkl]}`$ (157)
$`\lambda _i\lambda _{[jklpq]}\lambda _i`$ $`=`$ $`4\mathrm{\Omega }_{jklpq}\frac{2}{n}\lambda _{[jklpq]}`$ (158)
$`\lambda _i\lambda _{[i_1j_1\mathrm{}i_mj_mk]}\lambda _i`$ $`=`$ $`4i^m\mathrm{\Omega }_{}^{(2m+1)}{}_{i_1j_1\mathrm{}i_mj_mk}{}^{}\frac{2}{n}\lambda _{[i_1j_1\mathrm{}i_mj_mk]}`$ (159)
$`\lambda _i\lambda _{[i_1j_1\mathrm{}i_mj_m]}\lambda _i`$ $`=`$ $`\frac{2}{n}\lambda _{[i_1j_1\mathrm{}i_mj_m]}.`$ (160)
The last result follows from (154) because of (118).
We note here further simple results that may help streamline larger tasks, for example one approach to the proof of (149)–(153):
$`\lambda _i\lambda _{[ij]}`$ $`=`$ $`n\lambda _j`$ (161)
$`\lambda _i\lambda _{[ijk]}`$ $`=`$ $`\frac{2}{3}\frac{(n^21)}{n}\lambda _{[jk]}`$ (162)
$`\lambda _{[ij]}\lambda _{[ijk]}`$ $`=`$ $`\frac{2}{3}(n^21)\lambda _k`$ (163)
$`\lambda _i\lambda _{[ijkl]}`$ $`=`$ $`\frac{n}{2}\lambda _{[jkl]}if_{jkl}`$ (164)
$`\lambda _{[ij]}\lambda _{[ijkl]}`$ $`=`$ $`\frac{1}{3}(n^24)\lambda _{[kl]}`$ (165)
$`\lambda _{[ijk]}\lambda _{[ijkl]}`$ $`=`$ $`\frac{n}{3}(n^24)\lambda _l`$ (166)
$`\lambda _i\lambda _{[ijklm]}`$ $`=`$ $`\frac{2}{5n}(n^21)\lambda _{[jklm]}\frac{4}{5}if_{[jkl}\lambda _{m]}`$ (167)
$`\lambda _{[ij]}\lambda _{[ijklm]}`$ $`=`$ $`\frac{1}{5}(n^24)\lambda _{[klm]}`$ (168)
$`\lambda _{[ijk]}\lambda _{[ijklm]}`$ $`=`$ $`\frac{2}{15n}(n^21)(n^24)\lambda _{[lm]}`$ (169)
$`\lambda _{[ijkl]}\lambda _{[ijklm]}`$ $`=`$ $`\frac{2}{15}(n^21)(n^24)\lambda _m.`$ (170)
Hermitian conjugation gives results such as
$$\lambda _{[ji]}\lambda _i=n\lambda _j.$$
(171)
Inspection of (161), (162), (164) and (167) suggests the general result
$$\mathrm{Tr}\lambda _i\lambda _{[ii_2\mathrm{}i_s]}=0(s\mathrm{even}\mathrm{or}\mathrm{odd}),$$
(172)
which is easily proved using the results of Sec. 5.2. If $`s`$ is even we find, using (129), $`\mathrm{Tr}(\lambda _i\lambda _{[ii_2\mathrm{}i_s]})\mathrm{Tr}(\lambda _i\mathrm{\Omega }_{ii_2\mathrm{}i_sk}\lambda _k)=2\mathrm{\Omega }_{ii_2\mathrm{}i_si}=0`$. If $`s`$ is odd, $`\mathrm{Tr}(\lambda _i\lambda _{[ii_2\mathrm{}i_s]})=\mathrm{Tr}(\lambda _i\lambda _{[i[i_2\mathrm{}i_s]]})\mathrm{Tr}(\lambda _i\lambda _{[i}\lambda _{k]})\mathrm{\Omega }_{i_2\mathrm{}i_sk}=0`$.
Some of the principal results to be derived below require as input more trace results. First, and in agreement with results displayed above, we expect
$$\mathrm{Tr}(\lambda _{[ij]}\lambda _{[iji_3\mathrm{}i_{2s}]})=0.$$
(173)
A typical proof here, using the methods of Sec. 5.2, is
$$\mathrm{Tr}(\lambda _{[ij]}\lambda _{[ijklpq]})=\mathrm{Tr}(\lambda _{[ij]}(i\mathrm{\Omega }_{ijklpqr}\lambda _r))=2\mathrm{\Omega }_{ijklpqr}f_{ijr}=0,$$
(174)
upon use of (68). The analogues to (173) for $`s=4`$ and $`s=5`$ of (173) however depend on (69) and (71), results which remain unproved until the methods of Sec. 8 can be called upon. Second
$$\mathrm{Tr}(\lambda _{[ijk]}\lambda _{[ijki_4\mathrm{}i_{2s}]})=0.$$
(175)
A typical proof, here for $`s=3`$, is
$`\mathrm{Tr}(\lambda _{[ijk]}\lambda _{[ijklpq]})`$ $`=`$ $`\mathrm{Tr}(\lambda _{[ijk]}(i\mathrm{\Omega }_{ijklpqr}\lambda _r))`$ (176)
$`=`$ $`\mathrm{\Omega }_{ijklpqr}\frac{4}{n}\mathrm{\Omega }_{ijkab}f_{abr}`$
$`=`$ $`2\mathrm{\Omega }_{ijklpqr}f_{u[ij}d_{k]ur}=0,`$
where (121), (115) and (57) have been used. The last equality follows from the total symmetry of $`d^{(3)}`$ and the antisymmentry of $`\mathrm{\Omega }^{(7)}`$.
### 5.4 Further trace results
We are interested here in trace results of the type
$$\mathrm{Tr}(\lambda _{[ij\mathrm{}s]}\lambda _{[ij\mathrm{}s]}).$$
(177)
Even traces of this sort are of primary interest in virtue of their relationship to $`(\mathrm{\Omega }^{(2s+1)})^2`$. For example, for $`s=2`$,
$`4(\mathrm{\Omega }^{(5)})^2`$ $`=`$ $`\mathrm{Tr}\lambda _{[ijkpq]}\mathrm{Tr}\lambda _{[ijkpq]}`$ (178)
$`=`$ $`\mathrm{Tr}\lambda _i\lambda _{[jkpq]}\mathrm{Tr}\lambda _i\lambda _{[jkpq]}`$
$`=`$ $`2\mathrm{T}\mathrm{r}\left(\lambda _{[jkpq]}\lambda _{[jkpq]}\right),`$
where (142) and (114) have been used. However, proceeding recursively for higher $`s`$ brings the odd traces into the picture. We have two approaches to either even or odd traces, and one works better for the odd and the other for the even traces. We begin with the even traces for which we have a nice general result. We note a generalisation of results embedded in the previous subsections:
$$\mathrm{Tr}(\lambda _{[i_1\mathrm{}i_{2s}k]}\lambda _{[i_1\mathrm{}i_{2s}k]})=\frac{2}{(2s+1)}\frac{(n^21)}{n}\mathrm{Tr}(\lambda _{[i_1\mathrm{}i_{2s}]}\lambda _{[i_1\mathrm{}i_{2s}]}).$$
(179)
A brief look at the case $`s=3`$ will indicate clearly that this result is valid in general. Thus we write
$`\mathrm{Tr}(\lambda _{[ijklpqr]}\lambda _i\lambda _j\lambda _k\lambda _l\lambda _p\lambda _q\lambda _r)`$ $`=`$ $`\frac{1}{7}\mathrm{Tr}(\lambda _i\lambda _{[jklpqr]}\lambda _i\lambda _{[jklpqr]})`$ (180)
$`+`$ $`\frac{1}{7}\mathrm{Tr}(\lambda _j\lambda _{[klpqri]}\lambda _i\lambda _j\lambda _{[klpqr]})+\mathrm{}.`$
Now we use the cyclic property of the trace to justify the use of results of the type (159) and (160) in the first six terms, and of (149) to the seventh. One can see from (118) the Omega tensor terms of (159) do not contribute (which is why this approach is better for the odd traces than for the even ones), and then it is easy to see that, after taking due care of the signs of the first six terms, everything cancels except the contribution of the seventh term of (180), which gives the right side of (179) at $`s=3`$.
Reduction of the right side of (179) is much harder because the same approach brings in the Omega tensor pieces of (157), (159) etc., non-trivially. This caused us to adopt a related but distinct approach to such traces in the next section, although the approach just followed does work, but rather less well. To say enough to allow a comparison of methods to be made, let us refer back to (178). We may drop the second set of square brackets and reinsert others judiciously in suitable places whenever this is allowed by existing antisymmetries. Then the development of the first set of square brackets yields
$`8(\mathrm{\Omega }^{(5)})^2`$ $`=`$ $`\mathrm{Tr}(\lambda _j\lambda _{[kpq]}\lambda _j\lambda _{[kpq]})`$ (181)
$``$ $`\mathrm{Tr}(\lambda _k\lambda _{[pqj]}\lambda _j\lambda _k\lambda _{[pq]})`$
$`+`$ $`\mathrm{Tr}(\lambda _p\lambda _{[qjk]}\lambda _j\lambda _k\lambda _p\lambda _q)`$
$``$ $`\mathrm{Tr}(\lambda _q\lambda _{[jkp]}\lambda _j\lambda _k\lambda _p\lambda _q).`$
The cyclic property of the trace now allows use of (157),(156), (155) and (149), in that order so that after cancellations, we obtain
$$8(\mathrm{\Omega }^{(5)})^2=\mathrm{Tr}\lambda _{[kpq]}4if_{kpq}+\left(\frac{2}{n}+\frac{2}{n}(n^21)\right)\mathrm{Tr}(\lambda _{[kpq]}\lambda _k\lambda _p\lambda _q).$$
(182)
Now use of (116) and (151) leads directly to the answer obtained before: (79) with (82). The higher order even traces get successively harder in this approach, but we will see a comparable increase in the price associated with passing to higher $`s`$ is present also in our favoured method of Sec. 6.
## 6 The recursion relations for the $`(\mathrm{\Omega }^{(2m1)})^2`$
We illustrate the general approach by reference to the case $`m=5`$. Since, by eq. (130),
$$2\mathrm{\Omega }_{ijklpqrst}=\mathrm{Tr}\lambda _{[ijklpqrs]t},$$
(183)
we may write
$`4(\mathrm{\Omega }^{(9)})^2=4\mathrm{\Omega }_{ijklpqrst}\mathrm{\Omega }_{ijklpqrst}`$ $`=`$ $`\mathrm{Tr}\left(\lambda _{[ijklpqrs]}\lambda _t\right)\mathrm{Tr}\left(\lambda _{[ijklpqrs]}\lambda _t\right)`$ (184)
$`=`$ $`2\mathrm{T}\mathrm{r}\left(\lambda _{[ijklpqrs]}\lambda _{[ijklpqrs]}\right).`$
Here we have used (142) and the trace result (118). The key steps now follow. We can remove the first set of square brackets completely and then reinsert them round the indices $`jklpqrs`$. Then we expand the second set of square brackets to expose, in each of the eight terms that thereby arise, the matrix $`\lambda _i`$:
$`16(\mathrm{\Omega }^{(9)})^2=\lambda _{iab}\lambda _{[jklpqrs]bc}(\delta _{cd}\lambda _{ide}\lambda _{[jklpqrs]ea}`$ $``$ $`\lambda _{scd}\lambda _{ide}\lambda _{[jklpqr]ea}`$
$`+\lambda _{[rs]cd}\lambda _{ide}\lambda _{[jklpq]ea}`$ $``$ $`\lambda _{[qrs]cd}\lambda _{ide}\lambda _{[jklp]ea}`$
$`+\lambda _{[pqrs]cd}\lambda _{ide}\lambda _{[jkl]ea}`$ $``$ $`\lambda _{[lpqrs]cd}\lambda _{ide}\lambda _{[jk]ea}`$
$`+\lambda _{[klpqrs]cd}\lambda _{ide}\lambda _{jea}`$ $``$ $`\lambda _{[jklpqrs]cd}\lambda _{ide}\delta _{ea}),`$ (185)
where $`a,b,\mathrm{},e=1,\mathrm{},n`$ are matrix element indices, $`\lambda _{iab}(\lambda _i)_{ab}`$. Now we may use (142) once more. The second term of (142) gives zero contribution, or rather its contributions to the eight terms of (185) cancel pairwise. Turning next to the contributions that come from the first term of (142), we see the second, fourth, sixth and seventh terms of (185) vanish because of trace results such as (118). The first term gives
$$2\mathrm{T}\mathrm{r}\left(\lambda _{[jklpqrs]}\right)\mathrm{Tr}\left(\lambda _{[jklpqrs]}\right)=8(\mathrm{\Omega }^{(7)})^2,$$
(186)
by steps like those that yielded (184). The eighth term gives
$`2\mathrm{T}\mathrm{r}\left(\lambda _{[jklpqrs]}\lambda _{[jklpqrs]}\right)\mathrm{Tr}I_n`$ $`=`$ $`(2)\frac{2}{7}\frac{(n^21)}{n}\mathrm{Tr}\left(\lambda _{[jklpqr]}\lambda _{[jklpqr]}\right)n`$ (187)
$`=`$ $`\frac{8}{7}(n^21)(\mathrm{\Omega }^{(7)})^2,`$
where the result (179) has been used. There thus remains to be treated a set of two terms, one each from the third and fifth lines of (185). We next display the two terms in question together with the results of evaluating them
$`2\mathrm{T}\mathrm{r}(\lambda _{[jklpqrs]}\lambda _{[rs]})\mathrm{Tr}(\lambda _{[jklpq]})`$ $`=`$ $`8.\frac{5}{7}(\mathrm{\Omega }^{(7)})^2,`$ (188)
$`2\mathrm{T}\mathrm{r}(\lambda _{[jklpqrs]}\lambda _{[pqrs]})\mathrm{Tr}(\lambda _{[jkl]})`$ $`=`$ $`8.\frac{3}{7}(\mathrm{\Omega }^{(7)})^2.`$ (189)
Proofs of (188) and (189) are given in below in Sec. 7. We may now collect the contributions (186)–(189) to produce the final answer
$$(\mathrm{\Omega }^{(9)})^2=\frac{1}{14}(n^216)(\mathrm{\Omega }^{(7)})^2.$$
(190)
We have presented this calculation in detail because every aspect of it works for higher cases in almost exactly the same fashion. The main difference for the case $`m=5`$ of $`(\mathrm{\Omega }^{(11)})^2`$ is that there are now three terms in the set of terms that arise in the same way as did (188) and (189), namely
$`2\mathrm{T}\mathrm{r}(\lambda _{[jklpqrsuv]}\lambda _{[uv]})\mathrm{Tr}\lambda _{[jklpqrs]}`$ $`=`$ $`8.\frac{7}{9}(\mathrm{\Omega }^{(9)})^2`$ (191)
$`2\mathrm{T}\mathrm{r}(\lambda _{[jklpqrsuv]}\lambda _{[rsuv]})\mathrm{Tr}\lambda _{[jklpq]}`$ $`=`$ $`8.\frac{5}{9}(\mathrm{\Omega }^{(9)})^2`$ (192)
$`2\mathrm{T}\mathrm{r}(\lambda _{[jklpqrsuv]}\lambda _{[pqrsuv]})\mathrm{Tr}\lambda _{[jkl]}`$ $`=`$ $`8.\frac{3}{9}(\mathrm{\Omega }^{(9)})^2.`$ (193)
With the aid of these results, proved below in Sec 7.1, we reach the $`m=5`$ analogue of (190)
$$(\mathrm{\Omega }^{(11)})^2=\frac{2}{45}(n^225)(\mathrm{\Omega }^{(7)})^2.$$
(194)
Indeed it is possible to infer the general result relating the squares of the $`(2s1)`$\- and $`(2s+1)`$-cocycles associated with the Racah-Casimir operators of order $`s`$ and $`(s+1)`$:
$$(\mathrm{\Omega }^{(2s+1)})^2=\frac{4}{2s(2s1)}(n^2s^2)(\mathrm{\Omega }^{(2s1)})^2,$$
(195)
and hence
$$(\mathrm{\Omega }^{(2m1)})^2=\frac{2^{2m3}n}{(2m2)!}\underset{r=1}{\overset{m1}{}}(n^2r^2).$$
(196)
The last two results show in full the expected factors that force the absence of the $`su(n)`$-algebra cocycle/Omega tensor $`\mathrm{\Omega }^{(2m1)}`$ whenever $`m>n`$. Indeed, the last factor in (196) is $`(n^2(m1)^2)`$ and hence $`(\mathrm{\Omega }^{(2m1)})^2=0`$ whenever $`n<m`$. These results are also crucial in the discussion of Racah-Casimir operators, their eigenvalues and of generalised Dynkin indices for $`su(n)`$.
## 7 Proof of results like (188)-(193)
We begin with the simplest result (188) for which we develop
$`2\mathrm{T}\mathrm{r}(\lambda _{[jklpqrs]}\lambda _{[rs]})\mathrm{Tr}\lambda _{[jklpq]}`$ $`=`$ $`4i\mathrm{Tr}[\left(\frac{2}{n}i\mathrm{\Omega }_{jklpqrs}i\mathrm{\Omega }_{t[jklpqr}d_{s]tx}\lambda _x\right)if_{rsy}\lambda _y]\mathrm{\Omega }_{jklpq}`$ (197)
$`=`$ $`8\mathrm{\Omega }_{t[jklpqr}d_{s]tx}f_{rsx}\mathrm{\Omega }_{jklpq}`$
$`=`$ $`8.\frac{5}{7}\mathrm{\Omega }_{tklpqrs}d_{jtx}f_{rsx}\mathrm{\Omega }_{jklpq}`$
$`=`$ $`\frac{40}{7}\mathrm{\Omega }_{tklpqrs}f_{kla}f_{pqb}d_{abj}f_{rsx}d_{jtx}`$
$`=`$ $`\frac{40}{7}\mathrm{\Omega }_{tklpqrs}f_{kla}f_{pqb}f_{rsx}d_{}^{(4)}{}_{(tabx)}{}^{}`$
$`=`$ $`\frac{40}{7}\mathrm{\Omega }_{tklpqrs}\mathrm{\Omega }_{tklpqrs},`$
where we used (124) and (119) in the first line. Next, by opening up the square brackets, we find seven terms of which two vanish upon use of (66), while the remaining five are seen to be equal after relabelling. This accounts for the fraction that appears in the third line. In the fourth line we have used the definition of $`\mathrm{\Omega }_{jklpq}`$, which sets the scene for using the definition (10) of $`\mathrm{\Omega }_{tklpqrs}`$ to reach the last line. It may be noted that it is the first Omega tensor which allows the required symmetries to be implied for the remaining factors in order to build the second Omega tensor.
It should suffice to illustrate things fully to sketch the proof of the most complicated member of the set of results (188)–(193). This requires the $`su(n)`$ relation:
$$\mathrm{Tr}(\lambda _t\lambda _{(abc)})=\mathrm{Tr}\lambda _{(tabc)}=\frac{4}{n}\delta _{t(a}\delta _{bc)}+2d_{st(a}d_{bc)s},$$
(198)
which also follows from (135). Then, putting in some brackets, we get
$`2\mathrm{T}\mathrm{r}(\lambda _{[jklpqrsuv]}\lambda _{[[pq][rs][uv]]})\mathrm{Tr}\lambda _{[jkl]}`$ $`=`$ $`4\mathrm{T}\mathrm{r}\mathrm{\Omega }_{w[jklpqrsu}d_{v]wt}\lambda _tf_{pqa}f_{rsb}f_{uvc}\lambda _a\lambda _b\lambda _cf_{jkl}`$ (199)
$`=`$ $`8\mathrm{\Omega }_{w[jklpqrsu}d_{v]wt}f_{pqa}f_{rsb}f_{uvc}f_{jkl}d_{ht(a}d_{bc)h}`$
$`=`$ $`8.\frac{3}{9}\mathrm{\Omega }_{wklpqrsuv}d_{jwt}f_{klj}f_{pqa}f_{rsb}f_{uvc}d_{ht(a}d_{bc)h}`$
$`=`$ $`\frac{24}{9}\mathrm{\Omega }_{wklpqrsuv}f_{klj}f_{pqa}f_{rsb}f_{uvc}d_{}^{(5)}{}_{(abcj)w}{}^{}`$
$`=`$ $`\frac{24}{9}\mathrm{\Omega }_{wklpqrsuv}\mathrm{\Omega }_{wklpqrsuv}.`$
Here the steps can be seen to be similar to steps already used. In virtue of the work of Sec. 3.2, the first Omega tensor provides enough antisymmetry properties to justify the use of (12) to identify the second one. The fraction in the third line follows opening up of square brackets to expose nine terms, of which six vanish because of (69), i.e.
$$\mathrm{\Omega }^{(9)}{}_{wjklpqrsu}{}^{}f_{jkl}^{}=0,$$
(200)
and the remaining three are equal. The first term of (198) fails to contribute to line two of (199) because of Jacobi identities that also rely on the antisymmetry properties provided by the first Omega tensor.
It can be seen that as one goes, notionally, to higher $`m`$ all the same patterns persist. Although this may require results beyond those explicitly provided here, no problems should be encountered in finding these, the generalisations of (69) being given in Sec.8. We remark also that the coefficients that appear on the right side of (188) and (189), and on the right side of (191)–(193) also conform to a rather obvious pattern, which affords a check on the work, and is instrumental in producing the crucial $`(n^2s^2)`$ factors of the recursion relations (195).
## 8 Proof of (71)
We first prove here that
$$\mathrm{\Omega }^{(11)}{}_{ijklpqrstuv}{}^{}f_{tuv}^{}=0.$$
(201)
This is a critical case because it is the simplest one of the type discussed in Sec.3.4 in which there is a non-trivial invariant totally antisymmetric tensor of the same order as the right side of the identity to be proved, namely the tensor $`\stackrel{~}{\mathrm{\Omega }}^{(8)}`$ of (92). In this case, moreover, the methods of Sec. 3.4 do not offer a viable approach. The method of proof to be given for (201) extends straightforwardly to its analogue for $`\mathrm{\Omega }^{(13)}`$. But then we meet a further critical case
$$\mathrm{\Omega }^{(15)}{}_{ijklpqrstuvwxyz}{}^{}f_{xyz}^{}=0.$$
(202)
This case is critical because it is the simplest one in which there is a non-trivial invariant totally antisymmetric tensor of the same order as $`\mathrm{\Omega }^{(15)}`$ itself, the tensor $`\stackrel{~}{\mathrm{\Omega }}^{(15)}`$ of (93). However there is no obstacle to extending to this case the method of proof to be given for (201).
To prove (201), we start with
$`4\mathrm{\Omega }^{(11)}{}_{ijklpqrstuv}{}^{}f_{tuv}^{}`$ $`=`$ $`\mathrm{Tr}\lambda _{[ijklpqrstuv]}\mathrm{Tr}\lambda _{[tuv]}`$ (203)
$`=`$ $`\mathrm{Tr}\lambda _{[ijklpqrstu]v}\mathrm{Tr}\lambda _{[tu]v}`$
$`=`$ $`2\mathrm{T}\mathrm{r}\lambda _{[ijklpqrstu]}\lambda _{[tu]},`$
using now familiar steps.
To make progress with computing the right side of (203), we drop the second set of square brackets and open out the other set to expose $`\lambda _u`$ in each of its ten terms. This serves to enable a second use of (142):
$`(\lambda _{[ijklpqrst]ab}\delta _{dc}`$ $``$ $`\lambda _{[jklpqrstab}\lambda _{i]dc}`$
$`+\lambda _{[klpqrstab}\lambda _{ij]dc}`$ $``$ $`\lambda _{[lpqrstab}\lambda _{ijk]dc}`$
$`+\lambda _{[pqrstab}\lambda _{ijkl]dc}`$ $``$ $`\lambda _{[qrstab}\lambda _{ijklp]dc}`$
$`+\lambda _{[rstab}\lambda _{ijklpq]dc}`$ $``$ $`\lambda _{[stab}\lambda _{ijklpqr]dc}`$
$`+\lambda _{[tab}\lambda _{ijklpqrs]dc}`$ $``$ $`\delta _{ab}\lambda _{[ijklpqrst]dc})\frac{1}{5}\lambda _{ubd}\lambda _{tce}\lambda _{uea},`$ (204)
where the labels $`a,\mathrm{},e=1,\mathrm{},n`$ are matrix element indices, hence unaffected by antisymmetrisation.
It is easy to check that all ten contributions from the second term of (142) cancel pairwise. The ten contributions from the first term of (142) then are $`\frac{2}{5}`$ times
$`\mathrm{Tr}\lambda _{[ijklpqrst]}\mathrm{Tr}\lambda _t`$ $``$ $`\mathrm{Tr}\lambda _{[jklpqrst}\mathrm{Tr}\lambda _{i]t}`$
$`+\mathrm{Tr}\lambda _{[klpqrst}\mathrm{Tr}\lambda _{ij]t}`$ $``$ $`\mathrm{Tr}\lambda _{[lpqrst]}\mathrm{Tr}\lambda _{ijk]t}`$
$`+\mathrm{Tr}\lambda _{[pqrst}\mathrm{Tr}\lambda _{ijkl]t}`$ $``$ $`\mathrm{Tr}\lambda _{[qrst]}\mathrm{Tr}\lambda _{ijklp]t}`$
$`+\mathrm{Tr}\lambda _{[rst}\mathrm{Tr}\lambda _{ijklpq]t}`$ $``$ $`\mathrm{Tr}\lambda _{[st}\mathrm{Tr}\lambda _{ijklpqr]t}`$
$`+\mathrm{Tr}\lambda _{[t}\mathrm{Tr}\lambda _{ijklpqrs]t}`$ $``$ $`n\mathrm{Tr}(\lambda _{[ijklpqrst]}\lambda _t).`$ (205)
Terms 1 and 9 here are zero trivially, terms 2, 4, 6, 8 are zero using (118). Also term 10 is zero by (172). This leaves terms 3, 5 and 7. Terms 3, 5 and 7 are, to within a common factor, given by
$`\mathrm{\Omega }^{(7)}{}_{[klpqrst}{}^{}f_{ij]t}^{}`$ (206)
$`\mathrm{\Omega }^{(5)}{}_{[pqrst}{}^{}\mathrm{\Omega }_{}^{(5)}_{ijkl]t}`$ (207)
$`f_{[rst}\mathrm{\Omega }^{(7)}{}_{ijklpq]t}{}^{}.`$ (208)
The terms (206), (208) are zero since they are the result of extending the antisymmetrisation of expressions that are already zero by $`Ad`$-invariance, cf. (38). Similarly the term (207) is zero by (42).
Acknowledgements. This work was partly supported by the DGICYT, Spain ($`\mathrm{\#}`$PB 96-0756) and PPARC, UK. One of the authors (JA) wishes to thank the theory group at Imperial College, London, for their hospitality during the last stages of this paper.
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# Is dynamic ultrametricity observable in spin glasses?
## I Introduction
A remarkable feature of mean-field models of spin glasses is the fact that pure states are organized in a hierarchical way. Given three equilibrium configurations (1,2,3) they determine a triangle whose sides are measured by the overlaps $`q_{12}`$, $`q_{13}`$ and $`q_{23}`$. If the configurations are chosen with the Gibbs weight then the two smallest overlaps are equal: all triangles are isosceles. This property, valid in the thermodynamic limit, defines an ultrametric space .
A first question to ask is whether ultrametricity exists in finite dimensional systems. Numerically, this is investigated by searching low temperature ground states of very small systems, and looking for an ultrametric organization between them . Experimentally, however, spin glasses are always out of equilibrium, and another relevant question may be to know what signatures of ultrametricity can be observable in a non-equilibrium system in the thermodynamic limit. In a relaxing (aging) mean-field system, it has been shown that ultrametricity would manifest itself in the dynamical behavior : the correlation functions at three times $`t_1>t_2>t_31`$, $`C(t_1,t_2)`$, $`C(t_2,t_3)`$ and $`C(t_1,t_3)`$ satisfy the relation
$$C(t_1,t_3)=\mathrm{min}\{C(t_1,t_2),C(t_2,t_3)\}.$$
(1)
This property is often termed ‘dynamic ultrametricity’. A link between dynamic and static ultrametricity has been established by Franz et al , who showed that the existence of an ultrametric solution from the dynamic point of view implies, under certain assumptions, the hierarchical organization of pure states in short-range systems. However, we are not aware of any report of dynamical ultrametricity (in the form (1)) in a simulated finite dimensional system, or in experiments . The reason may be that dynamic ultrametricity such as described by (1) is extremely difficult to observe numerically in an aging system . Hence, for the 3D Edwards-Anderson model, Franz and Ricci-Tersenghi resorted to an indirect argument involving the coupling of replicas. Their simulation (together with their argument) then suggests a situation for three dimensions similar to the one observed in mean-field models.
From the experimental point of view, one obviously cannot use the method of . Moreover, noise autocorrelations are extremely difficult to measure, but one can check the version of Eq. (1) involving the response functions :
$$M^{\text{TRM}}(t_1)|_{t_w=t_3}=\text{min}\{M^{\text{TRM}}(t_1)|_{t_w=t_2},M^{\text{TRM}}(t_2)|_{t_w=t_3}\},$$
(2)
where $`M^{\text{TRM}}(t_a)|_{t_w=t_b}`$ is the thermoremanent magnetization at time $`t_a`$ when the field was cut off at time $`t_b`$. Such a relation is not observed experimentally. Indeed, the fit
$$M^{\text{TRM}}(t_a)|_{t_w=t_b}\left(\frac{h(t_a)}{h(t_b)}\right)$$
(3)
works very well for suitable functions $`h(t)`$ , and it is incompatible with dynamic ultrametricity.
As has been often remarked, glassy dynamics can be studied in two complementary forms. One can study the relaxational aging system — and then the large parameter is the waiting time $`t_w`$ — or one can drive the system in such a way it becomes stationary . In this second ‘rheological’ form, the relevant parameter is the intensity of the driving force. Superconducting vortex systems with disorder have been studied this way : in the presence of current the vortices are driven by a Lorentz force. In a spin glass, though less obvious to implement in experiments, the role of a ‘stirring’ force can be played by non-symmetric couplings .
Consider a system to which some non-conservative force of strength $`\epsilon `$ is added. It is by now well established that in a quench experiment, after some transient, the system becomes stationary. If the system is glassy, the relaxation times in the stationary regime will diverge as $`\epsilon 0`$. (The time necessary to achieve stationarity is of the same order, and will diverge as well — we will not study this in the paper.) Indeed, we have that correlations decay for small $`\epsilon `$ as $`C(t)=C_f(t)+C_s(t,\epsilon )`$, where the subindices denote ‘fast’ and ‘slow’. \[Recall that since we are in a stationary state, all the functions depend on a single time argument.\] $`C_f(t)`$ represents the ‘cage’ motion, or the fluctuations inside a domain, and is not affected by a small perturbation. $`C_s(t,\epsilon )`$ is the slow motion, in which the system moves along the ‘almost flat’ directions that are responsible for aging (e.g. structural rearrangements in a glass, domain wall motions in coarsening). The Edwards-Anderson parameter is given by $`qlim_t\mathrm{}lim_{\epsilon 0}C(t)`$. A similar decomposition holds for $`M^{\text{TRM}}(t)`$.
The dynamic ultrametricity discussed above becomes extremely simple in terms of the scaling of $`C_s(t,\epsilon )`$ as $`\epsilon 0`$. One possibility is indeed simple scaling:
$$C_s(t,\epsilon )=f\left(\frac{t}{\overline{t}(\epsilon )}\right)$$
(4)
where $`\overline{t}(\epsilon )\mathrm{}`$ when $`\epsilon 0`$, and $`f`$ is a scaling function. Whenever Eq. (4) holds, we have that for two values of the correlation $`C_1<C_2<q`$:
$$\frac{t(C_1)}{t(C_2)}\underset{\epsilon 0}{}\text{constant},$$
(5)
where the relation $`C_s(t)`$ has been inverted to give $`t(C)`$. The behavior of Eqs. (4,5) is not the only possibility. An alternative would be that $`C_s(t)`$ is the sum of several terms of the type (4), each having a different scaling with $`\epsilon `$. An extreme example of this is the ultrametric law, which for all $`C_1<C_2<q`$ results in
$$\frac{t(C_1)}{t(C_2)}\underset{\epsilon 0}{}\mathrm{}.$$
(6)
The relationship of this property to Eq. (1) is immediate since, in the limit $`\epsilon 0`$, one has $`C(t(C_1)+t(C_2))C(t(C_1))=\mathrm{min}\{C(t(C_1)),C(t(C_2)\},`$ which is the application of Eq. (1) to a stationary situation. This point is discussed in more detail in section II C. The two following examples illustrate our discussion. Simple scaling is, for small $`\epsilon `$:
$$C_s(t)f\left(\frac{t}{\overline{t}(\epsilon )}\right)t(C)\overline{t}(\epsilon )j(C),$$
(7)
where $`j(C)`$ is the inverse function of $`f`$, while dynamic ultrametricity may be obtained with the small $`\epsilon `$ scaling
$$\mathrm{ln}[t(C)]\overline{t}(\epsilon )j(C)C_s(t)f\left(\frac{\mathrm{ln}(t)}{\overline{t}(\epsilon )}\right),$$
(8)
where $`\overline{t}(\epsilon )`$ may diverge with $`\epsilon 0`$ as, for example, $`\overline{t}(\epsilon )\epsilon ^a`$ and $`j(C)`$ is a positive decreasing function of $`C`$. Equation (8) clearly implies (6) and (1), and expresses the hierarchical scaling in a very straightforward way. In the stationary case, dynamic ultrametricity simply differs from simple scaling in that to make the correlation curves for different small $`\epsilon `$ collapse, one has to rescale $`\mathrm{ln}(t)`$ instead of $`t`$ with a function of $`\epsilon `$.
Let us emphasize that such a simple statement of ultrametricity is not possible in the aging case. Consider, as an example, an aging system with correlations evolving as $`C(\tau +t_w,t_w)=\mathrm{ln}(t_w)/\mathrm{ln}(t_w+\tau )`$, which is non-ultrametric since $`C(t_3,t_1)=C(t_3,t_2)C(t_2,t_1)`$ (as opposed to Eq. (1)). If we compare the time-differences $`\tau (C)`$ to reach two correlation values $`C_1<C_2`$, it is easy to find:
$$\frac{\tau (C_1)}{\tau (C_2)}(t_w)^{\frac{1}{C_1}\frac{1}{C_2}}\underset{t_w\mathrm{}}{}\mathrm{}.$$
(9)
Hence, a criterion like (6) with $`t_w`$ playing the role of the large parameter is not applicable in the aging case, and one has to go back to Eq. (1).
Our aims in this paper are the following:
(i) Within mean-field, the solution exhibiting dynamical ultrametricity relies on an asymptotic analysis of the dynamical equations involving two-point correlation and response functions, which may thus be questionable. We will show that this analysis is indeed correct by solving numerically the dynamical equations governing the driven dynamics of a mean-field spin glass in the stationary regime. This is done on a very wide range of time scales, making us confident that the asymptotic solution is the correct one.
(ii) Since ultrametricity should hold strictly in the asymptotic limit of zero driving force (or equivalently in the infinite waiting time limit in the aging case), it is important to study also the preasymptotic regime to understand how ultrametricity gradually develops, and how a full hierarchy of time scales appears. We analyze then this preasymptotic regime in detail, thus characterizing the onset of ultrametricity.
(iii) Having the mean-field dynamical behavior in mind, it is very tempting to see if something similar happens in a finite dimensional system. We have then performed a Monte-Carlo simulation of the 3D Edwards-Anderson model with asymmetrical couplings. As is (unfortunately) usual in 3D spin glass simulations, our numerical results may be interpreted in two different ways. They are indeed compatible with an extremely slow appearance of dynamic ultrametricity, i.e. much slower than in the mean-field case. But an alternative view is that there is asymptotically only one relevant time scale, or, in other words, that dynamic ultrametricity is not present at all. In both cases the conclusion is that even if one assumes that ultrametricity is present, it has not fully developed for experimentally accessible time windows.
The paper is organized as follows. In the next section, the mean-field model under study is presented, and its preasymptotic behavior is detailed in Section III. Section IV presents our numerical results in three dimensions and Section V contains our conclusions.
## II A mean-field driven spin glass
### A Model
We focus in the paper in a mean-field spin glass model first introduced in Ref. , where the statics was solved. The equilibrium dynamics was recently worked out . It consists in a slight modification of the spherical $`p`$-spin model. We consider indeed $`N`$ continuous variables $`s_i`$ ($`i=1,\mathrm{},N`$) interacting through the Hamiltonian
$$H=\underset{p=2}{\overset{\mathrm{}}{}}\underset{j_1<\mathrm{}<j_p}{}J_{j_1\mathrm{}j_p}s_{j_1}\mathrm{}s_{j_p}.$$
(10)
In this expression, the $`J`$’s are random Gaussian variables, symmetrical about the permutation of ($`j_1,\mathrm{},j_p`$) with mean zero and variance
$$\overline{J_{j_1\mathrm{}j_p}^2}=\frac{p!J_{p}^{}{}_{}{}^{2}}{2N^{p1}},$$
(11)
so that the thermodynamic limit is well-defined. A spherical constraint $`_{i=1}^Ns_i^2(t)=N`$ is moreover imposed to the spins. It is convenient to define also the function
$$g(C)\frac{1}{2}\underset{p=2}{\overset{\mathrm{}}{}}J_{p}^{}{}_{}{}^{2}C^p.$$
(12)
The interesting case is when quadratic couplings are present together with quartic and/or higher order interactions. It was shown that the model belongs then to the universality class of the Sherrington-Kirkpatrick model : it has a continuous transition between a paramagnetic phase and a spin glass phase characterized by a full replica symmetry pattern and a non-trivial probability distribution of overlaps $`P(q)`$. Quantitatively, below the transition, one has
$$x(C<q)_0^Cdq^{}P(q^{})=T\frac{g^{\prime \prime \prime }(C)}{2(g^{\prime \prime }(C))^{3/2}};x(C>q)=1,$$
(13)
where $`q`$ is the Edwards-Anderson parameter. This holds only if $`J_2`$ is big enough so that $`x^{}(C)=P(C)>0`$. We consider in the following the combination of second and sixth order terms, keeping then only $`J_2`$ and $`J_6`$ different from zero. In this particular case, one has $`T_c=J_2`$ (independent of $`J_6`$). The positivity of $`x^{}(C)`$ gives moreover the inequality $`J_{6}^{}{}_{}{}^{2}<J_{2}^{}{}_{}{}^{2}/15`$.
### B Driven dynamics
To study the driven dynamics of the model, it is usual to consider a Langevin equation
$$\frac{s_i(t)}{t}=\mu (t)s_i(t)\frac{\delta H}{\delta s_i(t)}+\epsilon f_i^{\text{nc}}(t)+\eta _i(t).$$
(14)
As discussed in the introduction, we take for the (non-conservative) driving force
$$f_i^{\text{nc}}=\underset{\begin{array}{c}j_1<\mathrm{}<j_{k1}\\ j_1,\mathrm{},j_{k1}i\end{array}}{}\stackrel{~}{J}_i^{j_1\mathrm{}j_{k1}}s_{j_1}\mathrm{}s_{j_{k1}}.$$
(15)
The parameter $`\epsilon `$ controls the strength of the force. The couplings $`\stackrel{~}{J}`$’s in the driving force are random Gaussian variables, symmetrical about the permutations of ($`j_1,\mathrm{},j_{k1}`$), with mean zero and variance
$$\overline{\stackrel{~}{J}_i^{j_1\mathrm{}j_{k1}}\stackrel{~}{J}_i^{j_1\mathrm{}j_{k1}}}=\frac{k!}{2N^{k1}};\overline{\stackrel{~}{J}_i^{j_1\mathrm{}j_{k1}}\stackrel{~}{J}_{j_r}^{j_1\mathrm{}i\mathrm{}j_{k1}}}=0.$$
(16)
These couplings are partially uncorrelated and contain thus an antisymmetrical part. This makes it impossible to write the driving force as the derivative of an energy. The parameter $`\mu (t)`$ ensures the spherical constraint and $`\eta _i(t)`$ ($`i=1,\mathrm{},N`$) are random Gaussian variables with mean 0 and variance $`2T`$, where $`T`$ is the temperature of the heat bath.
The dynamics of the model is better analyzed in terms of the autocorrelation function $`C(t,t^{})_is_i(t)s_i(t^{})/N`$ and the response function $`R(t,t^{})_i\delta s_i(t)/\delta \eta _i(t^{})/N`$, since in the thermodynamic limit, $`N\mathrm{}`$, $`C(t,t^{})`$ and $`R(t,t^{})`$ verify closed Dyson equations . The presence of the non-conservative force allows to replace two-time functions $`C(t,t^{})`$ by single arguments functions $`C(tt^{})`$, and the following equations are obtained:
$`{\displaystyle \frac{\mathrm{d}C(t)}{\mathrm{d}t}}=`$ $`\mu C(t)+{\displaystyle _0^t}dt^{}\mathrm{\Sigma }(tt^{})C(t^{})+{\displaystyle _0^{\mathrm{}}}dt^{}\left[\mathrm{\Sigma }(t+t^{})C(t^{})+D(t+t^{})R(t^{})\right],`$ (17)
$`{\displaystyle \frac{\mathrm{d}R(t)}{\mathrm{d}t}}=`$ $`\mu R(t)+{\displaystyle _0^t}dt^{}\mathrm{\Sigma }(tt^{})R(t^{}),`$
$`\mu =`$ $`T+{\displaystyle _0^{\mathrm{}}}dt^{}\left[D(t^{})R(t^{})+\mathrm{\Sigma }(t^{})C(t^{})\right],`$
$`D(t)`$ $`g^{}(C(t))+\epsilon ^2{\displaystyle \frac{k}{2}}C(t)^{k1},\mathrm{\Sigma }(t)g^{\prime \prime }(C(t))R(t).`$
These integro-differential equations are associated with initial conditions $`C(0)=1`$, $`R(0^+)=1`$ and with the condition $`R(t<0)=0`$ (causality). The use of $`C(t)=C(t)`$ in the derivation of Eqs. (17) made these equations non-causal in the time difference, as can be seen from the last integral in the equation for $`C(t)`$.
### C Asymptotic solution of the dynamical equations
In order to make the paper self-contained, and to set our notations, the asymptotic analysis of Eqs. (17) is now briefly recalled. A more detailed computation can be found in . ‘Asymptotic analysis’ means that the limit $`\epsilon 0`$ is taken.
The first step of the analysis consists in making the decomposition $`C(t)=C_s(t)+C_f(t)`$ and $`R(t)=R_s(t)+R_f(t)`$ between a fast and a slow part, and to derive the equations verified by each part. The equations for $`C_f`$ and $`R_f`$ are solved by making the ansatz that they satisfy the fluctuation-dissipation theorem (FDT) $`TR_f(t)=\mathrm{d}C_f(t)/\mathrm{d}t`$. Taking the limit $`t\mathrm{}`$ gives then the following relation:
$$\mu +\frac{g^{}(q)g^{}(1)}{T}=\frac{T}{1q},$$
(18)
where $`qlim_t\mathrm{}lim_{\epsilon 0}C(t)`$ is the Edwards-Anderson parameter. The slow parts verify:
$`{\displaystyle \frac{\mathrm{d}C_s(t)}{\mathrm{d}t}}=`$ $`\mu C_s(t)+{\displaystyle _0^t}dt^{}\mathrm{\Sigma }_s(tt^{})C_s(t^{})+{\displaystyle _0^{\mathrm{}}}dt^{}\left[\mathrm{\Sigma }_s(t+t^{})C_s(t^{})+D_s(t+t^{})R_s(t^{})\right]`$ (19)
$`+{\displaystyle \frac{g^{}(1)g^{}(q)}{T}}C_s(t)+{\displaystyle \frac{1q}{T}}D_s(t),`$
$`{\displaystyle \frac{\mathrm{d}R_s(t)}{\mathrm{d}t}}=`$ $`\mu R_s(t)+{\displaystyle _0^t}dt^{}\mathrm{\Sigma }_s(tt^{})R_s(t^{})+{\displaystyle \frac{g^{}(1)g^{}(q)}{T}}R_s(t)+{\displaystyle \frac{1q}{T}}\mathrm{\Sigma }_s(t),`$
$`\mu =`$ $`T+{\displaystyle _0^{\mathrm{}}}dt^{}\left[D_s(t^{})R_s(t^{})+\mathrm{\Sigma }_s(t^{})C_s(t^{})\right]+{\displaystyle \frac{g^{}(1)qg^{}(q)}{T}},`$
$`D_s(t)`$ $`g^{}(C_s(t)),\mathrm{\Sigma }_s(t)g^{\prime \prime }(C_s(t))R_s(t).`$
The second step stems from the observation that once the first order derivatives are dropped out (which is of course justified in the slow regime), the above equations become invariant under a reparametrization of time $`th(t)`$. This suggests that $`C_s`$ and $`R_s`$ are related through a reparametrization-invariant formula, namely
$$R_s(t)=\frac{X(C_s(t))}{T}\frac{\mathrm{d}C_s(t)}{\mathrm{d}t}.$$
(20)
This amounts to an extension of the FDT to this non-equilibrium situation by the introduction of an effective temperature $`T_{\text{eff}}(C)T/X(C)`$. In the same spirit, the function $`f`$ defining ‘triangles’ is introduced:
$$C_s(t_1+t_2)f(C_s(t_1),C_s(t_2)),$$
(21)
Let us emphasize that the main assumption of the analysis is that the functions $`X(C_s,\epsilon )`$ and $`f(C_1,C_2,\epsilon )`$ have a continuous limit when $`\epsilon 0`$. It is also convenient to define the function $`\overline{f}`$: $`C_s(t_1)\overline{f}(C_s(t_2),C_s(t_1+t_2))`$. This allows to rewrite Eqs. (19) in such a way that the time disappears:
$`0=`$ $`{\displaystyle \frac{T}{1q}}C_s+{\displaystyle \frac{1q}{T}}g^{}(C_s)+{\displaystyle \frac{1}{T}}{\displaystyle _{C_s}^q}dC_s^{}g^{\prime \prime }(C_s^{})X(C_s^{})\overline{f}(C_s^{},C_s)`$ (22)
$`+{\displaystyle \frac{1}{T}}{\displaystyle _0^{C_s}}dC_s^{}X(C_s^{})g^{\prime \prime }(C_s^{})\overline{f}(C_s,C_s^{}){\displaystyle \frac{1}{T}}{\displaystyle _0^{C_s}}dC_s^{}g^{\prime \prime }(C_s^{})F[\overline{f}(C_s,C_s^{})],`$
$`0=`$ $`{\displaystyle \frac{T}{1q}}F[C_s]+{\displaystyle \frac{1q}{T}}H[C_s]+{\displaystyle \frac{1}{T}}{\displaystyle _{C_s}^q}dC_s^{}g^{\prime \prime }(C_s^{})X(C_s^{})F[\overline{f}(C_s^{},C_s)],`$
$`F[C_s]`$ $`{\displaystyle _{C_s}^q}dC_s^{}X(C_s^{}),H[C_s]{\displaystyle _{C_s}^q}dC_s^{}X(C_s^{})g^{\prime \prime }(C_s^{}),`$
where Eq. (18) has been used.
The last step consists in computing explicitly $`X`$ and $`f`$ from Eqs. (22). This is done by introducing ‘fixed points’ of $`f`$ . A fixed point $`q^{}`$ of $`f`$ satisfies
$$f(q^{},q^{})=q^{}.$$
(23)
In the stationary context we are discussing, the physical meaning of such values of the correlation is very clear, since from the definition of $`f`$,
$$q^{}=C_s(t)C_s(2t)=f(C_s(t),C_s(t))=C_s(t),$$
(24)
which simply means that $`q^{}`$ is a plateau in the correlation function. Two trivial such fixed points are $`q=0`$ and $`q=1`$. It is then rather straightforward to solve Eqs. (22). This gives $`X(C)`$ for $`C[0,q]`$,
$$X(C)=T\frac{g^{\prime \prime \prime }(C)}{2(g^{\prime \prime }(C))^{3/2}},$$
(25)
and the matching with the FDT regime ($`X(C>q)=1`$) determines $`q`$ which satisfies $`T=(1q)\sqrt{g^{\prime \prime }(q)}`$. The fact that $`X(C)`$ coincides with $`x(C)`$ (Eq. (13)) is a remarkable property of this class of mean-field spin glass models . The properties of $`f`$ are also obtained. It is first shown that $`C[0,q]`$, $`f(C,C)=C`$, and there is hence a continuum of fixed points. Crudely speaking, it can be said that, in the asymptotic limit, each value of $`C<q`$ corresponds to a plateau value in the correlation function (in the sense that the correlation will not decay below this value in a finite time). This ultimately means that we will have a behavior as described by Eq. (6). It is also shown that
$$(C_1,C_2)[0,1]^2[q,1]^2f(C_1,C_2)=\mathrm{min}\{C_1,C_2\},$$
(26)
which is the ultrametric relation, as presented in the introduction.
The conclusion of this section, is that, apart from the trivial short time behavior where $`C(t)>q`$ and FDT holds $`X(C>q)=1`$, the relaxation below $`q`$ is characterized by a hierarchy of time scales and a non-trivial FDT $`X(C<q)<1`$.
## III Ultrametricity from dynamical correlations
As stated in the preceding section, the asymptotic solution of the dynamical equations exhibits a ‘many-plateau pattern’, which is an unusual feature. To see how this asymptotic solution is approached, Eqs. (17) were solved numerically, in the particular case where $`g(C)=C^2/2+C^6/30`$, and $`k=2`$. We shall mainly work at $`T=0.25=0.25T_c`$, where the Edwards-Anderson parameter is $`q0.787`$. To solve these equations, a combination of numerical methods of Refs. has been used.
### A Hierarchy of time scales
A first element that is missed by the above analysis is the functional form of the correlation functions. They are depicted in Fig. 1, for different values of $`\epsilon `$ and for $`T=0.25`$. As expected, two different regimes are clearly present. For $`C(t)>q`$, the relaxation depends very weakly on the asymmetry, while for $`C(t)<q`$, the smaller the asymmetry, the slower the relaxation.
The main question that cannot be answered analytically is the precise dependence of the relaxation times on the parameter $`\epsilon `$ that controls the strength of the asymmetry. The ‘many-plateau pattern’ discussed above means that, in the preasymptotic regime (non-zero $`\epsilon `$), the correlation will stay a large, albeit finite time around the same value. This characteristic time increases with decreasing $`C`$, as expressed by Eq. (6). In order to test this separation of time scales, we compute the ratio of relaxation times for two fixed values $`C_1`$ and $`C_2`$ of the correlation as a function of $`\epsilon `$. Two such ratios are represented in Fig. 2, where it is clearly seen that they indeed diverge in the small asymmetry limit. A consequence is that the relaxation cannot be represented by a single time scale, but involves a full hierarchy of them. In particular, a stretched exponential fit such as the one proposed in for the Sherrington-Kirkpatrick model with asymmetry can only be approximate.
The numerical solution of the dynamical equations suggests the following dependence for $`C<q`$:
$$t(C)\mathrm{exp}\left[\overline{t}(\epsilon )j(C)\right];\overline{t}(\epsilon )\epsilon ^{0.65}.$$
(27)
In this expression, $`j(C)`$ is a positive decreasing function of the correlation. It plays the same role as in the example of the introduction, Eq. (8). Numerically, $`j(C)`$ is consistent with a linear variation: $`j(C)j(C=0)bC`$. For a fixed value $`C`$ of the correlation, $`t(C)`$ grows faster than a power law of $`\epsilon `$, as was noted in Ref.. It also follows that
$$\frac{t(C_1)}{t(C_2)}\mathrm{exp}\left[\left(j(C_1)j(C_2)\right)\overline{t}(\epsilon )\right],$$
(28)
which means that the ratios of two time scales may be fitted by the same functional form (27) as the time scales themselves. This is also displayed in Fig. 2.
Let us note that the scaling (27) is very reminiscent of the ‘creep’ regime scaling for vortex glasses , with the role of $`\epsilon `$ played by the current and $`C`$ a measure of the average squared transverse displacements along the vortex.
The presence of this hierarchy of time scales in (27) implies that the correlation curves can not be superimposed by rescaling the time. This is illustrated in Fig. 3, where the time is rescaled so that the curves meet at the value $`C=0.3`$. It can be seen that the curves are more and more horizontal around $`C=0.3`$, when the asymmetry is decreased. On the contrary, if we rescale the logarithm of the time to make the curves meet at $`C=0.3`$ (Fig. 4) we find that the curves tend to collapse. In other words: shifting $`\mathrm{ln}(t)`$ (as in simple scaling) does not make the curves collapse, while stretching $`\mathrm{ln}(t)`$ does. A similar picture would have been obtained by choosing any value $`C[0,q]`$.
This rescaling (or absence thereof) of the correlation functions in the driven dynamics of a spin glass is a direct and very simple test of dynamic ultrametricity. Let us emphasize again that, as already mentioned in the introduction, this is not true for the aging regime.
### B Ultrametric relation and FDT
The function $`f(C_1,C_2)`$ introduced in Eq. (21) satisfies in the asymptotic limit the ultrametric relation (26). It is then natural to try to understand its preasymptotic behavior. A three dimensional view of this two-variable function is given in Fig. 5, for $`\epsilon =0.248`$ (the slowest relaxation in Fig. 1). Also plotted in the plane ($`C_1,C_2`$) are the constant-$`f`$ contours. These contours would be right angles in the limit of vanishing asymmetry.
To see how the function $`f`$ evolves towards its asymptotic value (the ultrametric relation (26)), the evolution of the contours for different $`\epsilon `$ is represented in Fig. 6. These curves are very clearly evolving towards right angles, as expected from the analysis of the preceding section.
It is interesting to compare these curves with the ones obtained from Ref. for a system ($`p`$-spin with $`p=3`$) with a single time scale, i.e. without ultrametricity. The result is shown in Fig. 7, for correlations that evolve on a similar range of time scales to make the comparison relevant. In this case, the asymptotic analysis reveals that there exists asymptotically a function $`f`$ defined as above, but it is not given by the ultrametric relation (26). The difference between the two systems is very clear from the comparison of Figs. 6 and 7: in the latter, the function $`f`$ rapidly saturates to its asymptotic (non-ultrametric) value.
To make the analysis of the mean-field dynamics complete, the usual plot of the integrated response function $`\chi (t)_0^tdt^{}R(t^{})`$ as a function of the correlation function $`C(t)`$, parameterized by the time $`t`$, is done in Fig. 8. As expected from the above analytical results, the FDT holds for $`C(t)>q`$ and it is strongly violated for smaller values of the correlation. It is clear that in the limit of zero asymmetry, the analytic expression for the function $`X(C)`$ that generalizes the FDT to non-equilibrium situations will be recovered. The fact that the same limiting FDT violations happen in gently driven and aging systems has been suggested in Ref. , and we have verified here that this (mean-field) result holds also for systems with many time scales.
## IV Simulation of the 3D Edwards-Anderson model
We turn now to the numerical results obtained for a 3D spin glass model, the results of the preceding sections being a guide to investigate its stationary driven dynamics. We use the same notations for quantities that play a similar role in the simulation and in the mean-field model, the distinction between the two cases being clear from the context.
### A Model and details of the simulation
The model under study is defined through its Hamiltonian
$$H=\underset{i,j}{}J_{ij}s_is_j,$$
(29)
where $`s_i`$ ($`i=1,\mathrm{},N`$) are $`N=L^3`$ Ising spins located on the sites of 3D cubic lattice of linear size $`L`$, with periodic boundary conditions. The sum $`i,j`$ runs over nearest neighbors and the $`J`$’s are chosen randomly from a bimodal distribution $`J_{ij}=\pm 1`$. The model has been extensively studied : it exhibits a second order phase transition at the critical temperature $`T_c=1.11\pm 0.04`$ from a paramagnetic to a spin glass phase .
To drive the system, a coupling $`\stackrel{~}{J}_{ij}`$ is added on each link. The $`\stackrel{~}{J}`$’s are chosen from a bimodal distribution $`\stackrel{~}{J}_{ij}=\pm \epsilon `$, and are antisymmetrical: $`\stackrel{~}{J}_{ij}=\stackrel{~}{J}_{ji}`$. The small parameter $`\epsilon `$ controls the strength of the driving force. The spins are randomly sequentially updated through a standard Metropolis algorithm, and one Monte Carlo step represents $`N`$ attempts to update a spin. Numerical results are presented for a linear size $`L=20`$ ($`N=8000`$ spins), where finite size effects are negligible for the time scales investigated here .
The temperature has been chosen so that the Edwards-Anderson parameter is comparable to the one of the mean-field case studied before. At $`T=0.60.54T_c`$, the Edwards-Anderson parameter was roughly estimated in an off-equilibrium simulation (with $`\epsilon =0`$) through its dynamical definition $`q=lim_t\mathrm{}lim_{t_w\mathrm{}}C(t,t_w)0.8`$ (we had $`q0.787`$ in the mean-field study). Similarly, 5 different values of the asymmetry were studied: $`\epsilon =0.5`$, 0.4, 0.3, 0.25 and 0.2, so that the range of time scales is comparable to the previous mean-field results. Remarkably, initial conditions are irrelevant, since a stationary state is reached after a time which depends on the intensity of the drive. As was already noted in Ref. , this is an unusual feature for glassy systems where cooling procedures are known to be crucial. Stationarity allows moreover to average the correlation functions over different initial times, so that very few averages over the disorder have been necessary (typically 5), provided the first steps of the simulation are discarded. Stationarity has been carefully checked throughout the simulation.
### B Looking for ultrametricity
The dynamical quantity of interest is the spin-spin autocorrelation function, which in the TTI regime reads
$$C(t)\frac{1}{N}\underset{i=1}{\overset{N}{}}\overline{s_i(t+t_0)s_i(t_0)}.$$
(30)
The overline means that an average over disorder is performed, while $`\mathrm{}`$ stands for an average over different initial times $`t_0`$, all chosen in the TTI regime.
The autocorrelation functions for different values of the asymmetry at $`T=0.6`$ are represented in Fig. 9. As for the mean-field case, two different regimes are present. The short time relaxation towards $`q0.8`$ is very weakly affected by the driving force, while the time to relax towards 0 dramatically increases when $`\epsilon `$ is lowered.
The crucial point is to look for the possible presence of the ultrametric relaxation pattern described in the preceding sections. We have seen above that a simple test of the presence of ultrametricity is the rescaling of the time. The same rescaling that was done in Fig. 3 for the mean-field case is now performed for the 3D case in Fig. 10. Although not perfect, this rescaling works remarkably well for the smallest values of $`C`$. It is important to note that even if ultrametricity is absent, the values of correlation $`1<Cq`$ (which do not depend on the asymmetry) will not scale together with the $`C<q`$ portion of the curves. This may be a source of errors if the value of $`q`$ is unknown. Thus, preasymptotic effects affect the quality of the rescaling in the region $`Cq`$, since it can be seen from Fig. 9 that the plateau is not completely developed. All these points may explain why the rescaling is good only for the smallest values of $`C`$.
However, a very slow flattening of the curves around the value $`C=0.3`$ cannot be completely excluded from this figure. But it is very clear that if there is a separation of time scales, then it will become evident only if a much larger time window is analyzed.
The second interesting quantity to study is the function $`f(C_1,C_2)`$ introduced in Eq. (21). Its evolution for different values of the asymmetry is depicted in Fig. 11, with the same construction as for Figs. 6 and 7. The function $`f`$ indeed evolves, but very slowly when compared to the mean-field case. The quantitative behavior of $`f`$ for the 3D model is more reminiscent of the $`p=3`$ case than of the mean-field spin glass case. Once again, this figure does not allow to decide clearly between a limiting smooth – as it is for the $`p=3`$ case – or ultrametric $`f`$-function, because interesting things may happen on time scales that are inaccessible in our simulation time.
In our opinion, the important point that clearly emerges from the figures is that, on a given time window (we have 4 decades here), the relaxation does not appear to be typical of an ultrametric system and a single-time-scale description is very accurate. This may explain why dynamic ultrametricity has not been observed in aging experiments, which span some six decades $``$ (100 Hz. - 10 hours).
### C FDT
We have also investigated the way the FDT is violated in the driven 3D Edwards-Anderson model, since the mean-field theory suggests that these violations are the same as for an aging system. To our knowledge, there are no numerical confirmation of this prediction for spin glasses in the literature. For this purpose, correlation functions have to be compared to susceptibility curves. The correlation functions are taken from Fig. 9. The way of computing the susceptibility is now standard . A stationary magnetic field $`h_i`$ is applied in each site at time $`t_0`$. It is random in space, and we take it from a Gaussian distribution with mean 0 and variance $`\overline{h_ih_j}=h_0^2`$. The staggered magnetization
$$m(t)\frac{1}{h_0N}\underset{i=1}{\overset{N}{}}h_is_i(t),$$
(31)
is recorded for all $`t>t_0`$. In the linear response regime (we work with $`h_0=0.1`$, as in previous studies ), the susceptibility is obtained from $`m(t)`$ as $`\chi (t)m(t)/h_0`$. The results are averaged over several (from 50 to 300) realizations of the field, and over different initial times $`t_0`$. They are presented in Fig. 12.
The results for the violations of the FDT are clearly very similar to the mean-field ones. These curves exhibit two different regimes: for short times ($`Cq`$), the FDT is well satisfied, whereas for longer times ($`Cq`$), it is violated. The curves also clearly saturate to a smooth limiting curve in the small asymmetry limit. The shape of this limiting curve is compatible with the limiting non-trivial parametric curves that have already been found in simulations of 3D spin glasses in the aging regime . The fact that both situations (aging and driven dynamics) exhibit the same kind of FDT violations deepens, in our opinion, the physical meaning of the quantity $`X(C)`$ (the so-called fluctuation-dissipation ratio) that can be extracted from this plot .
## V Conclusions
We have studied in this paper the behavior of spin glasses in a ‘rheological’ setting, in which the dynamics is stationary and the control parameter is the strength of the driving force.
We have checked the asymptotic analysis for mean-field models through the numerical integration of the equations governing the dynamics. This has confirmed the presence of a full hierarchy of time scales in the relaxation of the correlation, Eq. (26), and a ‘many-time-scale, many-temperature scenario’ . This scenario can be seen as the dynamic counterpart of static ultrametricity. Simple scaling time dependencies, such as assumed in Ref. could be justified for mean field models only as an approximation valid in some time-window, but do not hold strictly in the large time limit. Previous numerical studies of the aging dynamics of the Sherrington-Kirkpatrick model had pointed out the lack of simple $`t/t_w`$-scaling , but did not investigate the presence of ultrametricity. Since there exist now powerful algorithms to solve dynamical two-time equations for large times , it would be interesting to have an analysis of the aging dynamics of a mean-field ultrametric model following the lines introduced in this paper.
For the 3D spin glass, our simulation suggests that the stationary driven dynamics is accurately interpreted within a single time scale relaxation pattern. We cannot decide whether our results imply an extremely slow appearance of ultrametricity (such that, even at experimental times, it is not fully observable), or its absence.
Whether transient or permanent, the single time scale relaxation obtained in our 3D simulation is in complete accordance with all the known numerical and experimental studies of the aging regime of 3D spin glasses . In all known cases, two-time functions are indeed very well approximated by $`M(t,t_w)(t/t_w)`$, in contradiction with the dynamical mean-field theory, as we emphasized all along the paper. It is interesting to note that the version of the ‘droplet theory’ which predicts the scaling $`C(t,t_w)𝒞(\mathrm{ln}(t)/\mathrm{ln}(t_w))`$ fails also in reproducing experiments and simulations, since the logarithmic law is too slow.
Numerical results on 4D spin glasses also show that correlation functions are rather well represented by a $`t/t_w`$-scaling in the aging regime , as is found in 3D. It would be extremely interesting to go back to the 4D simulations with asymmetrical couplings to study precisely the stationary regime, and see if, like in our simulation, there is a single time scale. It would be a disappointing result if all the richness of the ultrametric construction in spin glasses were a pure $`D\mathrm{}`$ feature.
The situation is made even more puzzling by the fact that this construction and the associated separation of time scales are often invoked to explain the results of temperature cycling experiments . This is perhaps related to the difficulty in reproducing these experimental results with the 3D EA model with numerical simulations .
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# Testing Dispersion Relations of Quantum 𝜅–Poincaré Algebra on Cosmological Ground
(October 26, 2000)
## Abstract
Following the procedure proposed recently by Martin and Brandenberger we investigate the spectrum of the cosmological perturbations in the case when the “trans–Plackian” dispersion relations are motivated by the quantum $`\kappa `$-Poincaré algebra. We find that depending on the choice of initial conditions of the perturbations, the spectrum either differs from the flat one for for instantaneous Minkowski vacuum or, in the case of initial conditions minimizing energy density, leads to the observed scale-invariant Harrison–Zel’dovich spectrum in the Friedmann epoch.
In the recent years a growing mass of evidence appeared, indicating that on short (trans-Planckian) scale the usual space-time symmetries are drastically modified, eg., by existence of a fundamental length scale (see e.g. ). Some time ago Lukierski, Nowicki, Ruegg, and Tolstoy using formalism of quantum algebras derived a quantum deformation of the Poincaré algebra, which is a high energy extension of the standard low-energy Poincaré algebra. This algebra, called $`\kappa `$-Poincaré algebra includes the parameter $`\kappa `$ of dimension of (inverse) length, usually identified with the Planck length. The algebra (in the so-called bicrossproduct basis) ) takes the form:
$`[M_{\mu \nu },M_{\rho \tau }]=`$ $`i\left(\eta _{\mu \tau }M_{\nu \rho }\eta _{\nu \rho }M_{\nu \tau }+\eta _{\nu \rho }M_{\mu \tau }\eta _{\nu \tau }M_{\mu \rho }\right),`$ (1)
$`[M_i,P_j]=`$ $`iϵ_{ijk}P_k,[M_i,P_0]=iP_i,`$ (3)
$`[N_i,P_j]=`$ $`i\delta _{ij}\left({\displaystyle \frac{\kappa }{2}}\left(1e^{2P_0/\kappa }\right)+{\displaystyle \frac{1}{2\kappa }}\stackrel{}{P}^2\right)+{\displaystyle \frac{i}{\kappa }}P_iP_j,`$ (5)
$`[N_i,P_0]=`$ $`iP_i,`$ (7)
$`[P_\mu ,P_\nu ]=`$ $`0,`$ (9)
where $`P_\mu =(P_i,P_0)`$ are space and time components of four-momentum, and $`M_{\mu \nu }`$ are modified Lorentz generators with rotations $`M_k=\frac{1}{2}ϵ_{ijk}M_{ij}`$, and boosts $`N_i=M_{0i}`$. It can be readily found that the first Casimir of this algebra, defining the mass-shell is of the form<sup>1</sup><sup>1</sup>1There exists another realization of this quantum algebra, with slightly different co-product, for whose the first Casimir is of the form $`𝒞_1^{bcp^{}}=(\stackrel{}{P})^2e^{+P_0/\kappa }\left(2\kappa \mathrm{sinh}\left(\frac{P_0}{2\kappa }\right)\right)^2`$ and the dispersion relation $`(k)^2e^{+\omega /\kappa }\left(2\kappa \mathrm{sinh}\left(\frac{\omega }{2\kappa }\right)\right)^2=0`$. This case differs from the cases considered in this papers by the fact that there is a cut-off for three-momentum: $`k\kappa `$ when $`\omega \mathrm{}`$. For this reason we will consider this case in a separate paper.
$$𝒞_1^{bcp}=(\stackrel{}{P})^2e^{P_0/\kappa }\left(2\kappa \mathrm{sinh}\left(\frac{P_0}{2\kappa }\right)\right)^2.$$
(10)
Equations (1, 10), with additional structures of quantum algebra (coproducts and antipodes) form the algebraic framework for field theoretical constructions. In the first step in this construction one should identify the deformed $`\kappa `$-Minkowski space. This can be done by taking the set of generators dual to the $`\kappa `$-Poincaré algebra, $`(\widehat{x}^\mu ,\mathrm{\Lambda }^{\rho \sigma })`$, which form the $`\kappa `$-Poincaré group and dividing it by its Lorentz subgroup (observe that in algebra (1) the Lorentz generators are not deformed.) This has been done in the papers . As a results one obtains that the deformed $`\kappa `$-Minkowski space must be necessarily a non-commuting space:
$$[\widehat{x}^\mu ,\widehat{x}^\nu ]=\frac{i}{\kappa }\left(\delta _0^\mu \widehat{x}^\nu \delta _0^\nu \widehat{x}^\mu \right)$$
(11)
Now the question arises as to how one can define a quantum theory on such a space. Such a construction has been presented recently in . This construction consists of three major steps. First one $`\kappa `$-deforms the classical local field theory on standard Minkowski space obtaining as a result a local $`\kappa `$-deformed field theory on $`\kappa `$-deformed, non-commutative Minkowski space. Then one observes that momenta in (1) are commutative and uses an appropriately defined $`\kappa `$-deformed Fourier transform to obtain a $`\kappa `$ deformed field theory on commutative momentum space. In the last step one makes use of the standard inverse Fourier transform to get a non-local, $`\kappa `$-deformed field theory on standard Minkowski space. This non-locality exhibits itself in non-polynomial structure of the Casimir operator (10), but the fact that the Minkowski space is now standard makes it possible to use the standard realization of momentum operators as derivatives over commuting Minkowski space positions, $`P_\mu =i/x^\mu `$.
One can now turn to the free massless scalar field (see ). The invariant wave operator on $`\kappa `$-deformed Minkowski space $`\frac{}{\widehat{x}_\mu }\frac{}{\widehat{x}^\mu }`$ can be expressed, in the momentum space as $`𝒞_1^{bcp}\left(1\frac{𝒞_1^{bcp}}{4\kappa ^2}\right)`$. We see that the spectrum of the modified massless wave operator contains the deformed massless mode $`𝒞_1^{bcp}\varphi =0`$ and the tachyon $`\left(𝒞_1^{bcp}4\kappa ^2\right)\varphi =0`$. In this paper we will consider only the first branch, leaving the possible physical meaning of the second one to future investigations, and thus we are left with a field theory whose dynamics is governed by the Casimir operator, (10).
This Casimir operator leads to the following dispersion relation
$$(k)^2e^{\omega /\kappa }\left(2\kappa \mathrm{sinh}\left(\frac{\omega }{2\kappa }\right)\right)^2=0.$$
(12)
It should be noted that the wave equation leading to the above dispersion relation is non-local in time. Following the discussion in we solve this equation for $`\omega ^2`$ so that the resulting dispersion relation corresponds to an operator of second order in time derivatives, and non-local in space:
$$\omega ^2=\left[\kappa \mathrm{log}\left(1+\frac{k}{\kappa }\right)\right]^2.$$
(13)
Below we will use this expression to define a theory of modified dynamics of fluctuations in inflationary universe. It should be noted that except of the second Newton’s law there is no clear principle saying which relation is physically more fundamental. Therefore, strictly speaking one cannot say that the theory described by (13) and investigated below is derived from $`\kappa `$-Poincaré algebra, though it is not excluded that one could derive this theory from the first principles using some procedure different from that presented in .
It is worth pausing for a moment to show that the dispersion relation (13) indeed leads to modified wave equation covariant under $`\kappa `$-Poincaré symmetry. For the scalar field $`\varphi `$ this equation reads
$$P_0^2\varphi =\left[\kappa \mathrm{log}\left(1+\frac{|\stackrel{}{P}|}{\kappa }\right)\right]^2\varphi .$$
(14)
It is clear that this equation is invariant under three-dimensional rotations generated by $`M_i`$. Let us consider generalized boosts parameterized by infinitesimal parameter $`\epsilon `$ resulting from eq. (1).
$$\delta _\epsilon P_0=[\epsilon ^iN_i,P_0]=i\epsilon ^iP_i$$
(15)
$$\delta _\epsilon P_j=[\epsilon ^iN_i,P_j]=i\epsilon _j\left(\frac{\kappa }{2}\left(1e^{2P_0/\kappa }\right)+\frac{1}{2\kappa }\stackrel{}{P}^2\right)+\frac{i}{\kappa }\epsilon ^iP_iP_j$$
(16)
One can check by direct computation that eq. (14) is covariant under transformations (15), (16) <sup>2</sup><sup>2</sup>2In the course of this calculation one uses eq. (14) twice: first to express $`P_0\delta P_0`$ as $`\kappa \mathrm{log}\left(1+\frac{|\stackrel{}{P}|}{\kappa }\right)\delta P_0`$ and second to rewrite the $`e^{2P_0/\kappa }`$ factor in terms of $`P_i`$. .
Since the parameter $`\kappa `$ is assumed to be of order of Planck length, for a long time the $`\kappa `$-Poincaré algebra remained a nice mathematical construction without any direct, testable physical applications and consequences (see however ). Recently however some areas have been identified, where the behavior of matter at very high, “trans–Planckian” frequencies has a direct, low energy consequences. One of these areas is the black hole physics. The question whether modified high frequency behavior has an effect on Hawking radiation and Hawking–Beckenstein entropy formula has been considered (in the context of sonic black holes) by Unruh , Corley and Jacobson, and others (for recent review see ). These authors proposed a rather ad hoc dispersion relations, to wit
$$\omega =\kappa \mathrm{tanh}^{1/p}\left[\left(\frac{k}{\kappa }\right)^p\right],\text{ (Unruh) }$$
(17)
and
$$\omega ^2=k^2\frac{k^4}{\kappa ^2},\text{ (Corley–Jacobson)},$$
(18)
where we again denoted by $`\kappa `$ the characteristic length.
In the recent papers Niemeyer and Martin and Branderberger considered another setting in which there is a direct relation between trans–Planckian physics and low energy phenomenology. Namely, they considered quantum fluctuations during inflationary stage, giving rise to large scale formation in the Friedmann universe. These fluctuations are enormously red shifted in the course of inflation, and thus the wavelengths of the modes that correspond to the present large scale structure were, at the initial time, well in the realm of the trans–Planckian physics. It follows that it is justified to ask, to what extend the modified dispersion relations have their imprints on the present day large scale structure of the universe. Since we know that in order to agree with observations, the spectrum of fluctuations should be (almost) flat, this setting is perfect to test validity of modified dispersion relations, derived or motivated by other means. One should stress at this point that the relation (12) we are to investigate in this paper are motivated by considerations, having root in some fundamental physics (in the case at hands, from quantum groups).
To set the stage let us recall some fundamental facts concerning early cosmology and large scale structure formation. We consider spatially flat Friedmann universe with the conformally flat line element
$$ds^2=a^2(\eta )\left(d\eta ^2+\underset{i=1}{\overset{3}{}}dx_i^2\right).$$
(19)
The conformal time $`\eta `$ is related to the cosmic time $`t`$ by relation $`dt=a(\eta )d\eta `$; in particular the De–Sitter universe corresponds to $`\eta =H^1e^{Ht}`$, $`a(\eta )=l_H/\eta `$.
To compute power spectra of observable quantities in the standard case (i.e., with unmodified dispersion relations) one considers equation of evolution of modes of Fourier components of massless scalar field $`\mu _n`$
$$\mu _n^{\prime \prime }+\left(n^2\frac{a^{\prime \prime }}{a}\right)\mu _n=0,$$
(20)
where prime denotes derivative with respect to conformal time $`\eta `$. When the wavelength of fluctuations is much longer than the characteristic cosmological scale, the Hubble radius, $`\lambda (\eta )a(\eta )2\pi /nl_H`$, so that the first term in the parenthesis is small compared to $`a^{\prime \prime }/a`$, the solution of this equation is $`\mu _n(\eta )=C_na(\eta )`$. One can then calculate the resulting spectra, to obtain
$$n^3P=n^3|C_n|^2$$
(21)
The observed, flat, scale invariant spectrum corresponds to the case when $`C_nn^{3/2}`$.
This behaviour can be readily seen from the explicit general solution of eq. (20). One has in this case
$$\mu _n(\eta )=\alpha _n\left(1+\frac{i}{n\eta }\right)e^{in\eta }+\beta _n\left(1\frac{i}{n\eta }\right)e^{in\eta }.$$
(22)
For large $`\eta `$ (corresponding to early times) the second term in parentheses is small and the fluctuation behaves as a wave with slowly changing amplitude; starting from $`\eta _H=2\pi /n`$, corresponding to the moment when fluctuation crosses the horizon, there are no oscillations and the wave gets frozen. Thus to find the spectrum we have to compute
$$|C_n|=\frac{1}{a(\eta _H)}\left|\mu _n\left(\eta _H\right)\right|=\frac{2\pi }{l_H}\frac{1}{n}\left|\mu _n\left(\eta _H\right)\right|.$$
(23)
We see therefore that in order to lead to flat spectrum, the coefficients $`\alpha _n`$, $`\beta _n`$ should behave as $`1/\sqrt{n}`$.
Equation (20) corresponds to the standard dispersion relation
$$\omega ^2=k^2=\frac{n^2}{a^2}.$$
In the case of modified dispersion relation $`\omega ^2=\mathrm{{\rm Y}}^2(k)`$ one should replace $`n^2`$ with
$$n_{mod}^2=a^2(\eta )\mathrm{{\rm Y}}^2(k)=a^2(\eta )\mathrm{{\rm Y}}^2\left(\frac{n}{a(\eta )}\right).$$
(24)
Therefore the equation we are to consider is
$$\mu _n^{\prime \prime }+\left(a^2(\eta )\mathrm{{\rm Y}}^2\left(\frac{n}{a(\eta )}\right)\frac{a^{\prime \prime }}{a}\right)\mu _n=0.$$
(25)
This equation should be appended by initial conditions at for $`\mu _n`$ and $`\mu _n^{}`$. In this paper we will consider two natural sets of initial conditions: the “minimal energy condition” of ,
$$\mu _n(\eta _i)=\sqrt{\frac{a(\eta _i)}{2}}\mathrm{{\rm Y}}^{1/2}\left(\frac{n}{2\pi a(\eta _i)}\right),$$
(26)
$$\mu _n^{}(\eta _i)=\pm i\sqrt{\frac{1}{2a(\eta _i)}}\mathrm{{\rm Y}}^{1/2}\left(\frac{n}{2\pi a(\eta _i)}\right);$$
(27)
and the instantaneous Minkowski vacuum conditions, as in the case of standard dispersion relation
$$\mu _n(\eta _i)=\frac{1}{\sqrt{2n}},\mu _n^{}(\eta _i)=\pm i\sqrt{\frac{n}{2}}.$$
(28)
Observe that in the case of standard dispersion relation, $`\mathrm{{\rm Y}}^2=n^2`$ these initial condition are identical.
Our goal would be therefore to solve eq. (25) with initial conditions (26, 27) or (28), and by making use of expression (21) to find the power spectrum of cosmological perturbation. We will consider two regimes: (I) when modifications of dispersion relations are relevant and $`\mathrm{{\rm Y}}`$ term dominates in eq. (25); (II) when the perturbation is large as compared to the scale defined by $`\kappa `$, in which case we have to do with standard dispersion relation and eq. (22) holds. With this equation we can compute modulus of the coefficient $`C_n`$.
Let us turn to the dispersion relation (13) which is valid in regime (I)
$$\mathrm{{\rm Y}}(k)=\kappa \mathrm{log}\left(1+\frac{k}{\kappa }\right).$$
(29)
In this regime, assuming $`a(\eta )=l_H/\eta `$ eq. (25) takes the form
$$\frac{d^2\mu _n^{(I)}}{d\eta ^2}+\left[\frac{4\pi ^2ϵ^2}{\eta ^2}\mathrm{log}^2\left(1+\frac{n\eta }{2\pi ϵ}\right)\frac{2}{\eta ^2}\right]\mu _n^{(I)}=0,$$
(30)
where we introduced the parameter $`ϵ=(\kappa l_H)`$ which is a ratio of two relevant length scales $`l_H`$, the size of the cosmological horizon and $`1/\kappa `$. If we assume that $`1/\kappa `$ is of order of Planck length, then $`ϵ`$ is equal to the size of the horizon expressed in Planck unit. The numerical value of this parameter depends therefore on details of the dynamics of inflation. Following we assume that $`ϵ10^6`$.
To solve equation (30) we make use of the fact that for equation $`\mu ^{\prime \prime }+W^2(\eta )\mu `$ the solution is of the approximate form
$$\mu =\frac{1}{\sqrt{2W}}\mathrm{exp}\left(\pm iW(\eta )𝑑\eta \right).$$
This solution is valid in the adiabatic regime, where
$$\frac{1}{2}\left(\frac{W^{\prime \prime }}{W^3}\frac{3}{2}\frac{W^{}^2}{W^4}\right)1.$$
(31)
In our case we get the condition
$$\frac{3n^2\eta ^2+4n\pi ϵ\eta \mathrm{log}(1+\frac{n\eta }{2\pi ϵ})+(2\pi ϵ+n\eta )^2\mathrm{log}(1+\frac{n\eta }{2\pi ϵ})^2}{16\pi ^2ϵ^2(2\pi ϵ+n\eta )^2\mathrm{log}(1+\frac{n\eta }{2\pi ϵ})^4}1.$$
In regime (I) the fluctuations are described therefore by
$$\mu _n^{(I)}=A_n^{(I)}\sqrt{\frac{\eta }{4\pi ϵ\mathrm{log}\left(1+\frac{n\eta }{2\pi ϵ}\right)}}\mathrm{exp}\left[2\pi iϵLi_2\left(\frac{n\eta }{2\pi ϵ}\right)\right]$$
$$+B_n^{(I)}\sqrt{\frac{\eta }{4\pi ϵ\mathrm{log}\left(1+\frac{n\eta }{2\pi ϵ}\right)}}\mathrm{exp}\left[+2\pi iϵLi_2\left(\frac{n\eta }{2\pi ϵ}\right)\right],$$
(32)
where $`Li_2(x)𝑑x\mathrm{log}(1x)/x`$ is the polylogarithm function.
Let us now turn to initial conditions (26, 27). At the time $`\eta =\eta _i`$ we find
$$\mu _n^{(I)}(\eta _i)=A_n^{(I)}\sqrt{\frac{\eta _i}{4\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}}\mathrm{exp}\left[2\pi iϵLi_2\left(\frac{n\eta _i}{2\pi ϵ}\right)\right]$$
$$+B_n^{(I)}\sqrt{\frac{\eta _i}{4\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}}\mathrm{exp}\left[+2\pi iϵLi_2\left(\frac{n\eta _i}{2\pi ϵ}\right)\right]=\sqrt{\frac{\eta _i}{4\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}}$$
(33)
and
$$\frac{1}{i}\mu _n^{(I)}{}_{}{}^{}(\eta _i)=A_n^{(I)}\sqrt{\frac{2\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}{2\eta _i}}\mathrm{exp}[2\pi iϵLi_2(\frac{n\eta _i}{2\pi ϵ})]$$
$$B_n^{(I)}\sqrt{\frac{2\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}{2\eta _i}}\mathrm{exp}\left[+2\pi iϵLi_2\left(\frac{n\eta _i}{2\pi ϵ}\right)\right]=\sqrt{\frac{2\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}{2\eta _i}},$$
(34)
Which leads to
$$A_n^{(I)}=\mathrm{exp}\left[2\pi iϵLi_2\left(\frac{n\eta _i}{2\pi ϵ}\right)\right],B_n^{(I)}=0,$$
(35)
and
$$\mu _n^{(I)}=\sqrt{\frac{\eta }{4\pi ϵ\mathrm{log}\left(1+\frac{n\eta }{2\pi ϵ}\right)}}\mathrm{exp}\left[2\pi iϵLi_2\left(\frac{n\eta _i}{2\pi ϵ}\right)2\pi iϵLi_2\left(\frac{n\eta }{2\pi ϵ}\right)\right].$$
(36)
In the case of linear standard vacuum initial conditions, (28) we similarly obtain
$$A_n^{(I)st}=\left(\frac{1}{\sqrt{8n}}\sqrt{\frac{2\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}{2\eta _i}}+\sqrt{\frac{n}{8}}\sqrt{\frac{\eta _i}{4\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}}\right)e^{2\pi iϵLi_2\left(\frac{n\eta _i}{2\pi ϵ}\right)}$$
(37)
$$B_n^{(I)st}=\left(\frac{1}{\sqrt{8n}}\sqrt{\frac{2\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}{2\eta _i}}\sqrt{\frac{n}{8}}\sqrt{\frac{\eta _i}{4\pi ϵ\mathrm{log}\left(1+\frac{n\eta _i}{2\pi ϵ}\right)}}\right)e^{2\pi iϵLi_2\left(\frac{n\eta _i}{2\pi ϵ}\right)}$$
(38)
Now we turn to regime (II), and we have to match above solution with $`\mu _n^{(II)}`$ given by (22) at $`\eta =\eta _1`$. Before doing that, let us pause for a moment to find the characteristic value of the conformal time $`\eta _1`$. It corresponds to the moment when the wavelength of the perturbation equals the characteristic length $`\kappa ^1`$. Thus
$$\eta _1=\frac{2\pi ϵ}{n}.$$
(39)
It must be checked if at this value of conformal time we are still in the region of validity of adiabatic approximation, i.e., if condition (31) is satisfied. One finds the condition
$$\frac{3+4\mathrm{log}(2)^2+\mathrm{log}(4)}{32\pi ^2ϵ^2\mathrm{log}(2)^4}1$$
which indeed holds for any $`ϵ10^1`$ and larger. On the other hand we must also satisfy the condition
$$4\pi ^2ϵ^2\mathrm{log}(2)2$$
so that we can safely assume that fluctuations do not feel the cosmic expansion and neglecting the factor $`(2/\eta ^2)`$ in eq. (30) was indeed justified. This condition is well satisfied for $`ϵ`$ slightly larger than $`1`$. We see therefore that our approximate solution holds at the matching point $`\eta =\eta _1`$. Physically this means that for such values of the parameter $`ϵ`$, the fluctuation relaxes sufficiently slowly from the regime ruled by the trans-Planckian physics and at the same time reaches the regime described by standard physics well before the fluctuation length becomes comparable with the Hubble size. This analysis of the adiabaticity conditions is in perfect agreement with the results of papers , .
Let us thus match the solution (36) with the solution (22). In view of eq. (23) we will be mainly interested in the $`n`$ dependence of the coefficients. Computing $`\mu _n^{(I)}(\eta _1)`$ and its derivative and comparing with appropriate expressions obtained from (23) we easily find
$$\alpha _n\frac{1}{\sqrt{n}}e^{2\pi iϵLi_2\left(\frac{n\eta _i}{2\pi ϵ}\right)},\beta _n\frac{1}{\sqrt{n}}e^{2\pi iϵLi_2\left(\frac{n\eta _i}{2\pi ϵ}\right)},$$
(40)
both multiplied by complicated numerical ($`n`$-independent) coefficients. This means that
$$|C_n|n^{3/2}$$
and we recover the flat spectrum of the fluctuations.
For instantaneous Minkowski vacuum initial conditions (28), the final result changes. One finds
$$\alpha _n,\beta _nconst+\frac{1}{n},$$
and thus one cannot obtain the scale invariant spectrum. This corresponds to the result of Martin and Brandenberger who also concluded that these standard initial conditions do not lead to the correct spectrum of fluctuations.
This result is not very surprising. The initial conditions (28) are natural in the case of standard dispersion relation. However, it would be hardly understandable if they remain valid in the case of modified nonlinear dispersion. The choice of the minimal energy conditions of Martin and Brandenberger (26, 27), on the other hand, is justified by powerful physical principle that the fluctuations leading to large scale structure originate from vacuum fluctuations with minimal possible energy.
This result is the first test passed by $`\kappa `$ physics. This is important because till now $`\kappa `$-Poincaré symmetry was nothing but a nice algebraic construction. I hope that the result reported in this paper would encourages one to look for new grounds where predictions of modified dispersion relations might be possibly checked (for a very recent attempt along this line, see .)
Some open problems remain, of course. Probably the most important one is what is the source of “trans-Plackian frequencies reservoir”. In other words one should understand the initial conditions (26, 27) from the point of view of fundamental $`\kappa `$-quantum field theory. It is also interesting to see imprints of $`\kappa `$-physics on early stages of inflation, especially in the pictures where the inflation starts immediately after Planck era.
Acknowledgement. I would like to thank Professor Jerzy Lukierski for enlightening discussions concerning $`\kappa `$ physics, and to the anonymous referee, whose comments made it possible to improve this paper greatly.
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# Topology of Diophantine Sets: Remarks on Mazur’s Conjectures
## 1. Introduction
One of the main themes in model theory is to understand the structure of definable sets: given a first-order language $`L`$ and an $`L`$-structure $`M`$, describe the $`L`$-definable subsets of $`M^n`$ for various $`n𝐙_{>0}`$. Here, a set $`SM^n`$ is called $`L`$-definable if there exists an $`L`$-formula $`\psi (𝐱)`$ with free variables $`𝐱=(x_1,\mathrm{},x_n)`$ such that for any $`𝐚M^n`$, $`𝐚SM\psi (𝐚)`$. A set is called existentially definable (respectively, positive existentially, or diophantine) if $`\psi (𝐱)`$ can be taken to be $`𝐛\varphi (𝐱,𝐛)`$, with $`\varphi `$ quantifier-free (respectively, quantifier- and negation-free, or atomic).
The natural geometric examples of such structures arise as in the following definition:
###### Definition 1.1.
If $`R`$ is a commutative ring with unit, it admits a natural interpretation for any first order language $`L`$ of the form $`L_R=(+,,=,c_i)`$ where $`c_i`$ are primary predicates (“constants”), less in number than $`|R|`$. We call $`L_R`$ a *ring language*. We define $`L_𝐙=(+,,0,1)`$ and $`L_t=(+,,0,1,t)`$ for any $`tR`$.
###### Example 1.2.
(Tarski, cf. , pp. 202–206) (a) An algebraically closed field $`k`$ admits elimination of quantifiers in the language $`L_𝐙`$. Hence any $`L_𝐙`$-definable subset in $`k^n`$ is a boolean combination of sets defined by an equation. Thus, the definable sets for an algebraically closed field are exactly the classical sets of algebraic geometry – one deduces for example that the only definable subsets of $`k`$ are finite or cofinite, a fact which at first sight seems not so obvious.
(b) The field of real numbers $`𝐑`$ admits elimination of quantifiers in the language $`L_{}=(0,1,+,,)`$ of ordered fields. Hence every definable set in $`𝐑^n`$ is a boolean combination of semi-algebraic sets (i.e., solution sets to systems of equations of the form $`f(𝐱)=0g(𝐱)0`$). This gives a nice description of the definable subsets of $`𝐑`$: they are finite unions of intervals.
(c) More examples in the same vein exist, e.g., a description of definable sets over $`p`$-adic fields (, ), or generalization of $`(𝐑,L_{})`$ via o-minimal expansions.
(d) To give an example with a different language, existentially definable sets of $`𝐙`$ in the language $`(0,1,+,|)`$ are unions of arithmetic progressions (a result of Lipshitz ).
The moral is that if the (existentially) definable sets for such $`M`$ have a sufficiently easy description, then the first-order (respectively, existential) theory of $`M`$ is decidable – this is the case in the above examples. Conversely, if definable sets are combinatorially complicated, one expects the corresponding theory to be undecidable.
###### Example 1.3.
(a) Consider the rational integers $`(𝐙,L_𝐙)`$. It is impossible to describe the $`L_𝐙`$-definable sets of $`𝐙`$ in terms of “classical” sets (e.g., finite sets, arithmetic progressions, …). Eventually, this leads to the undecidability of the full theory of $`(𝐙,L_𝐙)`$.
(b) The celebrated theorem of Davis, Matijasevich, Putnam and Robinson describes the existentially definable sets of $`𝐙`$: they are exactly the recursively enumerable sets, whose complexity outranges by far that of decidable (hence, certainly, of computable) sets – and the undecidability of the existential theory of $`(𝐙,L_𝐙)`$ follows.
This maxim, the interplay between (un)decidability and definable sets, applies in particular to the field $`(𝐐,L_𝐙)`$ of rational numbers. The field structure of $`𝐐`$ admits the same kind of “wild” definable sets as the integers; this follows from J. Robinson’s theorem that $`𝐙`$ is a definable subset of $`𝐐`$ (, theorem 3.1). The question whether the same can happen if we restrict to the existentially definable sets is still open.
In the next paragraph, we will present a conjecture by Mazur, which – although it does not characterize the existentially definable sets of $`𝐐`$ – poses severe restrictions on their real topological structure. In the subsequent section, we prove that this conjecture implies there is no “diophantine model” (cf. infra) of $`(𝐙,L_𝐙)`$ in $`(𝐐,L_𝐙)`$ – this generalizes Mazur’s observation that his conjecture implies that $`𝐙`$ is not an $`L_𝐙`$-diophantine subset of $`𝐐`$. In particular, any proof of the diophantine undecidability of $`𝐐`$ “along traditional lines” fails if Mazur’s conjecture is true. In the final paragraph, we comment upon a non-archimedean version of this conjecture. Though most of these observations are folklore, they do not seen to have been written down previously.
## 2. Mazur’s conjectures
In , and , Barry Mazur has proposed and discussed several conjectures and questions about the behaviour of the set of $`𝐐`$-rational points of a variety over $`𝐐`$ under taking topological closure w.r.t. some metric induced by a valuation on $`𝐐`$. The conjecture that we will concentrate upon (the weakest) is the following:
###### Conjecture 2.1.
*(Mazur , Conjecture 3)* For any variety $`V`$ over $`𝐐`$, the (real) topological closure of $`V(𝐐)`$ in $`V(𝐑)`$ has only a finite number of real topological components.
There is some evidence for this conjecture, especially for such $`V`$ which possess special geometric properties (mostly related to the canonical class of $`V`$) – and no counterexample to it is known. Also observe that, with $`𝐐`$ replaced by $`𝐑`$ in 2.1, the “conjecture” says that a real variety has only finitely many real connected components. This holds true; it could be deduced from Tarski’s results – there is even an explicit bound on the betti numbers of $`V(𝐑)`$, the so-called Milnor-Thom theorem (cf. ).
###### Example 2.2.
(a) Conjecture 2.1 is true for curves $`V`$. One can assume $`V`$ to be projective and non-singular. The case where $`V`$ has genus $`g2`$ is settled by Faltings’s theorem, which says that $`V(𝐐)`$ is a finite set (). If $`V`$ has genus 0, then either $`V(𝐐)`$ is empty, or $`V`$ is $`𝐐`$-birational to $`𝐀^1`$, and $`𝐀^1(𝐐)`$ is topologically dense in $`𝐀^1(𝐑)`$. Finally, assume that $`V`$ has genus 1. It is known that $`V(𝐑)`$ is isomorphic to the “circle group” $`𝐑/𝐙`$ or to $`𝐑/𝐙\times 𝐙/2`$ (see , V). Every proper closed subgroup of the circle group is finite (see , theorem 1.34). Hence, if $`V(𝐐)`$ is not finite, then it is dense in every component of $`V(𝐑)`$ that it intersects.
(b) To provide a higher dimensional example, let $`V`$ be a variety satisfying weak approximation (i.e., such that $`V(𝐐)V(𝐐_p)`$ is dense). Then the conjecture holds true for $`V`$. This holds, e.g., if $`V`$ is a smooth complete intersection of two quadrics in projective space of dimension at least 5 (cf. ).
###### Remark 2.3.
Mazur has made even stronger conjectures, some of which had to be slightly modified, due to the construction of a counterexample by Colliot-Thélène, Skorobogatov and Swinnerton-Dyer (). For an extensive (unsurpassable) exposition and more examples, we refer to the original sources , and . We will concentrate on the model-theoretical aspects of the conjectures, which are already present in 2.1 – but let the reader be warned about making too bold generalizations of 2.1. A non-archimedean version will be considered in the last paragraph of this paper.
###### Remark 2.4.
The $`(𝐐,L_𝐙)`$-existentially definable subsets, in the sense of the introduction, are precisely images of projections from $`V(𝐐)`$ to affine space $`𝐀_𝐐^n`$ for various $`V`$ and $`n`$.
A more model-theoretic conjecture would be that *the real topological closure of a $`(𝐐,L_𝐙)`$-existentially definable set is an $`(𝐑,L_{})`$-definable set* (i.e., a semi-algebraic set). This implies 2.1, since a semi-algebraic set has only finitely many components.
We do not know whether conjecture 2.1 is equivalent to this statement. Note that J. Robinson’s argument (in ) shows that it is wrong when the word “existentially” is erased.
## 3. Diophantine models of $`𝐙`$ in $`𝐐`$
Mazur has observed that conjecture 2.1 implies that $`𝐙`$ is not diophantine in $`𝐐`$ in the language $`L_𝐙`$; indeed, if $`𝐙`$ would arise as the projection of $`V(𝐐)`$ for some variety $`V`$, then since $`𝐙`$ has infinitely many real components and the projection is continuous, the same would hold for $`V(𝐑)`$.
However, many proofs of the undecidability of the diophantine theory of structures $`(R,L_R)`$ as in (1.1) do not give that $`𝐙`$ is a diophantine *subset* of $`R`$, but rather produce a *diophantine model* of $`(𝐙,L_𝐙)`$ in $`(R,L_R)`$, in the sense of the following definition:
###### Definition 3.1.
A model $`(M,L,\varphi )`$ is a triple consisting of a first order language $`L`$ which consists of $`i`$-ary predicates $`\{P_{i,\alpha }\}`$, a set $`M`$ and an interpretation $`\varphi `$ of $`L`$ in $`M`$ (we will often leave out $`\varphi `$ of the notation). Note that any cartesian power $`M^k,(k1)`$ is likewise a model for $`L`$ via “diagonal interpretation”.
We say that a model $`(M^{},L^{}=\{P_{i,\alpha }^{}\},\varphi ^{})`$ admits a *diophantine model* in $`(M,L,\varphi )`$ if there exists a set-theoretical bijection between $`M^{}`$ and a subset of some cartesian power $`M^k(k1)`$, such that the image is diophantine, and such that the induced inclusions $`\varphi ^{}(P_{i,\alpha }^{})M^{ik}`$ are diophantine.
A similar notion of *(positive) existential model* exists.
###### Example 3.2.
(a) If $`(M^2,L)`$ admits a diophantine model in $`(M,L)`$, then the latter structure is said to admit *diophantine storing* (cf. ). This is true, for example, for $`(𝐙,L_𝐙)`$. For non-algebraically closed rings $`(R,L_R)`$ admitting diophantine storing, one can always choose $`k=1`$ in the above definition. For if $`(M^{},L^{},\varphi ^{})`$ admits a diophantine model in $`(R^2,L_R)`$ and $`(R,L_R)`$ admits diophantine storing, then $`(M^{},L^{},\varphi ^{})`$ admits a diophantine model in $`(R,L_R)`$ (since conjuntions of diophantine formulas are again diophantine if the quotient field of $`R`$ is not algebraically closed – cf. , §3).
(b) Typically, diophantine models of the integers $`(𝐙,L_𝐙)`$ in ring languages $`(R,L_R)`$ arise in the following way: a commutative algebraic group $`G`$ (e.g., the multiplicative group of a quadratic ring, or an elliptic curve) is assumed to have rank one over $`R`$, and the set $`𝐙`$ has a diophantine model as the $`R`$-rational points $`G(R)`$ of $`G`$ \- the relation “addition” is automatically mapped to a diophantine subset of $`G^3(R)`$, since the group law on $`G`$ is a morphism. The most problematic point is defining the relation “multiplication”. For an example, consider the proof that $`(𝐙,L_𝐙)`$ admits a diophantine model in $`(R:=S[t],L_t)`$ for any commutative unitary domain $`S`$ of characteristic zero, see Denef . He takes for $`G`$ the torus $`𝐆_{m,R[\sqrt{\mathrm{\Delta }}]}`$ of discriminant $`\mathrm{\Delta }=t^21`$, which is non-split over $`R`$; $`G(R)`$ has rank one: any $`R`$-point is given by a solution $`(x_n,y_n)`$ to the Pell-equation $`X^2\mathrm{\Delta }Y^2=1`$ (i.e., a power $`u^n=x_n+y_n\sqrt{\mathrm{\Delta }}`$ of the fundamental unit $`u=t+\sqrt{\mathrm{\Delta }}`$). Multiplication $`(x_r,y_r)(x_s,y_s)=(x_n,y_n)`$ is defined by saying that $`f:=y_ny_ry_s`$ has a zero at $`t=1`$, i.e., $`(hR)(f=(t1)h).`$
(c) It is not known whether, if the ring $`R`$ contains $`𝐙`$, the set $`𝐙`$ itself is a diophantine subset of $`R`$ whenever $`(𝐙,L_𝐙)`$ admits a diophantine model in $`(R,L_𝐙)`$.
The following result formalizes the technique of proof of many undecidability results:
###### Observation 3.3.
Assume that $`R`$ is as in 1.1, and there is a polynomial whose coefficients belong to $`\varphi (L)`$ but that has no zero in the fraction field of $`R`$. If $`(M^{},L^{})`$ has an undecidable diophantine theory, and admits a diophantine model in $`(R,L_R)`$, then the diophantine theory of $`(R,L_R)`$ is undecidable. ∎
###### Remark 3.4.
Without any restrictions on $`R`$, if $`(M^{},L^{})`$ has an undecidable (positive) existential theory and admits a (positive) existential model in $`(R,L_R)`$, then the (positive) existential theory of $`(R,L_R)`$ is undecidable.
The technique of many undecidability proofs for rings $`(R,L_R)`$ is based on the fact that, via a construction as in (3.2(b)), one can find a diophantine model of the integers $`(𝐙,L_𝐙)`$ in $`(R,L_R)`$, and then rely on the fact that the diophantine theory of the integers is undecidable (, ). It has been suggested that, with this more flexible definition, one would be able to find a diophantine model of the integers in the rationals:
###### Question 3.5.
Does $`(𝐙,L_𝐙)`$ admit a diophantine model in $`(𝐐,L_𝐙)`$?
However, even this is impossible if we assume Mazur’s conjecture:
###### Theorem 3.6.
Mazur’s conjecture 2.1 implies that there is no diophantine model of $`(𝐙,L_𝐙)`$ in $`(𝐐,L_𝐙)`$.
###### Proof.
Assume that there is such a diophantine model $`(D,L_D)`$, with $`D𝐐^k`$. Then there is an affine variety $`V`$ over $`𝐐`$ admitting a finite morphism $`f:V_𝐐𝐀_𝐐^k`$ defined over $`𝐐`$ such that $`f(V(𝐐))=D`$.
If $`D`$ is discrete (i.e., infinite and totally disconnected in the real topology), the traditional proof applies: the real topological closure of $`V(𝐐)`$ in $`V(𝐑)`$ is also mapped to $`D`$ by $`f`$, and hence it has infinitely many components.
If $`D`$ is not discrete (which seems to be the case for the typical infinite diophantine set in $`𝐐`$, say, the set of squares), then we show that one can select (in a computable way) a discrete subset $`\stackrel{~}{D}`$ of $`D`$. Then the above proof, applied to $`\stackrel{~}{D}`$, gives the result.
Here are the details of the construction of $`\stackrel{~}{D}`$. We only have to treat the case where the real topological closure $`\overline{V}`$ of $`V(𝐐)`$ has only finitely many connected components. Since $`f`$ is continuous, the mean value theorem implies that $`f(\overline{V})`$ is the union of finitely many closed subsets in $`𝐑^k`$. In particular, the topological closure $`\overline{D}`$ of $`D`$ contains finitely many closed subsets, and since $`D`$ is infinite, one of these subsets, say, $`D_0`$, is not a point. By composing $`f`$ with a suitable $`𝐐`$-rational projection $`\pi :𝐀_𝐐^k𝐀_𝐐^1`$ which does not map $`D_0`$ to a point, we may assume $`k=1`$. By composing with a fractional linear transformation defined over $`𝐐`$, we may assume $`\pi (D_0)`$ to be the unit interval $`I=[0,1]`$. Let $`d_n`$ be the element of $`d`$ corresponding to $`n𝐙`$. Let us consider the set
$$\stackrel{~}{Z}=\{n𝐙|\frac{1}{2j+1}\pi (d_n)\frac{1}{2j}\text{ for some }j𝐙_{>0}\}.$$
The set $`\stackrel{~}{Z}`$ is Turing computable (since $`D=\{\pi (d_n)\}`$ is a listable subset of $`𝐐`$, it is easy to write a Turing program to check the inequalities), hence it is recursively enumerable (by Kleene’s normal form theorem, cf. , 2.3-2.4), so by , it is diophantine in $`(𝐙,L_𝐙)`$. Also, $`\stackrel{~}{Z}`$ is infinite, since $`\pi (D)I`$ is dense in $`I`$. We now set
$$\stackrel{~}{D}=\{d_n|n\stackrel{~}{Z}\}.$$
By construction, the set $`\stackrel{~}{D}`$ is diophantine in $`(D,L_D)`$, and hence a fortiori in $`(𝐐,L_𝐙)`$. So there exists a variety $`\stackrel{~}{V}`$ over $`𝐐`$ and a $`𝐐`$-morphism $`\stackrel{~}{f}:\stackrel{~}{V}𝐀_𝐐^1`$ such that $`\stackrel{~}{f}(\stackrel{~}{V}(𝐐))=\stackrel{~}{D}`$. However, the real closure of the set $`\stackrel{~}{D}`$ has infinitely many connected components in the real topology by construction. Hence the same holds for $`\stackrel{~}{V}(𝐐)`$, contradicting Mazur’s conjecture. ∎
## 4. Non-archimedean aspects of Mazur’s conjectures
In (, II.2), Mazur has devised a conjecture of the above type which applies to any completion of a number field, not just an archimedean one. As it makes sense for any global field, we formulate it as follows:
###### Question 4.1.
Let $`V`$ be a variety over a global field $`K`$, $`v`$ a valuation on $`K`$, and $`K_v`$ the completion of $`K`$ w.r.t. $`v`$. For every point $`pV(K_v)`$, let $`W(p)`$ be the Zariski closure of $`(V(K)U)`$, where $`U`$ ranges over all $`v`$-open neighbourhoods of $`p`$ in $`V(K_v)`$. Is the set $`\{W(p):pV(K_v)\}`$ finite?
In our next theorem, we observe that the answer to this question is negative in positive characteristic:
###### Theorem 4.2.
Let $`K=𝐅_q(t)`$ be the rational function field over a finite field $`𝐅_q`$ of positive characteristic $`p>0`$, and let $`v`$ be the valuation corresponding to the place $`t^1`$ of $`K`$ (i.e., $`v(a)=q^{\mathrm{deg}(a)}`$ for $`a𝐅_q[t]`$). Then there is a variety $`V`$ over $`K`$ for which the answer to question 4.1 is negative.
###### Proof.
In (lemma 1) Pheidas proves that, for $`p>2`$, projection onto the $`x`$-coordinate of the $`K`$-rational points of the space curve $`V_p`$ given by
$$V_p:xt=u^pu,x^1t^1=v^pv$$
gives the set $`D_p=\{t^{p^s}|s𝐙_0\}.`$ For $`p=2`$, Videla () proved that the set $`D_2`$ is the projection onto the $`x`$-coordinate of
$$V_2:x+t=u^2+u,u=w^2+t,x^1+t^1=v^2+v,v=s^2+t^1.$$
Already the sets $`W(p)`$ for $`pV(K)`$ are disjoint, since their $`x`$-coordinates are separated ($`v(t^{p^s}t^{p^r})>1`$ for all $`rs`$). This gives a negative answer to question 4.1. ∎
Thinking of the analogy between function fields and number fields, one can ask for the strict analogue of question 3.5 for global function fields.The answer to it is *positive*:
###### Theorem 4.3.
For any prime power $`q`$, $`q=p^n`$, $`p>0`$, the polynomial ring $`(𝐅_q[t],L_t)`$ admits a diophantine model in the ring of rational functions $`(𝐅_q(t),L_t)`$.
###### Proof.
The proof is a bit indirect: we show that the polynomial ring has a diophantine model in the positive rational integers, and the latter has a diophantine model in the field of rational functions.
More precisely, $`𝐅_q[t]`$ is a recursive ring (cf. Rabin ), because $`𝐅_q`$ is recursive (since finite), and hence the same holds for the polynomial ring over $`𝐅_q`$ (cf. Fröhlich and Sheperdson ). So there exists an injective map $`\theta :𝐅_q[t]𝐙_0`$ such that the graphs of addition and multiplication are recursive on $`𝐙_0`$, and hence $`(\theta (𝐅_q[t]),\theta (L_t))`$ is a diophantine model of $`(𝐅_q[t],L_t)`$ in $`(𝐙_0,L_𝐙)`$.
For the second step, we first recall a construction of Pheidas and Videla (, ). Let $`v`$ denote the $`t`$-valuation on $`𝐅_q(t)`$, i.e., $`v(x)`$ is the order of $`x`$ at zero. For any $`k𝐙_0`$, let $`[k]`$ denote the equivalence class of elements $`x𝐅_q(t)`$ with $`v(x)=k`$. For positive integers $`a`$ and $`b`$, let $`a|_pb`$ denote the relation $`(n𝐙_0)(a=bp^n)`$.
Consider the structure $`S=(𝐙_0,(+,|_p,0,1))`$. Firstly, multiplication is diophantine in $`S`$ (, corollary on p. 529). Secondly, the set of equivalence classes $`[k]`$ as above is a model for $`S`$ in which the relations in $`S`$ can be defined by diophantine formulas in $`(𝐅_q(t),L_t)`$ between arbitrary representatives of the equivalence classes in $`𝐅_q(t)`$. We conclude that for arbitrary elements $`x,y,z𝐅_q(t)`$ the relations $`[v(x)]=[v(y)+v(z)]`$ and $`[v(x)]=[v(y)v(z)]`$ are diophantine in $`𝐅_q(t)`$.
The problem with this encoding is that we do not know the existence of a diophantine set in $`𝐅_q(t)`$ which contains exactly one representative for each such equivalence class. We fix this problem as follows. We know from the proof of theorem 4.2 that the set $`D_p=\{t^{p^k},k𝐙_0\}`$ is diophantine in $`(𝐅_q(t),L_t)`$, and this will be our model.
To define addition and multiplication on elements of this set, we introduce the following switching between $`t^{p^k}`$ and $`[k]`$: the set $`\{(k,p^k),k𝐙_0\}`$ is recursively enumerable in $`𝐙_0^2`$, so by Matijasevich’s theorem, it is diophantine in $`𝐙_0`$. Then, by the aforementioned results, the set $`=\{([k],[p^k]),k𝐙_0\}`$ is diophantine over $`(𝐅_q(t),L_t)`$.
For the function symbols $`R\{+,\}`$ on $`𝐙_0`$, we let the corresponding symbol $`\stackrel{~}{R}`$ for $`x,y,zD_p`$ be defined by
$`z=x\stackrel{~}{R}y`$ $`(x_1,y_1,z_1𝐅_q(t))(((x_1,x),(y_1,y),(z_1,z))^3[R(x_1,y_1)]=[z_1]).`$
For $`R\{+,\}`$, the righthand side of the equivalence is diophantine in $`(𝐅_q(t),L_t)`$ by what we have said before and the fact that for any two elements $`w_1,w_2𝐅_q(t)`$, the statement $`[w_1]=[w_2]`$ is equivalent with $`(v(w_1/w_2)0)(v(w_2/w_1)0)`$, which is diophantine by and .
Finally, $`(D_p,\stackrel{~}{+},\stackrel{~}{},t,t^p)`$ is a diophantine model of $`(𝐙,L_𝐙)`$ in $`(𝐅_q(t),L_t)`$. This finishes the proof of the theorem. ∎
Of course, the above theorem still does not settle the following problem:
###### Question 4.4.
Is the polynomial ring $`𝐅_q[t]`$ a diophantine *subset* of the field of rational functions $`𝐅_q(t)`$?
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# Schemata as Building Blocks: Does Size Matter?
## 1 Introduction
A very large proportion of scientific endeavour has been associated with the question: What are “things” made of? the reason being that this is an indispensable requirement for understanding how and why a “thing” functions. The answer has always tended to be: “things” are made of other, more elementary “things”. In the physical sciences this is obvious: a table is made up of atoms, which in turn are made up of electrons and a nucleus, which in its turn … In biology a living organism (generally) is composed of various organs and tissues, which in their turn are made of various types of cells, which in their turn are composed of various constituents such as nuclei, protoplasm, mitochondria etc. From the nucleus we pass to chromosomes, genes, DNA and RNA etc. The latter, along with other important ingredients, forming an elaborate “computer programme” for the construction of the organism. In computer science high level languages are composed of more elementary languages until we arrive at the most basic machine level language recognized by the computer itself.
What have all these examples in common? They show that all things are made out of “building blocks”, whether they be tables, giraffes or computer programmes. Inevitably there exists a hierarchy of building blocks, the hierarchy being ordered more often than not according to scale and complexity. One can think of building blocks as the “effective degrees of freedom” (EDOF) of a system, which in their turn are composed of more “fine grained”, elementary degrees of freedom. For complex systems the former are composed of very large numbers of the latter. What building blocks one uses to describe a system depends very much on what one wants to say about it. In particular, on how fine grained a description one requires. Almost always a more coarse grained description will suffice.
So what has the above to do with genetic algorithms (GAs)? At the most basic level all the above can be coded as bit strings (of course) and in some way or other be associated with the notion of adaptation and optimization of some “fitness” function in a complicated environment. Trying to understand these problems in adaptation and optimization at the level of the fundamental, microscopic degrees of freedom is prohibitively difficult: there are simply too many and, more often than not, they interact in a highly non-linear fashion. In order to describe these systems both qualitatively and quantitatively one needs to know how EDOF emerge. In particular, in the context of GAs, if we coded the above problems as such how would some of the known EDOF emerge? In GAs one can universally represent EDOF as schemata. Of course, not all are of equal utility. EDOF if they are to be useful in the description of a system must display a certain degree of integrity.
The question is: what schemata are utilized by a GA? Or rather, what are the typical properties possessed by successful schemata. Naturally, the answer to this question will depend on the fitness landscape of interest. However, one might enquire as to if or not such properties exist in generic classes of landscapes. In fact theory has tried to be even more ambitious. The schema theorem and the related building block hypothesis , propose that the EDOF in GAs are fit, short schemata irrespective of the landscape! This is an extremely strong statement. The fact that fit schemata are preferred is intuitively understandable, though we will see some counterexamples to this later, whilst the purported preference for short schemata is a supposed consequence of the destructive effect of crossover.
In this paper we will investigate theoretically and experimentally the evolution of schemata and in particular how this evolution depends on the defining length of the schemata. Our theoretical analysis will be based on an exact evolution equation for schemata for the case of proportional selection and $`1`$-point crossover. In section 2 we will give a brief overview of the equation. In section 3 we will discuss some of its more important theoretical ramifications and in section 4 we will show how numerical experiments confirm the theoretical predictions. Finally in section 5 we will draw some conclusions.
## 2 Schema Equation
In this section we give without proof (see , for more details) the schema evolution equation for the relative proportion, $`P(\xi ,t)n(\xi ,t)/n`$, of the schema $`\xi `$ of defining length $`l`$ and order $`N_2`$ at time $`t`$ in a canonical GA of population $`n`$ consisting of chromosomes of $`N`$ bits evolving with respect to proportional selection, point mutation and $`1`$-point crossover. In the limit $`n\mathrm{}`$ $`P(\xi ,t)`$ gives the probability of finding the schema $`\xi `$ at time $`t`$. Explicitly
$`P(\xi ,t+1)=𝒫(\xi )P_c(\xi ,t)+{\displaystyle \underset{\xi /_i}{}}𝒫(\xi /_i\xi )P_c(\xi /_i,t)`$ (1)
where the sum is over all schemata, $`\xi /_i`$, that differ by at least one bit in one of the defining bits of the schema $`\xi `$, and the effective mutation coefficients $`𝒫(\xi )`$ and $`𝒫(\xi /_i\xi )`$ represent the probabilities that the schema $`\xi `$ remains unmutated and the probability that the schema $`\xi /_i`$ mutates to the schema $`\xi `$ respectively. $`P_c(\xi ,t)`$ is the mean proportion of schemata $`\xi `$ at time $`t`$ after selection and crossover. Explicitly
$`P_c(\xi ,t)=P^{}(\xi ,t){\displaystyle \frac{p_c}{N1}}{\displaystyle \underset{k=1}{\overset{l1}{}}}(P^{}(\xi ,t)P^{}(\xi _l(k),t)P^{}(\xi _R(lk))`$ (2)
where $`P^{}(\xi ,t)=(\overline{f}(\xi ,t)/\overline{f}(t))P(\xi ,t)`$, $`\overline{f}(\xi ,t)`$ being the mean fitness of the schema $`\xi `$ and $`\overline{f}(t)`$ the average population fitness. $`p_c`$ is the crossover probability and $`k`$ the crossover point. The quantities $`P^{}(\xi _L,t)`$ and $`P^{}(\xi _R,t)`$ are defined analogously to $`P^{}(\xi ,t)`$ but refer to the schemata $`\xi _L`$ and $`\xi _R`$ which are the parts of $`\xi `$ to the left and right of $`k`$ respectively. One can illustrate the content of the equation with a simple diagramatic example: $`1|01`$ is a schema with $`l=7`$ and $`N_2=3`$. The crossover point is at $`k=7`$ hence $`\xi _L`$ has $`N_2=1`$ and $`l=1`$ while $`\xi _R`$ has $`N_2=2`$ and $`l=4`$. Note that the equation takes into account exactly both the effects of schema destruction and schema reconstruction. At the level of strings the equation will be equivalent to other exact formulations , and is most closely related to the analogous equation for a canonical GA evolving with respect to proportional selection and recombination derived by Altenberg based on earlier work in genetics by Karlin and Liberman .
There are several notable features of the above equation: first of all it implies that crossover as an operator imposes the idea of a schema. This can easily be seen by considering the above equation for the case where $`\xi `$ is the entire string, $`C_i`$. The reconstruction probability depends on the relative fitness of strings that contain the constituent elements, $`C_i^L`$ and $`C_i^R`$ of $`C_i`$, but given that there can be many strings that contain $`C_i^L`$ or $`C_i^R`$ one must take an average over these strings. In this sense we are averaging over the “degrees of freedom” represented by the bits that are not contained in $`C_i^L`$ or $`C_i^R`$. The equation also shows that the effects of reconstruction will outweigh destruction if the parts of a string are more selected than the whole. This is closely related to the notion of linkage disequilibrium from population genetics. However, linkage disequilibrium there is measured by the covariance, $`C(\xi _L,\xi _R)P(\xi ,t)P(\xi _L,t)P(\xi _R,t)`$, of the relative frequencies of $`\xi _L`$ and $`\xi _R`$. Here, we see that the relevant measure of schemata growth is the covariance, $`C^{}(\xi _L,\xi _R)P^{}(\xi ,t)P^{}(\xi _L,t)P^{}(\xi _R,t)`$, of the fitness-weighted relative frequencies of $`\xi _L`$ and $`\xi _R`$. Clearly $`C^{}`$ does not have to be of the same sign as $`C`$.
The equation also shows an hierarchy in structure both in complexity and size and also in time. This is obvious in the very nature of the equation which relates a schema $`\xi `$ to its “building blocks” of lower order and smaller defining length, $`\xi _L`$ and $`\xi _R`$. These in turn are related to even smaller, lower order building blocks $`\xi _{LL}`$, $`\xi _{LR}`$, $`\xi _{RL}`$ and $`\xi _{RR}`$ which are associated with reconstruction of the schemata $`\xi _L`$ and $`\xi _R`$. The hierarchy terminates at $`1`$-schemata, i.e. schemata with $`N_2=1`$, which are immune to the effects of crossover. The hierarchical nature of the evolution in terms of complexity and size is manifest in (2) as schemata of a certain “size”’ (defining length and order) are related to schemata of smaller size, which in their turn are related to yet smaller schemata etc. The structure is also temporal as schemata at time $`t`$ are related to smaller schemata at time $`t1`$, which in turn are related to yet smaller schemata at time $`t2`$ etc.
Note the “form invariance” of the equation under a coarse graining. What does this mean? Consider the equation at the level of a complete string. As we have pointed out, the very notion of recombination introduces a coarse graining in that we can write the recombination contribution in terms of $`P(C_i^L,t)`$ and $`P(C_i^R,t)`$, which involve summing over the microscopic (bit) degrees of freedom of $`\{C_iC_i^L\}`$ and $`\{C_iC_i^R\}`$ respectively, where $`\{\}`$ denotes set difference. Coarse graining here simply means that we have forfeited detailed knowledge about the microscopic degrees of freedom not contained in $`C_i^L`$ or $`C_i^R`$. The corresponding evolution equations for $`C_i^L`$ and $`C_i^R`$ involve recombination terms where a further coarse graining must be carried out via a summation, for example, over the microscopic degrees of freedom of $`\{C_iC_i^{LL}\}`$ in the case of the schema $`C_i^{LL}`$. However, irrespective of the degree of coarse graining the form of the evolution equation remains exactly the same.
We see then that a general form of the building block hypothesis is inherent in the very structure of the evolution equation for a canonical GA. Recombination builds complex schemata from more primitive consituents which in turn are constructed from yet more elementary building blocks until we arrive at the ultimate building blocks — $`1`$-schemata. The question of whether a GA utilizes building blocks to find a good solution can be seen to be related to whether schema reconstruction or destruction is the most important effect. However, this is clearly not the only criterion. We have said that the evolution equation contains a generalized form of the block hypothesis: simply that larger, more complex schemata are constructed from more primitive building blocks irrespective of whether they are fit schemata utilized by the GA in finding fit chromosomes. This is not however the only way to grow a schemata. Perhaps the schemata was in the initial population and was of high fitness. To see whether indeed a GA uses fit, short building blocks as the standard building block hypothesis and Schema theorem purport we must examine more closely the idea of fitness.
## 3 What do we mean by fit?
Why are certain schemata preferred over others? Because they are fitter of course. But what does one really mean by this statement? Consider the following contrived but instructive example: consider a $`2`$-schemata, i.e. $`N_2=2`$, problem with crossover but neglecting mutation, and with a fitness landscape where $`f(01)=f(10)=0`$ and $`f(11)=f(00)=1`$. The steady state solution of the schema evolution equation is
$`P(11)=P(00)={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{p_c}{2}}{\displaystyle \frac{(l1)}{(N1)}}\right)P(01)=P(10)={\displaystyle \frac{p_c}{4}}{\displaystyle \frac{(l1)}{(N1)}}`$ (3)
For $`l=N`$ and $`p_c=1`$ we see that half the steady state population is composed of strings that have zero fitness! Note that this fixed point is in fact a stable one. Such results lead one to doubt whether the concept of fitness is the most relevant one in gauging the growth of a schema.
Another, more relevant, example is associated with a GA of binary alleles with mutation and proportional selection but without crossover (the Eigen model ). Consider a “needle-in-a- haystack” type fitness landscape where there exists one string, $`C_m`$ — the “master sequence”, of high fitness all the rest being of equal low fitness. When the mutation rate, $`\mu `$, is zero then the steady state population, assuming $`n`$ is large, is such that $`P(C_m)=1`$. When $`\mu >0`$ but small, $`P(C_m,t\mathrm{})1`$. However, the population is clustered around the master sequence in that the Hamming distance between the master sequence and the large majority of other strings in the steady state population is small. Increasing $`\mu `$ one reaches a critical value, $`\mu _c`$, beyond which $`P(C_m,t\mathrm{})1/2^N`$, which is exactly the proportion expected in a completely random population. The sharp phase transition at $`\mu _c`$ is familiar from thermodynamics being due to the competition between “energy” (selection) and “entropy” (mutation). In fact, the quantity $`(1/2)\mathrm{ln}(\mu /(1\mu ))`$ is the precise analog of the thermodynamic temperature. For $`\mu \mu _c`$ the evolution of the GA cannot be well understood by thinking of evolution on the needle-in-a-haystack landscape. Thermodynamically this corresponds to considering the energy landscape as opposed to the more physically relevant free energy landscape. In the entropy dominated regime for $`\mu \mu _c`$ every state has the same free energy hence there is effectively no selection acting. Once again we are led to call into question the usefulness of the standard notion of fitness.
As a third simple example consider the effect of mutation without crossover in the context of a model that consists of $`2`$-schemata, $`11`$, $`01`$, $`10`$, $`00`$, where each schema can mutate to the two adjacent ones when the states $`11`$, $`10`$, $`00`$, $`01`$ are placed clockwise on a circle. For example, $`11`$ can mutate to $`10`$ or $`01`$ but not to $`00`$. We assume a simple degenerate fitness landscape: $`f(11)=f(01)=f(10)=2`$, $`f(00)=1`$. Clearly there is no selective advantage for any one of the three degenerate schemata over the others. In a random population, $`P(11)=\mathrm{}=P(00)=\frac{1}{4}`$. If there is uniform probability, $`\mu `$, for each schema to mutate to an adjacent one then the evolution equation that describes this system is
$`P(i,t+1)=(12\mu )P^{}(i,t)+\mu (P^{}(i1,t)+P^{}(i+1,t))`$ (4)
For $`\mu =0`$ the steady state population is $`P(11)=P(01)=P(10)=1/3`$, $`P(00)=0`$. Thus we see the synonym symmetry of the landscape associated with the degeneracy of the states $`11`$, $`10`$ and $`01`$ is unbroken. However, for $`\mu >0`$, the schema distribution at $`t=1`$ starting from a random distribution at $`t=0`$ is $`P(11)=2/7`$, $`P(01)=P(10)=(2\mu )/7`$, $`P(00)=(1+2\mu )/7`$. Thus, we see that there is an induced breaking of the landscape synonym symmetry due to the effects of mutation. In other words the population is induced to flow along what in the fitness landscape is a flat direction.
Why is it that the examples above lead us to reconsider the idea of fitness? Fitness is intrinsically associated in the standard picture with a particular genetic operator — selection. In the above we have two examples that exhibit regimes where other genetic operators, crossover and mutation respectively, can dominate and another example wherein populations are forced to flow along directions with zero selection gradient. In such regimes intuition gleaned from the normal fitness landscape is of little value. Clearly what is required is a generalization of fitness, a type of “effective fitness”, that treats the various genetic operators on a more democratic footing and where population flows take place in an “effective fitness” landscape. In the case of selection and mutation the problem can be reformulated into a problem in equilibrium thermodynamics hence the standard thermodynamic free energy may be utilized. In more general circumstances one must find a more general effective fitness.
A natural candidate for an effective fitness has been given in . Specifically,
$`P(\xi ,t+1)={\displaystyle \frac{f_{\mathrm{eff}}(\xi ,t)}{\overline{f}(t)}}P(\xi ,t)`$ (5)
Comparing with equation (1) one finds
$`f_{\mathrm{eff}}(\xi ,t)=𝒫(\xi )\left({\displaystyle \frac{P_c(\xi ,t)}{P(\xi ,t)}}\right)\overline{f}(\xi ,t)+{\displaystyle \underset{\xi /_i}{}}𝒫(\xi /_i\xi )\left({\displaystyle \frac{P_c(\xi ,t)}{P(\xi ,t)}}\right)\overline{f}(t)`$ (6)
This definition is very natural from an evolutionary viewpoint as it gives a direct measure of the reproductive success of a given schema. In the limit $`\mu 0`$, $`p_c0`$ one finds that $`f_{\mathrm{eff}}(\xi ,t)\overline{f}(\xi ,t)`$.
Using the concept of effective fitness one can much better understand the three examples given earlier. For instance, in the first example, although the schemata $`01`$ and $`10`$ have zero fitness their effective fitness is non-zero. Similarly, in the third example although the schemata $`11`$, $`10`$ and $`01`$ are degenerate in terms of fitness this degeneracy is lifted by the effect of mutation and this is manifest at the level of the effective fitness. The above also leads to the idea of an effective selection coefficient, $`s_{\mathrm{eff}}=f_{\mathrm{eff}}(\xi ,t)/\overline{f}(t)1`$, that measures directly selective pressure including the effects of genetic operators other than selection. If we think of $`s_{\mathrm{eff}}`$ as being approximately constant in the vicinity of time $`t_0`$, then $`s_{\mathrm{eff}}(t_0)`$ gives us the exponential rate of increase or decrease of growth of the schema $`\xi `$ at time $`t_0`$. In the limit of a continuous time evolution the solution of (5) is
$`P(\xi ,t)=P(\xi ,0)\mathrm{e}^{_0^ts_{\mathrm{eff}}𝑑t^{}}`$ (7)
Using the evolution equation and the concept of effective fitness one can formulate a new Schema theorem that unlike the standard schema theorem is an equality rather than an inequality Schema Theorem
$`P(\xi ,t+1)={\displaystyle \frac{f_{\mathrm{eff}}(\xi ,t)}{\overline{f}(t)}}P(\xi ,t)`$ (8)
The interpretation of this equation is clear and analogous to the old schema theorem: schemata of higher than average effective fitness will be allocated an “exponentially” increasing number of trials over time. We put the word exponentially in quotes as the real exponent, $`^ts_{\mathrm{eff}}𝑑t^{}`$, is not, except for very simple cases such as a flat fitness landscape, of the form $`\alpha t`$, where $`\alpha `$ is a constant. The above says much more than the standard schema theorem: first of all it is an exact equation not just a lower bound; secondly it gives a deeper insight into the role of crossover. This comes about because the equation takes into account schema creation. The standard schema theorem emphasizes only the destructive effect of crossover.
Armed with the above we can much more readily investigate the reconstructive aspect of crossover. In fact we will see generically that it is a more important effect. This is of great relevance in the link between the schema theorem and the building block hypothesis. The former gives a quantitative bite to the latter by showing that the destructive effects of crossover are greater the longer the defining length of the schema; thus leading to the notion that short, fit schemata are favoured. We can readily see that this is not generally true. If schema reconstruction outweighs that of destruction then crossover is more positive the longer the schema. Under such circumstances long, fit schemata are favoured. We will see this confirmed experimentally in the next section. Finally, we may remark that unless the fitness landscape warrants it there is no reason to think of a building block as a “local” object, as seems to be the case in the work on Royal Road functions where attempts were made to validate the building block hypothesis by giving high fitness to a very small number of states associated with localized blocks of $`1`$s.
Another natural definition of effective fitness which leads us to a very simple proof of Geiringer’s theorem follows from splitting the evolution equation into those terms that are linear in $`P(\xi ,t)`$ and those that are independent of it. For instance, in the case of selection and crossover we have
$`P(\xi ,t+1)={\displaystyle \frac{f_{\mathrm{eff}}^{}(\xi ,t)}{\overline{f}(t)}}P(\xi ,t)+j(t)`$ (9)
where $`f_{\mathrm{eff}}^{}(\xi ,t)=(1p_c\frac{(l1)}{N1})\frac{\overline{f}(\xi ,t)}{\overline{f}(t)}`$ and $`j(t)=\frac{p_c}{N1}_{k=1}^{l1}P^{}(\xi _L,t)P^{}(\xi _R,t)`$. The corresponding effective selection coefficient is $`s_{\mathrm{eff}}^{}=((1p_c\frac{(l1)}{N1})\frac{\overline{f}(\xi ,t)}{\overline{f}(t)}1)`$. In the limit of a continuous time evolution (9) may be formally integrated to yield
$`P(\xi ,t)=\mathrm{e}^{_0^ts_{\mathrm{eff}}^{}(t^{})𝑑t^{}}P(\xi ,0)+\mathrm{e}^{_0^ts_{\mathrm{eff}}^{}(t^{})𝑑t^{}}{\displaystyle _0^t}j(t^{})\mathrm{e}^{_0^t^{}s_{\mathrm{eff}}^{}(t^{\prime \prime })𝑑t^{\prime \prime }}𝑑t^{}`$ (10)
In a flat landscape $`s_{\mathrm{eff}}^{}=p_c(l1)/(N1)`$, hence
$`P(\xi ,t)=\mathrm{e}^{p_c\frac{(l1)}{(N1)}t}P(\xi ,0)+{\displaystyle \frac{p_c}{(N1)}}\mathrm{e}^{p_c\frac{(l1)}{(N1)}t}\times `$
$`{\displaystyle \underset{k=1}{\overset{l1}{}}}{\displaystyle _0^t}P^{}(\xi _L,t^{})P^{}(\xi _R,t^{})\mathrm{e}^{p_c\frac{(l1)}{(N1)}t^{}}𝑑t^{}`$ (11)
Notice that dependence on the initial condition, $`P(\xi ,0)`$, is exponentially damped unless $`\xi `$ happens to be a $`1`$-schema, the solution of the $`1`$-schemata equation being $`P(i,t)=P(i,0)`$. An immediate consequence is that when considering the source term describing reconstruction the only non-zero terms that need to be taken into account are those which arise from $`1`$-schemata, as any higher order term will always have an accompanying exponential damping factor. Thus, we see that the fixed point distribution for a GA with crossover evolving in a flat fitness landscape is
$`P^{}(\xi )=Lt_t\mathrm{}P(\xi ,t)={\displaystyle \underset{i=1}{\overset{N_2}{}}}P(i,0)`$ (12)
which is basically Geiringer’s Theorem in the context of schema distributions and simple crossover. We see here that the theorem appears in an extremely simple way as a consequence of the solution of the evolution equation. Note that the fixed point distribution arises purely due to the effects of recombination, the long time behaviour depending only on the initial distribution of the most elementary building blocks — the $`1`$-schemata.
Geiringer’s theorem will also hold in a non-flat landscape if selection is only very weak, where what we mean by weak is that $`\frac{\overline{f}(\xi ,t)}{\overline{f}(t)}(1+ϵ)`$ and $`ϵ<(p_c((l1)/(N1))/1p_c((l1)/(N1)))l>1`$. In this case anything other than a $`1`$-schema will once again be exponentially damped. In this case however due to the non-trivial landscape certain $`1`$-schemata are preferred over others. A concrete example of such a landscape is $`f_i=1+\alpha _i`$ where $`_i|\alpha _i|(ϵ/(2+ϵ))`$ and $`f_i`$ is the fitness of the ith bit. Note there is no need to restrict to a linear fitness function here, arbitrary epistasis is allowed as long as it does not lead to large fitness deviations away from the mean. In this case $`\frac{\overline{f}(\xi ,t)}{\overline{f}(t)}<1+ϵ`$.
So what can we glean from our evolution equation in terms of schema size vis a vis the building block hypothesis? For a flat fitness landscape the effective selection coefficient for a $`2`$-schema $`ij`$ is
$`s_{\mathrm{eff}}=p_c\left({\displaystyle \frac{l1}{N1}}\right)+p_c\left({\displaystyle \frac{l1}{N1}}\right){\displaystyle \frac{P(i,0)P(j,0)}{P(ij,t)}}`$ (13)
One sees that schema reconstruction will exceed that of schema destruction if $`i`$ and $`j`$ are negatively correlated, i.e. if the linkage disequilibrium is negative. When will this be the case? The fitness landscape itself can quite easily induce negative correlations between $`\xi _L`$ and $`\xi _R`$. In this case there is a competition between the (anti-) correlating effect of the landscape and the mixing effect of crossover. For instance, in the neutral case of a Kaufmann $`k=0`$ landscape when $`\delta f_{\xi _L},\delta f_{\xi _R}>0`$, where $`\delta f_{\xi _L}=\overline{f}(\xi ,t)\overline{f}(t)`$, then $`1+\frac{2N_2}{N}\delta f_\xi <(1+\frac{2N_L}{N}\delta f_{\xi _L})(1+\frac{2N_R}{N}\delta f_{\xi _R})`$, so selection induces an anti-correlation, hence in an uncorrelated initial population, $`P^{}(\xi ,t)<P^{}(\xi _L,t)P^{}(\xi _R,t)`$. In this case there will clearly be a preference for large schemata rather than small ones. In the case of deceptive problems, for instance the minimal deceptive problem , typically there is a positive correlation. Under such circumstances schema destruction will be the dominant effect and one should see a preference for short, fit schemata as hypothesized by the canonical Schema theorem and Block Hypothesis. In a fitness landscape that has both a deceptive and a non-deceptive component a given schema may typically be reached via deceptive or non-deceptive channels. For instance, for a $`3`$-schema, $`ijk`$, perhaps the channel $`ij+kijk`$ is deceptive whereas the channel $`i+jkijk`$ is not. If both channels are equally likely in terms of selection then crossover will favour the non-deceptive channel and this will be manifest in the fact that one channel contributes positively to the effective fitness whereas the other contributes negatively. It has been hypothesized that generically populations will evolve via non-deceptive channels. Further theoretical results for $`k=0,2`$ can be found in .
## 4 Experimental Results
In this section we present experimental evidence for many of the statements and theoretical results we have presented. Given our desire to investigate the effects of crossover vis a vis its effect on schema length we considered a GA with $`\mu =0`$ as mutation, being a local operator, can have no direct effect on schema length. We will first consider the case of a non-epistatic landscape — the well known counting ones, or unitation problem. We considered a population of $`5000`$ $`8`$-bit strings, thus the maximum schema length is $`8`$. We chose a large population so as to be able to ignore finite size effects. Figures 1 and 2 are plots of $`M(l)`$ versus time where $`M(l)(n_{opt}(l)n_{opt}(8))/n_{opt}(8)`$. Here, $`n_{opt}(l)`$ is the number of optimal $`2`$-schemata of defining length $`l`$ normalized by the total number of length $`l`$ 2-schemata per string, i.e. $`9l`$. By optimal $`2`$-schemata we mean schemata containing the global optimum $`11`$. $`n_{opt}(8)`$ is the number of optimal $`2`$-schemata of defining length $`8`$. Figure 1 is with $`p_c=0`$ and Figure 2 with $`p_c=1`$. We show averages over $`30`$ different runs. Without crossover there is essentially no preference for schemata of a given length. Adding in crossover leads to a remarkable change: Schemata prevalence is ordered monotonically with respect to length but with the larger schemata being favoured. This is in complete accord with the theoretical prediction of the previous section based on considerations of the evolution equation. We emphasize that this is purely an effect of the crossover operator and shows clearly that schema reconstruction is more important than schema destruction.
This effect is measured nicely by the effective fitness function in Figures 3 and 4. What we in fact graph is $`F(l)(f_{\mathrm{eff}}(l)f_{\mathrm{eff}}(8))/f_{\mathrm{eff}}(8)`$. where $`f_{\mathrm{eff}}(l)`$ is the effective fitness of optimal $`2`$-schemata of size $`l`$ and $`f_{\mathrm{eff}}(8)`$ is the analogous quantity for optimal $`2`$-schemata of size $`8`$. Note that the effective fitness of larger schemata is greater than that of shorter ones for the first $`6`$ or so generations, in fact significantly so given that the fluctuations in Figure 3 are of the order of $`0.0005`$ whereas the relative effective fitness advantage of schemata of size $`8`$ relative to size $`2`$ is, from Figure 4, about $`0.01`$, i.e. $`20`$ times larger! As can be seen from (7) a positive selection coefficient is associated with a schema that is growing in number relative to another. After $`6`$ generations the curves in Figure 2 start to converge again which coincides with the effective fitness now being larger for the smaller schemata. Roughly speaking one can think of the effective fitness as being a measure of the gradient of the curves in Figures 1 and 2. If one repeats the experiments for schemata of order $`3`$ or higher one will once again find a preference for long, “non-local” schemata; non-local in the sense that there is no preference whatsoever for the three bits to be found together.
One might say that this is all fine and good but we have only shown what happens for non-epistatic landscapes. An important aspect of the counting ones landscape is not so much that it is non-epistatic as rather it is neutral in the absence of crossover, in that there is no preference in the landscape itself for schemata of a certain length. This means we can study directly the geometrical effect of crossover without having to worry about the complicating effects of selection. Having established generic properties of crossover we can then turn to various classes of fitness landscape to investigate the intricate relationship between the two operators. We will first introduce epistasis by considering what happens in the case of landscapes of the form
$`f(C_i)={\displaystyle \underset{j}{}}1_j+{\displaystyle \frac{ϵ}{N^\pm }}{\displaystyle \underset{jkC_i}{}}l_{jk}^{\pm 1}`$ (14)
where the first sum is over all the $`1`$s of the string $`C_i`$ and the second is over all pairs of $`1`$s, $`l_{jk}`$ being the defining length of the schema with $`1`$s at the points $`j`$ and $`k`$. $`N^\pm =_{jk}l_{jk}^{\pm 1}`$ is a normalization constant, being a sum over all optimal $`2`$-schemata of the form $`11`$ of the optimum string $`11111111`$. $`ϵ`$ simply controls the size of the length-dependent epistatic term relative to the counting ones term. Summing over the lengths of the schemata gives a landscape where there is a preference for long schemata. On the contrary, summing over the inverses of the lengths gives a landscape where there is a bias for short schemata. In the former case the epistasis can be thought of as giving rise to an effective repulsion between pairs of $`1`$s and in the latter an effective attraction. In both cases the epistasis between string bits depends on their distance apart.
The results can be seen in Figures 5-8. In Figure 5 we see the effect of a bias for large schemata of magnitude $`ϵ=0.3`$ with $`p_c=0`$. The graph is quite similar to that of Figure 2 hence one can see that crossover leads to an effective bias for larger schemata that is similar in many respects to an effective repulsion between schema bits. In Figure 6 we see what happens when one includes a bias for small schemata of magnitude $`ϵ=0.75`$ with $`p_c=1`$. Note that the effect of crossover is to completely annul the effect of the landscape bias. In Figures 7 and 8 we see the evolution of $`F(l)`$ in both cases. In Figures 9 and 10 we see what happens in a deceptive landscape. The landscape we chose was one where for each of the 28 different pairs of the $`8`$-bit string we have two possible sets of conditions on the fitness of each pair: i) $`f(11)=3`$, $`f(01)=f(10)=1`$, $`f(00)=2`$; and ii) $`f(11)=3`$, $`f(01)=f(10)=2`$, $`f(00)=1`$. The first set clearly is deceptive, the second clearly not. As a function of the total number of deceptive pairs, $`n_d`$, we can vary how deceptive the total landscape is. For $`n_d=28`$ the landscape is totally deceptive, whilst for $`n_d=0`$ there is no deception. In Figure 10 we see the effects of crossover on a totally deceptive landscape. Note that the effect of crossover in this case is to increase the bias for short schemata due to the fact that schema destruction is more important than schema reconstruction.
## 5 Conclusions
In this paper we have analyzed the consequences of an exact evolution equation for GAs which applies, in a very natural way, directly to schemata thus allowing for a critical analysis of the Schema theorem and the Building Block hypothesis. We saw that the very structure of the equation, taking into account as it does schema reconstruction, contains a general form of the building block hypothesis; longer, higher order schemata being constructed from smaller, lower order schemata when schema reconstruction dominates. The ultimate building blocks were shown to be $`1`$-schemata as they are immune to the effects of crossover.
We noted that a building block interpretation was natural in the case where schema reconstruction dominates, irrespective of whether the blocks were fit or not. In order to investigate under what conditions fit, small schemata were combined into fit, large schemata we found it useful to critically examine the concept of fitness. We used explicit examples to demonstrate that selective fitness and the corresponding fitness landscape were inadequate to intuitively understand GA evolution. We therefore introduced the notion of effective fitness, showing that it was a more relevant concept than pure selective fitness in governing the reproductive success of a schema. Based on this concept and the evolution equation we introduced a new schema theorem that showed that schemata of high effective fitness receive an exponentially increasing number of trials as a function of time. We also showed that generically there is no preference for short, low-order schemata. In fact, if schema reconstruction dominates the opposite is true. Only in deceptive problems does it seem that short schemata will be favoured, and then only in totally deceptive problems, as the system will tend to seek out the non-deceptive channels if they exist.
We performed various experiments to verify our theoretical results in both epistatic and non-epistatic landscapes. For non-epistatic landscapes we confirmed that there is indeed a preference for large schemata. In fact we showed that schema prevalence is a monotonically increasing function of schema defining length. In a class of epistatic landscapes designed to give an effective repulsion or attraction between pairs of bits we showed that crossover in its action was analogous to a bit-bit repulsion, thus favouring long schemata. For a model deceptive landscape we showed that in the case of total deception contrary to all the previous cases, and as predicted on theoretical grounds, short schemata were favoured.
It would naturally be very interesting to see if other exact results besides Geiringer’s theorem follow very simply from the evolution equation. A pressing matter is the search for approximation schemes within which the equations can be solved, as for a general landscape an exact solution will be impossible. In this respect techniques familiar from statistical mechanics, such as the renormalization group might well prove very useful. Of course, much more experimental analysis is needed on a wider set of test landsacpes. In particular it will be of interest to test the hypothesis that GAs will seek out non-deceptive trajectories if possible.
### Acknowledgements
This work was partially supported through DGAPA-UNAM grant number IN105197. CRS is grateful to an anonymous referee for bringing the work of Lee Altenberg to his attention and to Adam Prüghel-Bennett for useful comments on the manuscript.
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# 1. Introduction
## 1. Introduction
We are interested in the continuous dependence of entropy solutions to hyperbolic conservation laws
$$_tu+_xf(u)=0,u(x,t)RI,xRI,t>0,$$
$`1.1`$
where the flux $`f:RIRI`$ is a smooth and convex function. After works by Liu and Yang and Dafermos , we aim at deriving sharp $`L^1`$ estimates of the form
$$u^{II}(t)u^I(t)_{w(t)}+_s^tM(\tau ;u^I,u^{II})𝑑\tau u^{II}(s)u^I(s)_{w(s)},0st,$$
$`1.2`$
for any two entropy solutions of bounded variation $`u^I`$ and $`u^{II}`$ of (1.1), where $`._{w(t)}`$ is a weighted $`L^1`$-norm equivalent to the standard $`L^1`$ norm on the real line. In (1.2), the positive term $`M(\tau ;u^I,u^{II})`$ is intended to provide a sharp bound on the strict decrease of the $`L^1`$ norm. The estimate with $`w1`$ and $`M0`$ is of course well-known.
Recall that the fundamental issue of the uniqueness and continuous dependence for hyperbolic systems of conservation laws was initiated by Bressan and his collaborators (see and the references therein). A major contribution came from Liu and Yang who introduced a decreasing $`L^1`$ functional ensuring (1.2) for scalar conservation laws and systems of two equations. This research culminated in papers published simultaneously by Bressan, Liu and Yang , Hu and LeFloch , and Liu and Yang , which contain particularly simple proofs of the continuous dependence of entropy solutions for systems.
In the present paper, we restrict attention to scalar conservation laws and, following the approach developed by the second author (Hu and LeFloch and LeFloch ), we investigate the stability issue from the standpoint of Holmgren’s and Haar’s methods ( and the references therein). The problem under consideration is (essentially) equivalent to showing the uniqueness and $`L^1`$ stability for the following hyperbolic equation with discontinuous coefficient:
$$_t\psi +_x\left(a\psi \right)=0,\psi (x,t)RI,xRI,t>0.$$
$`1.3`$
That is, for solutions with bounded variation we aim at deriving an estimate like
$$\psi (t)_{w(t)}+_s^t\stackrel{~}{M}(\tau ;a,\psi )𝑑\tau \psi (s)_{w(s)},0st.$$
$`1.4`$
For the application to (1.1) one should define $`a`$ by
$$a=a(u^I,u^{II})=\frac{f(u^{II})f(u^I)}{u^{II}u^I}.$$
$`1.5`$
One may also consider the equation (1.3) for more general coefficients $`a`$.
Recall that the existence and uniqueness of solutions to the Cauchy problem associated with (1.3) was established in LeFloch in the class of bounded measures, under the assumption $`a_xE`$ for some constant $`E`$. The latter holds when $`a`$ is given by (1.5) (at least when $`u^I`$ and $`u^{II}`$ contain no rarefaction center on the line $`t=0`$ which holds “generically”). See also Crasta and LeFloch for further existence results.
It must be observed that we restrict attention here to more regular solutions, having bounded total variation, as this is natural in view of the application to the conservation law (1.1). In this direction, recall that an $`L^1`$ stability result like (1.4) was established in (see, therein, Section 5, and our Theorem 2.2 below) in the class of piecewise Lipschitz continuous solutions, with $`\stackrel{~}{M}0`$ however. This uniqueness and stability result was achieved under the assumption that the coefficient $`a`$ does not contain any rarefaction shock (see Section 2 below for the definition). In , the following essential observation was made:
The linearized equation $`(1.3)`$-$`(1.5)`$ based on two entropy $`1.6`$
solutions of $`(1.1)`$ does not exhibit rarefaction shocks.
(This is also true for systems of conservation laws, as far as solutions with small amplitude are concerned.) One of our aims here is to extend the $`L^1`$ stability result for (1.3) in to arbitrary solutions of bounded variation.
The present paper relies also heavily on the contribution by Liu and Yang who, for approximate solutions constructed by the Glimm scheme, discovered a weighted norm having a sharp decay of the form (1.2). Subsequently, the Liu-Yang’s functional was extended by Dafermos (, Chapter 11) to arbitrary functions of bounded variation (BV) and, using the notion of generalized characteristics, Dafermos derived precisely an estimate of the form (1.2) valid for BV solutions.
The aim of this paper is to provide a new derivation and some generalization of this $`L^1`$ functional. Toward the derivation of bounds like (1.2) or (1.4), we make the following preliminary observations:
The content of this paper is as follows.
In Section 2, we consider piecewise constant solutions of (1.3) and introduce a class of weighted norms satisfying a sharp bound of the form (1.4). See Theorem 2.3 below. All undercompressive and Lax discontinuities contribute to the decrease of the $`L^1`$ norm. For the sake of comparison, we also consider the $`L^1`$ norm without weight; see Theorem 2.2.
In Section 3, we point out that the setting of Section 2 covers the case of the conservation law (1.1). Passing to the limit in wave front tracking approximations, in Theorem 3.5 we arrive to the sharp bound (1.2) for general BV solutions. The proof is based on fine convergence properties established earlier by Bressan and LeFloch and on a technique of stability of nonconservative products developed by DalMaso, LeFloch, and Murat and LeFloch and Liu .
Next, in Sections 4 and 5 we return to the equation (1.3) studied in Section 2 but, now, we deal with general BV solutions. We follow closely ideas developed by Dafermos for solutions of (1.1), and extend them to the linear equation (1.3). Using generalized characteristics we establish first a maximum principle in Theorem 4.5. Finally, in Theorem 5.1 using the technique of generalized characteristics, we establish the sharp $`L^1`$ stability property (1.4) directly, for general BV solutions of (1.3). The result applies in particular to the conservation law (1.1) and allows us to recover (1.2).
Throughout the paper, we always assume that all functions of bounded variation under consideration are normalized to be defined everywhere as right-continuous function.
## 2. Decreasing Weighted Norms for Piecewise Constant Solutions
Given a piecewise constant function $`a:RI\times RI_+RI`$, let us consider the linear hyperbolic equation
$$_t\psi +_x\left(a\psi \right)=0,\psi (x,t)RI,$$
$`2.1`$
and restrict attention to piecewise constant solutions. By definition, the function $`a`$ admits a set of jump points $`𝒥(a)`$, consisting of finitely many straightlines defined on open time intervals, together with a finite set of interaction points $`(a)`$, consisting of the end points of the lines in $`𝒥(a)`$. The function $`a`$ is constant in each connected component of the complement $`𝒞(a)`$ of $`(a)𝒥(a)`$. At a point $`(x,t)𝒥(a)`$ we denote by $`\lambda ^a=\lambda ^a(x,t)`$ the speed of the discontinuity and $`a_\pm =a_\pm (x,t)=a(x\pm ,t)`$ the left- and right-hand traces. It is tacitly assumed that the discontinuity speeds $`\lambda ^a`$ remain uniformly bounded. Finally the function is normalized to be right-continuous. A similar notation is used for the function $`\psi `$.
The geometrical property of the coefficient $`a`$ play a central role for the analysis of (2.1), so we recall the following terminology :
###### Definition 2.1
A point $`(x,t)𝒥(a)`$ is called a Lax discontinuity iff
$$a_{}(x,t)>\lambda ^a(x,t)>a_+(x,t),$$
a slow undercompressive discontinuity iff
$$\lambda ^a(x,t)\mathrm{min}(a_{}(x,t),a_+(x,t)),$$
a fast undercompressive discontinuity iff
$$\lambda ^a(x,t)\mathrm{max}(a_{}(x,t),a_+(x,t)),$$
and a rarefaction-shock discontinuity iff
$$a_{}(x,t)<\lambda ^a(x,t)<a_+(x,t).$$
For each $`t>0`$, we denote by $`(a),𝒮(a),(a)`$, and $`(a)`$ the set of points $`(x,t)𝒥(a)`$ corresponding to Lax, slow undercompressive, fast undercompressive, and rarefaction-shock discontinuities, respectively.
###### Theorem 2.2
Consider a piecewise constant speed $`a=a(x,t)`$. Let $`\psi `$ be any piecewise constant solution of $`(2.1)`$. Then we have for all $`0st`$
$`\psi (t)_{L^1(RI)}+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}2\left(a_{}(x,\tau )\lambda ^a(x,\tau )\right)|\psi _{}(x,\tau )|d\tau `$ $`2.2`$
$`=\psi (s)_{L^1(RI)}+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}2\left(\lambda ^a(x,\tau )a_{}(x,\tau )\right)|\psi _{}(x,\tau )|d\tau .`$
In (2.2), the left-hand traces are chosen for definiteness only. Indeed it will be noticed in the proof below that for all $`(x,\tau )(a)(a)`$
$$\left(\lambda ^a(x,\tau )a_{}(x,\tau )\right)|\psi _{}(x,\tau )|=\left(\lambda ^a(x,\tau )a_+(x,\tau )\right)|\psi _+(x,\tau )|$$
Observe that the Lax discontinuities contribute to the decrease of the $`L^1`$ norm, while the rarefaction-shocks increase it. On the other hand, the undercompressive discontinuities don’t modify the $`L^1`$ norm. When $`a`$ contains no rarefaction shocks (this is the case when (2.1) is a linearized equation derived from entropy solutions of a conservation law, as discovered in Hu and LeFloch ), Theorem 5.1 yields
$$\psi (t)_{L^1(RI)}\psi (s)_{L^1(RI)},0st,$$
$`2.3`$
where we neglected the favorable contribution of the Lax discontinuities appearing in the left-hand side of (2.2). In particular, (2.3) implies that the Cauchy problem for (2.1) admits a unique solution (in the class of piecewise constant functions at this stage), provided $`a`$ has no rarefaction-shock discontinuities.
On the other hand, it is clear that the sign of the function $`\psi `$ is important for the sake of deriving the $`L^1`$ stability of the solutions $`\psi `$ of (2.1). For instance, if $`\psi `$ has a constant sign for all $`(x,t)`$, then (2.3) holds as an equality
$$\psi (t)_{L^1(RI)}=\psi (s)_{L^1(RI)},0st,$$
which implies that the Cauchy problem for (2.1) admits at most one solution $`\psi `$ of a given sign.
###### Demonstration Proof
Denote by $`(E)`$ the projection of a subset $`E`$ of the $`(x,t)`$-plane on the $`t`$-axis. By definition, any piecewise Lipschitz continuous solution $`\psi `$ is also Lipschitz continuous in time with values in $`L^1(RI)`$. So, it is enough to derive (2.2) for all $`tE:=\left((a)(\psi )\right)`$. The latter is just a finite set. The following is valid in each open interval $`I`$ such that $`IE=\mathrm{}`$.
We denote by $`x_j(t)`$ for $`tI`$ and $`j=1,\mathrm{},m`$ the discontinuity lines where the function $`\psi (.,t)`$ changes sign, with the convention that
$$(1)^j\psi (x,t)0\text{ for }x[x_j(t),x_{j+1}(t)].$$
$`2.4`$
Set $`\psi _j^\pm (t)=\psi _\pm (x_j(t),t)`$, $`\lambda _j(t)=\lambda ^a(x_j(t),t)`$, etc. Then by using that $`\psi `$ solves (2.1) we find (for all $`t`$ in the interval $`I`$)
$`{\displaystyle \frac{d}{dt}}{\displaystyle _{RI}}|\psi (x,t)|𝑑x`$
$`={\displaystyle \frac{d}{dt}}{\displaystyle \underset{j=1}{\overset{m}{}}}(1)^j{\displaystyle _{x_j(t)}^{x_{j+1}(t)}}\psi (x,t)𝑑x`$
$`={\displaystyle \underset{j=1}{\overset{m}{}}}(1)^j\left({\displaystyle _{x_j(t)}^{x_{j+1}(t)}}_t\psi (x,t)dx+\lambda _{j+1}(t)\psi _{j+1}^{}(t)\lambda _j(t)\psi _j^+(t)\right)`$
$`={\displaystyle \underset{j=1}{\overset{m}{}}}(1)^j\left({\displaystyle _{x_j(t)}^{x_{j+1}(t)}}_x(a(x,t)\psi (x,t))dx+\lambda _{j+1}(t)\psi _{j+1}^{}(t)\lambda _j(t)\psi _j^+(t)\right)`$
$`={\displaystyle \underset{j=1}{\overset{m}{}}}(1)^j\left((a_j^+(t)\lambda _j(t))\psi _j^+(t)+(a_j^{}(t)\lambda _j(t))\psi _j^{}(t)\right).`$
The Rankine-Hugoniot relation associated with (2.1) reads
$$(a_j^+(t)\lambda _j(t))\psi _j^+(t)=(a_j^{}(t)\lambda _j(t))\psi _j^{}(t),$$
$`2.5`$
therefore by (2.4)
$$\frac{d}{dt}_{RI}|\psi (x,t)|𝑑x=2\underset{j=1}{\overset{m}{}}\pm (a_j^\pm (t)\lambda _j(t))|\psi _j^\pm (t)|.$$
$`2.6`$
Consider each point $`x_j(t)`$ successively. If $`x_j(t)`$ is a Lax discontinuity, then $`a_j^{}(t)>\lambda _j(t)>a_+(t)`$ and both coefficients $`\pm (a_j^\pm (t)\lambda _j(t))`$ are negative. If $`x_j(t)`$ is a rarefaction-shock discontinuity, then $`a_j^{}(t)<\lambda _j(t)<a_+(t)`$ and the coefficients $`\pm (a_j^\pm (t)\lambda _j(t))`$ are positive. These two cases lead us to the two sums in (2.2). Indeed one just needs to observe the following: if $`(x,\tau )`$ correspond to a Lax or rarefaction-shock discontinuity of the speed $`a`$, but $`\psi `$ does not change sign at $`(x,\tau )`$ (so it is not counted in (2.6)), then actually by the Rankine-Hugoniot relation (see (2.5)) we conclude easily that
$$\psi _{}(x,\tau )=\psi _+(x,\tau ))=0,$$
and so it does not matter to include the point $`(x,\tau )`$ in the sums (2.2).
Suppose next that $`x_j(t)`$ is an undercompressive discontinuity. Then the two sides of (2.5) have different sign, therefore
$$(a_j^+(t)\lambda _j(t))\psi _j^+(t)=(a_j^{}(t)\lambda _j(t))\psi _j^{}(t)=0,$$
and the corresponding term in (2.6) vanishes. $`\mathit{}`$
Our objective now is to derive an improved version of Theorem 2.2, based on a weighted $`L^1`$ norm adapted to the equation (2.1). For piecewise constant functions, we set
$$\psi (t)_{w(t)}:=_{RI}|\psi (x,t)|w(x,t)𝑑x,$$
$`2.7`$
where $`w=w(x,t)>0`$ is a piecewise constant and uniformly bounded function. We determine this function based on the following constrain on its jumps, at each discontinuity of the speed $`a`$,
$$w_+(x,t)w_{}(x,t)=\{\begin{array}{cc}0\hfill & \text{ if }(x,t)𝒮(a),\hfill \\ & \\ 0\hfill & \text{ if }(x,t)(a).\hfill \end{array}$$
$`2.8`$
The weight is chosen so that the left-hand trace of a slow undercompressive discontinuity and the right-hand trace of a fast one are weighted more. This is consistent with the immediate observation that the terms $`\left(\lambda _j(t)a_j^{}(t)\right)|\psi _j^{}(t)|`$ and $`\left(a_j^+(t)\lambda _j(t)\right)|\psi _j^+(t)|`$ have a favorable (negative) sign for slow and fast undercompressive discontinuities, respectively. On the other hand, the jumps of $`w`$ at Lax or rarefaction-shock discontinuities will remain unconstrained. This choice is motivated by the two observations: (i) Lax shocks already provide us with a good contribution in (2.2), and (ii) rarefaction shocks are the source of instability and non-uniqueness and cannot be “fixed up”.
The constrain in (2.8) is different for slow and for fast undercompressive discontinuities. To actually exhibit a (uniformly bounded) weight satisfying (2.8), we put a restriction on how the nature of the discontinuities change in time as wave interactions take place. (An incoming wave may be a slow undercompressive one and become a fast one after the interaction, etc. A different constrain is placed before and after the interaction.)
Precisely, we suppose that, to the speed $`a=a(x,t)`$, we can associate on one hand a function $`\kappa :RI\times RI_+RI`$ having bounded total variation and such that $`𝒥(\kappa )𝒥(a)`$ and $`(\kappa )(a)`$, and on the other hand a partition of the discontinuities
$$𝒥(a)=𝒥^I(a)𝒥^{II}(a),$$
$`2.9`$
so that, for each $`(x,t)𝒥(a)`$, the limits $`\kappa _\pm =\kappa _\pm (x,t)`$ determine if the wave is slow or fast on its left or right side, as follows:
$$\text{sgn }\left(a_\pm (x,t)\lambda (x,t)\right)=\{\begin{array}{cc}\text{sgn }\kappa _{}\hfill & \text{ if }(x,t)𝒥^I(a),\hfill \\ \text{sgn }\kappa _{}\hfill & \text{ if }(x,t)𝒥^{II}(a).\hfill \end{array}$$
$`2.10`$
Here we use $`\text{sgn }(y)=1,0,1`$ iff $`y<0,y=0,y>0`$, respectively. Therefore a discontinuity $`(x,t)𝒥^I(a)`$ (for instance) is
a Lax one iff $`\kappa _{}<0\text{ and }\kappa _+>0,`$
a slow undercompressive one iff $`\kappa _{}0\text{ and }\kappa _+0,`$
a fast undercompressive one iff $`\kappa _{}0\text{ and }\kappa _+0,`$
a rarefaction-shock iff $`\kappa _{}>0\text{ and }\kappa _+<0.`$
Furthermore, to measure the strength of the jumps, we introduce a piecewise constant function, $`b=b(x,t)`$, having the same jump points as the function $`a`$. For instance, we could assume that there exist constants $`C_1,C_2>0`$ such that at each discontinuity of $`a`$
$$C_1|a_+(y,t)a_{}(y,t)||b_+(y,t)b_{}(y,t)|C_2|a_+(y,t)a_{}(y,t)|.$$
$`2.11`$
However, strictly speaking, this condition will not be used, in the present section at least.
Based on the functions $`\kappa `$ and $`b`$ and for $`t`$ except wave interaction times, we can set
$`V^I(x,t)`$ $`={\displaystyle \underset{\genfrac{}{}{0pt}{}{(y,t)𝒥^I(a),}{y<x}}{}}|b_+(y,t)b_{}(y,t)|,`$ $`2.12`$
$`V^{II}(x,t)`$ $`={\displaystyle \underset{\genfrac{}{}{0pt}{}{(y,t)𝒥^{II}(a),}{y<x}}{}}|b_+(y,t)b_{}(y,t)|,`$
so that the total variation of $`b(t)`$ on the interval $`(\mathrm{},x)`$ decomposes into
$$TV_{\mathrm{}}^x(b(t))=V^I(x,t)+V^{II}(x,t).$$
$`2.13`$
Fix some parameter $`m0`$. Consider now the weight-function defined for each $`(x,t)𝒞(a)`$ by
$$w(x,t)=\{\begin{array}{cc}m+V^I(\mathrm{},t)V^I(x,t)+V^{II}(x,t)\hfill & \text{if }\kappa (x,t)>0,\hfill \\ & \\ m+V^I(x,t)+V^{II}(\mathrm{},t)V^{II}(x,t)\hfill & \text{if }\kappa (x,t)0.\hfill \end{array}$$
$`2.14`$
It is immediate to see that indeed (2.8) holds and that with (2.11)
$$mw(x,t)m+TV(b(t))m+C_2TV(a(t)),xRI.$$
$`2.15`$
Note also that the weight depends on $`b`$ and $`a`$, but not on the solution.
###### Theorem 2.3
Consider a piecewise constant speed $`a=a(x,t)`$ admitting a decomposition $`(2.9)`$-$`(2.10)`$ and satisfying the total variation estimate $`(2.15)`$. Consider the weight function $`w=w(x,t)`$ defined by $`(2.13)`$. Let $`\psi `$ be any piecewise constant solution of the linear hyperbolic equation $`(2.1)`$. Then the weighted norm $`(2.7)`$ satisfies for all $`0st`$
$`\psi (t)_{w(t)}`$ $`2.16`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}\left(2m+TV(b)|b_+(x,\tau )b_{}(x,\tau )|\right)\left|a_{}(x,\tau )\lambda (x,\tau )\right||\psi _{}(x,\tau )|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )𝒮(a)(a)}{}}|b_+(x,\tau )b_{}(x,\tau )|\left|a_{}(x,\tau )\lambda (x,\tau )\right||\psi _{}(x,\tau )|d\tau `$
$`=\psi (s)_{w(s)}+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}\left(2m+TV(b)\right)\left|a_{}(x,\tau )\lambda (x,\tau )\right||\psi _{}(x,\tau )|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}|b_+(x,\tau )b_{}(x,\tau )|\left|a_{}(x,\tau )\lambda (x,\tau )\right||\psi _{}(x,\tau )|d\tau .`$
The statement (2.16) is sharper than (2.2), as all discontinuities contribute now to the decrease of the weighted $`L^1`$ norm. Note that as $`m\mathrm{}`$, we recover exactly (2.2) from (2.16).
###### Demonstration Proof
We proceed similarly as in the proof of Theorem 2.2. However, $`x_j(t)`$ for $`tI`$ (some open interval avoiding the interaction points in $`a`$ or $`\psi `$) denote now all the jump points in either $`a`$ or $`\psi `$. We obtain as before the identity
$`{\displaystyle \frac{d}{dt}}{\displaystyle _{RI}}|\psi (x,t)|w(x,t)𝑑x`$ $`2.17`$
$`={\displaystyle \underset{j=1}{\overset{m}{}}}\left((\lambda _j(t)a_j^{}(t))|\psi _j^{}(t)|w_j^{}(t)+(a_j^+(t)\lambda _j(t))|\psi _j^+(t)|w_j^+(t)\right)`$
$`={\displaystyle \underset{j=1}{\overset{m}{}}}\left(\text{sgn }\left(\lambda _j(t)a_j^{}(t)\right)w_j^{}(t)+\text{sgn }\left(a_j^+(t)\lambda _j(t)\right)w_j^+(t)\right)\left|\lambda _j(t)a_j^{}(t)\right||\psi _j^{}(t)|,`$
where we used the Rankine-Hugoniot relation (2.5).
If $`x_j(t)`$ is a Lax discontinuity in $`𝒥^I(a)`$, then by (2.11) we have $`\kappa _{}<0`$ and $`\kappa _+>0`$. So by (2.14) we find
$`w_j^{}`$ $`=m+V^I(x_j(t))+V^{II}(\mathrm{})V^{II}(x_j(t)),`$
$`w_j^+`$ $`=m+V^I(\mathrm{})V^I(x_j(t)+)+V^{II}(x_j(t)+),`$
and so
$`\text{sgn }\left(\lambda _j(t)a_j^{}(t)\right)w_j^{}(t)+\text{sgn }\left(a_j^+(t)\lambda _j(t)\right)w_j^+(t)`$ $`2.18a`$
$`=w_j^{}(t)w_j^+(t)`$
$`=2mTV(b)+|b_j^+(t)b_j^{}(t)|.`$
If $`x_j(t)`$ is a rarefaction-shock discontinuity in $`𝒥^I(a)`$, then by (2.11) we have $`\kappa _{}>0`$ and $`\kappa _+<0`$. By (2.13) we find
$`w_j^{}`$ $`=m+V^I(\mathrm{})V^I(x_j(t))+V^{II}(x_j(t)),`$
$`w_j^+`$ $`=m+V^I(x_j(t)+)+V^{II}(\mathrm{})V^{II}(x_j(t)+),`$
and so
$`\text{sgn }\left(\lambda _j(t)a_j^{}(t)\right)w_j^{}(t)+\text{sgn }\left(a_j^+(t)\lambda _j(t)\right)w_j^+(t)`$ $`2.18b`$
$`=w_j^{}(t)+w_j^+(t)`$
$`=2m+TV(b)+|b_j^+(t)b_j^{}(t)|.`$
If $`x_j(t)`$ is a fast undercompressive discontinuity in $`𝒥^I(a)`$, then by (2.11) we have $`\kappa _{}0`$ and $`\kappa _+0`$. By (2.13) we find
$`\text{sgn }\left(\lambda _j(t)a_j^{}(t)\right)w_j^{}(t)+\text{sgn }\left(a_j^+(t)\lambda _j(t)\right)w_j^+(t)`$ $`2.18c`$
$`=w_j^{}(t)w_j^+(t)`$
$`=m+V^I(x_j(t))+V^{II}(\mathrm{})V^{II}(x_j(t))`$
$`mV^I(x_j(t)+)V^{II}(\mathrm{})+V^{II}(x_j(t)+)`$
$`=|b_j^+(t)b_j^{}(t)|.`$
Similarly for slow undercompressive discontinuities in $`𝒥^I(a)`$ we obtain
$$\text{sgn }\left(\lambda _j(t)a_j^{}(t)\right)w_j^{}(t)+\text{sgn }\left(a_j^+(t)\lambda _j(t)\right)w_j^+(t)=|b_j^+(t)b_j^{}(t)|.$$
$`2.18d`$
Using (2.18) in (2.17) we conclude that
$`\psi (t)_{w(t)}`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}\left(2m+TV(b)|b_+(y,\tau )b_{}(y,\tau )|\right)\left|a_{}(x,\tau )\lambda (x,\tau )\right||\psi _{}(x,\tau )|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )𝒮(a)(a)}{}}|b_+(y,\tau )b_{}(y,\tau )|\left|a_{}(x,\tau )\lambda (x,\tau )\right||\psi _{}(x,\tau )|d\tau `$
$`=\psi (s)_{w(s)}`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}\left(2m+TV(b)+|b_+(y,\tau )b_{}(y,\tau )|\right)\left|a_{}(x,\tau )\lambda (x,\tau )\right||\psi _{}(x,\tau )|d\tau ,`$
which is equivalent to (2.16). $`\mathit{}`$
Using that $`(a)`$ is included in the set of points where $`\psi `$ changes sign, it is easy to deduce from (2.16) that:
###### Corollary 2.4
Under the assumptions and notations in Theorem 2.3, we have for all $`0st`$
$`\psi (t)_{w(t)}`$
$`\psi (s)_{w(s)}+(2m+TV(b))\underset{\genfrac{}{}{0pt}{}{(x,\tau )(a)}{s\tau t}}{sup}|b_+(x,\tau )b_{}(x,\tau )|{\displaystyle _s^t}TV(\psi (\tau ))d\tau `$
and, in particular, letting $`m\mathrm{}`$
$$\psi (t)_{L^1(RI)}\psi (s)_{L^1(RI)}+2\underset{\genfrac{}{}{0pt}{}{(x,\tau )(a)}{s\tau t}}{sup}|b_+(x,\tau )b_{}(x,\tau )|_s^tTV(\psi (\tau ))d\tau .$$
$`2.19`$
Finally, in view of Corollary 2.4, in case the function $`a`$ contains no rarefaction shocks, we deduce that
$$\psi (t)_{w(t)}\psi (s)_{w(s)},0st.$$
Observe that this result is achieved, based on a weight that depends on an arbitrary function, $`b`$, and on the sole assumption that a decomposition (2.9)-(2.10) of the jumps of $`a`$ is available. However, our result in this section covers only piecewise constant solutions. We will see in Section 5 that a stronger structure assumption on the coefficients $`a`$ is necessary to handle general solutions of bounded variation.
## 3. Sharp $`L^1`$ Estimate for Hyperbolic Conservation Laws
In this section, we apply Theorem 2.3 to the case that $`a`$ is the averaging coefficient (1.5) based on two entropy solutions of (1.1). First, we check that the assumptions required in Section 2 on the coefficient $`a`$ do hold in this situation. Therefore Theorem 2.3 applies to the piecewise constant solutions defined by the wave-front traking (also called polygonal approximation) algorithm proposed by Dafermos in . Next, we observe that, with a suitable choice of the definition of the wave strengths, the weighted norm in Section 2 reduces to Liu-Yang’s functional. Finally we rigorously justify the passage to the limit in the estimate of Theorem 2.3 when the number of wave fronts tends to infinity and exact entropy solutions of (1.1) are recovered.
Consider the nonlinear scalar conservation law:
$$_tu+_xf(u)=0,u(x,t)RI,$$
$`3.1`$
where the flux $`f:RIRI`$ is a smooth function. Let $`u^I`$ and $`u^{II}`$ be two bounded entropy solutions of (3.1) having bounded total variation. Given $`h>0`$ let us approximate the data $`u^I(0)`$ and $`u^{II}(0)`$ by piecewise constant functions $`u^{I,h}(0)`$, $`u^{II,h}(0)`$, having finitely many jumps and such that as $`h0`$
$$u^{I,h}(0)u^I(0),u^{II,h}(0)u^{II}(0)\text{ in the }L^1\text{ norm,}$$
$`3.2`$
$$TV(u^{I,h}(0))TV(u^I(0)),TV(u^{II,h}(0))TV(u^{II}(0)).$$
$`3.3`$
Applying Dafermos’ scheme , we can construct corresponding, piecewise constant, approximate solutions $`u^{I,h}`$ and $`u^{II,h}`$ having finitely many jump lines and for $`ts0`$ and $`p[1,\mathrm{}]`$
$$u^{I,h}(t)_{L^p(RI)}u^{I,h}(s)_{L^p(RI)},u^{II,h}(t)_{L^p(RI)}u^{II,h}(s)_{L^p(RI)},$$
$`3.4`$
and for all $`\mathrm{}A+M(ts)BM(ts)`$
$`TV_{A+M(ts)}^{BM(ts)}\left(u^{I,h}(t)\right)TV_A^B\left(u^{I,h}(s)\right),`$ $`3.5`$
$`TV_{A+M(ts)}^{BM(ts)}\left(u^{II,h}(t)\right)TV_A^B\left(u^{II,h}(s)\right).`$
More precisely, the functions $`u^{I,h}`$ and $`u^{II,h}`$ are exact solutions of (3.1) satisfying therefore the Rankine-Hugoniot relation at every jump. They contain two kinds of jump discontinuities: Lax shocks satisfy the so-called Oleinik entropy inequalities, while rarefaction jumps do not, but have small strength, that is
$$|u^{I,h}(x+,t)u^{I,h}(x,t)|h,|u^{II,h}(x+,t)u^{II,h}(x,t)|h.$$
$`3.6`$
Furthermore, for a subsequence $`h0`$ at least, we have for each time $`t0`$
$$u^{I,h}(t)u^I(t),u^{II,h}(t)u^{II}(t)\text{ in the }L^1\text{ norm.}$$
To study the $`L^1`$ distance between these approximate solutions, we set
$$\psi :=u^{II,h}u^{I,h},$$
which is one solution of the linear hyperbolic equation
$$_t\psi +_x\left(a^h\psi \right)=0,a^h(x,t):=\frac{f(u^{II,h}(x,t))f(u^{I,h}(x,t))}{u^{II,h}(x,t)u^{I,h}(x,t)}.$$
$`3.7`$
First of all, based on Theorem 2.2 and (3.5)-(3.6), we obtain immediately:
###### Theorem 3.1
The approximate solutions $`u^{I,h}`$ and $`u^{II,h}`$ satisfy the following $`L^1`$ stability estimate for all $`0st`$
$`u^{II,h}(t)u^{I,h}(t)_{L^1(RI)}`$ $`3.8`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}2\left(a^h(x,\tau )\lambda ^{a^h}(x,\tau )\right)|u^{II,h}(x,\tau )u^{I,h}(x,\tau )|d\tau `$
$`u^{II,h}(s)u^{I,h}(s)_{L^1(RI)}+2h(ts)f^{\prime \prime }_{\mathrm{}}\left(TV(u^{I,h}(0))+TV(u^{II,h}(0))\right).`$
From the functions $`u^I`$ and $`u^{II}`$ we define the function $`a`$ as in (3.7). Recall that the wave front tracking scheme converge locally uniformly (see the proof of Theorem 3.5 below for a the definition), so that the BV solutions $`u^I`$ and $`u^{II}`$ are endowed with additional regularity properties. Consider for instance the function $`u^I`$. In particular, for all but countably many times $`t`$ and for each $`x`$, either $`x`$ is a point of continuity of $`u^I`$ in the classical sense (say $`(x,t)𝒞(u^I)`$) or else it is a point of jump in the classical sense (say $`(x,t)𝒥(u^I)`$) and, to the discontinuity, one can also associate a shock speed, denoted by $`\lambda ^I(x,t)`$.
From the properties shared by $`u^I`$ and $`u^{II}`$, one deduces immediately a similar property for the coefficient $`a`$. Excluding countably many times at most, at each point of jump of $`a`$ we can define the propagation speed $`\lambda ^a(x,t)`$ of the discontinuity located at the point $`(x,t)`$. Namely, we have
$$\lambda ^a(x,t)=\{\begin{array}{cc}\lambda ^I(x,t)\hfill & \text{ if }(x,t)𝒥(u^I),\hfill \\ \lambda ^{II}(x,t)\hfill & \text{ if }(x,t)𝒥(u^{II}).\hfill \end{array}$$
In the limit $`h0`$ we deduce from $`(3.8)`$ that:
###### Corollary 3.2
For all $`0st`$ we have
$`u^{II}(t)u^I(t)_{L^1(RI)}`$ $`3.9`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}2\left(a(x,\tau )\lambda ^a(x,\tau )\right)|u^{II}(x,\tau )u^I(x,\tau )|d\tau `$
$`u^{II}(s)u^I(s)_{L^1(RI)}.`$
We omit the proof of Corollary 3.2 as (3.9) is a consequence of a stronger estimate proven in Theorem 3.5 below (by taking $`m\mathrm{}`$ in (3.15)). Note that (3.9) is a stronger statement than the standard $`L^1`$ contraction estimate
$$u^{II}(t)u^I(t)_{L^1(RI)}u^{II}(s)u^I(s)_{L^1(RI)}.$$
###### Demonstration Proof
We apply the estimate (2.2) with $`\psi `$ replaced with $`u^{II,h}u^{I,h}`$. We just need to observe (see ) that all the rarefaction-shock discontinuities in $`a^h`$ are due to rarefaction fronts in $`u^{I,h}`$ or in $`u^{II,h}`$, which have small strength according to (3.6). In other words we have
$`{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}2\left(\lambda ^a(x,\tau )a_{}(x,\tau )\right)|\psi _{}(x,\tau )|d\tau `$
$`\underset{\genfrac{}{}{0pt}{}{(x,\tau )(a)}{s\tau t}}{sup}2|a_+(x,\tau )a_{}(x,\tau )|{\displaystyle _s^t}TV(\psi (\tau ))d\tau `$
$`2f^{\prime \prime }_{\mathrm{}}h{\displaystyle _s^t}TV(\psi (\tau ))𝑑\tau `$
$`2f^{\prime \prime }_{\mathrm{}}h(ts)\left(TV(u^{I,h}(0))+TV(u^{II,h}(0))\right).`$
This establishes (3.8). $`\mathit{}`$
We now want to apply Theorem 2.3 and control a weighted norm of $`u^{II,h}u^{I,h}`$. In this direction, our main observation is:
###### Lemma 3.3
When the function $`f`$ is strictly convex, the coefficient $`a^h`$ satisfies all of the assumptions $`(2.9)`$-$`(2.10)`$.
###### Demonstration Proof
The function $`a^h`$ is piecewise constant, and we can associate to this function an obvious decomposition of the form (2.9). To establish (2.10), consider for instance a jump point $`(x,t)𝒥(u^{I,h})𝒞(u^{II,h})`$, together with its left- and right-hand traces $`u_{}^I`$ and $`u_+^I`$. Since $`u^{I,h}`$ is a solution of (3.1), the corresponding speed $`\lambda =\lambda (x,t)`$ satisfies the Rankine-Hugoniot relation:
$$\lambda \left(u_+^Iu_{}^I\right)+f(u_+^I)f(u_{}^I)=0.$$
Thus the term in the left-hand side of (2.10) takes the form
$`a_\pm (x,t)\lambda (x,t)`$ $`={\displaystyle \frac{f(u^{II})f(u_\pm ^I)}{u^{II}u_\pm ^I}}{\displaystyle \frac{f(u_+^I)f(u_{}^I)}{u_+^Iu_{}^I}}`$
$`={\displaystyle _0^1}\left(f^{}\left(\theta u^{II}+(1\theta )u_\pm ^I\right)f^{}\left(\theta u_{}^I+(1\theta )u_\pm ^I\right)\right)𝑑\theta .`$
Thus we obtain
$`a_\pm (x,t)\lambda (x,t)=\mu \left(u^{II}u_{}^I\right),`$ $`3.10`$
$`\mu :={\displaystyle _0^1}{\displaystyle _0^1}f^{\prime \prime }\left(\rho \left(\theta u^{II}+(1\theta )u_\pm ^I\right)+(1\rho )\left(\theta u_{}^I+(1\theta )u_\pm ^I\right)\right)\theta 𝑑\theta 𝑑\rho .`$
Since $`f`$ is strictly convex, the coefficient is bounded away from zero. In view of (3.10), if we now choose $`\kappa (x,t):=u^{II,h}u^{I,h}`$, the desired property (2.10) holds true. $`\mathit{}`$
Next, we define the weight $`w^h=w^h(x,t)`$ associated with the function $`a^h`$, by the formula (2.14) in which we specify
$$\kappa ^h(x,t):=u^{II,h}u^{I,h}.$$
$`3.11`$
It follows immediately from Theorem 2.3 that:
###### Theorem 3.4
Suppose that the function $`f`$ is strictly convex. The approximate solutions constructed by Dafermos scheme satisfy the $`L^1`$ stability estimate for all $`0st`$
$`u^{II,h}(t)u^{I,h}(t)_{w^h(t)}`$ $`3.12`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a^h)}{}}\left(2m+TV(b^h)|b^h(x+,\tau )b^h(x,\tau )|\right)`$
$`\left|a^h(x,\tau )\lambda ^h(x,\tau )\right||u^{II,h}(x,\tau )u^{I,h}(x,\tau )|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )𝒮(a^h)(a^h)}{}}|b^h(x+,\tau )b^h(x,\tau )|\left|a^h(x,\tau )\lambda ^h(x,\tau )\right|`$
$`|u^{II,h}(x,\tau )u^{I,h}(x,\tau )|d\tau `$
$`=u^{II,h}(s)u^{I,h}(s)_{w^h(s)}`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a^h)}{}}\left(2m+TV(b^h)+|b^h(x+,\tau )b^h(x,\tau )|\right)`$
$`\left|a^h(x,\tau )\lambda ^h(x,\tau )\right||u^{II,h}(x,\tau )u^{I,h}(x,\tau )|d\tau ,`$
where $`a^h`$ is the averaging coefficient defined in $`(3.7)`$ and $`\lambda ^h(x,\tau )`$ represents the speed of the discontinuity located at $`(x,\tau )𝒥(a^h)`$.
We emphasize that (3.12) is an equality in which the contribution to the $`L^1`$ norm of each type of wave appears clearly. The coefficient $`a^h`$ exhibits three types of waves: the Lax and undercompressive discontinuities in $`a^h`$ contribute to the decay of the $`L^1`$ weighted distance. The statement (3.12) quantifies sharply this effect. On the other hand, the rarefaction-shocks appearing in the right-hand side of (3.12) increase the $`L^1`$ norm.
In the rest of this section, we assume that the function $`b=b^h`$ is chosen to be specifically
$$b^h(x+,t)b^h(x,t)=\{\begin{array}{cc}u^{I,h}(x+,t)u^{I,h}(x,t)\hfill & \text{ if }(x,t)𝒥(u^{I,h}),\hfill \\ u^{II,h}(x+,t)u^{II,h}(x,t)\hfill & \text{ if }(x,t)𝒥(u^{II,h}),\hfill \end{array}$$
$`3.13`$
but a more general definition is possible.
Our next purpose is to pass to the limit ($`h0`$) in the statement established in Theorem 3.4 for piecewise constant approximate solutions. We recover here a result derived by Dafermos via a different approach. Recall the notation $`𝒞(u^I)`$, $`𝒮(u^I)`$, etc introduced earlier. Denote by $`(u^I)`$ the countable set of interactions times. Let $`V^I(t)`$ be the total variation function associated with $`u^I(t)`$. Based on the functions $`V^I(t)`$ and $`V^{II}(t)`$, we then define the weight $`w`$ as in (2.14) but with (2.12) replaced by the total variation functions of $`u^I(t)`$ and $`u^{II}(t)`$, with $`\kappa :=u^{II}u^I`$ and
$$b(x+,t)b(x,t)=\{\begin{array}{cc}u^I(x+,t)u^I(x,t)\hfill & \text{ if }(x,t)𝒥(u^I),\hfill \\ u^{II}(x+,t)u^{II}(x,t)\hfill & \text{ if }(x,t)𝒥(u^{II}).\hfill \end{array}$$
$`3.14`$
Furthermore, to any functions of bounded variation $`u,v,w`$ in the space variable $`x`$ (the time variable being fixed) we associate the measure on $`RI`$
$$\mu =\left(a(u,v)f^{}(u)\right)(vu)dw$$
understood as the nonconservative product in the sense of Dal Maso, LeFloch and Murat and characterized by the following two conditions:
Note that, if $`u=u^I`$ and $`v=u^{II}`$, the two terms $`\left(a(u_\pm ,v_\pm )a(u_{},u_+)\right)(v_\pm u_\pm )`$ in fact coincide.
###### Theorem 3.5
Let the function $`f`$ be strictly convex and let $`u^I`$ and $`u^{II}`$ be two entropy solutions of bounded variation of the conservation law $`(1.1)`$. For all $`0st`$ we have
$`u^{II}(t)u^I(t)_{w(t)}`$ $`3.16`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)𝒥(u^I)}{}}q\left|a(u^I(x),u^{II}(x))a(u^I(x+),u^I(x))\right||u^{II}(x)u^I(x)|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)𝒥(u^{II})}{}}q\left|a(u^I(x),u^{II}(x))a(u^{II}(x+),u^{II}(x))\right||u^{II}(x)u^I(x)|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle _{RI}}\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)𝑑V^I𝑑\tau `$
$`+{\displaystyle _s^t}{\displaystyle _{RI}}\left(a(u^I,u^{II})f^{}(u^{II})\right)(u^Iu^{II})𝑑V^{II}𝑑\tau `$
$`u^{II}(s)u^I(s)_{w(s)}.`$
where $`q=q(\tau )=2m+TV(u^I(\tau ))+TV(u^{II}(\tau ))`$.
Observe that the terms in integrals in (3.16) globally contribute to the decrease of weighted norm, as is better seen rewriting the formula as follows ($`V_c^I`$ and $`V_c^{II}`$ being the continuous parts of the measures $`V^I`$ and $`V^{II}`$):
$`u^{II}(t)u^I(t)_{w(t)}`$ $`3.16^{}`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)𝒥(u^I)}{}}\left(q|u_+^Iu_{}^I|\right)\left|a(u_{}^I,u_+^{II})a(u_+^I,u_{}^I)\right||u^{II}u^I|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)𝒥(u^{II})}{}}\left(q|u_+^{II}u_{}^{II}|\right)\left|a(u_{}^I,u_{}^{II})a(u_+^{II},u_{}^{II})\right||u^{II}u^I|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )\left(𝒮(a)(a)\right)𝒥(u^I)}{}}\left|a(u_{}^I,u_{}^{II})a(u_+^I,u_{}^I)\right||u^{II}u^I||u_+^Iu_{}^I|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )\left(𝒮(a)(a)\right)𝒥(u^{II})}{}}\left|a(u_{}^I,u_{}^{II})a(u_+^{II},u_{}^{II})\right||u^{II}u^I||u_+^{II}u_{}^{II}|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle _{RI}}\left|a(u^I,u^{II})f^{}(u^I)\right||u^{II}u^I|𝑑V_c^I𝑑\tau `$
$`+{\displaystyle _s^t}{\displaystyle _{RI}}\left|a(u^I,u^{II})f^{}(u^{II})\right||u^Iu^{II}|𝑑V_c^{II}𝑑\tau `$
$`u^{II}(s)u^I(s)_{w(s)}.`$
The following estimate is a direct consequence of the definition (3.15):
###### Lemma 3.6
There exists a constant $`C>0`$ such that for all functions of bounded variation $`u,\stackrel{~}{u},v,\stackrel{~}{v},w`$ defined on some interval $`[\alpha ,\beta ]`$
$`\left|{\displaystyle _\alpha ^\beta }\left(a(u,v)f^{}(u)\right)(vu)𝑑w{\displaystyle _\alpha ^\beta }\left(a(\stackrel{~}{u},\stackrel{~}{v})f^{}(\stackrel{~}{u})\right)(\stackrel{~}{v}\stackrel{~}{u})𝑑w\right|`$ $`3.17`$
$`C\left(\stackrel{~}{u}u_{L^{\mathrm{}}(\alpha ,\beta )}+\stackrel{~}{v}v_{L^{\mathrm{}}(\alpha ,\beta )}\right)TV_{[\alpha ,\beta ]}(w).`$
###### Demonstration Proof of Theorem 3.5
Step 1 : Preliminaries.
For each $`t0`$, the functions $`V^{I,h}(t)`$ and $`V^{II,h}(t)`$ associated with the wave front tracking approximations $`u^{I,h}(t)`$ and $`u^{II,h}(t)`$ are of uniformly bounded variation as $`h0`$. The measures $`dV^{I,h}`$ and $`dV^{II,h}`$ are also Lipschitz continuous in time (with constant independent of $`h`$) for the weak convergence, except at interaction points. On the other hand, interaction times in the limiting solutions are at most countable. Therefore, extracting subsequences if necessary, the measures $`dV^{I,h}`$ and $`dV^{II,h}`$ converge to some limiting (non-negative) measures, say:
$$dV^{I,h}(t)d\overline{V}^I(t)dV^{II,h}(t)d\overline{V}^{II}(t).$$
$`3.18`$
By lower semi-continuity, we have at each time $`t`$
$$dV^I(t)d\overline{V}^I(t),dV^{II}(t)d\overline{V}^{II}(t),$$
$`3.19`$
and, in particular, at each $`(x,t)`$
$$V^I(x,t)\overline{V}^I(x,t),V^{II}(x,t)\overline{V}^{II}(x,t).$$
$`3.20a`$
$`V^I(+\mathrm{},t)V^I(x,t)\overline{V}^I(+\mathrm{},t)\overline{V}^I(x,t),`$ $`3.20b`$
$`V^{II}(+\mathrm{},t)V^{II}(x,t)\overline{V}^{II}(+\mathrm{},t)\overline{V}^{II}(x,t).`$
Based on the functions $`\overline{V}^I(t)`$ and $`\overline{V}^{II}(t)`$, on the coefficient $`\kappa :=u^{II}u^I`$ and on the function in (3.14), we can define a weight denoted by $`\overline{w}`$, along the same lines as in (2.14). We will show that the left-hand side of (3.16) is bounded above by
$`u^{II}(t)u^I(t)_{\overline{w}(t)}`$ $`3.21`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)𝒥(u^I)}{}}\overline{q}\left|a(u^I(x),u^{II}(x))a(u^I(x+),u^I(x))\right||u^{II}(x)u^I(x)|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)𝒥(u^{II})}{}}\overline{q}\left|a(u^I(x),u^{II}(x))a(u^{II}(x+),u^{II}(x))\right||u^{II}(x)u^I(x)|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle _{RI}}\left(a\right(u^I,u^{II})f^{}(u^I))(u^{II}u^I)d\overline{V}^I(y,\tau )d\tau `$
$`+{\displaystyle _s^t}{\displaystyle _{RI}}\left(a(u^I,u^{II})f^{}(u^{II})\right)(u^Iu^{II})𝑑\overline{V}^{II}(y,\tau )𝑑\tau `$
where $`\overline{q}:=2m+\overline{V}^I(+\mathrm{})+\overline{V}^{II}(+\mathrm{})`$, and that (3.21) coincides with the desired upper bound $`u^{II}(s)u^I(s)_{w(s)}`$. The former statement is postponed to Step 5 below and we focus now on the latter.
Fix some $`ts0`$ and rewrite (3.12) in the equivalent form
$`u^{II,h}(t)u^{I,h}(t)_{w^h(t)}`$ $`3.22`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a^h)}{}}\left(2m+TV(b^h)\right)\left|a^h(x,\tau )\lambda ^h(x,\tau )\right||u^{II,h}(x,\tau )u^{I,h}(x,\tau )|d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )𝒥^I(a^h)}{}}|b^h(x+,\tau )b^h(x,\tau )|\left(a^h(x,\tau )\lambda ^h(x,\tau )\right)(u^{II,h}(x,\tau )u^{I,h}(x,\tau ))d\tau `$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )𝒥^{II}(a^h)}{}}|b^h(x+,\tau )b^h(x,\tau )|\left(a^h(x,\tau )\lambda ^h(x,\tau )\right)(u^{I,h}(x,\tau )u^{II,h}(x,\tau ))d\tau `$
$`=u^{II,h}(s)u^{I,h}(s)_{w^h(s)}`$
$`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a^h)}{}}\left(2m+TV(b^h)\right)\left|a^h(x,\tau )\lambda ^h(x,\tau )\right||u^{II,h}(x,\tau )u^{I,h}(x,\tau )|d\tau ,`$
or, with obvious notations,
$$u^{II,h}(t)u^{I,h}(t)_{w^h(t)}+\mathrm{\Omega }_1^h+\mathrm{\Omega }_2^h=u^{II,h}(s)u^{I,h}(s)_{w^h(s)}+\mathrm{\Omega }_3^h.$$
$`3.23`$
As the maximum strength of rarefaction fronts in $`u^{I,h}`$ and $`u^{II,h}`$ vanishes with $`h`$ (see (3.6)) and rarefaction shocks in $`a^h`$ arise only from these rarefaction fronts (see (1.6)), we have
$$\mathrm{\Omega }_3^h0\text{ as }h0.$$
$`3.24`$
On the other hand, we can always choose the (initial) approximations at time $`s`$ in such a way that
$$\overline{w}(s)=w(s)$$
$`3.25`$
and
$$\underset{h0}{lim}u^{II,h}(s)u^{I,h}(s)_{w^h(s)}=u^{II}(s)u^I(s)_{w(s)}.$$
$`3.26`$
It remains to prove that the limit of the left-hand side of (3.22) is exactly (3.21). This will be established in the following three steps.
Step 2 : We will rely on the local uniform convergence of the front tracking approximations (see Bressan and LeFloch ). For all but countably many times $`\tau `$ we have the following properties for $`u^I`$ (as well as for $`u^{II}`$):
We also recall from that, for all but countably many times $`t`$, the atomic parts of the measures $`\overline{V}^I`$ and $`\overline{V}^{II}`$ coincide with the one of $`V^I`$ and $`V^{II}`$, that is for each $`yRI`$
$`\overline{V}^I(y+,t)\overline{V}^I(y,t)=V^I(y+,t)V^I(y,t),`$ $`3.29`$
$`\overline{V}^{II}(y+,t)\overline{V}^{II}(y,t)=V^{II}(y+,t)V^{II}(y,t).`$
Following LeFloch and Liu who established the weak stability of nonconservative products under local uniform convergence, we want to show that
$`\mathrm{\Omega }_2^h(\tau ):=`$ $`{\displaystyle _{RI}}\left(a(u^{I,h}(y,\tau ),u^{II,h}(y,\tau ))f^{}(u^{I,h}(y,\tau ))\right)(u^{II,h}(y,\tau )u^{I,h}(y,\tau ))𝑑V^{I,h}(y)`$ $`3.30`$
$`+{\displaystyle _{RI}}\left(a(u^{I,h}(y,\tau ),u^{II,h}(y,\tau ))f^{}(u^{II,h}(y,\tau ))\right)(u^{I,h}(y,\tau )u^{II,h}(y,\tau ))𝑑V^{II,h}(y)`$
$``$ $`{\displaystyle _{RI}}\left(a(u^I(y,\tau ),u^{II}(y,\tau ))f^{}(u^I(y,\tau ))\right)(u^{II}(y,\tau )u^I(y,\tau ))𝑑\overline{V}^I(y)`$
$`+{\displaystyle _{RI}}\left(a(u^I(y,\tau ),u^{II}(y,\tau ))f^{}(u^{II}(y,\tau ))\right)(u^I(y,\tau )u^{II}(y,\tau ))𝑑\overline{V}^{II}(y).`$
By Lebesgue dominated convergence theorem and since a uniform bound in $`\tau `$ and $`h`$ is available, it will follow from (3.29) that
$`\mathrm{\Omega }_2^h={\displaystyle _s^t}\mathrm{\Omega }_2^h(\tau )𝑑\tau `$ $`{\displaystyle _s^t}{\displaystyle _{RI}}\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)𝑑\overline{V}^I(y,\tau )𝑑\tau `$ $`3.30^{}`$
$`+{\displaystyle _s^t}{\displaystyle _{RI}}\left(a(u^I,u^{II})f^{}(u^{II})\right)(u^Iu^{II})𝑑\overline{V}^{II}(y,\tau )𝑑\tau .`$
Given $`ϵ>0`$, select finitely many (large) jumps in $`u^I`$ or $`u^{II}`$, located at $`y_1,y_2,\mathrm{}y_n`$, so that
$$\underset{\genfrac{}{}{0pt}{}{xy_j}{j=1,2,\mathrm{},n}}{}|u^I(x+)u^I(x)|+|u^{II}(x+)u^{II}(x)|<ϵ.$$
$`3.31`$
To each $`y_j`$ we associate the corresponding discontinuity point $`y_j^h`$ in $`u^{I,h}`$ or $`u^{II,h}`$. To simplify the presentation we will focus on the case where $`y_j<y_j^h<y_{j+1}<y_{j+1}^h`$ for all $`j`$. The other cases can be treated similarly. In view of the local convergence property (3.27)–(3.28) and by extracting a covering of the interval $`[y_0,y_n]`$, we have also
$$|u^{I,h}(x)u^I(x)|+|u^{II,h}(x)u^{II}(x)|2ϵ,x(y_j^h,y_{j+1})(y_j,y_{j+1}).$$
$`3.32`$
In view of (3.30) we can construct functions $`u_ϵ^I`$ and $`u_ϵ^{II}`$ that are continuous everywhere except possibly at the points $`y_j`$ and such that the following conditions hold with $`u`$ replaced by either $`u^I`$ or $`u^{II}`$:
$`TV(u_ϵ;RI\{y_1,\mathrm{},y_n\})CTV(u;RI\{y_1,\mathrm{},y_n\}),`$ $`3.33`$
$`uu_ϵ_{\mathrm{}}Cϵ,TV(uu_ϵ;RI\{y_1,\mathrm{},y_n\})Cϵ,`$
where $`C`$ is independent of $`ϵ`$.
Consider the decompositions
$$_{RI}\left(a(u^{I,h},u^{II,h})f^{}(u^{I,h})\right)(u^{II,h}u^{I,h})𝑑V^{I,h}=\underset{j=0}{\overset{n}{}}_{(y_j^h,y_{j+1}^h)}\mathrm{}+\underset{j=1}{\overset{n}{}}_{\{y_j^h\}}\mathrm{}$$
and
$$_{RI}\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)𝑑\overline{V}^I=\underset{j=0}{\overset{n}{}}_{(y_j,y_{j+1})}\mathrm{}+\underset{j=1}{\overset{n}{}}_{\{y_j\}}\mathrm{}.$$
Here $`y_0^h=y_0=\mathrm{}`$ and $`y_{n+1}^h=y_{n+1}=+\mathrm{}`$. Thus in (3.30) we have to estimate
$`\mathrm{\Omega }_2^h(\tau )`$ $`={\displaystyle _{RI}}\left(a(u^{I,h},u^{II,h})f^{}(u^{I,h})\right)(u^{II,h}u^{I,h})𝑑V^{I,h}`$ $`3.34`$
$`{\displaystyle _{RI}}\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)𝑑\overline{V}^I`$
$`=T_1^h+T_2^h`$
with
$`T_1^h:=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{\{y_j^h\}}}\left(a(u^{I,h},u^{II,h})f^{}(u^{I,h})\right)(u^{II,h}u^{I,h})𝑑V^{I,h}`$
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{\{y_j\}}}\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)𝑑\overline{V}^I`$
and
$`T_2^h:=`$ $`{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{(y_j^h,y_{j+1}^h)}}\left(a(u^I,u^{II})f^{}(u^{I,h})\right)(u^{II,h}u^{I,h})𝑑V^{I,h}`$
$`{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{(y_j,y_{j+1})}}\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)𝑑\overline{V}^I.`$
First, relying on the convergence property (3.29) we have immediately
$`T_1^h`$ $`={\displaystyle \underset{j=1}{\overset{n}{}}}\left(a(u^{I,h}(y_j^h),u^{II}(y_j^h))\lambda ^{I,h}(y_j^h)\right)\left(u^{II,h}(y_j^h)u^{I,h}(y_j^h)\right)|u^{I,h}(y_j^h+)u^{I,h}(y_j^h)|`$
$`\left(a(u^I(y_j),u^{II}(y_j))\lambda ^I(y_j)\right)\left(u^{II}(y_j)u^I(y_j)\right)|u^I(y_j+)u^I(y_j)|,`$
so that
$$\left|T_1^h\right|C\underset{j=1}{\overset{n}{}}\underset{\pm }{}|u^{I,h}(y_j^h\pm )u^I(y_j\pm )|+|u^{II,h}(y_j^h\pm )u^{II}(y_j\pm )|.$$
Thus, in view of the local convergence at jump points (3.27a), for $`h`$ small enough we obtain
$$\left|T_1^h\right|Cϵ.$$
$`3.35`$
Relying on the simplifying assumption $`y_j<y_j^h<y_{j+1}<y_{j+1}^h`$ for all $`j`$, we can decompose $`T_2^h`$ as follows:
$`T_2^h`$ $`={\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{(y_j^h,y_{j+1})}}\left(a(u^{I,h},u^{II,h})f^{}(u^{I,h})\right)(u^{II,h}u^{I,h})𝑑V^{I,h}`$ $`3.36`$
$`\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)d\overline{V}^I`$
$`{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{(y_j,y_j^h]}}\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)𝑑\overline{V}^I`$
$`+{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{[y_{j+1},y_{j+1}^h)}}\left(a(u^{I,h},u^{II,h})f^{}(u^{I,h})\right)(u^{II,h}u^{I,h})𝑑V^{I,h}`$
$`=:T_{2,1}^h+T_{2,2}^h+T_{2,3}^h.`$
We first consider $`T_{2,2}^h`$:
$`T_{2,2}^h={\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{(y_j,y_j^h]}}`$ $`\left(a(u^I(y_j+),u^{II}(y_j+))f^{}(u^I(y_j+))\right)\left(u^{II}(y_j+)u^I(y_j+)\right)d\overline{V}^I(y)`$
$`+{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{(y_j,y_j^h]}}`$ $`\{(a(u^I(y),u^{II}(y))f^{}(u^I(y)))(u^{II}(y)u^I(y))`$
$`(a(u^I(y_j+),u^{II}(y_j+))f^{}(u^I(y_j+)))(u^{II}(y_j+)u^I(y_j+))\}d\overline{V}^I(y).`$
Therefore, with (3.17), we obtain
$`\left|T_{2,2}^h\right|`$ $`C{\displaystyle \underset{j}{}}\left|\overline{V}^I(y_j+)\overline{V}^I(y_j^h+)\right|`$
$`+CV^I(+\mathrm{})\left(\underset{y(y_j,y_j^h]}{sup}|u^I(y)u^I(y_j+)|+\underset{x(y_j,y_j^h]}{sup}|u^{II}(y)u^{II}(y_j+)|\right).`$
Since $`y_j^hy_j`$, we have $`\left|\overline{V}^I(y_j+)\overline{V}^I(y_j^h+)\right|0`$, therefore for $`h`$ sufficiently small
$$\left|T_{2,2}^h\right|Cϵ.$$
$`3.37`$
A similar argument for $`T_{2,3}^h`$ shows that
$$\left|T_{2,3}^h\right|Cϵ.$$
$`3.38`$
Next consider the decomposition
$`(a(`$ $`u^{I,h},u^{II,h})f^{}(u^{I,h}))(u^{II,h}u^{I,h})dV^{I,h}(a(u^I,u^{II})f^{}(u^I))(u^{II}u^I)d\overline{V}^I`$
$`=(a`$ $`(u^{I,h},u^{II,h})f^{}(u^{I,h}))(u^{II,h}u^{I,h})dV^{I,h}(a(u^I,u^{II})f^{}(u^I))(u^{II}u^I)dV^{I,h}`$
$`+\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)dV^{I,h}\left(a(u_ϵ^I,u_ϵ^{II})f^{}(u_ϵ^I)\right)(u_ϵ^{II}u_ϵ^I)dV^{I,h}`$
$`+\left(a(u_ϵ^I,u_ϵ^{II})f^{}(u_ϵ^I)\right)(u_ϵ^{II}u_ϵ^I)dV^{I,h}\left(a(u_ϵ^I,u_ϵ^{II})f^{}(u_ϵ^I)\right)(u_ϵ^{II}u_ϵ^I)d\overline{V}^I`$
$`+\left(a(u_ϵ^I,u_ϵ^{II})f^{}(u_ϵ^I)\right)(u_ϵ^{II}u_ϵ^I)d\overline{V}^I\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)d\overline{V}^I,`$
which, with obvious notation, yields a decomposition for $`T_{2,1}^h`$
$$T_{2,1}^h=M_1^h+M_2^h+M_3^h+M_4^h.$$
$`3.39`$
Using (3.17) and the local convergence property (3.31), we obtain
$`|M_1^h|`$ $`C{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{(y_j^h,y_{j+1})}}\left|dV^{I,h}\right|\left(\underset{(y_j^h,y_{j+1})}{sup}|u^{I,h}u^I|+\underset{(y_j^h,y_{j+1})}{sup}|u^{II,h}u^{II}|\right)`$ $`3.40`$
$`Cϵ.`$
Similarly using (3.17) and (3.33) we obtain
$`|M_2^h|`$ $`C{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{(y_j^h,y_{j+1})}}\left|dV^{I,h}\right|\left(\underset{(y_j^h,y_{j+1})}{sup}|u^Iu_ϵ^I|+\underset{(y_j^h,y_{j+1})}{sup}|u^{II}u_ϵ^{II}|\right)`$ $`3.41`$
$`Cϵ.`$
Dealing with $`M_4^h`$ is similar:
$`|M_4^h|`$ $`C{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle _{(y_j^h,y_{j+1})}}\left|d\overline{V}^I\right|\left(\underset{(y_j^h,y_{j+1})}{sup}|u^Iu_ϵ^I|+\underset{(y_j^h,y_{j+1})}{sup}|u^{II}u_ϵ^{II}|\right)`$ $`3.42`$
$`Cϵ.`$
Finally to treat $`M_3^h`$ we observe that, since $`u_ϵ^I`$ and $`u_ϵ^{II}`$ are continuous functions on each interval $`(y_j^h,y_{j+1})`$ and since $`dV^{I,h}`$ is sequence of bounded measures converging weakly-star toward $`d\overline{V}^I`$, we have for all $`h`$ sufficiently small
$$|M_3^h|ϵ.$$
$`3.43`$
Combining (3.39)–(3.43) we get
$$\left|T_{2,1}^h\right|Cϵ.$$
$`3.44`$
Combining (3.36)–(3.38) and (3.44) we obtain
$$\left|T_2^h\right|Cϵ$$
and thus with (3.34)-(3.35)
$$\left|\mathrm{\Omega }_2^h(\tau )\right|Cϵ\text{ for all }h\text{ sufficiently small.}$$
Since $`ϵ`$ is arbitrary, this completes the proof of (3.30).
Step 3 : Consider now the term
$$\mathrm{\Omega }_1(\tau )=\underset{(x,\tau )(a^h)𝒥(u^{I,h})}{}\left(2m+TV(b^h)\right)\left|a^h(x,\tau )\lambda ^h(x,\tau )\right||u^{II,h}(x,\tau )u^{I,h}(x,\tau )|.$$
$`3.45`$
On one hand, observe that
$$TV(b^h(\tau ))=TV(u^{I,h}(\tau ))+TV(u^{II,h}(\tau ))\overline{V}^I(+\mathrm{},\tau )+\overline{V}^{II}(+\mathrm{},\tau ).$$
$`3.46`$
For all but countably many $`\tau `$ the following holds. Extracting a subsequence if necessary we can always assume that for each $`j`$ either $`(y_j^h,\tau )(a^h)`$ for all $`h`$, or else $`(y_j^h,\tau )(a^h)`$ for all $`h`$. Then consider the following three sets: denote by $`J_1`$ the set of indices $`j`$ such that $`(y_j^h,\tau )(a^h)`$ and $`(y_j,\tau )(a)`$. Let $`J_2`$ the set of indices $`j`$ such that $`(y_j^h,\tau )(a^h)`$ and $`(y_j,\tau )(a)`$. Finally $`J_3`$ is the set of indices $`j`$ such that $`(y_j^h,\tau )(a^h)`$ and $`(y_j,\tau )(a)`$.
The local convergence property (3.27) implies
$`{\displaystyle \underset{jJ_1}{}}\left|a(u^{I,h}(y_j^h),u^{II,h}(y_j^h))a(u^{I,h}(y_j^h),u^{I,h}(y_j^h+))\right||u^{II,h}(y_j^h)u^{I,h}(y_j^h)|`$ $`3.47`$
$`{\displaystyle \underset{jJ_1}{}}\left|a(u^I(y_j),u^{II}(y_j))a(u^I(y_j),u^I(y_j+))\right||u^{II}(y_j)u^I(y_j)|.`$
(Indeed, given $`ϵ>0`$, choose finitely many jump points as in (3.31) and use (3.27) with $`ϵ`$ replaced with $`ϵ|u^I(z+)u^I(z+)|`$).
On the other hand for indices in $`J_2`$ or $`J_3`$ we have
$$\underset{jJ_2J_3}{}\left|a(u^{I,h}(y_j^h),u^{II,h}(y_j^h))a(u^{I,h}(y_j^h),u^{I,h}(y_j^h+))\right||u^{II,h}(y_j^h)u^{I,h}(y_j^h)|0$$
$`3.48`$
but
$$\underset{jJ_2J_3}{}\left|a(u^I(y_j),u^{II}(y_j))a(u^I(y_j),u^I(y_j+))\right||u^{II}(y_j)u^I(y_j)|=0.$$
$`3.49`$
Indeed, for each $`jJ_2`$, $`y_j`$ is a Lax shock but $`y_j^h`$ is not. Extracting a subsequence if necessary, it must be that the Lax inequalities are violated on the left or on the right side of $`y_j^h`$ for all $`h`$. So it must be that, assuming that it is the case on the left side, $`a(u^I(y_j),u^{II}(y_j))a(u^I(y_j),u^I(y_j+)0`$ while $`a(u^I(y_j^h),u^{II}(y_j^h))a(u^I(y_j^h),u^I(y_j^h+)0`$ for all $`h`$. But the latter converges toward the former by the local uniform convergence, which proves that $`a(u^I(y_j),u^{II}(y_j))a(u^I(y_j),u^I(y_j+)=0`$.
Combining (3.45)–(3.49) yields
$$\mathrm{\Omega }_1^h_0^t\underset{(x,\tau )(a)}{}\overline{q}(\tau )\left|a(x,\tau )\lambda (x,\tau )\right||u^{II}(x,\tau )u^I(x,\tau )|d\tau ,$$
$`3.50`$
where $`\overline{q}:=2m+\overline{V}^I(+\mathrm{})+\overline{V}^{II}(+\mathrm{})`$.
Step 4 : Continuity of the weighted norm.
Fix some time $`t`$. Recall that the weight $`\overline{w}(t)`$ is defined based on the total variation functions $`\overline{V}^{II}`$ and $`\overline{V}^I`$ and on the function $`u^{II}(t)u^I(t)`$. The weight $`w^h(t)`$ is defined based on the total variation functions $`V^{II,h}`$ and $`V^{I,h}`$ and on the function $`u^{II,h}(t)u^{I,h}(t)`$. On the other hand, $`u^{II,h}u^{I,h}(t)u^{II}u^I(t)`$, $`V^{II,h}\overline{V}^{II}`$ and $`V^{I,h}\overline{V}^I`$. Therefore we have
$$w(x,t)=\overline{w}(x,t)\text{ whenever }u^{II}(x,t)u^I(x,t)0.$$
$`3.51`$
Combining (3.51) and the $`L^1`$ convergence $`u^{II,h}u^{I,h}(t)u^{II}u^I(t)`$, we have
$$u^{II}(t)u^I(t)_{\overline{w}(t)}=\underset{h0}{lim}u^{II,h}(t)u^{I,h}(t)_{w^h(t)}.$$
$`3.52`$
Step 5 : The left-hand side of $`(3.16)`$ is bounded above by $`(3.21)`$.
First of all, the inequality
$`{\displaystyle _0^t}{\displaystyle _{RI}}\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)𝑑V^I(y,\tau )𝑑\tau `$ $`3.53`$
$`{\displaystyle _0^t}{\displaystyle _{RI}}\left(a(u^I,u^{II})f^{}(u^I)\right)(u^{II}u^I)𝑑\overline{V}^I(y,\tau )𝑑\tau `$
is a direct consequence of (3.19) and the definition of the nonconservative product in (3.15).
On the other hand, by the definition of the weighted norm and because of (3.20), similarly to (3.51) we have the inequality
$$w(x,t)\overline{w}(x,t)\text{ whenever }u^{II}(x,t)u^I(x,t)0.$$
$`3.54`$
Hence, by (3.53) and (3.54) the left-hand side of (3.16) is bounded above by (3.21). This completes the proof of Theorem 3.5. $`\mathit{}`$
## 4. Generalized Characteristics and Maximum Principle
We now return to the setting in Section 2 and aim at extending the analysis therein to arbitrary functions of bounded variation. For exact solutions of the hyperbolic equation
$$_t\psi +_x\left(a\psi \right)=0,$$
$`4.1`$
we will establish a maximum principle: Any solution of $`(4.1)`$ remains non-negative for all times if it is so initially. For a more precise (local) statement, our proof will make use of Dafermos-Filippov theory of generalized characteristics.
Our main assumption throughout this section is the following:
$$\text{There exists a constant }E\text{ such that }_xa\frac{E}{t}.$$
$`4.2`$
This is nothing but a generalization of the well-known Oleinik’s entropy inequality. To motivate (4.2), let us recall the following result.
Let $`f`$ be a strictly convex function and $`u`$ be an entropy solution (with bounded variation for all times) of the conservation law
$$_tu+_xf(u)=0,u(x,t)RI.$$
$`4.3`$
Then is is known that there exists a constant $`C=C(u)`$ such that
$$_xu\frac{C}{t}.$$
$`4.4`$
###### Lemma 4.1
If $`u^I`$ and $`u^{II}`$ are two entropy solutions of the conservation law $`(4.3)`$, then the averaging speed
$$a=a(u^I,u^{II}):=\frac{f(u^{II})f(u^I)}{u^{II}u^I}.$$
$`4.5`$
satisfies our assumption $`(4.2)`$, with $`E=supf^{\prime \prime }\left(C(u^I)+C(u^{II})\right)/2`$.
###### Demonstration Proof
Let us fix some time $`t>0`$. On each Borel set consisting of points of continuity of both $`u^I`$ and $`u^{II}`$, the following holds:
$`_xa`$ $`=_x{\displaystyle _0^1}f^{}\left(\theta u^I+(1\theta )u^{II}\right)𝑑\theta `$
$`={\displaystyle _0^1}f^{\prime \prime }\left(\theta u^I+(1\theta )u^{II}\right)\left(\theta _xu^I+(1\theta )_xu^{II}\right)𝑑\theta `$
$`{\displaystyle _0^1}supf^{\prime \prime }\left(\theta {\displaystyle \frac{C(u^I)}{t}}+(1\theta ){\displaystyle \frac{C(u^{II})}{t}}\right)d\theta `$
$`supf^{\prime \prime }{\displaystyle \frac{C(u^I)+C(u^{II})}{2t}}.`$
On the other hand, at a point $`x`$ where one of $`u^I`$ or $`u^{II}`$ is discontinuous, we have with an obvious notation
$$a_+a_{}=_0^1f^{}\left(\theta u_+^I+(1\theta )u_+^{II}\right)𝑑\theta _0^1f^{}\left(\theta u_{}^I+(1\theta )u_{}^{II}\right)𝑑\theta 0,$$
since $`f^{}`$ is an increasing function and (for instance by (4.4)) both $`u^I`$ and $`u^{II}`$ satisfy $`u_+^Iu_{}^I`$ and $`u_+^{II}u_{}^{II}`$. $`\mathit{}`$
By definition, a generalized characteristic $`y=y(t)`$ associated with the coefficient $`a`$ must satisfy for almost every $`t`$ (in its domain of definition)
$$a_+(y(t),t)y^{}(t)a_{}(y(t),t).$$
$`4.6`$
According to Filippov’s theory of differential equations , through each point $`(\overline{x},\overline{t})`$ there pass a maximal and a minimal generalized characteristic.
###### Definition 4.2
A generalized characteristic is said to be genuine iff for almost every $`t`$ it satisfies
$$y^{}(t)\{a_{}(y(t),t),a_+(y(t),t)\}.$$
$`4.7`$
###### Proposition 4.3
Any minimal backward generalized characteristic is genuine and for almost every $`t`$ satisfies
$$y^{}(t)=a_{}(y(t),t).$$
$`4.8`$
Similarly, for a maximal backward generalized characteristic we have $`y^{}(t)=a_+(y(t),t)`$.
###### Demonstration Proof
Here we only rely on the following consequence of (4.2): $`a_+a_{}`$ at each discontinuity point of the function $`a`$. Geometrically, this condition prevents the existence of rarefaction-shocks in $`a`$. On the other hand, rarefaction centers (also prevented by (4.2) for $`t>0`$) could still be allowed for the present purpose.
Consider $`(\overline{x},\overline{t})(\mathrm{},+\mathrm{})\times (0,\mathrm{})`$, and let $`y(t):=y(t;\overline{x},\overline{t})`$ be the minimal backward characteristic through $`(\overline{x},\overline{t})`$. We prove that it is genuine on its domain $`(s,\overline{t}]`$. We proceed as in and assume by contradiction that there is a measurable set $`J,\overline{J}(s,\overline{t}]`$ of positive Lebesgue measure, and $`\epsilon >0`$ such that
$$a_{}(y(t),t)y^{}(t)>2\epsilon ,tJ.$$
$`4.9`$
For each $`tJ`$ there exists $`\delta (t)>0`$ with the property
$$a_+(x,t)a_{}(y(t),t)\epsilon ,x(y(t)\delta (t),y(t)).$$
$`4.10`$
Finally, there is a subset $`IJ`$ with $`\mu ^{}(I)>0`$ (here $`\mu ^{}`$ denotes the outer measure) and $`\overline{\delta }>0`$ such that $`\delta (t)>\overline{\delta }`$ for $`tI`$.
Let $`\tau `$ be a density point of $`I`$, with respect to $`\mu ^{}`$. Thus there exists $`\overline{r},0<\overline{r}<\overline{t}\tau `$, so that
$$\frac{\mu ^{}(I[\tau ,\tau +r])}{r}>\frac{2|\alpha |+\epsilon }{2|\alpha |+2\epsilon },0<r\overline{r},$$
$`4.11`$
where
$$\alpha :=inf\{a_+(x,t)a_{}(y(t),t):s<t\overline{t},y(t)\overline{\delta }x<y(t)\}.$$
Now take a point $`y(y(\tau )\overline{\delta },y(\tau ))`$ with the property $`y>y(\tau )\frac{1}{2}\epsilon \overline{r}`$, and consider a forward characteristic $`z()`$ through $`(y,\tau )`$. We first observe that
$$z(t)<y(t),t>\tau ,$$
since $`y(t)`$ is the minimal backward characteristic through $`(\overline{x},\overline{t})`$.
In addition, we have
$$z(t)>y(t)\overline{\delta },t>[\tau ,\tau +\overline{r}].$$
Indeed, suppose by contradiction that for some $`r(0,\overline{r}]`$, $`z(t)>y(t)\overline{\delta }`$ for $`t>[\tau ,\tau +r)`$, but $`z(\tau +r)=y(\tau +r)\overline{\delta }`$. Then
$`0`$ $`=z(\tau +r)y(\tau +r)\overline{\delta }=y+{\displaystyle _\tau ^{\tau +r}}z^{}(t)𝑑ty(t){\displaystyle _\tau ^{\tau +r}}y^{}(t)𝑑t+\overline{\delta }`$
$`>{\displaystyle _\tau ^{\tau +r}}(z^{}(t)y^{}(t))𝑑t`$
$`={\displaystyle _{I[\tau ,\tau +r]}}(z^{}(t)a_{}(y(t),t)+a_{}(y(t),t)y^{}(t))𝑑t`$
$`+{\displaystyle _{[\tau ,\tau +r]I}}(z^{}(t)a_{}(y(t),t)+a_{}(y(t),t)y^{}(t))𝑑t`$
$`\epsilon \mu ^{}(I[\tau ,\tau +r])+\alpha (r\mu ^{}(I[\tau ,\tau +r]))>0,`$
by (4.9)-(4.11), which leads to a contradiction. In the same way one obtains
$`0`$ $`>z(\tau +\overline{r})y(\tau +\overline{r})=y+{\displaystyle _\tau ^{\tau +\overline{r}}}z^{}(t)𝑑ty(t){\displaystyle _\tau ^{\tau +\overline{r}}}y^{}(t)𝑑t`$
$`>\epsilon \mu ^{}(I[\tau ,\tau +\overline{r}])+\alpha (\overline{r}\mu ^{}(I[\tau ,\tau +r])){\displaystyle \frac{1}{2}}\epsilon \overline{r}>0,`$
which gives another contradiction. For the maximal backward characteristic the proof is similar. $`\mathit{}`$
###### Proposition 4.4
Forward characteristics leaving from some $`(\overline{x},\overline{t})`$ are unique when $`\overline{t}>0`$.
###### Demonstration Proof
Suppose there were two forward characteristics $`y()`$ and $`z()`$ through $`(\overline{x},\overline{t})`$ with $`y(\tau )<z(\tau )`$ for some $`\tau >\overline{t}`$. By (4.2) we have
$$z^{}(\tau )y^{}(\tau )a_{}(z(t),t)a_+(y(t),t)C_{\overline{t}}(z(\tau )y(\tau )).$$
$`4.12`$
Integrating (4.12) from $`\overline{t}`$ to $`\tau `$ one gets $`z(\tau )y(\tau )=0`$, which gives a contradiction. $`\mathit{}`$
###### Theorem 4.5
Let $`\psi =\psi (x,t)`$ be a solution of $`(4.1)`$ such that on some interval $`[\xi _0,\zeta _0]`$ we have
$$\psi (x,0)0,x[\xi _0,\zeta _0].$$
$`4.13`$
Let $`\xi =\xi (t)`$ be any forward generalized characteristic leaving from $`(\xi _0,0)`$, and $`\zeta =\zeta (t)`$ be any forward generalized characteristic leaving from $`(\zeta _0,0)`$.
Then we have for all $`t0`$
$$\psi (x,t)0,x(\xi (t),\zeta (t)).$$
$`4.14`$
Note that it may happen that $`\xi (t)=\zeta (t)`$ for $`t`$ large enough.
###### Demonstration Proof
Observe that the two characteristics cannot cross and fix any time $`t>0`$ such that $`\xi (t)<\zeta (t)`$. Fix also any two points such that $`\xi (t)<\overline{y}<\overline{z}<\zeta (t)`$. Let $`y(t)`$ and $`z(t)`$ be the maximal and minimal backward characteristics emanating from $`\overline{y}`$ and $`\overline{z}`$, respectively. These characteristics can not leave the region limited by $`\xi (t)`$ and $`\zeta (t)`$.
Integrating (4.1) in the domain bounded by the characteristics $`y(t)`$ and $`z(t)`$, and using that these characteristics are genuine, so that the flux terms along the vertical boundaries vanish identically, we arrive at
$$_{\overline{y}}^{\overline{z}}\psi (x,t)𝑑x=_{y(0)}^{z(0)}\psi (x,0)𝑑x0.$$
$`4.15`$
The last inequality is due to the fact that $`\psi (.,0)0`$ and the inequalities $`\xi _0=\xi (0)y(0)z(0)\zeta (0)=\zeta _0`$. Since $`\overline{y}`$ and $`\overline{z}`$ are arbitrary, we obtain (4.14). $`\mathit{}`$
## 5. A Sharp $`L^1`$ Estimate for Hyperbolic Linear Equations
Based on the maximum principle established in Section 4, we now derive a sharp estimate for the weighted norm introduced in Section 2. We restrict attention again to the situation where $`u^I`$ and $`u^{II}`$ are two entropy solutions of the conservation law $`(4.3)`$ and $`a`$ is the averaging speed given in (4.5). We define a weight by analogy with what was done in Section 2 in the special case of piecewise constant solutions.
Given a solution $`\psi `$ of the equation (4.1), we introduce weighted $`L^1`$ norm in the following way. Set
$$V^I(x,t)=TV_{\mathrm{}}^x(u^I(t)),V^{II}(x,t)=TV_{\mathrm{}}^x(u^{II}(t))$$
$`5.1`$
and fix some parameter $`m0`$. Then consider the weight-function defined, for each $`t0`$ and each point of continuity $`x`$ for $`u^I(t)`$ and $`u^{II}(t)`$, by
$$w(x,t)=\{\begin{array}{cc}m+V^I(\mathrm{},t)V^I(x,t)+V^{II}(x,t)\hfill & \text{if }\psi (x,t)>0,\hfill \\ & \\ m+V^I(x,t)+V^{II}(\mathrm{},t)V^{II}(x,t)\hfill & \text{if }\psi (x,t)0.\hfill \end{array}$$
$`5.2`$
It is immediate to see that
$$mw(x,t)m+TV(u^I(t))+TV(u^{II}(t)),xRI.$$
$`5.3`$
Finally the weighted norm on the solutions $`\psi `$ of (4.1) is defined by
$$\psi (t)_{w(t)}:=_{RI}|\psi (x,t)|w(x,t)𝑑x.$$
Note that the weight depends on the fixed solutions $`u^I`$ and $`u^{II}`$, but also on the solution $`\psi `$.
Our sharp estimate will involve the nonconservative product
$$\mu _\psi ^I(t)=\left(af^{}(u^I(t))\right)\psi (t)dV^I(t)$$
defined for all almost every $`t0`$ by
Here $`\lambda ^I(x,t)`$ is a the shock speed of the discontinuity in $`u^I`$ located at $`(x,t)`$. The measure $`\mu _\psi ^{II}(t)`$ is defined similarly. Regarding the expression (5.4b), it is worth noting that if $`(x,t)`$ is a point of approximate jump of $`u^I`$ and $`\psi `$, then the jump relation for the equation (4.1) reads
$$\left(a(x,t)\lambda ^I(x,t)\right)\psi (x,t)=\left(a(x+,t)\lambda ^I(x,t)\right)\psi (x+,t).$$
$`5.5`$
In the same way we define
$$\mu _\psi ^{II}(t)=\left(f^{}(u^{II}(t))a\right)\psi (t)dV^{II}(t).$$
We now prove:
###### Theorem 5.1
Let $`u^I`$ and $`u^{II}`$ be two entropy solutions of $`(1.1)`$ such that $`u^{II}u^I`$ admits finitely many changes of sign. Let $`\psi `$ be any solution of bounded variation of the hyperbolic equation $`(4.1)`$ satisfying the constrain
$$\psi \left(u^{II}u^I\right)0.$$
$`5.6`$
Then for all $`0st`$
$`\psi (t)_{w(t)}`$ $`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}\left(2m+TV(a)\right)\left|a_{}(x,\tau )\lambda (x,\tau )\right||\psi _{}(x,\tau )|d\tau `$ $`5.7`$
$`+{\displaystyle _s^t}{\displaystyle _{RI}}\left(a(\tau )f^{}(u^I(\tau ))\right)\psi (\tau )𝑑V^I(\tau )𝑑\tau +{\displaystyle _s^t}{\displaystyle _{RI}}\left(a(\tau )f^{}(u^{II}(\tau ))\right)\psi (\tau )𝑑V^{II}(\tau )𝑑\tau `$
$`\psi (s)_{w(s)}.`$
The assumption (5.6) is clearly satisfied with the choice $`\psi =u^{II}u^I`$. Therefore our previous result in Theorem 3.5 (derived via a completely different proof) can be regarded as a corollary of Theorem 5.1.
It is interesting to observe that, when $`u^{II}=u^I`$, the weight (5.2) becomes constant, and therefore (5.7) reduces to the $`L^1`$ estimate.
$`\psi (t)_{L^1(RI)}`$ $`+{\displaystyle _s^t}{\displaystyle \underset{(x,\tau )(a)}{}}\left(2m+TV(a)\right)\left|a_{}(x,\tau )\lambda (x,\tau )\right||\psi _{}(x,\tau )|d\tau `$
$`\psi (s)_{L^1(RI)}.`$
Also, note that under the assumption (5.6) $`\mu _\psi ^I(t)`$ and $`\mu _\psi ^{II}(t)`$ are positive except at points $`(x,t)(a)(a)`$. However, these negative terms are offset in (5.7) by the positve terms under the first integral.
###### Demonstration Proof
Fix any positive time $`t`$. By assumption we have finitely many points $`\mathrm{}=y_0<y_1<\mathrm{}<y_n<y_{n+1}=+\mathrm{}`$ such that, on each interval $`(y_i,y_{i+1})`$, we have $`\psi (t)0`$ when $`i`$ is odd and $`\psi (t)0`$ when $`i`$ is even. For every $`i=1,\mathrm{},n`$, consider the (unique by Proposition 4.4) forward characteristic $`y_i()`$ associated with the coefficient $`a`$ and issuing from the initial point $`(y_i,t)`$.
We will focus attention on some interval $`(y_i,y_{i+1})`$ with $`i`$ odd, say, and with $`\mathrm{}<y_i<y_{i+1}<+\mathrm{}`$. Except when specified differently, all of the characteristics to be considered from now on are associated with the solution $`u^{II}`$. For definiteness we will first study the case that the forward characteristic $`\chi _0()`$ (associated with $`u^{II}`$ and) issuing from the point $`(y_i,t)`$ is located on the right-side of the curve $`y_i`$, that is,
$$y_i(\tau )\chi _0(\tau ),t\tau t+\delta $$
for some $`\delta >0`$ sufficiently small.
Fix some (sufficiently small) $`ϵ>0`$ and denote by $`y_i<z_1<\mathrm{}<z_N<y_{i+1}`$ the points where $`u^I`$ has a jump larger or equal to $`ϵ`$, that is,
$$u_{}^{II}(z_I,t)u_+^{II}(z_I,t)ϵ,I=1,\mathrm{},N.$$
$`5.8`$
For each $`I=1,\mathrm{},N`$, consider also the forward characteristic $`\chi _I()`$ issuing from the point $`(z_I,t)`$. For definiteness, we will also assume that the forward characteristic $`\chi _{N+1}()`$ issuing from $`(y_{i+1},t)`$ satisfies
$$\chi _{N+1}(\tau )y_{i+1}(\tau ),t\tau t+\delta $$
for some $`\delta >0`$ sufficiently small.
Next, let us select a time $`s>t`$ with $`st`$ so small that the following properties hold:
For $`I=0,\mathrm{},N`$, and some integer $`k`$ to be fixed later, consider a mesh of the form
$$\chi _I(s)=x_I^0<x_I^1<\mathrm{}<x_I^k<x_I^{k+1}=\chi _{I+1}(s).$$
$`5.9`$
For $`I=0,\mathrm{},N`$ and $`j=1,\mathrm{},k`$, consider also the maximal backward characteristic $`\xi _I^j()`$ emanating from the point $`(x_I^j,s)`$ and identify its intercept $`z_I^j=\xi _I^j(t)`$ by the horizontal line at time $`t`$. Finally set also
$$z_0^0=y_i,z_N^{k+1}=y_{i+1},z_{I1}^{k+1}=z_I^0=z_I,I=1,\mathrm{},N.$$
To start the proof, we integrate the equation (4.1) satisfied by the function $`\psi `$, successively in each domain limited by the characteristics introduced above. Applying Green’s theorem, we arrive at the following five formulas:
$`(i)`$ Integrating (4.1) on the region
$$\left\{(x,\tau )/t<\tau <s,y_i(\tau )<x<\chi _0(\tau )\right\}$$
and multiplying by $`V^{II}(y_i,t)`$ one gets
$`{\displaystyle _{y_i(s)}^{\chi _0(s)}}\psi (x,s)V^{II}(y_i,t)𝑑x`$ $`+{\displaystyle _t^s}(y_i^{}a_+)\psi _+(y_i(\tau ),\tau )V^{II}(y_i,t)𝑑\tau `$ $`5.10i`$
$`+{\displaystyle _t^s}(a_{}\lambda _0)\psi _{}(\chi _0(\tau ),\tau )V^{II}(y_i,t)𝑑\tau =0.`$
$`(ii)`$ Integrating (4.1) on each of the regions
$$\left\{(x,\tau )/t<\tau <s,\xi _I^j(\tau )<x<\xi _I^{j+1}(\tau )\right\}$$
for $`I=0,\mathrm{},N`$ and $`j=1,\mathrm{},k`$, and then multiplying by $`V^{II}(z_I^j+,t)`$, one gets
$`{\displaystyle _{x_I^j}^{x_I^{j+1}}}\psi (x,s)V^{II}(z_I^j+,t)𝑑x{\displaystyle _{z_I^j}^{z_I^{j+1}}}\psi (x,t)V^{II}(z_I^j+,t)𝑑x`$ $`5.10ii`$
$`+{\displaystyle _t^s}(\lambda _I^ja_+)\psi _+(\xi _I^j(\tau ),\tau )V^{II}(z_I^j+,t)𝑑\tau +{\displaystyle _t^s}(a_{}\lambda _I^{j+1})\psi _{}(\xi _I^{j+1}(\tau ),\tau )V^{II}(z_I^j+,t)𝑑\tau =0.`$
$`(iii)`$ Integrating (4.1) on each of the regions
$$\left\{(x,\tau )/t<\tau <s,\chi _I(\tau )<x<\xi _I^1(\tau )\right\}$$
for $`I=0,\mathrm{},N`$, and multiplying by $`V^{II}(z_I+,t)`$ one gets
$`{\displaystyle _{\chi _I(s)}^{x_I^1}}\psi (x,s)V^{II}(z_I+,t)𝑑x{\displaystyle _{z_I}^{z_I^1}}\psi (x,t)V^{II}(z_I+,t)𝑑x`$ $`5.10iii`$
$`+{\displaystyle _t^s}(\lambda _Ia_+)\psi _+(\chi _I(\tau ),\tau )V^{II}(z_I+,t)𝑑\tau +{\displaystyle _t^s}(a_{}\lambda _I^1)\psi _{}(\xi _I^1(\tau ),\tau )V^{II}(z_I+,t)𝑑\tau =0.`$
$`(iv)`$ Integrating (4.1) on the regions
$$\left\{(x,\tau )/t<\tau <s,\xi _I^k(\tau )<x<\chi _{I+1}(\tau )\right\}$$
for $`I=0,\mathrm{},N`$, and multiplying by $`V^{II}(z_I^k+,t)`$ one gets
$`{\displaystyle _{x_I^k}^{\chi _{I+1}(s)}}\psi (x,s)V^{II}(z_I^k+,t)𝑑x{\displaystyle _{z_I^k}^{z_{I+1}}}\psi (x,t)V^{II}(z_I^k+,t)𝑑x`$ $`5.10iv`$
$`+{\displaystyle _t^s}(\lambda _I^ka_+)\psi _+(\xi _I^k(\tau ),\tau )V^{II}(z_I^k+,t)𝑑\tau +{\displaystyle _t^s}(a_{}\lambda _{I+1})\psi _{}(\chi _{I+1}(\tau ),\tau )V^{II}(z_I^k+,t)𝑑\tau =0.`$
$`(v)`$ Finally integrating (4.1) on the last region
$$\left\{(x,\tau )/t<\tau <s,\chi _{N+1}(\tau )<x<y_{i+1}(\tau )\right\}$$
and multiplying by $`V^{II}(y_{i+1},t)`$ one gets
$`{\displaystyle _{\chi _{N+1}(s)}^{y_{i+1}(s)}}\psi (x,s)V^{II}(y_{i+1},t)𝑑x+{\displaystyle _t^s}(\lambda _{N+1}a_+)\psi _+(\chi _{N+1}(\tau ),\tau )V^{II}(y_{i+1},t)𝑑\tau `$ $`5.10v`$
$`+{\displaystyle _t^s}(a_{}y_{i+1}^{})\psi _{}(y_{i+1}(\tau ),\tau )V^{II}(y_{i+1},t)𝑑\tau =0.`$
Next, summing all of the formulas (5.10) leads us to the general identity:
$`{\displaystyle _{y_i(s)}^{\chi _0(s)}}\psi (x,s)V^{II}(y_i,t)𝑑x+{\displaystyle \underset{I=0}{\overset{N}{}}}{\displaystyle \underset{j=0}{\overset{k}{}}}{\displaystyle _{x_I^j}^{x_I^{j+1}}}\psi (x,s)V^{II}(z_I^j+,t)𝑑x`$ $`5.11`$
$`+{\displaystyle _{\chi _{N+1}(s)}^{y_{i+1}(s)}}\psi (x,s)V^I(y_{i+1},t)𝑑x{\displaystyle \underset{I=0}{\overset{N}{}}}{\displaystyle \underset{j=0}{\overset{k}{}}}{\displaystyle _{z_I^j}^{z_I^{j+1}}}\psi (x,t)V^{II}(z_I^j+,t)𝑑x`$
$`={\displaystyle \underset{I=0}{\overset{N}{}}}{\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle _t^s}[V^{II}(z_I^j+,t)V^{II}(z_I^{j1}+,t)](\lambda _I^ja_{})\psi _{}(\xi _I^j(\tau ),\tau )𝑑\tau `$
$`{\displaystyle \underset{I=0}{\overset{N}{}}}{\displaystyle _t^s}[V^{II}(z_I+,t)V^{II}(z_{I1}^k+,t)](\lambda _I^ja_{})\psi _{}(\chi _I(\tau ),\tau )𝑑\tau `$
$`{\displaystyle _t^s}[V^{II}(y_i+,t)V^{II}(y_i,t)](\lambda _0a_{})\psi _{}(\chi _0(\tau ),\tau )𝑑\tau `$
$`{\displaystyle _t^s}[V^{II}(y_{i+1}+,t)V^{II}(z_N^k,t)](\lambda _{N+1}a_{})\psi _{}(\chi _{N+1}(\tau ),\tau )𝑑\tau `$
$`{\displaystyle _t^s}(y_i^{}a_+)\psi _+(y_i(\tau ),\tau )V^{II}(y_i,t)𝑑\tau {\displaystyle _t^s}(a_{}y_{i+1}^{})\psi _{}(y_{i+1}(\tau ),\tau )V^{II}(y_{i+1},t)𝑑\tau .`$
To estimate the right-hand side of (5.11), we recall that the solution $`u^I`$ of a scalar conservation laws satisfies
$$V^{II}(y_i,t)V^{II}(\chi _0(s),s),V^{II}(z_I^j+,t)V^{II}(x_I^j+,s),$$
for $`I=0,\mathrm{},N`$ and $`j=0,\mathrm{},k`$. Hence, choosing the difference $`x_I^{j+1}x_I^j`$ in (5.9) sufficiently small and since the function $`V^{II}(,t)`$ is nondecreasing, we conclude that the left-hand side of (5.11) can be bounded from below, as follows:
$$\text{ L.H.S. }_{y_i(s)}^{y_{i+1}(s)}\psi (x,s)V^{II}(x,s)𝑑x(st)\epsilon _{y_i(t)}^{y_{i+1}(t)}\psi (x,t)V^{II}(x,t)𝑑x.$$
$`5.12`$
Estimating the right-hand side of (5.11) is more involved. First note that each term arising in the left-hand side of (5.11) is non-positive. This follows from our condition (5.6). Indeed, consider a point $`(x,s)`$ of approximate jump or approximate continuity of $`u^I`$, $`u^{II}`$ and $`\psi `$. If all of these functions are continuous, the result is trivial. Call $`\lambda `$ the discontinuity speed. Based on the jump relation (5.5), we see that either $`\psi _{}(\lambda a_{})=\psi _+(\lambda a_+)=0`$, or else all of the terms $`\psi _{}`$, $`\lambda a_{}`$, $`\psi _+`$, and $`\lambda a_+`$ are distinct from zero.
Suppose first that $`(x,s)`$ is a point in the interior of the region limited by the two curves $`y_i(.)`$ and $`y_{i+1}(.)`$. In the latter case, since $`\psi 0`$ in the region under consideration, we deduce that $`\psi _{}>0`$ and $`\psi _+>0`$, while the terms $`\lambda a_{}`$ and $`\lambda a_+`$ are either both negative or both positive. Actually, in view of the sign condition (5.6), we have $`u_\pm ^{II}u_\pm ^I0`$ and, therefore, $`\lambda a_\pm `$ as follows from (3.10) (here we are dealing with a jump of $`u^{II}`$).
Consider next a point of the boundary $`y_i`$, for instance. So we now have $`\psi _{}<0`$ and $`\psi _+>0`$, while the terms $`\lambda a_{}`$ and $`\lambda a_+`$ opposite sign. Since no rarefaction-shock can arise, the discontinuity must be a Lax shock and so $`\lambda a_{}<0`$ and $`\lambda a_+>0`$. Again the corresponding term in (5.11) has a favorable sign. (Observe that the condition (5.6) was not used in this second case.)
Then, for all $`I=0,\mathrm{},N`$ and $`j=1,\mathrm{},k`$, let $`\theta _I^j()`$ be the (maximal, for definiteness) backward characteristic associated with $`u^I`$ and issuing from the point $`(\xi _I^j(\tau ),\tau )`$. Denote also by $`\theta (z_I^j;\tau )`$ its intercept with the horizontal line at time $`t`$. Setting
$$\stackrel{~}{a}(x,t;\tau ):=\frac{f(u^{II}(x,t))f(u^I(\theta (x,\tau ),t))}{u^{II}(x,t)u^I(\theta (x;\tau ),t)}$$
and using that the solution $`u^I`$ remains constant along the characteristic $`\theta _I^j()`$, we obtain
$$(\lambda _I^ja_{})(\xi _I^j(\tau ))=\lambda _I^j(z_I^j)\stackrel{~}{a}(z_I^j,t;\tau ).$$
$`5.13`$
Then consider the (maximum, for definiteness) backward characteristic $`y_I^j()`$ associated with $`a`$ and issuing from the point $`(\xi _I^j(\tau ),\tau )`$. By integrating $`\psi `$ along the characteristic $`y_I^j()`$ and using the inequality (4.4), we arrive at a lower bound for $`\psi `$
$$\psi (\xi _I^j(\tau ),\tau )\psi (y_I^j(t),t)\left(\frac{t}{\tau }\right)^E,t<\tau <s.$$
$`5.14`$
Upon choosing $`x_I^{j+1}x_I^j`$ in (5.9) so small that the oscillation of $`V_c^{II}()`$ over each interval $`(z_I^jz_I^{j+1})`$ does not exceed $`\epsilon `$ and recalling the standard estimates on Stieltjes integrals we deduce from (5.11)-(5.13) that
$`{\displaystyle \underset{I=0}{\overset{N}{}}}{\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle _t^s}[V^{II}(z_I^j+,t)V^{II}(z_I^{j1}+,t)]\left((\lambda _I^ja_{})\psi _{}(\xi _I^j(\tau ),\tau )\right)𝑑\tau `$ $`5.15`$
$`{\displaystyle \underset{I=0}{\overset{N}{}}}{\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle _t^s}[V_c^{II}(z_I^j,t)V_c^{II}(z_I^{j1},t)]\left(\lambda _I^j(z_I^j)\stackrel{~}{a}(z_I^j,t,\tau )\right)\psi (y_I^j(t),t)\left({\displaystyle \frac{t}{\tau }}\right)^E𝑑\tau `$
$`{\displaystyle _t^s}{\displaystyle \underset{I=0}{\overset{N}{}}}{\displaystyle \underset{j=1}{\overset{k}{}}}\left({\displaystyle _{z_I^{j1}}^{z_I^j}}\left(\lambda _I^j(x)\stackrel{~}{a}(x,t,\tau )\right)\psi (x,t)𝑑V_c^{II}(x,t)c\epsilon \right)\left({\displaystyle \frac{t}{\tau }}\right)^Ed\tau `$
$`={\displaystyle _t^s}\left({\displaystyle _{y_i}^{y_{i+1}}}\left(\lambda _I^j(x)\stackrel{~}{a}(x,t,\tau )\right)\psi (x,t)𝑑V_c^{II}(x,t)c(y_{i+1}y_i)\epsilon \right)\left({\displaystyle \frac{t}{\tau }}\right)^E𝑑\tau .`$
We now combine (5.10), (5.11) and (5.15), divide the resulting inequality by $`st`$, and let $`st`$, $`ϵ0`$, obtaining the following inequality:
$`{\displaystyle \frac{d^+}{dt}}{\displaystyle _{y_i(t)}^{y_{i+1}(t)}}\psi (x,t)V^{II}(x,t)𝑑x`$ $`{\displaystyle _{y_i}^{y_{i+1}}}\left(\lambda _I^j(x)a(x,t)\right)\psi (x,t)𝑑V_c^{II}(x,t)`$ $`5.16`$
$`{\displaystyle \underset{(x,t)𝒥(u^{II})}{}}\left(u_{}^{II}(x,t)u_+^{II}(x,t)\right)(\lambda ^Ia_{})(x,t)\psi _{}(x,t)`$
$`\left(u_{}^{II}(y_i,t)u_+^{II}(y_i,t)\right)(\lambda ^Ia_{})(y_i,t)\psi _{}(y_i,t)`$
$`\left(u_{}^{II}(y_{i+1},t)u_+^{II}(y_{i+1},t)\right)(\lambda ^Ia_{})(y_{i+1},t)\psi _{}(y_{i+1},t)`$
$`(y_i^{}a_+)\psi _+(y_i,t)V^{II}(y_i,t)`$
$`(a_{}y_{i+1}^{})\psi _{}(y_{i+1},t)V^{II}(y_{i+1},t)`$
The third and fourth terms in the right-hand side of (5.16) are due to the fact that $`\chi _0`$ and $`\chi _{N+1}`$ lie inside the region limited by $`y_i`$ and $`y_{i+1}`$.
We can next focus on the intervals $`(y_i,y_{i+1})`$ with $`i`$ even. Based on a completely symmetric argument and using now the weight $`m+V^{II}(\mathrm{},t)V^{II}(,t)`$ instead of $`V^{II}(,t)`$, we obtain
$`{\displaystyle \frac{d^+}{dt}}`$ $`{\displaystyle _{y_i(t)}^{y_{i+1}(t)}}\left(\psi (x,t)\right)\left(m+V^{II}(\mathrm{},t)V^{II}(x,t)\right)𝑑x`$ $`5.17`$
$``$ $`{\displaystyle _{y_i}^{y_{i+1}}}\left(\lambda _I^j(x)a(x,t)\right)\left(\psi (x,t)\right)𝑑V_c^{II}(x,t)`$
$`+{\displaystyle \underset{(x,t)𝒥(u^{II})}{}}\left(u_{}^{II}(x,t)u_+^{II}(x,t)\right)(\lambda ^Ia_{})(\psi _{})(x,t)`$
$`+\left(u_{}^{II}(y_i,t)u_+^{II}(y_i,t)\right)(\lambda ^Ia_{})(\psi _{})(y_i,t)`$
$`+\left(u_{}^{II}(y_{i+1},t)u_+^{II}(y_{i+1},t)\right)(\lambda ^Ia_{})(\psi _{})(y_{i+1},t)`$
$`(y_i^{}a_+)(\psi _+)(y_i,t)\left(m+V^{II}(\mathrm{},t)V^{II}(y_i,t)\right)`$
$`(a_{}y_{i+1}^{})(\psi _{})(y_{i+1},t)\left(m+V^{II}(\mathrm{},t)V^{II}(y_{i+1},t)\right).`$
By summation over $`i=1,\mathrm{},n`$ in (5.16) for $`i`$ odd and in (5.17) for $`i`$ even respectively, we obtain
$`{\displaystyle \frac{d^+}{dt}}`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}\left[\psi (x,t)\right]^+V^{II}(x,t)+\left[\psi (x,t)\right]^+\left(m+V^{II}(\mathrm{},t)V^{II}(x,t)\right)dx`$ $`5.18`$
$``$ $`{\displaystyle \underset{(x,t)(a)𝒥(u^{II})}{}}\left(m+V^{II}(\mathrm{},t)\right)\left|\lambda (x,t)a_{}(x,t)\right||\psi _{}(x,t)|`$
$`{\displaystyle \underset{(x,t)𝒥(u^{II})}{}}\left(u_{}^{II}(x,t)u_+^{II}(x,t)\right)\left(\lambda ^I(x,t)a_{}(x,t)\right)\psi _{}(x,t)`$
$`{\displaystyle _{RI}}\left(f^{}(u^{II}(y,t))a(y,t)\right)\psi (y,t)𝑑V_c^{II}(y,t),`$
where the superscript $`+`$ denotes the positive part of the functions $`\psi `$ and $`\psi `$ respectively.
Consider now the case where
$$\chi _0(\tau )y_i(\tau ),t\tau t+\delta ,$$
and
$$y_{i+1}(\tau )\chi _{N+1}(\tau ),t\tau t+\delta .$$
Assume that there exists a time $`\overline{\tau }>t`$ such that
$$\chi _0(\overline{\tau })<y_i(\overline{\tau }),y_{i+1}(\overline{\tau })<\chi _{N+1}(\overline{\tau })$$
(otherwise the curves of the two pairs will coincide, and we can reduce to the previous case). Let now $`\xi _0()`$ be the maximal backward characteristic emanating from $`(y_i(\overline{\tau }),\overline{\tau })`$, and $`\zeta _{N+1}()`$ be the minimal backward characteristic emanating from the point $`(y_{i+1}(\overline{\tau }),\overline{\tau })`$. Since characteristics cannot cross, we have that
$$y_i(t)<\xi _0(t),\zeta _{N+1}(t)<y_{i+1}(t),$$
Then, by finite propagation speed, there exists a time $`s>t`$ such that
$$y_i(\tau )<\xi _0(\tau ),\zeta _{N+1}(\tau )<y_{i+1}(\tau ),t\tau <s,$$
$$y_i(s)=\xi _0(s),\zeta _{N+1}(s)=y_{i+1}(s).$$
Instead of properties (c), (d), we will require that $`s`$ satisfies the following:
From then on we can proceed as before. Finally we write the inequality in (5.18) exchanging the roles of $`u^I`$ and $`u^{II}`$, and combining it with (5.17) we arrive exactly at the desired inequality (5.7) and the proof of Theorem 5.1 is completed. $`\mathit{}`$
## Acknowledgements
The authors are very grateful to C. Dafermos who communicated to them his lecture notes on the Liu-Yang functional in the context of general functions with bounded variation.
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# Freezing of dynamical exponents in low dimensional random media
## Abstract
A particle in a random potential with logarithmic correlations in dimensions $`d=1,2`$ is shown to undergo a dynamical transition at $`T_{dyn}>0`$. In $`d=1`$ exact results demonstrate that $`T_{dyn}=T_c`$, the static glass transition temperature, and that the dynamical exponent changes from $`z(T)=2+2(T_c/T)^2`$ at high temperature to $`z(T)=4T_c/T`$ in the glass phase. The same formulae are argued to hold in $`d=2`$. Dynamical freezing is also predicted in the 2D random gauge XY model and related systems. In $`d=1`$ a mapping between dynamics and statics is unveiled and freezing involves barriers as well as valleys. Anomalous scaling occurs in the creep dynamics.
The motion of topological defects by thermal activation over pinning barriers determines the slow glassy dynamics in numerous disordered systems, e.g. domain walls in dirty magnets, vortices in superconductors, dislocations in pinned lattices . While treating many interacting extended defects in the presence of disorder remains a major challenge, progress can be made on the simpler, already non trivial problem of a single point defect. A model of great interest is a particle diffusing in a random potential with log-correlations, i.e barriers growing logarithmically with scale. In 2D it precisely describes a single vortex in a XY spin model with gaussian random gauge disorder , and is relevant to a host of related systems e.g. vacancies in pancake lattices of layered 3D superconductors , dislocations in 2D lattices with smooth disorder , electrons on helium . As a prototype model of diffusion in a complex phase space or in random media, it is of broader interest to relaxation in glasses , transport in solids population biology , non hermitian quantum mechanics and vortex glass dynamical scaling . Similar models where used to study the dynamical generation of broad (e.g. power law) distributions of relaxation times and its relation to aging . Its static limit also appears in Integer Quantum Hall (QH) transition studies and its quantum extension in QH bilayer systems .
Early studies of this model, in the context of tracer diffusion in 2d potential flows, used RG methods perturbative in the disorder strength $`\sigma `$ . The position $`x(t)`$ of a particle (in any $`d`$) satisfies a Langevin equation:
$`\dot{x}(t)=V(x(t))+\eta (t)`$ (1)
$`\eta (t)\eta (t^{})=2T\delta (tt^{})`$ being a thermal noise ($`\overline{X}`$ and $`X`$ respectively denote disorder and thermal averages). It was found that if correlations grow logarithmically:
$`\mathrm{\Delta }(xx^{})=\overline{(V(x)V(x^{}))^2}4\sigma \mathrm{ln}|xx^{}|`$ (2)
the diffusion is anomalous with a temperature $`T`$ dependent “dynamical exponent” $`z`$ given by :
$`\overline{x(t)^2}t^{2/z},z=2+2{\displaystyle \frac{\sigma }{dT^2}}`$ (3)
Strikingly, this result for $`z`$ was conjectured to hold to all orders in $`g=\sigma /dT^2`$. This was confirmed in $`d=1`$ and $`d=2`$ by general arguments and a three loop calculation (which also found $`O(g^3)`$ corrections to $`z`$ in $`d3`$). Further support came from excellent agreement with simulations , performed in $`d=2`$ for $`0<g<0.8`$, and exact results in $`d=1`$ for the velocity $`vf^{z1}`$ at small applied force $`f`$, with $`z`$ again as in (3).
Recently, however, the statics of this model has been investigated in several works in $`d=2`$ and any $`d`$ . As a result we now know that there is in fact a transition at $`T=T_c=\sqrt{\sigma /d}`$ to a strong disorder, low $`T`$ glass phase. This glass phase appears non trivial, as dominated by a few states, reminiscent of replica symmetry breaking (RSB). It is related, approximately, to the Random Energy Model (REM) and, more closely, to the directed polymer on a Cayley tree (DPCT) . It is thus an outstanding problem to investigate whether this equilibrium transition has a dynamical counterpart, and how it can be compatible with (3).
In this Letter we solve the apparent paradox. In $`d=1`$ we demonstrate that there is indeed a dynamical transition at the static $`T_c`$ and obtain the dynamical exponent $`z(T)`$ at all $`T`$. This is achieved through exact results and a real space RG (RSRG) method. Simple arguments and bounds also show that a dynamical transition occurs in $`d=2`$, with analogous behaviour of $`z(T)`$. A similar dynamical freezing is also predicted in the 2d random gauge XY model, since the XY phase and its boundary are dominated by the single (i.e dilute) vortex limit .
We first give a simple argument indicating that the $`1/T^2`$ divergence in (3) cannot hold at low temperature. First, in any given sample (in finite $`d`$), the characteristic time $`t`$ for a particle to escape a region of size $`L`$ satisfies:
$`\mathrm{ln}tB/T[V_{max}(L)V_{min}(L)]/TB_{max}/T`$ (4)
i.e is given as $`T0`$ by the Arrhenius estimate where $`B`$ is the energy barrier encountered by the particle, obviously bounded by the difference $`B_{max}`$ between the absolute maximum $`V_{max}(L)`$ and minimum $`V_{min}(L)`$ of the potential in the region. Second, we now know that these extremal values satisfy $`V_{min}(L)=2\sqrt{\sigma /d}(d\mathrm{ln}L\frac{1}{2}\gamma \mathrm{ln}\mathrm{ln}L)+\delta V`$ where $`\gamma `$ is a universal number and $`\delta V`$ has $`O(1)`$ sample to sample fluctuations. Thus, defining $`z\mathrm{ln}t/\mathrm{ln}L`$ at large $`L`$ this yields, in any dimension $`d`$, the bound:
$`z4\sqrt{\sigma d}/TT0.`$ (5)
Since (3) is exact perturbatively to all orders in $`g=\sigma /T^2`$ in $`d=1,2`$, (5) implies that a dynamical transition must occur in $`d=1,2`$ at a finite temperature $`T_{dyn}(d)>0`$.
In $`d=1`$, which we now study, it is natural to guess that the upper bound in (5) gives the exact $`z(T)`$ for all $`T<T_c`$. Indeed one expects that the Arrhenius law (4) should hold in the energy dominated glass phase and furthermore $`B`$ should equal its upper bound in (4) since in $`d=1`$ there is only a single path. We now confirm this with analytical results at all $`T`$ using the first passage time approach . We also show that (i) in $`d=1`$ there is a direct correspondence between dynamical (e.g. exponents) and static quantities (ii) the dynamical exponent $`z`$ can be unambiguously defined. Our conclusions being independent of boundary conditions we choose them reflecting at $`x=0`$ and define $`t`$ as the first passage time at site $`x=L`$ of a particle which starts at $`x=0`$. Some of our results are valid for any potential landscape but are mostly applied to gaussian log-correlated potentials (2).
As shown below it is sufficient to focus, in any given disorder realization $`V(x)`$, on the (thermal) mean first passage time $`t_1(\{V\})t\tau `$. We first obtain its typical behaviour (i.e in any sample), its large deviations being studied later. The exact formula for $`t_1(\{V\})`$ is :
$$t_1(\{V\})=\frac{1}{T}_0^L𝑑y_0^L𝑑x\theta (yx)\mathrm{e}^{[V(y)V(x)]/T},$$
(6)
where $`\theta (x)`$ is the step function. Since the canonical partition function for the same disorder realization is $`Z_L(\{V\})_0^L𝑑x\mathrm{exp}(V(x)/T)`$, (6) strongly resembles two copies of the statics, one with disorder $`+V(x)`$ (dominated by minima at low $`T`$), the other with $`V(x)`$ (dominated by maxima of $`V`$ at low $`T`$). Dynamics and statics are thus directly connected by: $`\mathrm{ln}\{T[t_1(\{V\})+t_1(\{V\})]\}=\mathrm{ln}Z_L(\{V\})+\mathrm{ln}Z_L(\{V\})`$ and also
$`\mathrm{ln}Z_{L/2}(\{V\})`$ $`+`$ $`\mathrm{ln}Z_{L/2}(\{V\})\mathrm{ln}\{Tt_1(\{V\})\}`$ (7)
$``$ $`\mathrm{ln}Z_L(\{V\})+\mathrm{ln}Z_L(\{V\}).`$ (8)
If the ratio of upper to lower bound converges to unity in the thermodynamic limit, they uniquely determine the long time behavior in terms of the statics. This is the case for gaussian potentials with correlations growing at the most as a power of the logarithm of the scale. If the correlations do not grow faster than logarithmically with scale, then $`\mathrm{ln}Z_L(\{V\})\mathrm{ln}L`$. If, additionally, the intensive free energy $`f=F_L(\{V\})/\mathrm{ln}L=T\mathrm{ln}Z_L(\{V\})/\mathrm{ln}L`$ is self-averaging, the dynamical exponent $`z`$ is obtained as:
$`z(\{V\})\underset{L\mathrm{}}{lim}\mathrm{ln}t_1(\{V\})/\mathrm{ln}L=2f/T,`$ (9)
and is also self averaging. This is the case both for uncorrelated and log-correlated gaussian potentials. For the former (9) gives the normal diffusion value $`z=2`$. For the latter, using the known results for $`f`$ we obtain our main result for the dynamical exponent:
$`z(T)=z_A(T)2\left(1+\sigma /T^2\right)(\text{for}TT_\mathrm{c})`$ (10)
$`z(T)=4\sqrt{\sigma }/T(\text{for}TT_\mathrm{c}).`$ (11)
Thus a dynamical transition, away from the “annealed” value $`z_A(T)`$ given by (3), occurs at the same temperature $`T_{dyn}=T_\mathrm{c}=\sqrt{\sigma }`$ as the equilibrium transition. At $`T_c`$ freezing occurs in the thermal configurations which dominate $`t`$, simultaneously around minima and maxima of the potential (i.e in the two copies), thus in the effective barrier. Interestingly, this transition coincides with the onset of logarithmic corrections. Indeed one can also characterize typical finite size fluctuations. Using that $`F_L=f(\mathrm{ln}L\frac{1}{2}\gamma (T)\mathrm{ln}(\mathrm{ln}L))+\delta F`$ where $`\delta F`$ has $`O(1)`$ sample to sample fluctuations we find that $`t_1(\{V\})/\tau _{typ}`$ has a well defined $`O(1)`$ limit distribution at large $`L`$, the typical mean first passage time being $`\tau _{typ}=L^{z(T)}(\mathrm{ln}L)^{\alpha (T)}`$ with $`\alpha (T)=2\gamma (T)\sqrt{\sigma }/T`$, $`\gamma (T)=0`$ for $`T>T_c`$ but $`\gamma (T_c)=\frac{1}{2}`$ and $`\gamma (T)=\frac{3}{2}`$ for $`T<T_c`$. For faster growing (e.g. power law) correlations, the ratio between the bounds in Eq. (8) does not converge to one, and even the leading order of $`\mathrm{ln}t_1(\{V\})`$ still fluctuates at large $`L`$ as in the unbiased Sinai model .
We now show as promised that $`z`$ can be defined unambiguously from the mean first passage time alone. In principle one could define a full set of dynamical exponents $`z_p(\{V\})=lim_L\mathrm{}\mathrm{ln}t_p(\{V\})/p\mathrm{ln}L`$ from higher thermal moments $`t_p(\{V\})t^p`$. Using the expression :
$`t_p(\{V\})={\displaystyle \frac{p!}{T^p}}{\displaystyle _{0<x_i,y_i<L}}{\displaystyle \underset{i=1}{\overset{p}{}}}\theta _{y_i,x_i}\theta _{y_i,x_{i1}}\mathrm{e}^{\frac{V(y_i)V(x_{i1})}{T}}`$ (12)
and $`\theta _{x,y}\theta (xy)1`$ we obtain, comparing with (6), $`\mathrm{ln}t_p(\{V\})/p\mathrm{ln}L\mathrm{ln}t_1(\{V\})/\mathrm{ln}L+\mathrm{ln}p!/p\mathrm{ln}L`$, which together with the general inequality $`t^pt^p`$ leads to:
$$z_p(\{V\})=z_1(\{V\})\text{and}\overline{z_p}=\overline{z_1}.$$
(13)
for any integer $`p1`$. This is because the thermal distribution of $`t`$ has exponential decay and all moments $`p1`$ are controlled at large $`L`$ by the largest relaxation time in a given sample. Thus $`z=z_1`$ characterizes the dynamics.
By contrast, the distribution of escape times $`P_L(\tau )`$ with disorder realization has broad tails extending in the region $`\tau \tau _{typ}L^{z(T)}`$, where we obtained the estimate:
$`P_L(\tau )d\tau {\displaystyle \frac{d\tau }{\tau }}({\displaystyle \frac{\tau _{typ}}{\tau }})^\mu \mathrm{exp}(\mu ^2{\displaystyle \frac{\mathrm{ln}^2(\tau /\tau _{typ})}{8\mathrm{ln}L}}).`$ (14)
Here $`\mu =T/T_c`$ and (14) is valid for (i) $`T<T_c`$ and (ii) for $`T>T_c`$ and $`\stackrel{~}{z}\mathrm{ln}\tau /\mathrm{ln}L4`$. For $`T>T_c`$ and $`z(T)<\stackrel{~}{z}<4`$ one has simply $`P_L(\tau )\tau _{typ}^1(\tau _{typ}/\tau )^{1+\mu ^2}`$. This corresponds to a quadratic (plus linear) multifractal spectrum for rare occurrences of $`\stackrel{~}{z}`$. (14) can be obtained by a Kosterlitz RG analysis of (6), as in . The result is a non linear Kolmogorov equation for $`P_L(\tau )`$ as a function of $`\mathrm{ln}L`$, identical to the one describing the partition sum of two directed polymers on a Cayley tree, seeing opposite disorder, with constrained endpoints $`x<y`$. It yields (14) up to log corrections, neglected here. More empirically (14) can be obtained from the moments:
$`\overline{t^p}`$ $`=`$ $`{\displaystyle \frac{1}{T^p}}{\displaystyle _0^L}{\displaystyle \underset{i=1}{\overset{p}{}}}\left[dy_idx_i\theta (y_ix_i)\right]\mathrm{exp}\{{\displaystyle \underset{i,j=1}{\overset{p}{}}}[\mathrm{\Delta }(y_ix_j)..`$ (16)
$`..{\displaystyle \frac{1}{2}}\mathrm{\Delta }(y_iy_j){\displaystyle \frac{1}{2}}\mathrm{\Delta }(x_ix_j)]/2T^2\}.`$
which reads as a partition sum for $`2p`$ particles, $`p`$ of type $`y`$ (representing hills in the potential landscape) and $`p`$ of type $`x`$ (representing valleys). Same type particles attract via a potential $`\mathrm{\Delta }(r)/2T`$ while those of opposite type repel via $`\mathrm{\Delta }(r)/T`$. This is similar to estimating $`\overline{Z^p}`$ in the statics, except that there only one kind of particles (representing valleys) appears. Here hills and valleys play symmetric roles. (16) can be estimated for log-correlated potentials and for integer $`p1`$ as follows. At high $`T`$ an ”entropic” variational saddle point dominates (with all $`y`$’s and $`x`$’s far away $`O(L)`$ from each other). At low $`T`$ an ”energetic” saddle point dominates (all $`y`$’s close together within $`O(1)`$, all $`x`$’s close together, $`y`$’s and $`x`$’s far away). This yields the large $`L`$ behaviour (universal since unaffected by changes in $`\mathrm{\Delta }(r)`$ for small $`r`$):
$`\overline{t^p}`$ $`L^{2p(1+\sigma /T^2)}`$ $`\text{for}TT_{\mathrm{c},\mathrm{p}},`$ (17)
$``$ $`L^{2p(1/p+p\sigma /T^2)}`$ $`\text{for}T<T_{\mathrm{c},\mathrm{p}}.`$ (18)
i.e, as for $`\overline{Z^p}`$ in the statics (of this model , the REM and the DPCT) there is a sequence of transition temperatures $`T_{\mathrm{c},\mathrm{p}}=\sqrt{p}T_\mathrm{c}`$ for the moments. One can check via a saddle point calculation of $`\overline{t^p}=𝑑\tau \tau ^pP_L(\tau )`$ that (18) is consistent with (14) and that the $`T_{\mathrm{c},\mathrm{p}}`$ correspond to a change of behaviour from rare events to typical events dominance as the saddle point crosses $`\stackrel{~}{z}_{sp}=4`$. By analytically continuing $`\overline{t^p}`$ as $`p0`$ one can recover (11) and the transition at $`T_c`$ in typical behaviour. The ”entropic” saddle point still dominates for $`T>T_c`$, but the ”energetic” saddle point is replaced at $`T<T_c`$ by a one-step RSB ansatz. Each kind of ”particle” is arranged in $`p/m`$ groups of $`m`$ particles close by in space, while different groups are far away, with $`m=T/T_c`$ at the optimum, extending the static and the REM and DPCT replica calculations . Compared to conventional static RSB which only involves “valleys” the interesting feature here is a nonequilibrium RSB which also involves “hills”: indeed, near degeneracies of distant barriers result, in the glass phase, in a nonequilibrium splitting of the thermal distribution of the diffusing particle into a few packets in a single environment . These features are absent for weaker correlations (high $`T`$ ”entropic” saddle) or stronger ones , where an ”energetic” saddle without RSB dominates (degeneracies are subdominant in the Sinai landscape except in the presence of a bias ).
If an external force is applied, the creep velocity $`v`$ relates to the disorder averaged mean escape time of regions of size $`L_01/f`$ and is thus controlled by the annealed exponent $`z_A(T)`$, distinct from $`z(T)`$ at low $`T`$, a striking breakdown of naive dynamical scaling. Freezing manifests itself in large finite size corrections to $`v`$ due to undersampling of the disorder average. From , $`v^1=\frac{1}{L}t`$ holds for $`fL1`$, where $`t`$ is given by (6) in the tilted landscape $`V(x)fx`$. It can be approximated as the average over $`N`$ samples of sizes $`L_0`$ of the escape time in each sample, each distributed with $`P_{L_0}(\tau )`$. A saddle point estimate shows that below $`T_c`$, $`vf^{z_A1}`$ for $`y=\mathrm{ln}L/\mathrm{ln}f>y^{}=2/\mu ^21`$, but that for smaller sizes $`1<y<y^{}`$, the typical $`vf^{z_m1}`$ where the exponent $`z_m=\frac{4T_c}{T}\sqrt{(1+y)/2}y+1`$ smoothly interpolates between the annealed and quenched one.
The RSRG method, previously devised to describe diffusion in the Sinai landscape (for details see ) and in a broader class , allows to obtain complementary information, e.g. about distribution of positions. Here, in the log-correlated landscape, it is implemented numerically. From the original set of (i) alternating local extrema $`V(x_i)`$, (ii) their energy differences (“barriers” of heights $`F_i=|V(x_i)V(x_{i+1})|`$) and (iii) the segments between them (“bonds” of lengths $`\mathrm{}_i=x_{i+1}x_i`$), one constructs iteratively the “renormalized landscape at scale $`\mathrm{\Gamma }`$” by removing as $`\mathrm{\Gamma }`$ increases all barriers between $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }+d\mathrm{\Gamma }`$ and merging the corresponding bonds. This decimation retains only the large barriers and deep valleys. The Arrhenius dynamics of a particle starting at $`x_0`$ is
approximated by putting it at time $`t`$ at the bottom of the bond which contains $`x_0`$ in the renormalized landscape at $`\mathrm{\Gamma }=T\mathrm{ln}t`$. The errors are small if the distribution of renormalized barriers is broad compared to $`T`$ (infinitely broad in Sinai). Since here barriers remain finite, the method should be exact only as $`T/T_c0`$, i.e when the thermal packet is concentrated in a single well. Perturbative corrections in $`T/T_c`$ can be computed by considering the small occupation probability of neighboring secondary wells. In practice, the method works surprisingly better and gives the exact $`z(T)`$ up to $`T_c`$.
We have found that under renormalization the probability distributions for rescaled variables reach fixed point forms, namely $`\text{Prob}_\mathrm{\Gamma }((F\mathrm{\Gamma })/F_\mathrm{\Gamma }^{\mathrm{typ}})P^{}((F\mathrm{\Gamma })/F^{\mathrm{typ}})`$ for rescaled barriers (Fig. 1) and $`\text{Prob}_\mathrm{\Gamma }(\mathrm{}/\overline{\mathrm{}}_\mathrm{\Gamma })Q^{}(\mathrm{}/\overline{\mathrm{}}_\mathrm{\Gamma })`$. Here $`F_\mathrm{\Gamma }^{\mathrm{typ}}`$ flows to a constant $`F^{\mathrm{typ}}4.4\sqrt{\sigma }`$ and $`\overline{\mathrm{}}_\mathrm{\Gamma }=L/N_\mathrm{\Gamma }`$ is the average bond length. Since two barriers are decimated at each step the number of remaining barriers $`N_\mathrm{\Gamma }`$ satisfies $`_\mathrm{\Gamma }N_\mathrm{\Gamma }=2\alpha _\mathrm{\Gamma }N_\mathrm{\Gamma }`$ where $`\alpha _\mathrm{\Gamma }Prob_\mathrm{\Gamma }(F=\mathrm{\Gamma })\alpha ^{}=P^{}(0)/F^{\mathrm{typ}}`$, a constant. Thus the bond length grows as $`\overline{\mathrm{}}_\mathrm{\Gamma }\mathrm{exp}(2\alpha ^{}\mathrm{\Gamma })t^{1/z}`$ and using $`\mathrm{\Gamma }=T\mathrm{ln}t`$ one recovers the dynamical exponent $`z=1/(2\alpha ^{}T)`$. Numerically we find $`1/(2\alpha _\mathrm{\Gamma })4\sqrt{\sigma }(10.25\mathrm{ln}(4\pi \mathrm{ln}\overline{\mathrm{}}_\mathrm{\Gamma })/\mathrm{ln}\overline{\mathrm{}}_\mathrm{\Gamma })`$ and thus a value of $`z`$ consistent with (11). The diffusion front, computed as in , converges to a scaling form as $`\overline{Prob_\mathrm{\Gamma }(xt|00)}\overline{\mathrm{}}_\mathrm{\Gamma }^1q(X=x/\overline{\mathrm{}}_\mathrm{\Gamma })`$, represented in Fig. 1. $`\overline{\mathrm{}}_\mathrm{\Gamma }`$ is thus the only relevant lengthscale, all moments of the displacement scaling as $`\overline{x(t)^k}\overline{\mathrm{}}_\mathrm{\Gamma }^kt^{k/z}`$. Finally we obtained good numerical evidence (see Fig. 1) that all above asymptotic scaling functions, as well as $`\alpha _0\sqrt{\sigma }\alpha ^{}`$ and $`F_0F^{\mathrm{typ}}/\sqrt{\sigma }`$, do not change upon adding short-range disorder and are thus universal ($`\overline{\mathrm{}}_\mathrm{\Gamma }`$ does change by a constant factor.)
We now address $`d2`$. First we note that the bound (5) can be improved to $`z(T)2\sqrt{d\sigma }/T`$ as $`T0`$ for any $`d2`$. Indeed, in order to escape the particle now only needs to find a set of saddles which connects to the boundary. A percolation and counting argument shows that it can do so by remaining within sites such that $`V(x)/\mathrm{ln}L0`$. Thus the relevant barrier is bounded as $`BV_{min}`$. In $`d=2`$ this bound is likely to be saturated since the particle still finds deepest minima, yielding $`z=2\sqrt{2\sigma }/T=4T_c/T`$. Since this expression matches (3) at the static $`T_c`$ in $`d=2`$, a likely scenario is that $`T_{dyn}=T_c`$ and that the expression holds for all $`T<T_c`$ . A single vortex in a random gauge XY model will experience a similar dynamical freezing.
To conclude we demonstrated dynamical transitions in $`d=1,2`$. In $`d=1`$ we found anomalous scaling of the creep velocity, novel freezing phenomena involving barriers and a finite $`T`$ generalization of Arrhenius law $`\tau e^{2F_L/T}`$. Extensions will appear elsewhere.
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# 1 Introduction
## 1 Introduction
Recently, much attention has been given to noncommutative Yang-Mills theory (NCYM), after it is realized that NCYM arises naturally in a specific compactification of the matrix theory and in string theories as decoupling limits of the worldvolume theories on D-branes in nonzero $`B`$-field background . It is argued that NCYM and the ordinary Yang-Mills theory (CYM) arise from the same field theory regularized in different ways. The gauge fields of NCYM and CYM are related by requiring the equivalence of the gauge transformations of the two theories. The evidence for the equivalence between the large $`N`$ limit NCYM and CYM was also given in Refs. .
In the spirit of the holographic principle , the bulk supergravity dual description of NCYM was studied in Refs. . In accordance with the conjectured equivalence between NCYM and CYM, it is found out that thermodynamics of near-extremal D$`p`$-branes with nonzero $`B`$-field coincides exactly with that of the corresponding D$`p`$-brane without $`B`$-field . By observing that the near horizon geometry and thermodynamics of the supergravity solution for the D$`p`$-brane with nonzero $`B`$-field are identical to those of the supergravity solution for the D$`(p2)`$-brane with vanishing $`B`$-field and two of its transverse coordinates delocalized, it is further argued that NCYM in $`p+1`$ dimensions and CYM in $`p1`$ dimensions are equivalent.
It is the purpose of this paper to give a microscopic explanation for the equivalence of thermodynamics of various supergravity solutions for the D-brane systems mentioned in the above. We relate such supergravity solutions to microscopic D-brane systems through the boundary state formalism and show that such D-brane systems are related to one another through $`T`$-duality transformations, implying that the number of microscopic degrees of freedoms of the above mentioned different D-brane configurations are mapped one-to-one to one another under the $`T`$-duality transformations. We begin by reviewing the relevant ideas and elaborating on their connection in section 2, for the purpose of preparing for the discussion of the main result of this paper. In section 3, we show that various D-brane systems whose thermodynamics are shown to coincide with one another are related under the $`T`$-duality transformations.
## 2 General Properties
The bosonic action for an open string ending on a D$`p`$-brane in the constant NS $`B`$-field background can be written in the following form:
$$S=\frac{1}{4\pi \alpha ^{}}d^2\xi \left[\sqrt{h}h^{\alpha \beta }g_{\mu \nu }_\alpha X^\mu _\beta X^\nu ϵ^{\alpha \beta }_{ij}_\alpha X^i_\beta X^j\right],$$
(1)
by going to the gauge in which the NS $`B`$-field takes the form $`B=B_{ij}dX^idX^j`$ ($`i,j=0,\mathrm{},p`$) . Here, $`_{ij}B_{ij}+_{[i}A_{j]}`$, $`A_i`$ is the $`U(1)`$ gauge field living on the D$`p`$-brane worldvolume and $`\xi ^\alpha =(\tau ,\sigma )`$ is the string worldsheet coordinates. By applying the gauge transformation $`B_{ij}B_{ij}+_i\mathrm{\Lambda }_j_j\mathrm{\Lambda }_i`$ and $`A_iA_i\mathrm{\Lambda }_i`$, under which the action (1) is invariant, one can set $`A_i`$ equal to zero, which we assume from now on. In the $`h_{\alpha \beta }=\eta _{\alpha \beta }`$ gauge with constant $`g_{\mu \nu }`$ and $`B_{ij}`$, the variation of the bosonic open string action (1) with respect to the string coordinates $`X^\mu `$ yields the the following boundary conditions at the ends $`\sigma =0,\pi `$ of the open string:
$`g_{ij}_\sigma X^j+B_{ij}_\tau X^j=0,i,j=0,1,\mathrm{},p,`$ (2)
$`\delta X^a=0,a=p+1,\mathrm{},9.`$ (3)
The string propagator with this boundary conditions takes the following form :
$$X^i(z)X^j(z^{})=\alpha ^{}\left[g^{ij}\mathrm{ln}\frac{|zz^{}|}{|z\overline{z}^{}|}+G^{ij}\mathrm{ln}|z\overline{z}^{}|^2+\frac{1}{2\pi \alpha ^{}}\theta ^{ij}\mathrm{ln}\frac{z\overline{z}^{}}{\overline{z}z^{}}\right],$$
(4)
where $`z=\tau +i\sigma `$ and
$`G^{ij}`$ $`=`$ $`\left({\displaystyle \frac{1}{g+B}}\right)_S^{ij}=\left({\displaystyle \frac{1}{g+B}}g{\displaystyle \frac{1}{gB}}\right)^{ij},`$ (5)
$`\theta ^{ij}`$ $`=`$ $`2\pi \alpha ^{}\left({\displaystyle \frac{1}{g+B}}\right)_A^{ij}=2\pi \alpha ^{}\left({\displaystyle \frac{1}{g+B}}B{\displaystyle \frac{1}{gB}}\right)^{ij},`$ (6)
where the subscripts $`S`$ and $`A`$ respectively denote the symmetric and the antisymmetric parts of the matrix. $`G^{ij}`$ is interpreted as the effective metric seen by the open strings and $`\theta ^{ij}`$ has the interpretation as the noncommutativity parameter as can be seen from the following time ordered commutation relation that follows from Eq. (4):
$$[X^i(\tau ),X^j(\tau )]=T\left(X^i(\tau )X^j(\tau ^{})X^i(\tau )X^j(\tau ^+)\right)=i\theta ^{ij}.$$
(7)
So, the end-points of the open string target space coordinates associated with non-zero components of the constant $`B`$ field live in noncommutative space. Namely, in the presence of nonzero $`B`$-field, the D-brane worldvolume becomes noncommutative. This result was also obtained by the quantization of $`X^i`$ through the analysis of the time-averaged symplectic form on the phase space or through the Dirac bracket quantization procedure . (See also Refs. .)
D-branes can be alternatively described by the “boundary states” of the “close-string channel” description. The “open-string channel” and the “closed-string channel” descriptions of D-branes are mapped to one another through $`\tau \sigma `$, under which an open string one-loop diagram and a closed string tree diagram are interchanged. The boundary state is given by the product of a matter and a ghost parts, each of which is expressed as the product of the bosonic and the fermionic parts. The GSO projection selects specific linear combinations of such boundary states separately for the NS-NS and the R-R sectors. In this paper, we will be mainly concerned with the bosonic matter part $`|B_X`$ of the boundary state. By applying the transformation $`\tau \sigma `$ (followed by an appropriate rescaling of $`\tau `$ and $`\sigma `$) to the boundary conditions (3) on an open string, one obtains the following conditions on $`|B_X`$ at $`\tau =0`$:
$`(_\tau X^i+B_j^i_\sigma X^j)|_{\tau =0}|B_X=0,i,j=0,1,\mathrm{},p,`$ (8)
$`(X^a|_{\tau =0}x^a)|B_X=0,a=p+1,\mathrm{},9.`$ (9)
In the nonzero $`B`$-field background with some of coordinates compactified on a torus, the oscillator expansion for the closed string target space coordinates is
$$X^\mu =x^\mu +w^\mu \sigma +\tau g^{\mu \nu }(p_\nu B_{\nu \rho }w^\rho )+\frac{i}{\sqrt{2}}\underset{n0}{}\left[\frac{\alpha _n^\mu }{n}e^{in(\tau \sigma )}+\frac{\stackrel{~}{\alpha }_n^\mu }{n}e^{in(\tau +\sigma )}\right],$$
(10)
where $`w^\mu `$ is zero for noncompact directions and we assume that only the longitudinal components $`B_{ij}`$ of the two-form potential are nonzero. So, the conditions (9) on $`|B_X`$ at $`\tau =0`$ in terms of the oscillator modes take the following forms:
$`\widehat{p}^i|B_X=0,(\widehat{x}^ax^a)|B_X=0,`$ (11)
$`\left[(\mathrm{𝟏}+B)_j^i\alpha _n^j+(\mathrm{𝟏}B)_j^i\stackrel{~}{\alpha }_n^j\right]|B_X=0,`$ (12)
$`(\alpha _n^a\stackrel{~}{\alpha }_n^a)|B_X=0,\widehat{w}^a|B_X=0,`$ (13)
where $`\mathrm{𝟏}`$ is the $`(p+1)\times (p+1)`$ identity matrix.
We now discuss the $`T`$-duality transformation of closed string theory on $`T^d`$ and its effect on D-brane configurations. The $`T^d`$ part of the canonical Hamiltonian of the closed string is
$$H=\frac{1}{4\pi \alpha ^{}}_0^{2\pi }𝑑\sigma \left(\begin{array}{cc}X^{}& 2\pi \alpha ^{}P\end{array}\right)M(E)\left(\begin{array}{c}X^{}\\ 2\pi \alpha ^{}P\end{array}\right),$$
(14)
where the matrix $`M(E)`$ determined by $`E=g+B`$ is given by
$$M(E)=\left(\begin{array}{cc}gBg^1B& Bg^1\\ g^1B& g^1\end{array}\right),$$
(15)
$`X^{}=_\sigma X`$ and $`P=(g_\tau X+B_\sigma X)/(2\pi \alpha ^{})`$ is the conjugate momentum. It appears from Eq. (14) that the Hamiltonian has the $`O(d,d,𝐑)`$ symmetry, but since the eigenvalues of the operators $`\widehat{w}^i`$ and $`\widehat{p}^i`$ in the mode expansion (10) take integer values due to the periodicity condition $`X^iX^i+2\pi `$ of the compactified coordinates actually the Hamiltonian is invariant only under the $`O(d,d,𝐙)`$ subset <sup>2</sup><sup>2</sup>2The Hamiltonian for the open string has the same form (14) in terms of the string coordinates and the conjugate momentum and therefore appears to have the same symmetry as the closed string case. However, the $`O(d,d,𝐙)`$ target space duality symmetry is not a symmetry of the open string due to the absence of the winding modes..
Under the $`O(d,d,𝐙)`$ transformation with the transformation matrix $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)O(d,d,𝐙)`$, the background fields and the oscillator modes transform as
$`E`$ $``$ $`E^{}=(aE+b)(cE+d)^1,`$ (16)
$`\alpha _n(E)`$ $``$ $`(dcE^T)^1\alpha _n(E^{}),\stackrel{~}{\alpha }_n(E)(d+cE)^1\stackrel{~}{\alpha }_n(E^{}).`$ (17)
Note, this transformation is valid also for the $`n=0`$ case, where the oscillator modes are defined as
$$\alpha _0(E)\frac{1}{\sqrt{2}}G^1(pEw),\stackrel{~}{\alpha }_0(E)\frac{1}{\sqrt{2}}G^1(p+E^Tw).$$
(18)
By making use of the equivalence between the closed-string and open-string channel descriptions of D-branes through boundary states, we study the effect of $`T`$-duality symmetry of closed string theory on D-branes. For this purpose, it is convenient to divide the closed string coordinate mode expansion of the form (10) into the left-moving and the right-moving parts as $`X=\frac{1}{2}(X_{}+X_+)`$ with
$`X_{}`$ $`=`$ $`x+\sqrt{2}(\tau \sigma )\alpha _0+i\sqrt{2}{\displaystyle \underset{n0}{}}{\displaystyle \frac{\alpha _n}{n}}e^{in(\tau \sigma )},`$ (19)
$`X_+`$ $`=`$ $`x+\sqrt{2}(\tau +\sigma )\stackrel{~}{\alpha }_0+i\sqrt{2}{\displaystyle \underset{n0}{}}{\displaystyle \frac{\stackrel{~}{\alpha }_n}{n}}e^{in(\tau +\sigma )},`$ (20)
where the zero modes $`\alpha _0`$ and $`\stackrel{~}{\alpha }_0`$ are defined in Eq. (18). The effect of the $`O(d,d,𝐙)`$ symmetry transformation of the closed string theory on the open-string channel can be inferred by noting the map $`\tau \sigma `$ (and therefore $`_\tau _\sigma `$) that connects closed and open string channel descriptions of D-branes. The $`O(d,d,𝐙)`$ $`T`$-duality transformation (17) is generated by the following transformations:
* Factorized dualities $`D_i`$:
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)=\left(\begin{array}{cc}I_de_i& e_i\\ e_i& I_de_i\end{array}\right),$$
(21)
where $`I_d`$ is the $`d\times d`$ identity matrix and the $`d\times d`$ matrix $`e_i`$ has zero entries except for the $`(i,i)`$-component which is 1. When the $`B`$-field is zero and $`g`$ is diagonal, i.e. $`g=\mathrm{diag}(R_1^2,\mathrm{},R_d^2)`$, the factorized duality $`D_i`$ acts as the well-known large and small radii duality on the $`i`$-th coordinate, i.e. $`R_i1/R_i`$ while the remaining radii unchanged. The oscillator modes transform <sup>3</sup><sup>3</sup>3The extra factor of $`1/R_i^2`$ in the oscillator mode transformation is due to our choice of gauge in which the periodicity of the coordinates $`X`$ is fixed to be $`2\pi `$ and all the information on the size and the shape of the torus is encoded in the background field $`E=g+B`$. We have chosen such gauge because we wish to consider the general constant background fields $`g`$ and $`B`$. as $`\alpha _n^i\frac{1}{R_i^2}\alpha _n^i`$ and $`\stackrel{~}{\alpha }_n^i\frac{1}{R_i^2}\stackrel{~}{\alpha }_n^i`$, meanwhile the remaining modes $`\alpha _n^j`$ and $`\stackrel{~}{\alpha }_n^j`$ ($`ji`$) remain unchanged. This implies the transformation $`_\tau X^i_\sigma X^i`$ on the $`i`$-th closed string coordinate $`X^i`$. From the correspondence between the open string boundary conditions (3) and the conditions (9) on the boundary states, one can see therefore that the Dirichlet and the Neumann boundary conditions of the $`i`$-th open string coordinate are interchanged. For a general background $`E=g+B`$, although $`E`$ transforms in more complicated way, the oscillator modes transform similarly as the diagonal $`E`$ case, namely $`\alpha _n^i\frac{1}{g_{ii}}\alpha _n^i`$ and $`\stackrel{~}{\alpha }_n^i\frac{1}{g_{ii}}\stackrel{~}{\alpha }_n^i`$ with the remaining modes unchanged, and therefore $`_\tau X^i_\sigma X^i`$. So, the factorized dualities generally correspond to the usual $`T`$-duality transformations that transform D$`p`$-brane into D$`(p\pm 1)`$-brane.
* Basis change of the compactification lattice $`\mathrm{\Lambda }`$, i.e. $`EAEA^T`$ with $`AGL(d,𝐙)`$:
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)=\left(\begin{array}{cc}A& 0\\ 0& (A^T)^1\end{array}\right)\mathrm{s}.\mathrm{t}.AGL(d,𝐙).$$
(22)
Under this transformation, the background fields and the oscillator modes transform as
$$gAgA^T,BABA^T,\alpha _n(E)A^T\alpha _n(E^{}),\stackrel{~}{\alpha }_n(E)A^T\stackrel{~}{\alpha }_n(E^{}).$$
(23)
So, the derivatives of closed string coordinates transform as
$$_\tau XA^T_\tau X,_\sigma XA^T_\sigma X,$$
(24)
implying that the worldvolume dimensionality of a D-brane does not change under the $`GL(d,𝐙)`$ transformation. We see therefore that the D-brane system with the constant $`g`$ and $`B`$ and the associated NCYM are equivalent to those with $`AgA^T`$ and $`ABA^T`$, where $`AGL(d,𝐙)`$.
* Integer “$`\mathrm{\Theta }`$”-parameter shift of $`E`$, i.e. $`E_{ij}E_{ij}+\mathrm{\Theta }_{ij}`$ with $`\mathrm{\Theta }_{ij}=\mathrm{\Theta }_{ji}𝐙`$:
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)=\left(\begin{array}{cc}I_d& \mathrm{\Theta }\\ 0& I_d\end{array}\right)\mathrm{s}.\mathrm{t}.\mathrm{\Theta }^T=\mathrm{\Theta }.$$
(25)
Under this symmetry, the background fields and the oscillator modes transform as
$$gg,BB+\mathrm{\Theta },\alpha _n(E)\alpha _n(E^{}),\stackrel{~}{\alpha }_n(E)\stackrel{~}{\alpha }_n(E^{}).$$
(26)
implying that closed string coordinates remain unchanged in the form $`X=\frac{1}{2}(X_{}+X_+)`$, meanwhile the $`B`$-field shifts by the integer-valued antisymmetric matrix $`\mathrm{\Theta }`$ as $`BB+\mathrm{\Theta }`$. So, NCYM associated with D$`p`$-brane with constant $`B`$ is equivalent to NCYM associated with D$`p`$-brane with $`B+\mathrm{\Theta }`$. In particular, CYM associated with D$`p`$-brane without $`B`$-field is equivalent to NCYM associated with D$`p`$-brane in the integer-valued $`B`$-field background.
## 3 D$`p`$-Brane with Rank 2 $`B`$ Field and (D$`(p2)`$,D$`p`$) Bound State
In this section, we restrict our attention to the case of D$`p`$-brane with the rank 2 NS $`B`$-field and $`g_{\mu \nu }=\eta _{\mu \nu }`$. We choose the non-zero component of the $`B`$-field to be $`B_{p1,p}`$. In this case, the boundary condition (3) for the open string takes the following form:
$`_\sigma X^i`$ $`=`$ $`0,i=0,1,\mathrm{},p2,`$ (27)
$`_\sigma X^i^{}+B_{i^{}j^{}}_\tau X^j^{}`$ $`=`$ $`0,i^{},j^{}=p1,p,`$ (28)
$`\delta X^a`$ $`=`$ $`0,a=p+1,\mathrm{},9.`$ (29)
In the corresponding closed-string channel description of such D-brane configuration, the condition on the bosonic matter boundary state at $`\tau =0`$ is
$`_\tau X^i|_{\tau =0}|B_X`$ $`=`$ $`0,i=0,1,\mathrm{},p2,`$ (30)
$`(_\tau X^i^{}+B_{i^{}j^{}}_\sigma X^j^{})|_{\tau =0}|B_X`$ $`=`$ $`0,i^{},j^{}=p1,p,`$ (31)
$`(X^a|_{\tau =0}x^a)|B_X`$ $`=`$ $`0,a=p+1,\mathrm{},9.`$ (32)
The boundary condition of the form (29) can also be achieved by rotating a D$`(p1)`$-brane in the plane defined by one of its longitudinal direction (say, $`X^{p1}`$) and one of its transverse direction (say, $`X^p`$) and then applying the $`T`$-duality transformation along the $`X^p`$-direction. The rotation angle $`\theta `$ is then related to the $`B`$-field as $`B_{p1,p}=\mathrm{tan}\theta `$.
To obtain the long distance behavior of the massless bosonic fields in closed string states interacting with the D-brane, we project the boundary state onto the massless bosonic states . Namely, the long distance fluctuation of a field $`\mathrm{\Psi }`$ in the closed string spectrum is given by $`\delta \mathrm{\Psi }=P^{(\mathrm{\Psi })}|D|B_{\mathrm{NS},\mathrm{RR}}`$, where $`P^{(\mathrm{\Psi })}`$ is the projector for the field $`\mathrm{\Psi }`$, $`D=\frac{\alpha ^{}}{4\pi }_{|z|1}\frac{d^2z}{|z|^2}z^{L_0a}\stackrel{~}{z}^{\stackrel{~}{L}_0a}`$ ($`a=1/2`$ in the NS-NS sector and $`a=0`$ in the R-R sector) is the closed string propagator and $`|B_{\mathrm{NS},\mathrm{RR}}`$ is the boundary state for the NS-NS or the R-R sector of the closed string. The resulting long distance fluctuation behavior of the massless bosonic fields is
$`\delta \varphi `$ $`=`$ $`{\displaystyle \frac{3p+2\mathrm{sin}^2\theta }{2\sqrt{2}}}nT_p{\displaystyle \frac{V_{p+1}}{k_{}^2}},`$ (33)
$`\delta h_{\mu \nu }`$ $`=`$ $`nT_p{\displaystyle \frac{V_{p+1}}{k_{}^2}}\mathrm{diag}(𝖠,𝖠,\mathrm{},𝖠,𝖡,𝖡,𝖢,\mathrm{},𝖢),`$ (34)
$`\delta B_{p1,p}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\theta }{\sqrt{2}}}nT_p{\displaystyle \frac{V_{p+1}}{k_{}^2}},`$ (35)
$`\delta A_{01\mathrm{}p2}^{(p1)}`$ $`=`$ $`\pm \sqrt{2}nT_p\mathrm{sin}\theta {\displaystyle \frac{V_{p+1}}{k_{}^2}},`$ (36)
$`\delta A_{01\mathrm{}p}^{(p+1)}`$ $`=`$ $`\pm \sqrt{2}nT_p\mathrm{cos}\theta {\displaystyle \frac{V_{p+1}}{k_{}^2}},`$ (37)
where $`T_p=\sqrt{\pi }(2\pi \sqrt{\alpha ^{}})^{6p}`$ is the tension of a D$`p`$-brane with the unit brane charge, $`V_{p+1}`$ is the volume of the D$`p`$-brane worldvolume, $`k_{}^2=_{a=p+1}^9k_a^2`$ is the square of the transverse momentum and
$$𝖠=(p72\mathrm{sin}^2\theta )/8,𝖡=(p7+6\mathrm{sin}^2\theta )/8,𝖢=(p+12\mathrm{sin}^2\theta )/8.$$
(38)
Here, we multiplied the entire boundary state by an overall factor $`n`$, which has to be an integer due to the Dirac quantization condition, so that the D-brane can take arbitrary R-R charge. To express the long distance behavior of the massless fields in ordinary space, rather than in momentum space, we apply the following Fourier transformation, valid for $`p<D3`$:
$$d^{p+1}xd^{Dp1}y\frac{e^{ik_{}y}}{(Dp3)r^{Dp3}\mathrm{\Omega }_{Dp2}}=\frac{V_{p+1}}{k_{}^2},$$
(39)
where $`r=\sqrt{y^ay^a}`$ is the radial coordinate of the (overall) transverse space and $`\mathrm{\Omega }_n=2\pi ^{(n+1)/2}/\mathrm{\Gamma }((n+1)/2)`$ denotes the area of a unit $`n`$-sphere $`S^n`$. And then one has to rescale the fields in the following way so that the fields can be canonically normalized:
$$\phi =\sqrt{2}\kappa \varphi ,g_{\mu \nu }=2\kappa h_{\mu \nu },_{\mu \nu }=\sqrt{2}\kappa e^{\phi /2}B_{\mu \nu },𝒜=\sqrt{2}\kappa A,$$
(40)
where $`\kappa `$ is the ten-dimensional gravitational constant and $`A`$ denotes $`A^{(p1)}`$ or $`A^{(p+1)}`$. The resulting long distance behavior of the massless fields exactly reproduces the asymptotic behavior of the following Einstein-frame supergravity solution for the non-threshold D$`p`$\- and D$`(p2)`$-brane bound state constructed in Ref. :
$`g_{\mu \nu }dx^\mu dx^\nu `$ $`=`$ $`H^{\frac{7p}{8}}h^{\frac{1}{4}}\left[dt^2+\mathrm{}+dx_{p2}^2+h(dx_{p1}^2+dx_p^2)\right]`$ (42)
$`+H^{\frac{p+1}{8}}h^{\frac{1}{4}}\left[dy_1^2+\mathrm{}+dy_{7p}^2\right],`$
$`e^{2\phi }`$ $`=`$ $`H^{\frac{3p}{2}}h,=(1H^1)h\mathrm{cos}\theta \mathrm{sin}\theta dx^{p1}dx^p,`$ (43)
$`𝒜^{(p1)}`$ $`=`$ $`\pm (1H^1)\mathrm{sin}\theta dt\mathrm{}dx^{p2},`$ (44)
$`𝒜^{(p+1)}`$ $`=`$ $`\pm (1H^1)h\mathrm{cos}\theta dt\mathrm{}dx^p,`$ (45)
where $`H=1+\frac{2\kappa nT_p}{(7p)\mathrm{\Omega }_{8p}}\frac{1}{r^{7p}}`$ and $`h^1=\mathrm{cos}^2\theta +H^1\mathrm{sin}^2\theta `$. So, we see that the D$`p`$-brane with the rank 2 constant $`B`$-field is described in long distance region by the D$`p`$\- and D$`(p2)`$-brane bound state.
Just by considering the boundary condition (29) on the open string coordinates, it appears that the $`B`$-field component $`B_{p1,p}`$ or the rotation angle $`\theta `$ can take an arbitrary value. However, it turns out that $`B_{p1,p}=\mathrm{tan}\theta `$ can take only discrete values determined by the Dirac quantization condition for the brane charge (density). From Eq. (37), one can see that the charge densities for the D$`(p2)`$\- and the D$`p`$-branes are $`q_{p2}=\pm \sqrt{2}V_2nT_p\mathrm{sin}\theta `$ and $`q_p=\pm \sqrt{2}nT_p\mathrm{cos}\theta `$, where $`V_2`$ is the volume of the $`(x_{p1},x_p)`$-plane, which is transverse to the D$`(p2)`$-branes but is longitudinal to the D$`p`$-branes. Since the charge density of D$`p`$-brane takes only values which are integer multiples of the fundamental D$`p`$-brane charge density given by $`\mu _p=\sqrt{2}\sqrt{\pi }(2\pi \sqrt{\alpha ^{}})^{3p}`$, we see that $`q_p=N_p\sqrt{2}\sqrt{\pi }(2\pi \sqrt{\alpha ^{}})^{3p}`$, where $`N_p𝐙`$ is the total number of D$`p`$-branes. We note that the relative transverse directions $`x_{p1}`$ and $`x_p`$, which are transverse to the D$`(p2)`$-branes, are delocalized and therefore there are infinitely many fundamental D$`(p2)`$-branes (with the charge density $`\mu _{p2}=\sqrt{2}\sqrt{\pi }(2\pi \sqrt{\alpha ^{}})^{1p}`$) packed on the $`(x_{p1},x_p)`$-plane. However, there are finite numbers (say $`n_{p2}`$) of D$`(p2)`$-branes per $`(2\pi \sqrt{\alpha ^{}})^2`$ area of this 2-plane (Cf. Refs. ). So, the total number of D$`(p2)`$-branes is $`N_{p2}=n_{p2}V_2/(2\pi \sqrt{\alpha ^{}})^2`$ and the total charge density is $`q_{p2}=N_{p2}\mu _{p2}`$. (The total number $`N_{p2}`$ of D$`(p2)`$-branes is finite \[infinite\], if the volume $`V_2`$ of the 2-plane is finite \[infinite\].) Making use of these facts, we see that the allowed values of the $`B`$-field are restricted by the Dirac quantization condition as
$$B_{p1,p}=\mathrm{tan}\theta =\frac{q_{p2}}{q_p}\frac{1}{V_2}=\frac{n_{p2}}{N_p}(n_{p2},N_p𝐙).$$
(46)
Note, this is valid whether the volume $`V_2`$ is finite or infinite.
Since the (rank 2) $`B`$-field can take only rational values, one can always set its non-zero components equal to zero by applying the (integer-valued) $`T`$-duality transformations, as we explain in the following. We consider only the relevant 2-dimensional part of the D$`p`$-brane worldvolume associated with non-zero components $`B_{p1,p}=B_{p,p1}=n_{p1}/N_p`$ of the $`B`$-field, i.e. we consider the $`O(2,2,𝐙)`$ $`T`$-duality transformation. The first step is to make the $`B`$-field take an integer value by applying the $`T`$-duality transformation with the $`O(2,2,𝐙)`$ matrix of the form (22). We can achieve this, for example, by choosing the entries of the $`GL(2,𝐙)`$ matrix $`A`$ in Eq. (22) to be integer multiples of $`N_p`$. The second step is to transform away the resulting integer valued non-zero $`B`$-field components by applying the integer “$`\mathrm{\Theta }`$”-parameter shift $`T`$-duality transformation with the transformation matrix of the form (25). We choose the $`2\times 2`$ antisymmetric matrix $`\mathrm{\Theta }`$ in Eq. (25) to be the negative of the $`B`$-field transformed through the first step. We note that the $`T`$-duality transformations that we applied in the above steps do not change the worldvolume dimensionality of the D-brane. So, we have related a D$`p`$-brane system with constant rank 2 $`B`$-field to a D$`p`$-brane system without $`B`$-field. This result implies the equivalence between the NCYM and CYM associate with such D-brane systems. This result also explains the microscopic origin of the equivalence between the thermodynamics of the supergravity solutions for D$`p`$-branes with nonzero $`B`$-field and those with vanishing $`B`$-field . Namely, the microscopic degrees of freedom of such brane configurations responsible for the thermodynamics are mapped one-to-one under the $`T`$-duality transformation mentioned above.
We notice from the above that different D$`p`$-brane systems with the non-zero $`B`$-field component $`B_{p1,p}=\mathrm{tan}\theta =n_{p2}/N_p`$ with the fixed number $`N_p`$ of D$`p`$-branes but with the different number densities $`n_{p2}`$ of D$`(p2)`$-branes (and therefore different values of $`B_{p1,p}`$ or $`\theta `$) can be mapped under the $`T`$-duality transformations to the same number of D$`p`$-branes without the $`B`$-field. Therefore, the stringy microscopic degrees of freedom for these different D-brane systems are in one-to-one correspondence under the $`T`$-duality transformations. This gives the microscopic explanation for the $`\theta `$ independence of the thermodynamic quantities for the nonextreme supergravity solutions for the (D$`(p2)`$,D$`p`$) bound states, which was previously observed in Refs. . In particular, this implies the equivalence of the two extreme limits corresponding to the cases $`\theta =0`$ and $`\theta =\pi /2`$ to the case with a finite nonzero $`\theta `$. The $`\theta =0`$ case is just the D$`p`$-branes without $`B`$-field, i.e. $`B_{p1,p}=\mathrm{tan}\theta =0`$. In the $`\theta =\pi /2`$ case (i.e. $`B_{p1,p}=\mathrm{tan}\theta =\mathrm{}`$), there are infinitely many D$`(p2)`$-branes per $`(2\pi \sqrt{\alpha ^{}})^2`$ area of the $`(x_{p2},x_p)`$-plane, i.e. $`n_{p2}=\mathrm{}`$. In this case, the second term in the second line of the open string boundary condition in Eq. (29) dominates, thereby the boundary condition (29) becoming that of open strings attached to D$`(p2)`$-branes. From this, one can see the equivalence of the system of D$`p`$-branes with nonzero constant rank 2 $`B`$-field and the system of infinitely many D$`(p2)`$-branes densely stacked on the 2-dimensional plane in the transverse space (thereby the D$`(p2)`$-branes becoming delocalized on the 2-plane), which was previously conjectured .
We comment on the decoupling limit of the D-brane worldvolume theories. The NCYM decoupling limit is defined as the limit in which $`\alpha ^{}\epsilon ^{\frac{1}{2}}0`$ and $`g_{ij}\epsilon 0`$ such that $`G^{ij}`$ and $`\theta ^{ij}`$ in Eq. (6), including $`B_{\mathrm{SW}}=B/(2\pi \alpha ^{})`$, are held fixed . (Note, the difference in the convention of the $`B`$-field $`B_{ij}`$ in this paper from the one $`B_{\mathrm{SW}ij}`$ in Ref. .) So, in the NCYM decoupling limit, the $`B`$-field behaves as
$$B_{p1,p}=\mathrm{tan}\theta =\frac{n_{p2}}{N_p}=\frac{\stackrel{~}{b}}{\alpha ^{}},$$
(47)
with $`\alpha ^{}0`$ and the noncommutative parameter $`\stackrel{~}{b}`$ held fixed. So, in order for the noncommutative effect on the D-brane worldvolume to survive in the decoupling limit, $`B_{p1,p}`$ and therefore the number density $`n_{p2}`$ of D$`(p2)`$-branes have to go to infinity. The total number of the D$`(p2)`$-branes is given by $`N_{p2}=n_{p2}V_2/(2\pi \sqrt{\alpha ^{}})^2`$. Note, the NCYM decoupling limit condition on the coordinates $`x_{p1,p}`$ of the supergravity solution (45) is $`x_{p1,p}=\frac{\alpha ^{}}{\stackrel{~}{b}}\stackrel{~}{x}_{p1,p}`$ such that $`\stackrel{~}{x}_{p1,p}`$ are held fixed . So, the total number of D$`(p2)`$-branes is reexpressed as
$$N_{p2}=\alpha ^{}n_{n2}\stackrel{~}{V}_2/(4\pi ^2\stackrel{~}{b}^2),$$
(48)
where we used the relation $`\stackrel{~}{V}_2=\left(\frac{\alpha ^{}}{\stackrel{~}{b}}\right)^2V_2`$. From Eq. (47), we see that the D$`(p2)`$-brane number density goes to infinity as $`n_{n2}1/\alpha ^{}`$. So, even if the D$`(p2)`$-brane number density $`n_{p2}`$ diverges in the NCYM decoupling limit, the total number $`N_{p2}`$ of the D$`(p2)`$-branes can be finite, if $`\stackrel{~}{V}_2<\mathrm{}`$. As pointed out in the above, in the $`n_{p2}\mathrm{}`$ limit the open string boundary condition (29) reduces to that for the open string ending on D$`(p2)`$-brane, meaning that the long distance behavior of the massless bosonic fields in the closed string states interacting with such D-brane system, obtained from the boundary state through the projection, reproduces that of the supergravity solution for the D$`(p2)`$-branes. Therefore, in the NCYM decoupling limit, the system of D$`p`$-branes with rank 2 constant $`B`$-field reduces to the system of $`N_{p2}`$ numbers of D$`(p2)`$-branes, which are densely packed on the 2-plane associated with non-zero components of the $`B`$-field, implying the equivalence of NCYM in $`p+1`$ dimensions and CYM with the gauge group $`U(N_{p2})`$ in $`p1`$ dimensions . The rank $`N_{p2}`$ of the gauge group is determined by $`N_p`$, $`\stackrel{~}{b}`$ and $`\stackrel{~}{V}_2`$ through Eqs. (47) and (48). When $`\stackrel{~}{V}_2=\mathrm{}`$, the total number $`N_{p2}`$ of the D$`(p2)`$-branes is infinite and therefore NCYM in $`p+1`$ dimensions is equivalent to CYM with the gauge group $`U(\mathrm{})`$ in $`p1`$ dimensions , as pointed out in Ref. . The equivalence of the various D-brane configurations connected through the $`T`$-duality transformations, on which we elaborated in the previous paragraphs, therefore gives a microscopic explanation for the equivalence among the following gauge theories: $`(i)`$ CYM in $`p+1`$ dimensions associated with the decoupling limit of D$`p`$-branes without $`B`$-field, $`(ii)`$ NCYM in $`p+1`$ dimensions associated with the NCYM decoupling limit of D$`p`$-branes with nonzero constant rank 2 $`B`$-field, $`(iii)`$ CYM with the gauge group $`U(N_{p2})`$ or $`U(\mathrm{})`$ in $`p1`$ dimensions.
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# Regular dendritic patterns induced by non-local time-periodic forcing
## I Introduction
Complex patterns observed in nature attract a considerable interest recently . The complicated spatiotemporal behavior that leads to the formation of such patterns is usually associated with the instabilities of systems transforming under non-equilibrium conditions. A spectacular example of practical importance is dendritic growth that appears in anisotropic systems where the propagation of the transformation front is coupled with the diffusion of a conserved quantity. The respective diffusional instabilities lead to the formation of a traveling quasi-steady-state tip that emits surface undulations which evolve into side arms in directions determined by the anisotropy of the system . Many of the technologically important materials form by dendritic solidification . Analogous phenomena have been reported in biological systems , anisotropic Hele-Shaw cells , and in cosmology/particle physics . Although experiments on the freezing of transparent liquids clarified many of the essential features of dendrite formation, important questions remained open. For example, the role played by thermal fluctuations in side-branching is the subject of continuing investigations . It is anticipated that a selective amplification of thermal noise is responsible for the side-branching and for the observed irregularity of dendritic patterns .
While the steady-state behavior of dendritic growth is understood fairly well in the framework of recent theories (microscopic solvability and phase-field theories ), less is known of the dynamic response of the dendritic morphology to time-dependent external perturbations. The non-linear systems often develop regular patterns under periodic forcing . Considering the inherent non-linearity of the equations describing dendritic growth, it is reasonable to expect that such periodic perturbations lead to resonance patterns that might be used to influence the growth morphology. Besides its scientific interest, a detailed understanding of such phenomena could open novel routes for designing materials for specific applications.
Time-periodic forcing of dendritic side-branching has already been realized via heating locally the dendrite tip by laser beam , or by exposing it to an oscillatory flow field . Both methods lead to the formation of fairly regular morphologies in a suitable frequency range. However, these methods cannot be easily used for spatially homogeneous controlling of growth morphologies in large volumes. To circumvent this problem, we demonstrated the possibility for tuning the dendritic morphology by spatially homogeneous time periodic forcing via modulated pressure and heating, that change the undercooling instantaneously and homogeneously in large volumes.
In this paper we present a detailed analysis of the dynamic response of dendritic solidification to oscillating pressure and heating. In identifying the resonance conditions and the most important process parameters, we rely on phase-field modeling, one of the most potent methods for describing morphology evolution. The theoretical predictions will be confronted with experiments on quasi-2D liquid crystal layers, known as suitable model materials . The rest of this paper is structured as follows. In section II we briefly summarize the phase-field equations that incorporate external forcing, and introduce quantities for characterizing growth morphology. Section III is devoted to the experimental setup and the details of the measurement techniques. In section IV, we present the phase-field predictions and compare them with the experimental results. In section V, our findings are discussed in the light of theory and experiment on tip-localized forcing. Finally, a few concluding remarks will be made in section VI.
## II Phase-field theory with external forcing
We investigate the growth of the crystalline phase into an undercooled single component liquid at a temperature between the melting point and the hypercooling limit (“unit undercooling”) . In this temperature range, the thermal diffusion controls the growth process, as opposed to the (molecule attachment) kinetics controlled mechanism taking place below the hypercooling temperature . Owing to the heat release during solidification, the crystalline phase heats up to its melting point, while the rest of the latent heat is transported into the cold liquid. Under such conditions, the planar front is inherently unstable against thermal perturbations despite the stabilizing effect of the interfacial free energy .
The phase-field theory is a powerful tool to study moving boundary problems. It represents a dynamic extension of the Cahn-Hilliard theory of first-order phase transformations, in which the evolution of a non-conserved order-parameter is coupled to thermal or chemical diffusion . Here we use a thermodynamically consistent version (Model I of ) which ensures locally positive entropy production and was modified to incorporate anisotropies of the interfacial free energy and the kinetic coefficient . The local state of the matter is represented by the phase-field $`\varphi (`$r$`,t)`$ .
To incorporate external forcing, the equations that describe the evolution of the phase-field and the temperature field $`T(`$r$`,t)`$ in 2D
$`\tau (\theta ){\displaystyle \frac{\varphi }{t}}=Q(T)p^{}(\varphi )G^{}(\varphi ){\displaystyle \frac{\delta }{\delta \varphi }}\{{\displaystyle \frac{1}{2}}[\kappa (\theta )]^2|\varphi |^2\},`$ (1)
$$\{c_p+[p(\varphi )1]L^{}(T)\}\frac{T}{t}+L(T)p^{}(\varphi )\frac{\varphi }{t}=k^2T,$$
(2)
have to be modified. Here $`\tau `$ is an empirical relaxation coefficient, whose inverse is an intrinsic interfacial mobility, and $`\frac{\delta }{\delta \varphi }`$ stands for variation with respect to $`\varphi `$. Other notations are as follows: $`p(\varphi )=\varphi ^3(1015\varphi +6\varphi ^2)`$, $`\kappa ^2/2`$ the coefficient of the square-gradient term in the entropy functional, that depends on the orientation represented by the angle $`\theta `$ (defined as $`tan(\theta )=_y\varphi /_x\varphi `$), $`G(\varphi )=\varphi ^2(1\varphi )^2/4e`$, $`e`$ is a constant that determines the height of the intermediate maximum of the double well-potential, $`Q(T)=_{T_m}^T\frac{L(\eta )}{\eta ^2}𝑑\eta =\frac{\mathrm{\Delta }G(T)}{T}`$, $`\mathrm{\Delta }G(T)`$ the Gibbs free energy difference between the liquid and solid, while $`c_p`$, $`L(T)`$, and $`k`$ are the specific heat of the liquid, the enthalpy difference between the liquid and the solid, and the thermal conductivity, respectively. The primed quantities denote derivatives with respect to the arguments. Note that $`\varphi `$ and $`p`$ are dimensionless. The units of $`k`$ and $`\kappa ^2`$ are W m<sup>-1</sup> K<sup>-1</sup> and J m<sup>-1</sup> K<sup>-1</sup>, respectively; while $`Q,G,c_p`$ and $`1/e`$ are measured in J m<sup>-3</sup> K<sup>-1</sup>.
The dominant effect of the pressure modulation $`p(t)=p_0+\mathrm{\Delta }p(t)`$ is a variation of the melting point described by the Clausius-Clapeyron law $`T_m(p)=T_m(p_0)+\mathrm{\Delta }p\mathrm{\Delta }V/S_m`$, where $`\mathrm{\Delta }V`$ is the volume change upon solidification, and $`S_m`$ is the entropy of fusion. This enters into Eq. (1) via the Gibbs free energy difference which may be approximated as $`\mathrm{\Delta }G(T,p)S_m[T_m(p)T]`$ at small undercoolings ($`T_mTT_m`$). The other relevant quantities \[specific heat, heat of fusion, and entropy of fusion $`L^{}(T)`$\] are far less sensitive to the variation of pressure. For example, integrating the Clausius-Clapeyron relationship one obtains $`\mathrm{\Delta }L=\mathrm{\Delta }p\mathrm{\Delta }V[1+\mathrm{\Delta }c_p/S_m+T(\alpha _lV_l\alpha _sV_s)/\mathrm{\Delta }V]`$, where $`\mathrm{\Delta }c_p=c_{p,l}c_{p,s}`$ is the variation of the specific heat upon melting, $`\alpha _{l,s}`$ are the thermal expansion coefficients, while subscripts $`l`$ and $`s`$ refer to the liquid and the solid respectively. Under typical conditions of our experiments \[pressure amplitude $``$ 2 bar, undercooling of 1 K, and physical properties of CCH3 (Appendix A)\], the undercooling varies by $`6\%`$, while the relative change of $`L`$ (and $`S_m`$) is only $`10^3`$. We found that such pressure effects on $`L`$ and $`S_m`$ do not influence the results perceptibly; i.e. pressure modulation needs to be taken into account only via $`Q(T)`$ in Eq. (1). Modulated heating, in turn, appears as a time dependent source term in Eq. (2).
At this stage it is advantageous to introduce dimensionless variables. This is done similarly to that in . However, to avoid complications emerging from an oscillating melting point, the reduced temperature has been redefined in terms of a constant reference temperature $`T_r`$ as $`u(`$r$`,t)=(TT_r)/(T_rT_{\mathrm{}})`$, where $`T_{\mathrm{}}`$ is the initial temperature of the undercooled liquid. The mathematical problem is invariant to the choice of $`T_r`$ so far as $`A=u_m=(T_mT_r)/(T_rT_{\mathrm{}})`$ is inserted as in Eq. (3). (With $`T_r=T_m`$, the usual form of the phase-field model is recovered.)
The relevant physical properties are combined into dimensionless parameters $`\mathrm{\Delta }=c_p(T_rT_{\mathrm{}})/L`$, $`\alpha =\sqrt{2}\omega S_mL/(12c_p\sigma _o)`$, $`\tau _o=S_mD\beta _o/\sigma _o`$, and $`ϵ=\delta /\omega `$, $`\stackrel{~}{\sigma }^{^{}}=\mathrm{d}\stackrel{~}{\sigma }/\mathrm{d}\theta `$, where $`D`$ is the thermal diffusion coefficient, $`\delta `$ the interface thickness, and $`\omega `$ a reference length comparable to the size of a well developed crystallite. The anisotropies of the interfacial free energy $`\sigma `$ and the kinetic coefficient $`\beta `$ are represented by the dimensionless functions $`\stackrel{~}{\sigma }=1+\sigma _ncos(n\theta )/2`$ and $`\stackrel{~}{\beta }=1+\beta _ncos(n\theta )/2`$ of $`n`$-fold symmetry, that are related to the dimensional quantities via relationships $`\sigma (\theta )=\sigma _o\stackrel{~}{\sigma }(\theta )`$ and $`\beta (\theta )=\beta _o\stackrel{~}{\beta }(\theta )`$. Here $`\sigma _n`$ and $`\beta _n`$ are the anisotropy parameters. Lengths and time are scaled by $`\omega `$ and $`\omega ^2/D`$, respectively. The coefficient $`\kappa `$ and the constant $`e`$ are related to $`\delta `$ and the interfacial free energy .
With these notations Eqs. (1) and (2) transform to
$`ϵ^2\tau _o\stackrel{~}{\beta }\stackrel{~}{\sigma }{\displaystyle \frac{\varphi }{t}}=\varphi (1\varphi )\{\varphi {\displaystyle \frac{1}{2}}+30ϵ\alpha \mathrm{\Delta }[uA(t)]\varphi (1\varphi )\}`$ (3)
$`ϵ^2{\displaystyle \frac{}{x}}[\stackrel{~}{\sigma }\stackrel{~}{\sigma }^{^{}}{\displaystyle \frac{\varphi }{y}}]+ϵ^2{\displaystyle \frac{}{y}}[\stackrel{~}{\sigma }\stackrel{~}{\sigma }^{^{}}{\displaystyle \frac{\varphi }{x}}]`$ (4)
$`+ϵ^2[\stackrel{~}{\sigma }^2\varphi ]`$ (5)
$$\frac{u}{t}+\frac{1}{\mathrm{\Delta }}30\varphi ^2(1\varphi )^2\frac{\varphi }{t}=^2u+B(t),$$
(6)
where the modulated pressure and heating are incorporated via terms $`A(t)`$ and $`B(t)`$, respectively.
In this work Eqs. (3) - (6) have been solved numerically on an $`N\times N`$ rectangular grid, $`N=1000`$, that corresponds to an area $`l\times l`$ of dimensionless linear size $`l=5`$ and a grid spacing $`\mathrm{\Delta }x=0.005`$. An explicit finite difference scheme has been employed in the case of Eq. (3), while Eq. (6) has been handled by the alternating-direction implicit method which is unconditionally stable . At $`t=0`$ a crystalline nucleus ($`\varphi =0`$) is placed at the center of the model area filled by uniformly undercooled liquid ($`\varphi =1`$, $`u=1`$). The reduced temperature and the phase-field were kept constant ($`u=1`$ and $`\varphi =1`$) at the boundaries of the model space. To save CPU time, we used a primitive adaptive mesh technique: far from the solidification front Eq. (6) has been solved on a rough grid (of spacing $`10\mathrm{\Delta }x`$). To model the effect of thermal fluctuations and to suppress undesirable lattice effects that favor specific growth directions, a spatially and temporally uncorrelated noise of amplitude 0.01 and zero mean value has been added to the dimensionless temperature $`u`$ in every time step. If not stated otherwise, we use the dimensionless parameters: $`\alpha =350`$, $`\tau _o=20`$, $`ϵ=0.005`$, $`\mathrm{\Delta }x=0.005`$, and time step $`\mathrm{\Delta }t=10^4`$. Owing to the known limitations of phase-field modeling , we performed our calculations at relatively large undercoolings, $`\mathrm{\Delta }=0.400.68`$.
To characterize the spatiotemporal behavior of the evolving dendritic morphology the following quantities were determined:
1. We measured the width $`w`$ of the dendrite behind the tip at a distance of $`\zeta =0.75`$ (150 pixels). To investigate this feature for longer times, we performed separate simulations on a $`1200\times 600`$ grid, oriented so that the larger dimension be parallel with the direction of growth (here we used $`\mathrm{\Delta }t=810^5`$). The amplitude $`A_\nu `$ in the Fourier transform $`W(f)`$ of $`w(t)`$ characterizes the response synchronous with external forcing.
2. The symmetry of the growth patterns was characterized by the quantities a and b:
$$a=\frac{_{i,j=N/2}^{N/2}(\varphi _{i,j}\varphi _{i,j})^2}{M}$$
(7)
$$b=\frac{_{i,j=N/2}^0(\varphi _{i,j}\varphi _{j,i})^2}{K}$$
(8)
where $`M`$ and $`K`$ are the numbers of pixels in which $`0.4<\varphi <0.6`$ for the whole system and for the lower left quarter, respectively. Accordingly, $`a0`$ measures the symmetry of the whole domain with respect to the $`y`$axis, while $`b0`$ quantifies the axisymmetry of a main branch. Both parameters are zero for symmetric patterns.
## III Experimental setup and conditions
The nematic-smectic B ($`NS_B`$) phase transition of liquid crystals is recognized as an appropriate model of crystallization in liquids . A specialty of this phase transformation is a large anisotropy of the interfacial free energy , a property expected to be advantageous for the regularization of dendritic morphology on theoretical grounds (see section IV.A and ).
To test the predictions, experiments have been performed on thin CCH3 liquid crystal layers. Ready-made cells of E.H.C. Co. (Japan) KSRP-10 (of thickness 10 $`\mu `$m) and KSRP-02 (2 $`\mu `$m) have been filled with CCH3 (Merck, Darmstadt). The surface treatment of the bounding glass plates assured the planar alignment of both $`N`$ and $`S_B`$ phases \[the directors n$`(N)`$ and n$`(S_B)`$ that describe the average orientation of the elongated molecules are in the plane of the cell\], and the conducting layers on the bounding plates were used as electrodes.
For pressure modulation the liquid crystal cell was placed into a brass box (see Fig. 1) surrounded by a temperature controlled hot-stage of accuracy $`\pm `$ 3 mK. The gas pressure in the brass box has been regulated by a computer controlled solenoid valve system that switches on and off an excess pressure $`p_\mathrm{e}`$ preset between 0 to 2 bar with an accuracy of $`\pm `$ 0.03 bar. This allows square wave-like pressure modulations in the frequency range up to $`\stackrel{~}{\nu }2`$ Hz.
The modulated heat release in the bulk has been realized by periodically transmitting a high frequency (600 kHz) electric current through the liquid crystal layer produced by gating the signal of a function generator with $`\stackrel{~}{\nu }1`$ Hz. The local off–plane heat transport (a precondition of regulation with oscillatory heating, see section IV.A and ) is ensured by the quasi-2D sample geometry and by the heat transport through the bounding glass plates.
The growth patterns were monitored in transmitting mode via a polarizing microscope equipped with a CCD camera, a method relying on the different optical properties of the nematic and smectic B phases. The images were stored and processed by a PC. The spatial and time resolutions of the system were $`512\times 512`$ pixels and 0.04 s, respectively. The calibration procedure with 6.3 $`\times `$ objective and 3.2 $`\times `$ projector combination gave scale factors of $`1.35\pm 0.01`$ $`\mu `$m/pixel in $`x`$ direction and $`0.95\pm 0.01`$ $`\mu `$m/pixel in $`y`$ direction.
The same setup has been used to measure the pressure coefficient of the temperature of transformation between the smectic-B and nematic phases (see Appendix A).
## IV Results
### A Theoretical predictions
To find resonance patterns we investigated the parameter space defined by the frequency and amplitude of the modulation, the undercooling and the anisotropies of the interfacial free energy and the kinetic coefficient. The calculations were performed with four- and six-fold anisotropies ($`n=4`$ and $`6`$ in expressions for $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{\beta }`$).
Before presenting our results on non-local periodic forcing we characterize the unmodulated reference state.
#### 1 Free dendritic growth
Independently of the symmetry of the system, a weak uncorrelated (noise-induced) side-branching occurs in the unmodulated reference states shown in Figs. 2(a) and 3(a). The width $`w(t)`$ of a main dendrite arm measured at a dimensionless distance of $`\zeta =0.75`$ behind the tip \[Fig. 4(a)\] indicates surface undulations that can be decomposed into a spectrum of oscillations which covers the dimensionless frequency range of $`30`$ to $`170`$, centered around a characteristic frequency of $`f_0100`$ \[see insert in Fig. 4(b)\]. This behavior accords with that seen in the phase-field simulations of Karma and Rappel , and with experimental results . In agreement with , we find that far from the tip the characteristic wave length $`\lambda _0=v/f_0`$ depends much more weakly on the distance from the tip than predicted by the microscopic solvability theory , a behavior interpreted in terms of the “stretching” of the perturbations that travel along the curved perimeter of the dendrite.
Considering that the applied uncorrelated noise probes the dynamic response of the system to a broad variety of frequencies, it is reasonable to seek “resonance” in the vicinity of the characteristic frequency $`f_0`$ of spontaneous side branching.
#### 2 Pressure modulations
According to the Clausius-Clapeyron law, the modulated pressure translates into a time dependent melting point, and is represented by inserting $`A(t)=u_m(t)=a_0g(t)`$ and $`B(t)=0`$ into Eqs. (3) and (4). Since our experimental setup allows only square-wave modulations (characterized by the filling coefficient $`\xi =t_{\mathrm{on}}/t_0`$, where $`t_0`$ is the period of oscillations and $`t_{\mathrm{on}}`$ the pulse length), the simulations have been performed for this wave form (unless stated otherwise).
(a) The effect of forcing frequency: The frequency dependence of the growth patterns is shown in Figs. 2(b)–(f) and 3(b)–(f). At low forcing frequencies ($`\nu <15`$), the system alternates between two steady-state growth modes yielding uncorrelated side-branching that correspond to the actual undercooling. It is, however, remarkable that the switching transient between the constant pressure stages initiates more pronounced side arms than formed during the constant pressure periods. Regular growth morphologies \[Figs. 2(b)–(e) and 3(b)–(e)\] are observed in the lower half of dimensionless frequency range of spontaneous thermal side branching shown in the insert of Fig. 4(b)\]. At frequencies larger than these, the formation of side branches cannot follow the external forcing, and the uncorrelated thermal side-branching reappears \[Figs. 2(f) and 3(f)\].
In agreement with our earlier results , when regular morphologies are formed, the side-branching and the tip velocity correlate with the pressure modulation.
The parameters $`w(t)`$, $`a`$, and $`b`$ reflect the formation of regular morphologies. For example, a periodic variation of the width $`w(t)`$ of the dendrite is seen \[Fig. 4(a)\], which correlates with the external sinusoidal forcing and leads to the formation of regular side branches. Note that besides the forcing frequency (that appears with a far larger amplitude in $`W(f)`$ than the spontaneous undulations), its second and third harmonics ($`2\nu ,3\nu `$) are also present albeit with an amplitude that diminishes for higher order harmonics \[Fig. 4(b)\].
Apparently, periodic forcing with fixed frequency excites several surface modulations; one synchronous with the forcing frequency ($`\nu `$) and others oscillating with the higher harmonic frequencies ($`2\nu ,3\nu ,`$ …), which indicates a dynamic coupling among these modes associated with the non-linearity of the governing equations.
The relative amplitudes of higher harmonics vary with the wave form of the pressure modulation yielding different growth patterns (see Fig. 5). Fairly similar patterns \[cf. Figs. 5 (a) and (b)\] are observed for sinusoidal and square-wave modulations, provided that for the latter $`\xi =0.5`$ \[i.e., frequencies $`(2n+1)\nu `$, $`n=1,2,\mathrm{}`$ are present in the forcing spectrum\]. When the filling coefficient $`\xi `$ deviates from 0.5, the frequency/phase content of forcing varies (e.g., the second harmonic $`2\nu `$ appears) leading to such effects as the formation of side-branches of twice the number corresponding to the base frequency \[cf. Fig. 5(a)–(c)\]. To understand this behavior we analyze the frequency dependence of the dynamic response of the system to “asymmetric” forcing ($`\xi =0.3`$).
The Fourier amplitudes $`A_\nu `$ and $`A_{2\nu }`$ (corresponding to modes that oscillate with the forcing frequency and its double, respectively) and the symmetry parameters $`a`$ and $`b`$ are presented in Fig. 6 as a function of the forcing frequency $`\nu `$. $`A_\nu `$ and $`A_{2\nu }`$ show maxima if $`\nu `$ or $`2\nu `$ are close to $`f_{exc}70`$ \[Fig. 6(a)\], where $`f_{exc}`$ is the frequency where the synchronous response has the maximum amplitude. Note that $`f_{exc}`$ is somewhat below the characteristic frequency $`f_0100`$ of the unperturbed dendritic growth. A possible explanation of this deviation might be a different “stretching” of perturbations, related to differences seen between shapes of dendrite tips formed in free growth and under external forcing. It is also remarkable, that at low forcing frequencies the amplitude $`A_{2\nu }`$ of the second harmonic becomes larger than $`A_\nu `$. This inversion of the relative magnitudes of $`A_\nu `$ and $`A_{2\nu }`$ is responsible for the side-branch “doubling” shown in Fig. 5. A similar phenomenon is seen in Fig. 3(b), where the second harmonic is the dominant mode as indicated by the striking similarity of patterns corresponding to $`\nu =20`$ and 40 \[cf. Figs. 3(b) and (c)\]. Note the lack of alternating variation of the length of side-branches that appears when the synchronous mode is also present with a significant amplitude \[as in Fig. 5(c)\].
We find that the amplitude $`A_\nu `$ increases approximately exponentially with the distance $`\zeta `$ from the tip of the dendrite. In contrast, $`A_{2\nu }`$ shows a maximum as a function of $`\zeta `$ (Fig. 7).
Remarkably, the symmetry parameters $`a`$ and $`b`$ display minima at lower forcing frequencies \[around $`\nu 20`$; see Fig. 6(b)\] than the maximum of the Fourier-amplitude $`A_\nu `$. This difference in the positions of the extrema is a manifestation of the fact that axial symmetry and frequency content are different constituents of regularity. It appears that $`A_\nu `$ gives a closer representation of the regularity recognized by the human eye than the symmetry parameters.
The overall frequency dependence of the pattern formation supports our anticipation that regularization via periodic forcing with fixed frequencies is possible when the latter fall in the characteristic frequency range of noise-induced side branching.
(b) The effect of the amplitude of forcing: The Fourier amplitudes $`A_\nu `$ and $`A_{2\nu }`$ and the symmetry parameters show that with increasing forcing amplitude $`a_0`$ the regularity of the growth patterns increases (Fig. 8). This is manifested in monotonously increasing $`A_\nu `$ and $`A_{2\nu }`$ that saturate for large $`a_0`$. At the same time, $`a`$ and $`b`$ decrease with $`a_0`$, reflecting the increasing symmetry of the patterns.
(c) The effect of undercooling: The response to a given pressure modulation correlates with the tendency for “natural” side-branching (Fig. 9). While at large undercoolings (e.g., $`\mathrm{\Delta }=0.65`$), the formation of side-branches is rather intense, leading to well developed side-arms (both with or without forcing), at $`\mathrm{\Delta }=0.5`$ or below free growth produces essentially no side-branches, and external forcing is needed to trigger them, also reflected in their smaller size. Note the difference in spatial periodicity, that originates from the temperature dependence of the tip velocity.
The variation of the Fourier amplitudes $`A_\nu `$ and $`A_{2\nu }`$ with undercooling ($`\mathrm{\Delta }`$) is shown in Fig. 10. The maximum of $`A_\nu `$ might be understood as follows. We found that with increasing undercooling $`f_{exc}`$ increases. Accordingly, a fixed forcing frequency $`\nu `$ produces the maximum amplitude at the undercooling for which $`f_{exc}\nu `$. For the same reason, the maximum amplitude for the second harmonic $`2\nu `$ appears at a larger undercooling.
(d) The effect of anisotropy: Theoretical considerations and numerical simulations show that anisotropy plays a central role in the formation of dendritic morphology. It is, therefore, reasonable to expect that it has a similar importance in the formation of regular patterns. This expectation is verified by the dependence of growth forms on the magnitude of anisotropies in the system (Fig. 11). The growth patterns and variations of the Fourier amplitudes $`A_\nu `$ and $`A_{2\nu }`$, and of the symmetry parameters (Fig. 12) indicate that the regularizing effect of external forcing decreases dramatically with vanishing anisotropy, i.e. without well defined orientational preferences the external perturbations are unable to drive the amoebae-like growth forms into a regular pattern. Note that the tip velocity decreases with decreasing anisotropies (note the difference of simulation times the snapshots correspond to), the excitation with $`\nu =60`$ generates periodic side-arms with increasing spacing.
In these calculations we assumed that both the interfacial free energy and the kinetic coefficient are anisotropic. For the sake of simplicity, their ratio was usually kept constant $`\beta _4/\sigma _4=2`$. Our further investigations show, however, that resonance patterns form if any of these anisotropies (that of $`\sigma `$ or $`\beta `$) is sufficiently large.
#### 3 Modulated heating
(a) Alternating heating and cooling: They lead to essentially the same type of resonance patterns as pressure oscillations (Fig. 13), provided that that the net heat production in a period is negligible. A remarkable difference is, however, that while the average tip velocity $`v_0`$ is essentially independent of the amplitude of pressure modulations, it decreases with the amplitude of heat production/extraction (Fig. 14); a phenomenon that might be associated with different efficiencies of heating and cooling.
(b) Oscillatory heating: The introduction of a local off–plane thermal transport described by $`B(t)=b_0g(t)+h[u(`$r$`,t)u_{\mathrm{}}]`$ was necessary to prevent the melting of the crystal for the heating amplitudes needed to generate regular patterns. Here $`h<0`$ is a dimensionless heat transfer coefficient that imitates heat transfer perpendicular to the plane of the liquid crystal layer towards the surrounding of reduced temperature $`u_{\mathrm{}}`$. Since this term serves as a local heat sink, the in-plane thermal diffusion becomes less restrictive. This results in the formation of more compact objects (see Fig. 15) with thick main arms and less developed side-branches. While with an increasing amplitude of the heat pulses the side-arms become more regular, they become smaller due to the dissipated heat.
Much like pressure modulations, oscillatory heating leads to an oscillating tip velocity. Owing to a decreasing average undercooling accompanied with increasing heating amplitude $`b_0`$, the period averaged velocity $`v_0`$ of the tip decreases roughly linearly with $`b_0`$ (Fig. 16). At the same time, the amplitude of the oscillatory part of the velocity $`v_\nu `$ increases about linearly.
In summary, our investigations imply that under well defined conditions both types of non-local forcing can be used to control dendritic growth.
### B Experimental results
In defining the experimental conditions, we utilized the results of the computer simulations. For CCH3, the time and spatial resolution of our experimental setup is optimal at an undercooling of $``$ 1 K. Comparable variations of the undercooling ($``$ 0.1 – 0.2 K) are achievable by applying modulation amplitudes as high as 2 bar or 3 $`\times 10^4`$ W/cm<sup>2</sup>, as estimated on the basis of the pressure coefficient of the equilibrium temperature for the nematic–smectic B phases (see Appendix A) and other relevant properties.
The appropriate frequency range of forcing has been determined by Fourier analysis of the width $`w(t)`$ of a freely growing dendrite measured at 66.5 $`\mu `$m behind the tip. Without perturbation, the side-branching is essentially random \[Figs. 17(a) and 18(a)\]. The characteristic frequency, identified as the peak of the broad spectrum in the insert of Fig. 19, is roughly 1.8 Hz. Accordingly, our experimental setup was designed to cover the frequency range of 0 – 2 Hz.
#### 1 Pressure modulations
The experimental results for pattern formation in the 2 $`\mu `$m thick liquid crystal cell are summarized in Fig. 17. Resonance patterns of fairly regular side-branches are observed \[Figs. 17 (b)–(f)\] at forcing frequencies comparable with the characteristic frequency of the unperturbed state $`\stackrel{~}{\nu }1.8`$ Hz. The figure shows germs that nucleated heterogeneously at the same site with essentially the same orientation several times, yielding reproducible patterns. Although the respective dimensionless undercooling ($`\mathrm{\Delta }=0.06`$) is considerably smaller than in the simulations ($`\mathrm{\Delta }=0.40.68`$), the observed behavior follows closely the predictions of the phase-field theory. This similarity is especially striking when patterns formed at similar reduced frequencies ($`\stackrel{~}{\nu }/f_0`$) are compared \[cf. Figs. 2(c) and 17(b) of $`\stackrel{~}{\nu }/f_0=0.40`$ and 0.42; or Figs. 2(d) and 17(d) of $`\stackrel{~}{\nu }/f_0=0.6`$ and 0.65\]. An interesting observation is that far from the dendrite tip irregularities develop in the lengths of the side-branches, despite the even distance of their trunks. We observed an analogous phenomenon driven by the interaction with neighboring side-arms via thermal diffusion fields in long-time phase-field simulations .
The effect of the forcing amplitude is shown in Fig. 20. Pressure oscillations of amplitude of 1 bar efficiently regularize the side-branch formation. With increasing amplitude the phenomenon becomes more pronounced as predicted by the phase-field calculations. At the largest pressure amplitudes even the secondary side-branches correlate with the pressure oscillations \[see Fig. 20(c)\]. Note again that the competition of neighboring side-branches leads to irregularities in their lengths.
#### 2 Modulated heating
In analogy to pressure modulations, the periodic heating experiments on the 10 $`\mu `$m thick liquid crystal cell also reveal regular side-branching \[see Figs. 18(a)–(c)\]. The formation of side-branches correlates with the external forcing as illustrated in Fig. 18(b), where the black lines denote the position of the tip at the centers of the heating pulses. The correlation is also evident from the power spectrum of the width of the dendrite measured 66.5 $`\mu `$m behind the tip (Fig. 19). In full accord with the phase-field simulations for asymmetric square waves \[Fig. 4(b)\], the power spectrum indicates the presence of modes that are either synchronous with the forcing frequency $`\stackrel{~}{\nu }`$ = 0.46 Hz, or oscillate with doubled frequency ($`2\stackrel{~}{\nu }`$). Even a peak corresponding to the third harmonic ($`3\stackrel{~}{\nu }`$) may be identified, although with an amplitude that is close to the experimental uncertainty.
Remarkably, in the response to “symmetric” forcing $`\xi =0.5`$, the $`2\stackrel{~}{\nu }`$ mode is also present (see the Fourier-spectrum and the short side-arms of doubled frequency in Fig. 21), although this frequency is absent from the forcing spectrum. This finding confirms the non-linear behavior revealed by phase-field simulations, that higher harmonics missing from the forcing spectrum are also excited \[see $`W(f)`$ for sinusoidal forcing in Fig. 4\].
Increasing the heating amplitude or $`\xi `$ so that the period-averaged heating power reaches $`\overline{P}10^4`$ W/cm<sup>2</sup>, the formation of the side-branches is suppressed \[see Fig. 18(d)\], a phenomenon resembling that seen in the numerical simulations (Fig. 15). Further increase of the heating amplitude (and power) melts the dendrites back.
In line with our theoretical predictions (see Fig. 2. of ), a weak oscillation of the tip velocity has been observed that correlates with the forcing. However, its amplitude is just above the resolution of the present experimental setup.
A remarkable feature of the ”regularized” dendrites, not seen in the simulations, is a shift in the position of the side-branches on the two sides of the main tip \[see e.g. Figs. 18(b) and (c)\]. As a result, one cannot use the symmetry parameters $`a`$ and $`b`$ \[defined by eqs. (3) and (4)\] for characterization of the pattern’s regularity. For the same reason, the Fourier spectra shown in Figs. 19 and 21 have been evaluated from the half-width of the dendrite. Despite the shift of the side-branches, the power spectra on the two sides are fairly similar.
The shift in the position of the side-branches might be attributed to the asymmetry of the dendrite tip (due to the angular dependence $`\sigma (\theta )`$ ), shown by a high resolution snapshot of the tip region of a dendrite grown freely at $`\mathrm{\Delta }T=1.0`$ C (Fig. 22). Another remarkable feature is that the left side of the tip is faceted. The first observable surface undulation (marked with arrow in Fig. 22) that evolves later into a side-arm appears on the opposite side. In contrast, the faceted side remains smooth up to the same distance behind the tip.
Finally, one should mention that the electric heating in liquid crystals may have side-effects that are not incorporated into our phase-field model. Switching the electric field on, the orientation of the nematic director n$`(N)`$ changes from planar to homeotropic (perpendicular to the bounding plates), an effect that influences the magnitude and anisotropy of the interfacial free energy, and may induce local flow in the sample. However, these side-effects are of minor importance, since they are present in full strength much below the electric field needed for regularization.
## V Discussion
In this section we confront our results with those on the dynamic response of dendritic growth to local forcing with fixed frequency as emerging from the microscopic solvability theory (MST) and experiment. It is appropriate to mention, that some features of the MST derivation prevent a quantitative comparison. For example, the MST calculations were performed for a 2D symmetric dendrite, while the anisotropy, the stability coefficient, and the Péclet-number were assumed to be small; conditions that are not met in our simulations and experiments. Therefore, a qualitative comparison is only meaningful. An important further difference between previous work and ours is that we used non-local forcing, as opposed with the tip-localized forcing assumed in the MST.
Let us first recall some of the MST predictions on noise amplification (based on the Wentzel-Kramers-Brillouin approximation) that have been tested by comparison with numerical simulations and experiments for fixed frequency perturbations localized at the tip.
(i) Periodic forcing by a fixed frequency $`\nu `$ leads to surface oscillations of amplitude that increases exponentially with the distance from the dendrite tip up to a critical distance $`\zeta _c`$ proportional with $`1/\nu ^4`$. For larger distances the amplitude decreases and eventually dies away. If, in turn, the distance from the tip is fixed and $`\nu `$ is varied, a peak is observed in the amplitude.
(ii) The localized wave packets behave differently; they grow exponentially as they move to arbitrarily large distances from the tip, while the respective characteristic wave length increases.
Our simulations are in line with the previous MST and numerical predictions for tip-localized forcing in the following respects:
(a) In the vicinity of the dendrite tip, the amplitude of the mode synchronous with forcing increases roughly exponentially with the distance $`\zeta `$ from the tip. Owing to the excessive computation time and memory needed, we were unable to study the decay of this mode in detail. Nevertheless, as expected from the MST, at large frequencies \[Fig. 2(f) and 3(f)\] the synchronous mode can only be recognized in the vicinity of the dendrite tip. Another sign, that accords with the presence of a critical distance $`\zeta _c`$ beyond which the excited mode decays, is a maximum observed in the amplitude $`A_{2\nu }`$ of the second harmonic (Fig. 7) in simulations where the forcing spectrum contains $`2\nu `$ with a significant amplitude.
(b) The amplitude $`A_\nu `$ measured at a fixed distance behind the dendrite tip shows a maximum as a function of the forcing frequency \[Fig. 6(a)\].
An interesting feature, revealed by our phase-field simulations, is that besides the synchronous mode, the Fourier spectrum of the dendrite width contains the second and third harmonics with perceptible amplitudes. This finding is confirmed by our experiments. We believe that this is the first direct demonstration of such non-linear effects in connection with diffusional instability induced dendrite formation during a first-order phase transformation. It appears, that this non-linear behavior is not a peculiarity of spatially homogeneous forcing. Although not mentioned in the original works, traces of higher harmonics seem to be present in the respective power spectra for local forcing as well . It is worth mentioning, furthermore, that a similar behavior has been observed in the case of anomalous Saffman-Taylor fingering , suggesting that the appearance of higher harmonic modes in response to periodic external forcing is a common feature of dendrite formation whether driven by diffusional or mechanical instabilities.
Summarizing, in the present state of affairs it seems that the only specific feature of pattern formation under non-local forcing is a simultaneous triggering of side-branches (of first, second, and higher orders) on independent crystallites throughout the sample.
## VI Summary
Our computer simulations and experiments demonstrated that the dendritic morphology can be regularized by non-local time-periodic forcing realized by modulated pressure and Joule heating. These conditions lead to an oscillatory velocity of the dendrite tip, and yield side-branches at regular distances, provided that the frequency of the modulation is close to the “natural” frequency of free side-branch formation. The dynamic response of the system to such non-local forcing can be understood in general on the basis of previous theoretical and experimental results on tip-localized forcing. However, the non-linear effects, such as the appearance of higher harmonic modes, warrant further theoretical work.
###### Acknowledgements.
The authors express their thanks to Dr. György Szabó (Research Institute for Materials Science, Budapest) for calling their attention to the possibility of regulating dendritic morphology by pressure modulations, and to Dr. Nándor Éber (Research Institute for Solid State Physics and Optics, Budapest) for his help in solving technical problems. This work has been supported by research grants OTKA T025139, T031808, F022771 and EU HPMF-CT-1999-00132. Part of the computations were performed on computers donated by the Alexander von Humboldt Foundation.
## A Pressure dependence of the phase transition temperature
The pressure dependence of the $`NS_B`$ phase transition temperature has been measured by a procedure similar to that described in . For different $`p_e`$ a single $`S_B`$ monodomain has been kept at constant size by controlling the temperature for several hours (until the thermodynamic equilibrium state of the system has approached). Fig. 23 shows the change of the phase transition temperature depending on the pressure. Linear fit on the data gives a slope $`\frac{\mathrm{d}T_m}{\mathrm{d}p}=(0.032\pm 0.003)`$ K/bar for the Clapeyron coefficient.
The Clapeyron coefficient can be also calculated from the relation:
$$\frac{\mathrm{d}T_m}{\mathrm{d}p}=\frac{T_m^0\mathrm{\Delta }V_{NSB}}{\mathrm{\Delta }H}$$
(A1)
where $`T_m^0`$, $`\mathrm{\Delta }V_{NSB}`$ and $`\mathrm{\Delta }H`$ are the phase transition temperature at atmospheric pressure, the molar volume change on transition and the molar latent heat of fusion, respectively. From the relevant material parameters of CCH3, namely $`T_m^0=329.45`$ K , $`m=233`$ (molar mass), $`\rho =895.4`$ kg/m<sup>3</sup> (density at 80.6 C ), $`\mathrm{\Delta }H=6247`$ J/mol, $`\mathrm{\Delta }V_{NSB}/V=2.6`$ $`\%`$, and $`\alpha _N=25\times 10^4`$ 1/K (volumetric expansion coefficient in the nematic phase ) one obtains $`\frac{\mathrm{d}T_m}{\mathrm{d}p}=0.033`$ K/bar in excellent agreement with the experimental value.
Although this value is much larger than for metals, it is not unusual for liquid crystals. For example, $``$ 0.03 K/bar has been reported for the nematic – crystal transition in PAA, and $``$ 0.1 K/bar for the nematic smectic transition in p-methoxybenzoic acid . Comparable, or even larger coefficients have been measured on other organic substances such as camphene 0.214 K/atm, pivalic acid 0.0674 K/atm, and succinonitrile 0.0245 K/atm .
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