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# Evolution of the Spin Gap Upon Doping a 2-Leg Ladder
\[
## Abstract
The evolution of the spin gap of a 2-leg ladder upon doping depends upon the nature of the lowest triplet excitations in a ladder with two holes. Here we study this evolution using various numerical techniques for a $`t`$-$`t^{}`$-$`J`$ ladder as the next-near-neighbor hopping $`t^{}`$ is varied. We find that depending on the value of $`t^{}`$, the spin gap can evolve continuously or discontinuously and the lowest triplet state can correspond to a magnon, a bound magnon-hole-pair, or two separate quasi-particles. Previous experimental results on the superconducting two-leg ladder Sr<sub>14-x</sub>Ca<sub>x</sub>Cu<sub>24</sub>O<sub>41</sub> are discussed.
PACS: 71.27.+a, 75.50.Ee, 71.10.-w, 75.40.Mg
\]
Studies of strongly-correlated electrons confined to two-leg ladders and described by $`t`$-$`J`$ and Hubbard models have provided important insights into the high $`T_c`$ cuprate puzzle. These models are known to exhibit a gaped spin liquid state at half-filling and upon doping to evolve into a Luther-Emery state characterized by $`d_{x^2y^2}`$-like pairing and $`4k_F`$ CDW correlations . A key feature of the Luther-Emery state is the existence of a gap $`\mathrm{\Delta }_S`$ in the excitation energy of the spin degrees of freedom. At half-filling, the spin gap $`\mathrm{\Delta }_S`$ is set by the $`𝐊=(\pi ,\pi )`$ magnon excitation energy $`\mathrm{\Delta }_M`$ which is of order $`J/2`$ for an isotropic Heisenberg ladder with a near neighbor exchange interaction $`J`$. However, as discussed by Tsunetsugu et. al., there can be a discontinuous evolution of the spin gap upon doping. In particular, they note that a pair can be dissociated into two charge $`|e|`$ and spin $`S=1/2`$ quasi-particles, and the low-energy continuum for such scattering is set by the pair-binding energy $`\mathrm{\Delta }_P`$. Then, if the pair-binding energy is less than the half-filled spin gap ($`\mathrm{\Delta }_P<\mathrm{\Delta }_M`$), there will be a discontinuous decrease in the spin gap upon doping to a value equal to the pair-binding energy $`\mathrm{\Delta }_P`$. Thus, while there is still an $`S=1`$, $`𝐊=(\pi ,\pi )`$ magnon excitation with energy $`\mathrm{\Delta }_M`$ in the infinitesimally doped ladder, if $`\mathrm{\Delta }_M>\mathrm{\Delta }_P`$ a lower energy $`S=1`$ state exists in which a pair is dissociated into two quasi-particles.
There is a low energy continuum of excited states corresponding to two quasi-particles, each in an even parity $`k_y=0`$ state, which have a total momentum $`(k_x,k_y)=(0,0)`$. Here $`k_y=0`$ for a bonding and $`\pi `$ for an anti-bonding quasi-particle respectively. The singlet and triplet continua start at the same energy $`\mathrm{\Delta }_P`$. In addition to these scattering states, there can also be a bound $`S=1`$ state in which a bonding and an anti-bonding quasi-particle with momentum $`k_y=\pi `$ hybridize with a magnon excitation of the spin background . If there is such a bound magnon-pair with energy $`\mathrm{\Delta }_{MP}`$, then it will set the spin gap in the doped ladder provided $`\mathrm{\Delta }_{MP}<\mathrm{\Delta }_M`$ . Such a scenario occurs e.g. in ladders with anisotropic rung ($`J_{}`$) and leg ($`J_{}`$) couplings in the range $`0.4<J_{}/J_{}<1.4`$ . Here we combine exact diagonalization (ED) and density-matrix-renormalization-group (DMRG) techniques to investigate the evolution of the spin gap when one pair of holes is added to a $`t`$-$`t^{}`$-$`J`$ ladder. The next-near-neighbor one-electron hopping $`t^{}`$ provides a useful tuning parameter to study the interplay of the magnon gap $`\mathrm{\Delta }_M`$, the pair dissociation gap $`\mathrm{\Delta }_P`$, and the bound magnon-pair spin gap $`\mathrm{\Delta }_{MP}`$ in setting the spin gap $`\mathrm{\Delta }_S`$ of the lightly doped ladder.
The Hamiltonian for the $`t`$-$`t^{}`$-$`J`$ ladder is
$`H`$ $`=`$ $`J{\displaystyle \underset{i,\lambda }{}}(\stackrel{}{𝐒}_{i,\lambda }\stackrel{}{𝐒}_{i+1,\lambda }{\displaystyle \frac{1}{4}}n_{i,\lambda }n_{i+1,\lambda })`$ (1)
$`+`$ $`J{\displaystyle \underset{i}{}}(\stackrel{}{𝐒}_{i,1}\stackrel{}{𝐒}_{i,2}{\displaystyle \frac{1}{4}}n_{i,1}n_{i,2})`$ (2)
$`+`$ $`t`$ $`{\displaystyle \underset{i,\lambda ,s}{}}(c_{i,\lambda ,s}^{}c_{i+1,\lambda ,s}+h.c.)+t{\displaystyle \underset{i,s}{}}(c_{i,1,s}^{}c_{i,2,s}+h.c.)`$ (3)
$`+`$ $`t^{}{\displaystyle \underset{i,s}{}}(c_{i,1,s}^{}c_{i+1,2,s}+c_{i,2,s}^{}c_{i+1,1,s}+h.c.),`$ (4)
Here $`c_{i,\lambda ,s}^{}`$ creates an electron of spin $`s`$ on site $`i`$ of leg $`\lambda =1`$ or 2, $`\stackrel{}{𝐒}_{i,\lambda }=(c_{i,\lambda ,s}^{}\stackrel{}{\sigma }_{ss^{}}c_{i,\lambda ,s^{}})/2`$ and $`n_{i,\lambda }=\mathrm{\Sigma }_sc_{i,\lambda ,s}^{}c_{i,\lambda ,s}`$. We have taken both the near-neighbor leg and rung one-electron hopping matrix elements equal to $`t`$ and the diagonal next-near-neighbor term equal to $`t^{}`$. The exchange interaction $`J`$ is taken as isotropic between near-neighbor leg and rung sites and throughout this $`J/t=0.5`$.
We begin with our conclusions shown in Fig. 1(a) where we have plotted the excitation energies $`\mathrm{\Delta }E`$ of various triplet states versus $`t^{}`$. The spin gap $`\mathrm{\Delta }_S`$ of the two leg $`t`$-$`t^{}`$-$`J`$ ladder doped with two holes is defined as the difference between the ground state energies of the system with two holes and $`S=1`$ and $`S=0`$ respectively.
$$\mathrm{\Delta }_S=E_0\left(n_h=2,S=1\right)E_0\left(n_h=2,S=0\right).$$
(5)
The stars in Fig. 1(a) show $`\mathrm{\Delta }_S`$ versus $`t^{}/t`$ obtained from DMRG results on $`2\times L`$ ladders with $`L=32`$. The dashed line is the DMRG result for the magnon excitation of the undoped ladder obtained from
$$\mathrm{\Delta }_M=E_0\left(n_h=0,S=1\right)E_0\left(n_h=0,S=0\right).$$
(6)
That this difference in ground state energies corresponds to the $`(\pi ,\pi )`$ magnon is known from ED calculations in which the momentum of the excitation is specified. The open diamonds show the triplet excitation energy in the $`𝐊=(\pi ,\pi )`$ sector, obtained from a finite size scaling analysis using ED. Finally, the solid curve in Fig. 1 corresponds to the pair-binding energy calculated with DMRG from
$`\mathrm{\Delta }_P`$ $`=`$ $`E_0\left(n_h=2,S=0\right)+E_0\left(n_h=0,S=0\right)`$ (7)
$``$ $`2E_0\left(n_h=1,S={\displaystyle \frac{1}{2}}\right)`$ (8)
with $`n_h`$ the number of holes relative to the half-filled ladder (in agreement with the ED results for $`t^{}=0`$ in Ref. ). As shown in Fig. 1(b), $`\mathrm{\Delta }_P`$ sets the two quasi-particle continuum. Here infinite size extrapolated ED results for the lowest energy excited singlet and triplet states in the $`𝐊=(0,0)`$ sector are plotted as open symbols and the solid circles are DMRG data for the pair-binding energy $`\mathrm{\Delta }_P`$, Eq. (8). These energies are in good agreement, consistent with a picture in which a pair dissociates into two quasi-particles.
As discussed below, we have used ED, in which the momentum of the state can be specified, as well as DMRG calculations of the hole and spin correlations in order to interpret the results shown in Fig. 1(a). Here we summarize what these show. Basically, there are three different regimes set by $`t^{}/t`$. For $`0.5<t^{}/t<0.2`$, the discontinuous drop in the spin gap with doping reflects the fact that the pair binding energy $`\mathrm{\Delta }_P`$ is less than the $`(\pi ,\pi )`$ magnon energy $`\mathrm{\Delta }_M`$ of the undoped ladder. Thus, when the system is doped, a singlet pair can dissociate into two separate quasi-particles with total spin $`S=1`$, reducing the spin gap $`\mathrm{\Delta }_S`$ from $`\mathrm{\Delta }_M`$ to $`\mathrm{\Delta }_P`$. In this region, there is a bound magnon-hole pair with a minimum energy at $`(\pi ,\pi )`$ but its energy $`\mathrm{\Delta }_{MP}`$ is larger than $`\mathrm{\Delta }_P`$ so that the spin gap is set by $`\mathrm{\Delta }_P`$. In the region $`0.2t^{}/t0.35`$, the situation changes. The pair binding energy $`\mathrm{\Delta }_P`$ becomes greater than the energy to create a bound magnon-hole pair $`\mathrm{\Delta }_{MP}`$, but $`\mathrm{\Delta }_{MP}`$ is less than the energy to create a separate magnon $`\mathrm{\Delta }_M`$. Thus, in this parameter region the lowest energy triplet state of the 2-hole doped ladder has momentum $`(\pi ,\pi )`$ and corresponds to a bound magnon-hole pair so that $`\mathrm{\Delta }_S=\mathrm{\Delta }_{MP}`$. Finally, for $`0.35<t^{}/t<0.5`$, the energy of the triplet $`𝐊=(\pi ,\pi )`$ excitation becomes equal to the $`S=1`$ magnon energy of the undoped ladder. Here DMRG calculations of the spin and charge correlations show that the excitation corresponds to a magnon which is uncorrelated with the bound singlet pair. Thus, in this region, there is no discontinuity in the spin gap upon doping.
ED calculations were carried out on $`2\times L`$ ladders with $`L`$ an odd number of sites. Both periodic and anti-periodic boundary conditions for $`L`$ up to 13 were used . In Fig. 2 we show results for the triplet excitation energies in the $`𝐊=(\pi ,\pi )`$ sector for a sequence of $`2\times L`$ ladders with 2 holes for various values of $`t^{}`$. Here the excitation energy is measured relative to the 2-hole $`𝐊=(0,0)`$ ground state. The lowest triplet $`𝐊=(\pi ,\pi )`$ state is found to be separated from a quasi-continuum of higher energy states. In Fig. 2, the error bars mark the difference between the results obtained using periodic and anti-periodic boundary conditions with the open symbols marking the mean value. Since the actual longitudinal momentum for a finite ladder is $`\pi (11/L)`$, we have extrapolated these results using a scaling form $`A+B/L+c/L^2`$. The solid symbols denote the DMRG calculation of the spin gap $`\mathrm{\Delta }_S`$, Eq. (5), for an open $`2\times 32`$ ladder with 2 holes (larger lattices are also included for $`t^{}=0`$). For $`0.2t^{}0.5`$, the extrapolated ED results for the $`𝐊=(\pi ,\pi )`$ triplet pass through the DMRG spin gap $`\mathrm{\Delta }_S`$. However, for $`t^{}=0.5`$, the DMRG determined spin gap lays well below the extrapolated $`𝐊=(\pi ,\pi )`$ triplet. As discussed in the introduction, for $`0.2t^{}`$, the spin gap is set by the excitation in the triplet $`𝐊=(\pi ,\pi )`$ sector. However, for $`t^{}<0.2`$, the spin gap is set by the onset of the two quasi-particle continuum $`\mathrm{\Delta }_P`$ which goes to zero as $`t^{}`$ approaches $`0.5`$. Note that, even when $`\mathrm{\Delta }_S<\mathrm{\Delta }_{MP}`$, the magnon-hole pair state could still be locally stable if the decay process into 2 quasi-particles with the same momentum $`(\pi ,\pi )`$ is impossible. Although the ED results approach the $`𝐊=(\pi ,\pi )`$ magnon energy of the undoped ladder for $`t^{}>0.35`$, DMRG results show that the character of the triplet excitation changes from a bound magnon-hole-pair to a separate magnon and hole-pair state. Thus, in this regime, the spin gap is set by the excitation energy of the magnon $`\mathrm{\Delta }_M`$ and is therefore continuous upon doping.
In order to get a clearer picture of the nature of the triplet excitations which determine the spin gap, we have used DMRG results to study the spin and hole correlations in these states. In Fig. 3 a center section of a $`2\times 32`$ ladder with $`t^{}=0`$ is shown. The upper and middle diagrams show the probability of finding the second hole when the first hole is projected out at the center of the upper leg for the singlet and triplet state respectively. In both of these states, the two holes are bound and the most probable configuration for $`J/t=0.5`$ corresponds to having the holes on diagonal sites. Note that the triplet-bound state is more extended than the singlet bound state. The lowest diagram shows the spin distribution for the $`S_z=1`$ triplet state when the two holes are projected out at their most probable sites. It is clear that this state is a bound magnon-hole pair.
Similar calculations for $`t^{}=0.5`$ show that in the triplet state the two holes are unbound while for $`t^{}=0.5`$ the two holes are bound into a singlet and uncorrelated with the spin 1 excitation. This behavior is shown in Fig. 4, where we have plotted
$$S_z(\mathrm{}_x)S_z(\mathrm{}_x,1)P_h(i)P_h(j)/P_h(i)P_h(j)$$
(9)
versus $`\mathrm{}_x`$ for $`t^{}=0`$, $`0.5`$, and 0.5. Here $`P_h(i)`$ is the projection operator for a hole at the $`i^{\mathrm{th}}`$ site. For $`t^{}=0`$, we have set $`i=(16,2)`$ and $`j=(17,1)`$, corresponding to the most probable hole location. Here, as previously illustrated in Fig. 3, we see that the spin is bound to the hole-pair. For $`t^{}=0.5`$ we again have a situation where the holes are most likely to sit close to each other, and here we have projected them onto $`i=(16,1)`$ and $`j=(16,2)`$. However, in this case, the spin 1 is spread out corresponding to a magnon which is not bound to the hole-pair. Finally, for $`t^{}=0.5`$, one finds that the lowest energy triplet excitation corresponds to two separate quasi-particles.
We finish with a brief discussion of some experimental results for the superconducting two-leg doped ladder Sr<sub>2</sub>Ca<sub>12</sub>Cu<sub>24</sub>O<sub>41</sub> . Nuclear magnetic resonance measurements of the copper-63 Knight shift and relaxation time $`T_1`$ suggest a collapse of the spin gap with pressure. We believe this signals the appearance of new low lying triplet excitations upon doping the ladder planes and points towards a negative value of $`t^{}`$ . In this regime, due to the presence of a low-energy quasi-particle continum located predominantly around the zone center, momentum-resolved experiments like inelastic neutron scattering would be essential to search for sharp finite energy triplet excitations.
To summarize, using ED and DMRG calculations, we have found that the spin gap can evolve in different ways when two holes are doped into a 2-leg ladder. When the 2 holes are added, it is possible that the lowest energy triplet state simply remains the $`𝐊=(\pi ,\pi )`$ magnon so that there is no change in $`\mathrm{\Delta }_S`$. In this case the 2 added holes remain in a bound $`d_{x^2y^2}`$-like singlet state and a triplet magnon similar to that of an undoped ladder is created. As the length of the ladder increases, the interaction between these two entities becomes negligible. We see this happening for $`t^{}>0.35`$. It is also possible that the lowest energy triplet state has $`𝐊=(0,0)`$ and is set by the two-quasi-particle continuum corresponding to the pair-binding energy $`\mathrm{\Delta }_P`$. In this case, there is a discontinuous change in the spin gap upon doping and the lowest energy triplet state arises from the dissociation of a pair into two quasi-particles. We see this for the present model when $`t^{}<0.2`$. Finally, for the intermediate region $`0.2<t^{}<0.35`$ we find that the lowest energy triplet state has $`𝐊=(\pi ,\pi )`$ and corresponds to a bound magnon-hole pair with energy $`\mathrm{\Delta }_{MP}<\mathrm{\Delta }_M`$. In this case, there is again a discontinuous evolution in the spin gap from $`\mathrm{\Delta }_M`$ to $`\mathrm{\Delta }_{MP}`$ upon doping.
We would like to acknowledge useful discussions with Ian Affleck. S.R. White and D.J. Scalapino acknowledge support from the NSF under grant # DMR98-70930 and grant # DMR98-17242 respectively. D. Poilblanc thanks IDRIS (Paris) for allocation of CPU time on the NEC SX5 supercomputers.
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# The Identification of the Submillimeter Galaxy SMM J00266+1708
## 1. INTRODUCTION
Deep surveys of the submillimeter sky using the Submillimeter Common User Bolometer Array (SCUBA) camera (Holland et al. 1999) on the James Clerk Maxwell Telescope have uncovered a population of ultraluminous dusty galaxies at high-redshift (Smail, Ivison, & Blain 1997; Hughes et al. 1998; Barger et al. 1998; Eales et al. 1999; Blain et al. 1999a). This population accounts for a large fraction of the extragalactic background at mm/sub-mm wavelengths (Blain et al. 1999b) and hence is important to our understanding of the distant universe. The sub-mm population is thought to contribute significantly to both the total amount of star-formation (Blain et al. 1999b) and AGN activity (Almaini et al. 1999) at high-redshift. The sub-mm population will likely show a mixture of AGN and starburst properties given their apparent similarities to the local population of ultraluminous ($`L>10^{12}L_{\mathrm{}}`$) infrared galaxies (ULIGs, Sanders & Mirabel 1996). However, we could expect the majority ($`70`$–80%) of the sub-mm galaxies to be predominantly powered by starbursts since this has been found for the local ULIGs (Genzel et al. 1998). The early CO and X-ray data on the sub-mm population support the starburst nature of the population by showing the presence of sufficient molecular gas to fuel the star-formation activity (Frayer et al. 1998, 1999) and the lack of expected X-ray emission if mostly dominated by AGN (Fabian et al. 2000; Hornschemeier et al. 2000). Observations of the dust-rich sub-mm galaxies complement the studies of the ultraviolet-bright Lyman-break sources (Steidel et al. 1996, 1999) which tend to be much less luminous at infrared wavelengths (Chapman et al. 2000). Only by studying both the Lyman-break and the sub-mm populations of galaxies will a complete picture for the star-formation history of the universe emerge.
In order to understand the nature of the sub-mm population, we have been carrying out multi-wavelength observations of individual systems in the SCUBA Cluster Lens Survey (Smail et al. 1998). This survey represents sensitive sub-mm mapping of seven massive, lensing clusters which uncovered 15 background sub-mm sources. The advantage of this sample is that the amplification of the background sources allows for deeper source frame observations. Also, lensing by cluster potentials does not suffer from differential lensing so that the observed flux ratios will represent intrinsic values, despite the possible variation of source size at different wavelengths.
The most challenging aspect for follow-up observational studies of the sub-mm population is determining the proper counter-parts to the sub-mm emission and obtaining their redshifts (Ivison et al. 1998, 2000a). The large 15<sup>′′</sup> SCUBA beam leaves ambiguity in identifying the galaxy associated with the sub-mm emission. The early results based on optical imaging and spectroscopy were overly optimistic in the identification of the sub-mm counter-parts (Smail et al. 1998; Barger et al. 1999; Lilly et al. 1999). Radio data (Smail et al. 2000a) and initial near-infrared (NIR) imaging (Smail et al. 1999) suggest that several of the original candidate optical counter-parts (e.g., Barger et al. 1999) are incorrect. Despite their ultra-high luminosities, many sub-mm galaxies are nearly completely obscured by dust at ultraviolet/optical wavelengths. For these highly obscured galaxies, follow-up radio (Smail et al. 2000a) and/or mm interferometry (Downes et al. 1999; Bertoldi et al. 2000) as well as near-infrared observations are required in order to uncover the proper counter-part. The galaxy SMM J00266+1708 is an excellent example of such a source.
## 2. OBSERVATIONS
### 2.1. Background Data on SMM J00266+1708
The sub-mm galaxy SMM J00266+1708 is the second brightest galaxy in the SCUBA Cluster Lens Survey (Smail et al. 1998). Initially, the source was tentatively associated with a possible interacting pair of galaxies, M1 and M2, revealed by deep $`I`$-band imaging ($`3\sigma =26.1`$ mag) with the Hubble Space Telescope \[HST\] (Smail et al. 1998). Spectroscopy of the galaxy M2 showed a bright \[Oii\] emission line at $`z=1.226`$, consistent with a luminous star-forming galaxy (Barger et al. 1999). In the fall of 1998, we searched for redshifted CO(2$``$1) emission corresponding to the redshift of M2 at the Owens Valley Millimeter Array<sup>1</sup><sup>1</sup>1The Owens Valley Millimeter Array is a radio telescope facility operated by the California Institute of Technology and is supported by NSF grant AST 9981546. (OVRO). We failed to detect any CO emission at the redshift of M2; $`S(\mathrm{CO})<1.3`$ Jy$`\mathrm{km}\mathrm{s}^1`$ ($`3\sigma `$), assuming a standard $`300\mathrm{km}\mathrm{s}^1`$ line width. If M2 was the correct counter-part and the source was similar to the sub-mm galaxies with previous CO detections (Frayer et al. 1998, 1999), we would have expected a CO(2$``$1) line strength of approximately $`7\mathrm{Jy}\mathrm{km}\mathrm{s}^1`$. The nondetection of CO questions the association of M2 with the sub-mm galaxy. Additional optical spectroscopy showed that M1 is at the redshift of the foreground cluster ($`z=0.39`$), and hence, M1 and M2 are not an interacting pair after all (Barger et al. 1999). These results further bring into doubt the initial association of M1 and M2 with the sub-mm emission based on optical morphology alone.
Besides M1 and M2, the only other optically visible source spatially consistent with the SCUBA position is M8 (Fig. 1). However, the galaxy M8 shows no unusual properties that would suggest an association with a luminous sub-mm source. More significantly, the field has a weak radio source (Smail et al. 2000a) whose position is consistent with the SCUBA source, but is slightly offset from M8. Since the brightest sub-mm galaxies tend to have radio counter-parts (Smail et al. 2000a; Barger, Cowie, & Richards 2000), we could reasonably expect an association between the sub-mm galaxy and the radio source. If this is the case, the radio emission of SMM J00266+1708 is too weak to be consistent with a redshift of $`z=0.44`$ for M8 (Barger et al. 1999), based on the redshift dependency predicted for the sub-mm/radio spectral index (Carilli & Yun 1999). Although the sub-mm/radio spectral index only provides an estimate of the redshift given the uncertainties of source temperature and properties (Blain 1999), it does provide a powerful technique for discriminating between low-redshift ($`z0.5`$) and high-redshift ($`z2`$) galaxies. SMM J00266+1708 is expected to be at redshifts $`z>2`$ since its sub-mm/radio spectral index of $`\alpha =1`$ is much larger than any known galaxy at low redshift (Smail et al. 2000a). Therefore, it is unlikely that M8 is the sub-mm counter-part. Since none of the optically detected galaxies are plausibly associated with the sub-mm emission, we have carried out mm-continuum and near-infrared observations of SMM J00266+1708.
### 2.2. OVRO 1.3 mm Continuum Observations
We have taken sensitive 1.3 mm interferometric observations of SMM J00266+1708 in order to accurately constrain the position of the sub-mm source. SMM J00266+1708 was observed several times using the OVRO array in 1999. Approximately 10 hours of on source data were obtained during good conditions in the low resolution configurations of the array (baseline lengths ranging from 15m to 119m). We observed with four separate 1 GHz continuum bands centered at 229.0, 230.5, 233.5, and 235.0 GHz. The 4 GHz of total continuum bandwidth represents a factor of two increase in bandwidth over what was previously achievable at OVRO.
Fig. 1.— The OVRO 1.3 mm contour map overlaid on the HST $`I`$-band optical image (Smail et al. 1998). The rms level is 1.1 mJy/beam, and the contour levels are $`1\sigma \times (3,3,4,5)`$. The source is unresolved by the $`2\mathrm{}`$ OVRO beam. The position of the 21 cm radio source (Smail et al. 2000a) is shown by the cross labeled “VLA”, while the positional uncertainty of the sub-mm source is marked by the cross labeled “SCUBA”.
We observed the bright quasar 3C454.3 every 20 minutes for amplitude and phase calibration. Since 3C454.3 is 22 degrees from SMM J00266+1708, we also interweaved observations of the nearby quasar 0007+106 (B1950.0) in order to test the quality of the calibration. The systematic positional uncertainty of the data is estimated to be better than $`0\stackrel{}{\mathrm{.}}3`$. Observations of the planets Uranus and Neptune were used to calibrate the absolute flux scale. By using the flux history of 3C454.3, we estimate a flux calibration uncertainty of 20% for the data.
We combined the data for the four individual 1 GHz bands to produce the 1.3 mm detection at a mean frequency of 232.0 GHz. The contours in Figure 1 show the resultant natural–weighted image. No primary beam correction was required since the source was located within $`4\mathrm{}`$ (1/8 of the primary beam width) of the phase center. We made no correction for possible variations of source strength across the four individual bands. These variations are expected to be less than 10%, assuming dust emission, and were undetected within the uncertainties of the data.
### 2.3. Keck $`K`$-band Imaging
After the 1.3 mm detection, we obtained $`K`$-band ($`2.2\mu `$m) data to search for the galaxy responsible for the sub-mm emission at the position derived from the OVRO data. We observed SMM J00266+1708 using the Near Infrared Camera (NIRC) on Keck I<sup>2</sup><sup>2</sup>2The W. M. Keck Observatory is operated as a scientific partnership among the California Institute of Technology, the University of California, and the National Aeronautics and Space Administration. The Keck Observatory was made possible by the generous financial support of the W. M. Keck Foundation. on UT 1999 October 01. NIRC is a $`256\times 256`$ pixel InSb detector with a pixel scale of $`0\stackrel{}{\mathrm{.}}15`$ (Matthews & Soifer 1994). We observed using the standard $`K`$-band filter instead of the bluer $`K_\mathrm{s}`$ filter since the object is expected to be red. Integrations were taken using $`10\times 6`$ second coadds, and we randomly dithered the integrations within a $`8\mathrm{}\times 8\mathrm{}`$ box to provide uniformity across the image. We obtained a total of 4.3 hours of data on source. The seeing-disks of the stars observed throughout the night varied from $`0\stackrel{}{\mathrm{.}}3`$ to $`1\mathrm{}`$ (FWHM).
In reducing the data, a combined set of dark frames was subtracted from each individual exposure to remove the dark current as well as the bias level. The dark-subtracted exposures were divided by a normalized skyflat which was generated from the on object exposures themselves. Frames were sky-subtracted using the temporally–adjacent images to produce the reduced exposures. The individual reduced exposures were aligned to the nearest pixel using common objects in the frames. We observed a set of near-infrared standard stars (Persson et al. 1998) at a range of airmasses in order to correct the data for extinction. The uncertainty of the derived magnitude scale is estimated to be better than 0.04 magnitude, based on the dispersion in the zero-points derived throughout the night.
## 3. RESULTS
### 3.1. Images
Figure 1 shows the 1.3 mm continuum map. The source was detected at the $`5\sigma `$ level and was unresolved with a synthesized beam size of $`\theta _b=2\stackrel{}{\mathrm{.}}3\times 1\stackrel{}{\mathrm{.}}9`$. At the observed frequency of 232 GHz (1.29 mm), we derive a flux density of $`6.0\pm 1.1`$ mJy by fitting a Gaussian to the peak of the emission. The total uncertainty in the flux is 28% which includes the 20% calibration error and the rms error of 20% combined in quadrature. The strength of the 1.3 mm source is consistent with extrapolating the thermal dust spectrum expected from the sub-mm source (Fig. 3). We conclude that the 1.3 mm source is associated with the SCUBA source and is likely representative of the bulk of the sub-mm emission.
Figure 2 shows the $`K`$-band image. We detected a new galaxy $`1\stackrel{}{\mathrm{.}}5`$ south-east of the nucleus of M8 at the position of the OVRO 1.3 mm continuum source. Since the seeing varied significantly throughout the night, we combined only the best 2.4 hours of data to produce the $`K`$-band image. After convolving the data with a 2 pixel (FWHM) Gaussian to reduce pixel-to-pixel variations, the resolution of data is $`0\stackrel{}{\mathrm{.}}5`$ (FWHM), including seeing. The image has a $`1\sigma `$ limiting $`K`$-band surface brightness of $`\mu =24.8\mathrm{mag}/\mathrm{}\mathrm{}`$ or $`0.04\mu `$Jy/beam, adopting 646 Jy for the flux density equivalent for a zero $`K`$band magnitude (Neugebauer et al. 1987).
Fig. 2.— The Keck $`K`$-band ($`2.2\mu `$m) image. The arrow points to the new galaxy thought to be the counter-part of SMM J00266+1708. The rms of the image is 24.8 mag/$`\mathrm{}\mathrm{}`$ ($`0.04\mu `$Jy/beam), and the contours are $`1\sigma \times (3,3,4,5,6,8,10,15,20,30,50,80)`$. The seeing disk (beam) of the NIR data is shown in the lower left ($`0\stackrel{}{\mathrm{.}}5\times 0\stackrel{}{\mathrm{.}}5`$). The position of the 1.3 mm source is shown by the cross labeled “OVRO”.
Photometry on the galaxy was done using an aperture to sum up the emission from the regions around the south-eastern peak that are well separated from M8. The integrated emission for the galaxy is detected at about the $`10\sigma `$ level with a magnitude of $`K=22.45\pm 0.11`$ ($`0.68\pm 0.07\mu \mathrm{Jy}`$, $`K_{\mathrm{AB}}24.3`$). There is additional emission westward from the south-eastern peak and near M8. There is no evidence in the $`I`$band image for this extension. If we assume that this emission is associated with the south-eastern peak, we would derive a total magnitude of $`K=21.5`$ after subtracting the bright disk of M8 from the image. For the remainder of the paper, we neglect this additional emission since it is unclear what fraction is associated with M8, the south-eastern peak, or neither.
### 3.2. Astrometry
The $`K`$-band image was registered to the HST optical image on the APM coordinate system using the galaxies M1 and M8. The optical astrometry in the HST image has an rms uncertainty of only $`0\stackrel{}{\mathrm{.}}2`$. However, the dominant source of registration error in comparing the $`K`$-band data with the radio and 1.3 mm data is due to the systematic offsets between the APM system and the radio reference frame (Johnston et al. 1995). These offsets may be as large as $`1\mathrm{}`$ which we adopt as a conservative estimate of the total error in the $`K`$-band position. The positional error for the 1.3 mm source is $`0\stackrel{}{\mathrm{.}}5`$. This includes a systematic astrometry uncertainty of $`0\stackrel{}{\mathrm{.}}3`$ from calibration combined in quadrature with the statistical error related to the signal–to–noise (S/N) of approximately $`0\stackrel{}{\mathrm{.}}4`$ ($`\theta _b/[S/N]`$). The 21 cm radio data have a synthesize beam size of $`5\mathrm{}`$ and a positional error of about $`0\stackrel{}{\mathrm{.}}8`$ ($`\theta _b/[S/N]`$). The 850$`\mu `$m source detected by SCUBA has a positional accuracy of approximately $`3\mathrm{}`$ (Smail et al. 1998).
The positions of SMM J00266+1708 measured at 21 cm, 1.3 mm, 850$`\mu `$m, and $`2.2\mu `$m are all consistent with each other within the errors (Table 1). The positional coincidence alone suggests an association between the sources at different wavelengths. The probability of randomly finding a 22.5 mag $`K`$-band galaxy within a 1 square-arcsec box is only 1%, based on the observed $`K`$-band surface densities of galaxies (Djorgovski et al. 1995; Moustakas et al. 1997). By taking into account the red nature of the galaxy ($`IK>3.6`$), the likelihood of a random association decreases even more. Only about 10% of faint galaxies ($`22.0K22.9`$) have $`(IK)>3.5`$ (Moustakas et al. 1997). Hence, the probability of a chance association of such a faint red galaxy with the 1.3 mm source is only $`10^3`$. An even stronger case can be made on the association between the radio source and the 1.3 mm source. Radio source counts (Richards et al. 1999; Richards 2000) imply the probability of randomly finding such a strong radio source within 1 square arcsec of the 1.3 mm source is only about $`10^5`$. We conclude the sources detected at 21 cm, 1.3 mm, 850$`\mu `$m, and $`2.2\mu `$m are all likely associated with each other.
### 3.3. Spectral Energy Distribution and Redshift
Table 1 presents the flux density measurements and upper-limits observed for SMM J00266+1708. The 21 cm radio data are presented by Smail et al. (2000a). We measure a flux density of $`94\pm 15\mu \mathrm{Jy}`$ for the 21 cm radio source by fitting a Gaussian to the unresolved emission. We also report a 3.5 cm upper-limit for the galaxy based on sensitive observations of the cluster (A. R. Cooray, private communication). The sub-mm flux at 850$`\mu `$m has been previously tabulated by Barger et al. (1999), and the 450$`\mu `$m upper limit is discussed by Smail et al. (2000b). The optical $`I`$-band limit of 26.1 mag is derived from the observations presented by Smail et al. (1998). The 1.3 mm and $`K`$-band measurements are based on the observations presented in this paper.
As stated earlier (§2.1), SMM J00266+1708 is thought to lie at a high redshift of $`z2`$ based on the sub-mm/radio spectral index of the galaxy (Carilli & Yun 1999; Smail et al. 2000a). High redshifts are also required to account for the low 450$`\mu `$m/850$`\mu `$m ratio, assuming dust emission with properties similar to that found in low-redshift luminous starbursts (Dunne et al. 2000). The optical and NIR imaging data by themselves provide little constraint on the redshift of SMM J00266+1708 given the wide range of optical/NIR properties found for the sub-mm population (e.g., Ivison et al. 2000a). The best redshift constraints for SMM J00266+1708 come by fitting the spectral energy distribution (SED) of the galaxy from radio to sub-mm wavelengths. Figure 3 shows the observed SEDs of the low redshift ULIGs Arp 220 and Mrk 231 (Rigopoulou, Lawrence, & Rowan-Robinson 1996), as well as the $`z=1.44`$ extremely red galaxy (ERO) HR10 (Dey et al. 1999), the relatively blue sub-mm starburst SMM J14011+0252 (Ivison et al. 2000a), and the sub-mm galaxy SMM J02399-0136 which contains an AGN (Ivison et al. 1998). These ULIG systems are chosen for comparison since they represent the reasonable range of possible SEDs expected for the sub-mm population. Assuming a likely range of dust temperatures (30–50 K) and fitting all the data to the nearest 0.5 redshift unit, we estimate a redshift of $`z=3.5\pm 1.5`$ for SMM J00266+1708.
The wide range of optical flux densities and colors for ULIGs/sub-mm galaxies (Fig. 3) demonstrates the difficulty of attempting to estimate the far-infrared luminosities of ultraluminous galaxies based solely on optical observations (c.f., Adelberger & Steidel 2000). Bolometric luminosities estimated from sub-mm flux densities are significantly more robust since the observed emission from high-redshift galaxies arises near the peak of the SED.
Fig. 3.— The SED of SMM J00266+1708 (solid circles) adopting the best fit redshift of $`z=3.5\pm 1.5`$ compared to other ULIG systems. All data have been normalized to a rest wavelength of 350$`\mu `$m. The upper-limits for SMM J00266+1708 are marked with downward arrows. The solid line shows an example SED for dust emission assuming a typical temperature of $`T=45`$ K and $`\beta =1.5`$ for the power-law dependence of the dust absorption coefficient.
### 3.4. Lensing
In this section, we estimate the effect of gravitational lensing. SMM J00266+1708 is amplified by the potential of the $`z=0.39`$ foreground cluster Cl 0024+16. Assuming a redshift of $`z=1.23`$, Barger et al. (1999) find an amplification factor of 1.6 due to the cluster. Adopting a more realistic source redshift of $`z3.5`$, the amplification factor is increased to 2.1 using the relationships given by Schneider, Ehlers, & Falco (1992) to rescale the appropriate angular size distances. The amplification factor due to the cluster is fairly insensitive to redshift for $`z2`$.
Since SMM J00266+1708 is near the line of sight of the galaxy M8 ($`z=0.44`$), we could expect additional gravitational lensing due to M8. To estimate this effect, we approximate both M8 and SMM J00266+1708 as point sources and use the equations for a simple Schwarzschild lens (Schneider et al. 1992). M8 appears to be a normal spiral galaxy with a typical absolute $`V`$-band magnitude of $`M_V=20`$ (Smail et al. 1997). Assuming a total mass of $`2\times 10^{11}M_{\mathrm{}}`$ within the central 10 kpc of M8 (which corresponds to the angular separation of $`1\stackrel{}{\mathrm{.}}5`$ between SMM J00266+1708 and M8), we estimate that SMM J00266+1708 is amplified by a factor of 1.14 due to M8. Including lensing by both the cluster and M8, we expect a total magnification factor of about 2.4 ($`\pm 0.5`$) for SMM J00266+1708.
SMM J00266+1708 is the second sub-mm galaxy (ERO-N4 is the other example \[Smail et al. 1999\]) found in the Cluster Lens Survey thought to be lensed by a low redshift spiral galaxy. These results may suggest that lensing due to galaxies may play a significant role for the brightest sub-mm galaxies. The current observations of sub-mm counts indicate that future mm/sub-mm surveys covering large areas of the sky should detect many strongly lensed sources (Blain, Möller, & Maller 1999).
## 4. DISCUSSION
We have concentrated our follow-up studies on the nine brightest galaxies detected in the sub-mm Cluster Lens Survey (Smail et al. 1998). Currently, only three of the nine galaxies have optical counter-parts with redshifts. Two of these have been confirmed by CO observations at OVRO (Frayer et al. 1998, 1999), and the third is a Seyfert ring galaxy at $`z=1.06`$ recently detected in CO with the IRAM interferometer (Soucail et al. 1999; Kneib et al. 2000).
At least four of the nine systems are undetected at optical wavelengths ($`I>26`$–27 after correcting for lens amplification) and have only been detected with $`K`$-band imaging. Two of these are bright enough ($`K=19.1`$, 19.6 mag) to be classified as EROs (Smail et at. 1999). An additional faint ($`K=21`$, $`R>26`$) galaxy was found associated with a relatively bright (500 $`\mu `$Jy) radio counter-part (Ivison et al. 2000a, 2001). SMM J00266+1708 is the faintest sub-mm counter-part found to date at $`K=22.5`$. Only two of the nine sources still require deep $`K`$-band imaging and have uncertain counter-parts. Depending on the results for these last two unknown systems, the data suggest that approximately 40%–70% (4/9–6/9) of the sub-mm population as a whole are very faint/red galaxies which are undetected at optical wavelengths. It is still unclear what fraction of the sub-mm galaxies are EROs \[$`RK>6`$\] (Thompson et al. 1999). At least two out of nine ($`20`$%) galaxies in the Cluster Lens Survey are EROs (Smail et al. 1999), and other sub-mm surveys have uncovered additional faint EROs associated with SCUBA sources (Ivison et al. 2000b). Future, much deeper optical observations may show that a high fraction of sub-mm galaxies galaxies are EROs.
### 4.1. Intrinsic Properties of SMM J00266+1708
SMM J00266+1708 is an extremely faint galaxy. After correcting for lensing, the source-frame magnitude of the galaxy is $`K=23.4`$. We adopt a lensing magnification of 2.4, a redshift of $`z=3.5`$, and a cosmology of $`H_o=50\mathrm{km}\mathrm{s}^1`$ Mpc<sup>-1</sup> and $`q_o=0.5`$ in order to estimate the properties of SMM J00266+1708. Given the redshift uncertainty, we cannot accurately determine many of the intrinsic properties of SMM J00266+1708. One exception is the bolometric sub-mm/far-infrared (FIR) luminosity which is relatively insensitive to redshift for $`z1`$ (Blain & Longair 1993). From the 850$`\mu `$m measurement, the implied intrinsic FIR luminosity for SMM J00266+1708 is $`L(\mathrm{FIR})10^{13}L_{\mathrm{}}`$. This corresponds to a star-formation rate of massive stars ($`M>5M_{\mathrm{}}`$) of approximately 500–900 $`M_{\mathrm{}}`$ yr<sup>-1</sup> (Scoville et al. 1997; Condon 1992). These results are consistent with estimates for other sub-mm galaxies (Ivison et al. 1998, 2000a). It is still not known whether or not an AGN contributes significantly to the bolometric luminosity of SMM J00266+1708. If an AGN is present, it is not a strong radio source and must be heavily obscured at optical wavelengths (Fig. 3). Based on its SED, SMM J00266+1708 appears most similar to highly-reddened ULIG/starbursts such as Arp 220 or HR10 at a redshift of $`z3.5`$.
### 4.2. Comparison of Sub-mm and Lyman-break Populations
The relative importance of the sub-mm and Lyman-break galaxies to the global star-formation rate at high-redshift is an active area of discussion (Hughes et al. 1998; Guiderdoni et al. 1998; Blain et al. 1999b, 1999c; Trentham et al. 1999; Peacock et al. 2000). Depending on the exact contribution of AGN in the sub-mm population, the ultraluminous sub-mm galaxies ($`>10^{12}L_{\mathrm{}}`$; $`S[850\mu `$m$`]1`$ mJy) account for approximately 30–50% of the total amount of star-formation at high-redshift (Blain et al. 1999b). It is truly remarkable that such a high fraction of all star-formation at high-redshift occurs in ULIG/sub-mm galaxies. In contrast, ULIGs only contribute about 0.2% of the total amount of star formation in the local universe, based on the survey of Kim & Sanders (1998). Hence, the evolution in the amount of star formation occurring in ULIGs is about 100 times stronger than the global increase of the star-formation rate at high redshift seen for all galaxies (Madau et al. 1996; Steidel et al. 1999).
It has been suggested for both the sub-mm and Lyman-break populations that each represents the formative phases of massive $`L^{}`$ galaxies (Blain et al. 1999c; Giavalisco et al. 1998). In this scenario, the sub-mm systems may be associated with a very luminous, short-lived and heavily dust enshrouded starburst, while the Lyman-break galaxies would be associated with a longer-lived, less luminous phase of star formation (Blain et al. 1999c). If this scenario is correct, we could expect massive reservoirs of molecular gas associated with both populations. The detection of massive reservoirs of molecular gas (Frayer et al. 1998, 1999) suggests that the sub-mm population is associated with gas-rich massive galaxies ($`L^{}`$). The molecular gas masses of the sub-mm galaxies are 10–50 times greater than that found for the Milky Way. In contrast, the best studied Lyman-break galaxy cB58 (Yee et al. 1996; Pettini et al. 2000; Teplitz et al. 2000) has even less molecular gas than the Milky Way, after correcting the observed CO upper-limit for lensing and suspected metallicity effects (Frayer et al. 1997). Therefore, the Lyman-break sources may be sub-$`L^{}`$ systems representing the building blocks of more massive galaxies, as suggested by Lowenthal et al. (1997).
Most sub-mm galaxies are similar to SMM J00266+1708 in being too faint and/or too red to be included in the Lyman-break surveys (e.g., Steidel et al. 1999). The one notable exception is the sub-mm selected galaxy pair, SMM J14011+0252 J1/J2 (Ivison et al. 2000a; Adelberger & Steidel 2000), but for this system most of the blue light arises from J2, while the bulk of the bolometric luminosity is thought to be due to the much redder J1 component (Ivison et al. 2001). Due to the high levels of dust obscuration for the sub-mm galaxies, there is very little overlap between the sub-mm selected galaxies and optically-selected Lyman-break systems. Even for the Lyman-break galaxies, most of their star-formation ($``$80%) is obscured by dust at observed optical wavelengths (Peacock et al. 2000; Adelberger & Steidel 2000). Hence, it is clear that most of the star-formation activity at high-redshift is hidden from view at optical/ultraviolet wavelengths.
The results presented for SMM J00266+1708 have important implications on our general understanding of star-formation at high redshift. Roughly half of the total amount of star formation at high redshift is thought to occur in the sub-mm galaxies (Blain et al. 1999b), while the other half is inferred from the optically–selected Lyman-break sources (Peacock et al. 2000). Since about 50% of the sub-mm galaxies are undetected at optical wavelengths, a significant fraction ($`25`$%) of the total amount of high-redshift star formation may occur in the very faint/red sub-mm galaxies similar to SMM J00266+1708. Similar conclusions have also been reached for a radio-selected sample of sub-mm galaxies with faint optical counter-parts (Barger et al. 2000). The fact that many of the most luminous high-redshift starbursts/AGN are too faint to be studied at optical wavelengths highlights the importance that future sensitive mm–NIR wavelength instruments will have on our understanding of the distant universe.
## 5. CONCLUDING REMARKS
We report the identification of the sub-mm source SMM J00266+1708 with a faint red galaxy ($`K=22.5`$ mag) which is undetected at optical wavelengths, despite very deep observations. This source has an extremely high luminosity of approximately $`10^{13}L_{\mathrm{}}`$ even after correcting for lensing. The current data for the sub-mm Cluster Lens Survey suggest that 40%–70% of the sub-mm population as a whole are faint/red galaxies which are undetected at optical wavelengths. These faint/red sub-mm galaxies are thought to contribute significantly to the total amount of star formation at high redshift and are hence important to our understanding of the early evolution of galaxies.
The redshift of SMM J00266+1708 is currently unknown, but the galaxy is expected to be at a redshift $`z>2`$. Obtaining a redshift will be extremely challenging with current instrumentation. At $`K=22.5`$ mag, the galaxy pushes the capabilities of even the largest ground based telescopes. Since the galaxy is relatively red ($`IK>3.6`$), we expect H$`\alpha `$ to be the brightest optical emission line, and perhaps the only optical line currently detectable based on comparisons with the ERO HR10 (Dey et al. 1999). If the redshift is similar to that estimated from the SED of the galaxy ($`z3.5`$), H$`\alpha `$ would be shifted redward of $`K`$-band, making ground based observations extremely challenging.
We could expect much fainter $`K`$band magnitudes of 23–26 for similar ULIG/sub-mm galaxies ($`10^{12}L_{\mathrm{}}`$) which are unlensed. Therefore, obtaining optical/NIR redshifts for many of the sub-mm galaxies may have to wait for the Next Generation Space Telescope. Alternatively, redshifts could be directly measured from the CO lines themselves with future ground-based instruments, such as the Atacama Large Millimeter Array (ALMA), operating at mm-wavelengths (Blain et al. 2000). Sensitive interferometric observations at sub-mm/mm-wavelengths with the next generation of instruments will be crucial for our understanding of these dust-obscured systems.
We thank the staff at the Owens Valley Millimeter Array and the Keck Observatory who have made these observations possible. We thank Andrew Blain and Jean-Paul Kneib for useful discussion and their work on the SCUBA Lens Survey Sample. We thank Frazer Owen, Glenn Morrison, and Asantha Cooray for their work on the radio data. We thank Dave Thompson for his set of IRAF tasks called NIRCtools which were used to help reduce the NIR data.
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# I Introduction
## I Introduction
In mesonic semileptonic $`bc`$ transitions, the exclusive transitions to the ground state $`S`$-wave mesons $`BD,D^{}`$ make up approximately $`66\%`$ of the total semileptonic $`BX_c`$ rate . It would then be interesting to know what the corresponding semileptonic rate ratio $`\mathrm{\Gamma }_{\mathrm{\Lambda }_b\mathrm{\Lambda }_c}/\mathrm{\Gamma }_{\mathrm{\Lambda }_bX_c}`$ (termed $`R_E`$ in the following) is in semileptonic $`\mathrm{\Lambda }_b`$-decays. This is an important experimental issue since a knowledge of this ratio would greatly facilitate the analysis of semileptonic $`\mathrm{\Lambda }_b`$-decays. For example, if the semileptonic $`\mathrm{\Lambda }_b`$-decays were dominated by the quasi-elastic exclusive channel $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c+l^{}+\overline{\nu }_l`$ this would be of considerable help in the kinematical reconstruction of their decays in as much as the $`\mathrm{\Lambda }_c`$ baryon is easy to detect via its decay mode $`\mathrm{\Lambda }_cpK^{}\pi ^+`$. Unfortunately nothing is known experimentally about this ratio yet.
In this paper we attempt to address the problem of a determining the exclusive/inclusive ratio $`R_E`$ in semileptonic $`\mathrm{\Lambda }_b`$-decays from a theoretical point of view by consulting some model calculations which we critically scrutinize. We also attempt to extrapolate from the experimentally known results in the meson sector to the baryon sector.
As concerns the inclusive semileptonic rates of bottom mesons and bottom baryons one is now reasonably confident that they can be reliably calculated using the usual operator product expansion within HQET. The leading term in the OPE is given by the free heavy quark decay rate which clearly is the same for baryons and mesons. Radiative corrections to the free quark decay rate are quite large but again are identical for mesons and baryons. Differences in the inclusive semileptonic rates of mesons and the $`\mathrm{\Lambda }_b`$ baryon set in only at $`𝒪(1/m_b^2)`$. They affect the mesonic and $`\mathrm{\Lambda }_b`$ rates differently since there is no chromomagnetic $`𝒪(1/m_b^2)`$ correction in the $`\mathrm{\Lambda }_b`$ case. However, since the chromomagnetic term contributes only at the $`3.7\%`$ level, the difference in the inclusive semileptonic rates for mesons and baryons is predicted to be quite small.
A much more difficult task is to get a reliable theoretical handle on the quasi-elastic exclusive semileptonic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ rate. There exist a number of theoretical calculations on the exclusive decay $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c+l^{}+\overline{\nu }_l`$ using various model assumptions. They are of no great help since their predicted rate values may differ by factors of up to three and it is not easy to judge the reliability of the various model assumptions that enter the calculation. Ideally one would like to have model calculations that are valid both in the heavy meson and the heavy baryon sector. If these model calculations give sensible results in the heavy meson sector, where they can be checked against data, one would have more confidence in their predictions for the heavy baryon sector.
The paper is structured as follows. In Sec.II we take a first look at the leading order rate formula for the exclusive semileptonic decays of B mesons and $`\mathrm{\Lambda }_b`$ baryons in order to get a semiquantitative handle on the relative size of their rates. The analysis is refined in Sec.III for $`\mathrm{\Lambda }_b`$ baryons where $`1/m_b`$ and $`1/m_b^2`$ corrections and renormalization effects are included. In Sec.IV we recapitulate the calculation of the semileptonic inclusive decay rates. The results of the Sections III and IV are brought together in Sec.V where we discuss the exclusive/inclusive ratio $`R_E`$ of semileptonic $`\mathrm{\Lambda }_b`$ decays. We present numerical results on the exclusive/inclusive ratio for various models and give our best estimate of this ratio. In Sec.VI we classify the possible non-exclusive final states that will have to fill the gap between the exclusive and inclusive rates in semileptonic $`\mathrm{\Lambda }_b`$ decays. Sec.VII, finally, contains our conclusions.
## II The heavy quark limit
For a quick first appraisal of the question of how the exclusive semileptonic decays of mesons and baryons are related we turn to the heavy quark limit and list the leading order semileptonic rate formulae for the $`B(D+D^{})`$ and $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ transitions. One has
$$\frac{d\mathrm{\Gamma }\left\{\genfrac{}{}{0pt}{}{\mathrm{meson}}{\mathrm{baryon}}\right\}}{d\omega }=\frac{G_F^2|V_{bc}|^2M_1^5}{12\pi ^3}r^3\sqrt{\omega ^21}\left(3\omega (1+r^2)2r(2\omega ^2+1)\right)\left\{\genfrac{}{}{0pt}{}{\frac{\omega +1}{2}|F_{\mathrm{meson}}(\omega )|^2}{|F_{\mathrm{baryon}}(\omega )|^2}\right\},$$
(1)
where $`r=M_2/M_1`$ and $`\omega =(M_1^2+M_2^2q^2)/2M_1M_2`$. $`F_{\mathrm{meson}}(\omega )`$ and $`F_{\mathrm{baryon}}(\omega )`$ are the leading order Isgur-Wise transition form factors for the $`BD`$,$`D^{}`$ and $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ transitions, respectively. Throughout the paper we refer to $`M_1`$ and $`M_2`$ as the masses of the initial and final particles in the semileptonic decay process.
The free heavy quark decay rate (or leading order parton model rate) which we need later on is simply obtained by replacing the particle masses in (1) by the corresponding quark masses and setting the curly bracket in (1) to one, i.e. by taking the current coupling in the $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ case to be point-like. We shall encounter the integrated parton model rate for the $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ case again in Sec.IV. Finally, in the heavy quark limit one has to determine the final meson mass $`M_2`$ by taking the weighted average $`\overline{M}_D=1/4(M_D+3M_D^{})=1.973\mathrm{GeV}`$ with $`M_D=1.869\mathrm{GeV}`$ and $`M_D^{}=2.010\mathrm{GeV}`$. For the pseudoscalar bottom mass we take $`M_B=5.279\mathrm{GeV}`$. For the $`\mathrm{\Lambda }_Q`$-baryon masses we use $`M_{\mathrm{\Lambda }_b}=5.624\mathrm{GeV}`$ and $`M_{\mathrm{\Lambda }_c}=2.285\mathrm{GeV}`$.
When trying to compare the two rates in (1) one identifies two main determining factors which counteract each other. On the one hand one has the form factor expressions in the curly bracket which tend to enhance the mesonic rate due to the multiplicative factor $`(\omega +1)/2`$ in the mesonic case. Also, according to common prejudice the baryon form factor falls off more rapidly than the mesonic form factor. On the other hand one has the overall mass factor $`M_1^5r^3`$ which enhances the baryonic rate because $`M_1^5r^3=214.16\mathrm{GeV}^5`$ and $`M_1^5r^3=377.37\mathrm{GeV}^5`$ in the mesonic and baryonic cases, respectively.
It is evident that the choice of mesonic and baryonic Isgur-Wise functions plays a crucial role when comparing the two rates. As has been emphasized before, there exist some experimental knowledge on the mesonic Isgur-Wise function but nothing is known experimentally about the baryonic Isgur-Wise function yet <sup>3</sup><sup>3</sup>3The only available experimental result is from a preprint version of a DELPHI-analysis . This paper quotes a value of $`\rho ^2=1.81_{0.67}^{+0.70}\pm 0.32`$ for the slope of the baryonic Isgur-Wise function. However, since this paper has never been published, we shall not use their result in our analysis..
For quick reference it is sometimes convenient to characterize the fall-off behaviour of the Isgur-Wise functions by expanding it around the zero recoil point where one has the zero recoil normalization condition $`F(1)=1`$. Keeping terms up to second order in this expansion one has
$$F(\omega )=F(1)[1\rho ^2(\omega 1)+c(\omega 1)^2+\mathrm{}],$$
(2)
where the coefficients $`\rho ^2`$ and $`c`$ are called the slope parameter and the convexity parameter, respectively. The slope parameter is frequently used to characterize the fall-off behaviour of the Isgur-Wise function. The expansion (2) is useful if one studies the physics close to zero threshold but may give misleading results when calculating rates because the spectral weight function multiplying the form factor functions is essentially determined by the square root factor $`\sqrt{\omega ^21}`$ in (1) and is therefore strongly weighted towards the end of the spectrum. It goes without saying that the slope and convexity parameters are in general different for the mesonic and baryonic Isgur-Wise functions.
In order to proceed with our first appraisal of the magnitude of the exclusive mesonic and baryonic semileptonic rates we appeal to the spectator quark model where the mesonic and baryonic Isgur-Wise functions become related to one another . In the spectator quark model one has
$$F_{\mathrm{baryon}}(\omega )=\frac{\omega +1}{2}|F_{\mathrm{meson}}(\omega )|^2.$$
(3)
Explicit calculations show that the baryonic form factor is considerably underestimated by the spectator relation (3). Nevertheless, the spectator relation (3) may still serve as an effective lower bound on the baryonic form factor.
The physical picture behind the spectator quark model relation is quite simple. In the heavy baryon case there are two light spectator quarks that need to be accelerated in the current transition compared to the one spectator quark in the heavy meson transition. Thus the baryonic form factor is determined in terms of the square of the mesonic form factor. The factor $`(\frac{\omega +1}{2})`$ is a relativistic factor which insures the correct threshold behaviour of the baryonic form factor in the crossed $`e^+e^{}`$-channel .
In the relation between heavy meson and heavy baryon form factors was investigated in the context of a dynamical Bethe-Salpeter (BS) model. The above spectator quark model relation (3) in fact emerges when the interaction between the light quarks in the heavy baryon is switched off in the BS-interaction kernel. In the more realistic situation when the light quarks interact with each other, the heavy baryon form factor becomes flatter, i.e. the spectator quark model form factor may be used to bound the heavy baryon form factor from below. In Fig.1 we reproduce from the $`\omega `$ dependence of the spectator quark model form factor and that of two representative form factors with the interaction between the light quarks included. The starting point in is a mesonic form factor with a slope of $`\rho _{\mathrm{meson}}^2=1`$ which, according to (3), leads to a spectator form factor with a slope of $`\rho _{\mathrm{baryon}}^2=1.5`$. The spectator form factor is the lowest form factor shown in Fig.1. The interaction between the light quarks was introduced through a harmonic oscillator type kernel in the BS-equation. The two upper form factor curves in Fig.1 correspond to two different choices of the oscillator strength with which the light quarks interact or, equivalently, correspond to two different choices of the size parameter in the oscillator wave function. The interaction type form factors in Fig.1 have slopes of $`\rho _{\mathrm{baryon}}^2=0.81`$ (solid line) and 0.97 (short-dashed line) . They are considerably flatter than the spectator quark model form factor.
We shall now calculate exclusive rates for mesonic and baryonic transitions according to (1) using the spectator model relation (3). According to what was said before, the baryonic rate calculated in this way must subsequently be adjusted upward according to the analysis of . We went to considerable lengths in explaining the results of because we want to emphasize that the outcome of the baryonic rate estimate using the spectator quark model relation (3) must be viewed as providing only lower bounds on the true quasi-elastic baryonic rate. For the mesonic form factor we use the world average of the slope $`\rho _{\mathrm{meson}}^2`$ obtained by combining results from $`BD^{}`$ and $`BD`$ transitions, $`\rho _{\mathrm{meson}}^2=0.70`$ <sup>4</sup><sup>4</sup>4The CLEO Coll. also attempted a linear plus quadratic fit to the data, but the data was not good enough to determine the convexity parameter $`c`$ of the meson form factor with any accuracy.. Using $`V_{cb}=0.038`$, a linear meson form factor with the above slope, a baryonic form factor according to the spectator relation (3) and the rate formulae (1) one obtains $`\mathrm{\Gamma }_{\mathrm{meson}}=5.3010^{10}s^1`$ and $`\mathrm{\Gamma }_{\mathrm{baryon}}=5.0410^{10}s^1`$. As has been emphasized before the baryonic rate $`\mathrm{\Gamma }_{\mathrm{baryon}}=5.0410^{10}s^1`$ has to be adjusted upward in the more realistic situation of interacting light quarks. Looking at the model calculation for guidance, the increment in rate going from noninteracting (spectator) to interacting quarks is 1.28 and 1.37, respectively, for the two choices of oscillator strengths analyzed in . Adjusting the above baryonic rate accordingly our leading order estimate of the baryonic rate is thus $`\mathrm{\Gamma }_{\mathrm{baryon}}=(6.45÷6.90)10^{10}s^1`$. Starting from a mesonic exclusive/inclusive ratio of $`66\%`$ and assuming equal inclusive semileptonic rates for bottom baryons and mesons, which is sufficiently accurate for our semiquantitative calculation, our estimate for the exclusive/inclusive ratio in semileptonic $`\mathrm{\Lambda }_b`$ decays is $`R_E=(8086)\%`$. This is considerably larger than the mesonic exclusive/inclusive ratio $`R_E66\%`$.
Up to this point our semiquantitive analysis was done to leading order in HQET. How would finite mass effects affect our previous conclusions? One way of improving the previous analysis in the meson sector is to insert physical masses in the rate expression (1), thereby including part of the $`1/m_Q`$-corrections to (1). To do this we need to disentangle the $`BD`$ and $`BD^{}`$ rates in (1). One has :
$$\frac{d\mathrm{\Gamma }\left(BD\right)}{d\omega }=\frac{G_F^2|V_{bc}|^2M_1^5}{48\pi ^3}r^3(1+r)^2(\omega ^21)^{3/2}|F_{\mathrm{meson}}(\omega )|^2$$
(4)
and
$$\frac{d\mathrm{\Gamma }\left(BD^{}\right)}{d\omega }=\frac{G_F^2|V_{bc}|^2M_1^5}{48\pi ^3}r^3\sqrt{\omega ^21}(\omega +1)\left[(1r)^2(\omega +1)+4\omega (12\omega r+r^2)\right]|F_{\mathrm{meson}}(\omega )|^2.$$
(5)
Using now the world average of $`(\rho _{\mathrm{meson}}^2)^{BD}=0.66`$ and $`(\rho _{\mathrm{meson}}^2)^{BD^{}}=0.71`$ with linear form factors and taking physical $`D`$ and $`D^{}`$ masses one finds $`\mathrm{\Gamma }_{BD}=1.3910^{10}\mathrm{s}^1`$ and $`\mathrm{\Gamma }_{BD^{}}=3.9010^{10}\mathrm{s}^1`$ giving a total mesonic rate of $`\mathrm{\Gamma }_{BD+D^{}}=5.3010^{10}\mathrm{s}^1`$. It fully agrees with the above result and therefore leaves the aforegoing conclusions intact.
Continuing with our discussion on the contributions of nonleading effects in the $`1/m_Q`$-expansion we now turn to results of some model calculations in order to find out how nonleading effects may affect the above conclusions. In the mesonic sector Neubert and Rieckert analyzed an infinite momentum frame model and found that $`𝒪(1/m_Q)`$ effects raise the $`BD`$ and $`BD^{}`$ rates by $`15.7\%`$ and $`0.5\%`$, respectively, resulting in a rise of $`4.4\%`$ for the total $`D+D^{}`$ rate. Using a similar infinite momentum frame model König et al. find that the $`𝒪(1/m_Q)`$ effects raise the semileptonic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ rate by $`3\%`$ which is quite close to the $`4.4\%`$ found in in the bottom meson case. Judging from these model calculations our leading order comparison of the mesonic and baryonic rates and the conclusions drawn from it do not seem to be much affected by $`𝒪(1/m_Q)`$ corrections.
There also exist estimates of $`𝒪(1/m_Q^2)`$ corrections in the literature. Faustov and Galkin et al. use a relativistic quark model based on the quasipotential approach . They quote exclusive/inclusive branching ratios of $`(13.5+3.31.4)\%`$ and $`(39.1+6.53.9)\%`$ for semileptonic $`BD`$ and $`BD^{}`$ rates, where the second and third numbers refer to the $`𝒪(1/m_Q)`$ and $`𝒪(1/m_Q^2)`$ corrections, respectively. The $`𝒪(1/m_Q)`$ corrections in this model are considerably larger than in the infinite momentum frame models. Ivanov et al. investigated the role of finite mass effects in semileptonic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ decays without taking recourse to the heavy mass expansion. They found an overall rate reduction of $`9\%`$ relative to the infinite mass result . In the analysis of the present paper presented in Sec.III we obtain $`+5\%`$ and $`7\%`$ for the $`𝒪(1/m_Q)`$ and $`𝒪(1/m_Q^2)`$ corrections to the $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ decays, respectively. From all these model calculations one learns that the $`𝒪(1/m_Q)`$ corrections tend to increase the rates whereas the $`𝒪(1/m_Q^2)`$ corrections tend to decrease the rates, for both heavy meson and heavy baryon decays. Again, our leading order estimate of the relative size of the exclusive/inclusive ratios of bottom mesons and baryons is not likely to be affected much by including also $`𝒪(1/m_Q^2)`$ effects. The same holds true for renormalization effects of the weak current which affect the bottom baryon and bottom meson amplitudes equally and therefore drop out in the ratio of exclusive semileptonic bottom meson and baryon decays.
The conclusion drawn in this section on the predominance of the exclusive/inclusive ratio of semileptonic $`\mathrm{\Lambda }_b`$-decays over that of semileptonic B-decays carries over to the more sophisticated analysis of the next sections where we include $`1/m_Q`$ and $`1/m_Q^2`$ effects, and radiative corrections.
## III Exclusive Semileptonic Rate $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c+l^{}+\overline{\nu }_l`$
It is most convenient to represent the differential decay rate in terms of the helicity amplitudes of the process. One has (we take leptons to be massless)
$$\frac{d\mathrm{\Gamma }\left(\mathrm{\Lambda }_b\mathrm{\Lambda }_c\right)}{d\omega }=\frac{G_F^2|V_{bc}|^2}{96\pi ^3}\frac{q^2M_2^2\sqrt{\omega ^21}}{M_1}\left(|H_{\frac{1}{2}1}|^2+|H_{\frac{1}{2}1}|^2+|H_{\frac{1}{2}0}|^2+|H_{\frac{1}{2}0}|^2\right).$$
(6)
The helicity amplitudes are in turn related to the invariant amplitudes of the process via
$`\sqrt{q^2}H_{\frac{1}{2}0}^{V,A}`$ $`=`$ $`\sqrt{2M_1M_2(\omega 1)}\left((M_1\pm M_2)f_1^{V,A}\pm M_2(\omega \pm 1)f_2^{V,A}\pm M_1(\omega \pm 1)f_3^{V,A}\right),`$ (7)
$`H_{\frac{1}{2}1}^{V,A}`$ $`=`$ $`2\sqrt{M_1M_2(\omega 1)}f_1^{V,A},`$ (8)
where the invariant amplitudes are defined by
$`<\mathrm{\Lambda }_c(v_2)|J_\mu ^V|\mathrm{\Lambda }_b(v_1)>`$ $`=`$ $`\overline{u}_c(v_2)(f_1^V\gamma _\mu +f_2^Vv_{1\mu }+f_3^Vv_{2\mu })u_b(v_1),`$ (9)
$`<\mathrm{\Lambda }_c(v_2)|J_\mu ^A|\mathrm{\Lambda }_b(v_1)>`$ $`=`$ $`\overline{u}_c(v_2)(f_1^A\gamma _\mu +f_2^Av_{1\mu }+f_3^Av_{2\mu })\gamma _5u_b(v_1).`$ (10)
The total helicity amplitudes finally are given by
$$H_{\lambda _2\lambda _W}=H_{\lambda _2\lambda _W}^VH_{\lambda _2\lambda _W}^A,$$
(11)
where the choice of of the relative minus sign between the vector and axial vector helicity amplitudes reflects the $`(VA)`$ structure of the $`bc`$ current transition. The $`H_{\lambda _2\lambda _W}^{V,A}`$ are the helicity amplitudes for the vector ($`V`$) and axial vector ($`A`$) current induced transition in the decay $`1/2^+1/2^++W_{offshell}^{}`$ with $`\lambda _2`$ and $`\lambda _W`$ being the helicities of the final state baryon and the $`W`$ boson, respectively.
The remaining helicity amplitudes are related to the above helicity amplitudes (8) by parity. One has
$$H_{\lambda _2\lambda _W}^{V,A}=\pm H_{\lambda _2\lambda _W}^{V,A}.$$
(12)
It is well known that the complexity of the form factor structure exemplified by the set of six form factors $`f_i^{V,A}(i=1,2,3)`$ is considerably reduced in HQET. Working up to $`𝒪(1/m_Q)`$ in HQET and including also $`𝒪(\alpha _s)`$ corrections one finds
$`f_1^V(\omega )`$ $`=`$ $`F(\omega )+\left({\displaystyle \frac{1}{2M_1}}+{\displaystyle \frac{1}{2M_2}}\right)\left(\eta (\omega )+\overline{\mathrm{\Lambda }}F(\omega )\right)+{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}v_1(\omega ,\lambda )F(\omega ),`$ (13)
$`f_2^V(\omega )`$ $`=`$ $`F(\omega )\left({\displaystyle \frac{1}{M_2}}{\displaystyle \frac{1}{\omega +1}}\overline{\mathrm{\Lambda }}{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}v_2(\omega )\right),`$ (14)
$`f_3^V(\omega )`$ $`=`$ $`F(\omega )\left({\displaystyle \frac{1}{M_1}}{\displaystyle \frac{1}{\omega +1}}\overline{\mathrm{\Lambda }}{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}v_3(\omega )\right),`$ (15)
$`f_1^A(\omega )`$ $`=`$ $`F(\omega )+\left({\displaystyle \frac{1}{2M_1}}+{\displaystyle \frac{1}{2M_2}}\right)\left(\eta (\omega )+\overline{\mathrm{\Lambda }}F(\omega ){\displaystyle \frac{\omega 1}{\omega +1}}\right)+{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}a_1(\omega ,\lambda )F(\omega ),`$ (16)
$`f_2^A(\omega )`$ $`=`$ $`F(\omega )\left({\displaystyle \frac{1}{M_2}}{\displaystyle \frac{1}{\omega +1}}\overline{\mathrm{\Lambda }}{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}a_2(\omega )\right),`$ (17)
$`f_3^A(\omega )`$ $`=`$ $`F(\omega )\left({\displaystyle \frac{1}{M_1}}{\displaystyle \frac{1}{\omega +1}}\overline{\mathrm{\Lambda }}{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}a_3(\omega )\right).`$ (18)
The $`𝒪(\alpha _s)`$ corrections to the form factors have been taken from . They result from the $`𝒪(\alpha _s)`$ vertex correction to the current-induced $`bc`$ transition . The infrared singularity is regularized by the introduction of a fictitious gluon mass which is taken to be $`\lambda =0.2\mathrm{GeV}`$. At zero recoil, where the vertex correction is infrared finite, the renormalization is independent of the gluon mass regulator. However, away from zero recoil, the $`\alpha _s`$-correction functions $`v_1(\omega ,\lambda )`$ and $`a_1(\omega ,\lambda )`$ depend on the gluon mass regulator. This introduces a certain amount of model dependence in the renormalization procedure. The above value of the gluon mass was chosen according to the expectation that the exchange of virtual gluons in the vertex correction should be cut-off at frequencies $`k^01/R`$ with $`R1`$ fm being a typical hadronic scale. The argument of the $`\alpha _s`$ coupling, $`\overline{m}`$, is taken such that effects of higher order terms $`(\alpha _s\mathrm{ln}(m_b/m_c))^n`$ are minimized, $`\overline{m}=2m_bm_c/(m_b+m_c)2.31\mathrm{GeV}`$. In this way one avoids the use of the renormalization group improved summation of leading logarithms, which has been proven as inconsistent .
The binding energy of the $`\mathrm{\Lambda }_b`$ is denoted by $`\overline{\mathrm{\Lambda }}`$ which we take to be 0.6 GeV. The form factor function $`\eta (x)`$ results from the nonlocal contribution of the kinetic energy term of the $`1/m_Q`$ corrected HQET Lagrangian. It has been calculated in two different model approaches and has found to be negligibly small . Therefore, we can safely drop its contribution in the following. By neglecting the form factor $`\eta (x)`$ in the $`𝒪(1/m_Q)`$ result (18) the differential rate (6) is proportional to the square of the leading order Isgur-Wise function $`F(\omega )`$ with its zero recoil normalization $`F(1)=1`$. In this way we can meaningfully compare our results with the leading order results of other model calculations as will be done in Sec.V.
It is well known that $`𝒪(1/m_Q^2)`$ corrections to the unit zero recoil normalization of Isgur-Wise functions can be substantial. For example, by evaluating zero recoil sum rules, the authors of Refs. obtain
$`F(1)_{BD}=0.98\pm 0.07,`$ (19)
$`F(1)_{BD^{}}=0.91\pm 0.03.`$ (20)
in the mesonic case. In the $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ case the zero recoil sum rule gives a bound on the zero recoil value of the sole remaining form factor function $`f_1^A`$. The (unrenormalized) zero recoil sum rule leads to $`f_1^A(1)(10.165\mu _\pi ^2/\mathrm{GeV}^2)^{1/2}`$ . Using $`\mu _\pi ^2=0.5\mathrm{GeV}^2`$ , $`f_1^A(1)`$ must be smaller than 0.958. For definiteness we take a value close to the upper bound
$$f_1^A(1)=0.95.$$
(21)
This value is nicely corroborated by the finite heavy quark mass calculation of where one finds $`f_1^A(1)=0.97`$.
Nothing is known about the size of the $`𝒪(1/m_Q^2)`$ corrections to $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ away from zero recoil, except that they can be parametrized in terms of ten new $`\omega `$-dependent form factors and one new dimensionful constant , the magnitude and functional forms of which are not known. The lack of knowledge about the $`𝒪(1/m_Q^2)`$ corrections away from zero recoil prevents us from their exact treatment. On the other hand the size of the $`𝒪(1/m_Q^2)`$ correction at zero recoil is a clear indication that the $`𝒪(1/m_Q^2)`$ corrections cannot be neglected. We shall therefore adopt the following strategy. We smoothly extrapolate from the $`𝒪(1/m_Q^2)`$ information at zero recoil to the whole $`\omega `$-range. The appropiate amplitude for this extrapolation is the axial vector current $`S`$-wave amplitude the zero recoil value of which is determined by the zero recoil sum rules. We thus multiply the axial vector current $`S`$-wave amplitude everywhere by its zero recoil value $`f_1^A(1)=0.95`$. It is clear that the $`𝒪(1/m_Q^2)`$ corrections at zero recoil are exactly included in this approach. The lack of knowledge about the $`𝒪(1/m_Q^2)`$ corrections to the other partial wave amplitudes leaves us no choice but to leave them untreated. With this in mind it is gratifying to note that the $`S`$-wave contribution dominates the quasielastic rate. For example, using the standard form factor (47) with $`\rho _B^2=0.75`$ in a leading order calculation one finds that the $`S`$-wave contribution amounts to $``$ 66% of the total semileptonic rate.
In order to set up our procedure of how to incorporate the $`𝒪(1/m_Q^2)`$ corrections we define the relevant vector current ($`V`$) and axial vector current ($`A`$) partial wave amplitudes $`A_{LS}^{V,A}`$ in terms of the helicity amplitudes $`H_{\lambda _2\lambda _W}^{V,A}`$. $`L`$ denotes the orbital angular momentum of the final state and $`S=J_{\mathrm{current}}+S_{\mathrm{\Lambda }_c}`$ is the sum of the final state spin angular momenta where $`J_{\mathrm{current}}=1`$ in the zero lepton mass case that we are considering here. One has
$`A_{1\frac{1}{2}}^V`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}H_{\frac{1}{2}0}^V\sqrt{{\displaystyle \frac{4}{3}}}H_{\frac{1}{2}1}^V,`$ (22)
$`A_{2\frac{3}{2}}^V`$ $`=`$ $`\sqrt{{\displaystyle \frac{4}{3}}}H_{\frac{1}{2}0}^V+\sqrt{{\displaystyle \frac{2}{3}}}H_{\frac{1}{2}1}^V,`$ (23)
$`A_{0\frac{1}{2}}^A`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}H_{\frac{1}{2}0}^A\sqrt{{\displaystyle \frac{4}{3}}}H_{\frac{1}{2}1}^A,`$ (24)
$`A_{2\frac{3}{2}}^A`$ $`=`$ $`\sqrt{{\displaystyle \frac{4}{3}}}H_{\frac{1}{2}0}^A+\sqrt{{\displaystyle \frac{2}{3}}}H_{\frac{1}{2}1}^A.`$ (25)
By substituting invariant form factors according to (8) one can verify the correct threshold behaviour of the partial wave amplitudes, i.e. $`A_{1\frac{1}{2}}^V,A_{1\frac{3}{2}}^V(\omega 1)^{1/2}`$, $`A_{0\frac{1}{2}}^A(\omega 1)^0`$ and $`A_{2\frac{3}{2}}^A(\omega 1)^1`$.
According to the above strategy we now incorporate the $`𝒪(1/m_Q^2)`$ corrections by multiplying the $`S`$-wave amplitude $`A_{0\frac{1}{2}}^A`$ by the $`𝒪(1/m_Q^2)`$ zero recoil correction $`f_1^A(1)=0.95`$. Thus we write
$$A_{0\frac{1}{2}}^Af_1^A(1)A_{0\frac{1}{2}}^A=A_{0\frac{1}{2}}^A+(f_1^A(1)1)A_{0\frac{1}{2}}^A$$
(26)
For the first term on the $`r.h.s.`$ of (26) we substitute the $`𝒪(1/m_Q)`$ result according to (18). Contrary to this we use only the leading order result for $`A_{0\frac{1}{2}}^A`$ in the second term of (26) since it is already being multiplied by the $`𝒪(1/m_Q^2)`$ factor $`(f_1^A(1)1)`$. Including also the $`A_{2\frac{3}{2}}^A`$ partial wave amplitude the leading order expressions for the axial vector partial wave amplitudes read
$`A_{0\frac{1}{2}}^A`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{q^2}}}{\displaystyle \frac{2}{\sqrt{3}}}\sqrt{M_1M_2(\omega +1)}F(\omega )[(M_1M_2+2\sqrt{q^2})(1+{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}a_1(\omega ,\lambda ))`$ (28)
$`{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}(\omega 1)(M_2a_2(\omega )+M_1a_3(\omega ))]`$
$`A_{2\frac{3}{2}}^A`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{q^2}}}{\displaystyle \frac{2\sqrt{2}}{\sqrt{3}}}\sqrt{M_1M_2(\omega +1)}F(\omega )[(M_1M_2\sqrt{q^2})(1+{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}a_1(\omega ,\lambda ))`$ (30)
$`{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}(\omega 1)(M_2a_2(\omega )+M_1a_3(\omega ))]`$
Putting everything together we arrive at the differential rate. One obtains
$`{\displaystyle \frac{d\mathrm{\Gamma }\left(\mathrm{\Lambda }_b\mathrm{\Lambda }_c\right)}{d\omega }}`$ $`=`$ $`{\displaystyle \frac{G_F^2|V_{bc}|^2}{48\pi ^3}}{\displaystyle \frac{q^2M_2^2\sqrt{\omega ^21}}{M_1}}\{|H_{\frac{1}{2}1}^V|^2+|H_{\frac{1}{2}0}^V|^2+|H_{\frac{1}{2}1}^A|^2+|H_{\frac{1}{2}0}^A|^2`$ (33)
$`+\left((f_1^A(1))^21\right){\displaystyle \frac{2}{3}}{\displaystyle \frac{M_1M_2}{q^2}}(\omega +1)F^2(\omega )(M_1M_2+2\sqrt{q^2})`$
$`[(M_1M_2+2\sqrt{q^2})(1+2{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}a_1(\omega ,\lambda ))2{\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}(\omega 1)(M_2a_2(\omega )+M_1a_3(\omega ))]\}.`$
In Eq.(33) the first line incorporates the $`𝒪(1)`$ and $`𝒪(1/m_Q)`$ contributions including the radiative corrections as specified by (8) and by (18). The second and third lines comprise the $`𝒪(1/m_Q^2)`$ corrections (including radiative corrections) as described before. Since our aim is to compare our exclusive rate with the inclusive $`𝒪(\alpha _s)`$ rate written down in Sec.IV we have only retained radiative corrections up to $`𝒪(\alpha _s)`$ in the exclusive rate (33) for consistency reasons.
For the sake of completeness we separately list the $`𝒪(1/m_Q^2)`$ zero recoil corrections for the longitudinal and transverse pieces of the axial vector contributions. They are needed for the transverse/longitudinal separation shown in Fig.2. One has
$`|H_{\frac{1}{2}0}^A|^2|H_{\frac{1}{2}0}^A|^2+((f_1^A(1))^21){\displaystyle \frac{1}{6}}(A_{0\frac{1}{2}}^A)^2+\sqrt{{\displaystyle \frac{2}{9}}}(f_1^A(1)1)A_{0\frac{1}{2}}^AA_{2\frac{3}{2}}^A`$ (34)
$`|H_{\frac{1}{2}1}^A|^2|H_{\frac{1}{2}1}^A|^2+((f_1^A(1))^21){\displaystyle \frac{1}{3}}(A_{0\frac{1}{2}}^A)^2\sqrt{{\displaystyle \frac{2}{9}}}(f_1^A(1)1)A_{0\frac{1}{2}}^AA_{2\frac{3}{2}}^A`$ (35)
When summing the two contributions (35) in the rate formula the $`A_{0\frac{1}{2}}^AA_{2\frac{3}{2}}^A`$ interference contributions cancel out as is apparent in (33). As explained before we shall use the leading order results (30) for the second and third term in (35) since the factors $`((f_1^A(1))^21)`$ and $`(f_1^A(1)1)`$ multiplying them are already of $`𝒪(1/m_Q^2)`$.
Our numerical evaluation of (33) is based on the standard form factor (47) with $`\rho _B^2=0.75`$ using again $`V_{bc}=0.038`$. All parameters have been specified before. For the quasi-elastic rate we find $`\mathrm{\Gamma }^{\mathrm{excl}}=5.52\times 10^{10}s^1`$. The $`𝒪(1/m_Q)`$ and $`𝒪(1/m_Q^2)`$ corrections amount to $`+5.2\%`$ and $`6.6\%`$. The renormalization of the heavy quark current decreases the exclusive rate by $`8.8\%`$. In Fig.2 we show a plot of the $`\omega `$-spectrum of the quasi-elastic rate where we separately show the transverse ($`\lambda _W=\pm 1`$) and longitudinal contributions ($`\lambda _W=0`$) including $`𝒪(1/m_Q^2)`$ and $`𝒪(\alpha _s)`$ corrections calculated according to (35). The longitudinal rate dominates the spectrum except for a small region close to zero recoil. For the integrated rates we find $`\mathrm{\Gamma }_L^{\mathrm{excl}}`$/$`\mathrm{\Gamma }_T^{\mathrm{excl}}=1.89`$.
## IV Inclusive Semileptonic Rate $`\mathrm{\Lambda }_bX_c+l^{}+\overline{\nu }_l`$
To the leading order in the heavy mass expansion the inclusive rate is given by the free heavy quark decay rate which can be obtained from Eq.(1) using quark masses and setting the curly bracket equal to 1 . There are no $`𝒪(1/m_Q)`$ corrections to this result. Mass corrections come in at the order $`𝒪(1/m_Q^2)`$. In the case of $`\mathrm{\Lambda }_b`$-decay, where the light diquark system has spin $`0`$, the chromomagnetic contribution drops out and the mass corrections are determined by the nonperturbative kinetic energy parameter $`\mu _\pi ^2`$ alone.
Including also the $`\alpha _s`$-correction in the free quark decay rate , one has ($`x=(m_c/m_b)^2`$)
$$\mathrm{\Gamma }^{\mathrm{incl}}=\mathrm{\Gamma }_0\left(1\frac{2}{3}\frac{\alpha _s(m_b)}{\pi }g(x)\right)\left(1\frac{\mu _\pi ^2}{2m_b^2}\right),$$
(36)
where $`\mathrm{\Gamma }_0`$ is the lowest order (in $`\alpha _s`$) free quark decay rate
$$\mathrm{\Gamma }_0=\frac{G_F^2|V_{bc}|^2m_b^5}{192\pi ^3}I_0(x),I_0(x)=(1x^2)(18x+x^2)12x^2\mathrm{ln}x,$$
(37)
and the function $`g(x)`$ is determined by the $`𝒪(\alpha _s)`$ radiative corrections including all mass corrections as calculated in :
$`g(x)`$ $`=`$ $`h(x)/I_0(x),`$ (38)
$`h(x)`$ $`=`$ $`(1x^2)\left({\displaystyle \frac{25}{4}}{\displaystyle \frac{239}{3}}x+{\displaystyle \frac{25}{4}}x^2\right)+x\mathrm{ln}(x)\left(20+90x{\displaystyle \frac{4}{3}}x^2+{\displaystyle \frac{17}{3}}x^3\right)+x^2\mathrm{ln}^2(x)(36+x^2)`$ (39)
$`+`$ $`(1x^2)\left({\displaystyle \frac{17}{3}}{\displaystyle \frac{64}{3}}+{\displaystyle \frac{17}{3}}x^2\right)\mathrm{ln}(1x)4(1+30x^2+x^4)\mathrm{ln}(x)\mathrm{ln}(1x)`$ (40)
$``$ $`(1+16x^2+x^4)[6\mathrm{L}\mathrm{i}_2(x)\pi ^2]`$ (41)
$``$ $`32x^{3/2}(1+x)\left[\pi ^24\mathrm{L}\mathrm{i}_2(\sqrt{x})+4\mathrm{L}\mathrm{i}_2(\sqrt{x})2\mathrm{ln}(x)\mathrm{ln}\left({\displaystyle \frac{1\sqrt{x}}{1+\sqrt{x}}}\right)\right]`$ (42)
The $`𝒪(1/m_b^2)`$ corrections appear in the third factor of Eq.(36). For the value of the kinetic energy parameter $`\mu _\pi ^2`$ we take
$$\mu _\pi ^2=0.5\mathrm{GeV}^2,$$
(43)
where we assume equality of the kinetic energy parameter in the meson and baryon case.
It is well known and evident from Eq.(37) that the inclusive decay rate depends rather strongly on the exact value of the $`b`$-quark mass $`m_b`$ which is fraught with some uncertainties. We shall use the results of two recent theoretical analyses of the inclusive semileptonic decay rate. In the value of $`m_b`$ was determined from an analysis of $`\mathrm{{\rm Y}}`$ sum rules and $`B`$-meson semileptonic widths:
$`m_b=4.8\mathrm{GeV},m_c=1.325\mathrm{GeV},`$ (44)
where the charm quark mass was determined from the constraint
$$m_bm_c\mu _\pi ^2\left(\frac{1}{2m_c}\frac{1}{2m_b}\right)=\overline{M}_B\overline{M}_D.$$
(45)
The $`\overline{M}_{B,D}`$ are the spin-averaged masses $`\overline{M}_{B,D}=1/4(M_{B,D}+3M_{B^{},D^{}})`$ as before.
In the inclusive semileptonic $`B`$ decay rate was directly expressed in terms of the $`\mathrm{{\rm Y}}(1S)`$ meson mass instead of the $`b`$ quark mass. The authors of obtained:
$`\mathrm{\Gamma }^{\mathrm{incl}\mathrm{{\rm Y}}}={\displaystyle \frac{G_F^2|V_{bc}|^2}{192\pi ^3}}\left({\displaystyle \frac{m_\mathrm{{\rm Y}}}{2}}\right)^50.533\left[10.096ϵ0.029ϵ^2(0.28\lambda _2+0.12\lambda _1)/\mathrm{GeV}^2\right]`$ (46)
where $`ϵ=1`$ denotes the order of the expansion in $`m_\mathrm{{\rm Y}}`$. Mass corrections and radiative corrections are already taken into account. The parameters $`\lambda _1`$ and $`\lambda _2`$ in Eq.(46) are connected with the more familiar $`\mu _\pi ^2`$ and $`\mu _G^2`$ parameters by $`\mu _\pi ^2=\lambda _1`$ and $`\mu _G^2=3\lambda _2=0.36\mathrm{GeV}^2`$. For the $`\mathrm{\Lambda }_b`$ baryon we set $`\lambda _2=0`$ and assume equality of $`\lambda _1`$ in the meson and baryon case as before.
For the semileptonic inclusive $`bc`$ decay rate of the $`\mathrm{\Lambda }_b`$ we finally obtain $`\mathrm{\Gamma }^{\mathrm{incl}}=6.5010^{10}s^1`$ using the mass parameters from and $`\mathrm{\Gamma }^{\mathrm{incl}\mathrm{{\rm Y}}}=6.2310^{10}s^1`$ using the evaluation of . Again we have set $`V_{bc}=0.038`$. We mention that these two inclusive rate values include the $`𝒪(\alpha _s)`$ radiative corrections which lower the inclusive rates by about $`11\%`$. This value is not very far away from the $`8.8\%`$ by which the exclusive rate gets lowered by the same radiative corrections.
## V The Exclusive/Inclusive $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ Ratio
In this section we determine the exclusive/inclusive ratio $`R_E=\mathrm{\Gamma }^{\mathrm{excl}}/\mathrm{\Gamma }^{\mathrm{incl}}`$ in semileptonic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ decays based on our estimates for the inclusive rates derived in Sec.IV and on various phenomenological models for the baryonic Isgur-Wise function $`F(\omega )`$ entering in the exclusive differential rate Eq.(33). Of all the phenomenological models we shall mostly focus our attention on the sum rule calculation of .
We begin our discussion with the determination of the leading order $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ Isgur-Wise function by the QCD sum rule method given in . The shape of the Isgur-Wise function in can be very well reproduced by an exponential representation of the form
$$F(\omega )=\frac{2}{\omega +1}\mathrm{exp}\left((2\rho _B^21)\frac{\omega 1}{\omega +1}\right),$$
(47)
which has the correct zero recoil normalization $`F(1)=1`$ and a slope parameter given by $`\rho _B^2`$. The convexity parameter in this representation (proportional to $`(\omega 1)^2)`$) is given by $`c_B=1/8(1+4\rho _B^2+4\rho _B^4)`$ and is positive for $`\rho _B^20.207`$ as in most model calculations. We refer to this representation of the Isgur-Wise function as the standard form. For the $`\rho _B^2`$ parameter the authors of find $`\rho _B^2=0.85`$ and $`\rho _B^2=0.65`$ using diagonal and nondiagonal sum rules, respectively. As an average of these two values one obtains
$$\rho _B^2=0.75.$$
(48)
Using the average value of $`\rho _B^2`$, $`V_{bc}=0.038`$, the standard representation of the Isgur-Wise function (47) and the rate formula (33) from Sec.III one obtains the exclusive rate $`\mathrm{\Gamma }^{\mathrm{excl}}=5.5210^{10}s^1`$. From the inclusive rate calculated using the mass parameters given in $`\mathrm{\Gamma }^{\mathrm{incl}}=6.5010^{10}s^1`$ one finds $`R_E=0.85`$ for the exclusive/inclusive ratio. Note that the $`V_{bc}`$ dependence drops out in this ratio. The values of the slope parameter $`\rho _B^2`$ and the exclusive/inclusive ratio $`R_E`$ of the model of as well as those of other phenomenological models have been collected together in Table I. Table I uses the larger value of the two inclusive reference rates discussed in Sec.IV based on the mass parameters of . If one instead uses the inclusive rate of all $`R_E`$-values in Table I have to be increased by 4.3%. Radiative corrections do not affect the exclusive/inclusive ratios listed in Table I very much since they lower both the exclusive and inclusive rates (see Secs. III and IV). If they were left out the exclusive/inclusive ratio $`R_E`$ would be reduced by $`2\%`$.
It is clear from Table I that the form factor calculated in the quark confinement model is too flat to satisfy the bound $`\mathrm{\Gamma }^{\mathrm{excl}}\mathrm{\Gamma }^{\mathrm{incl}}`$. All other models in Table I satisfy this upper bound. Translating the upper bound $`R_E=1`$ into a lower bound on the slope parameter $`\rho _B^2`$ one obtains
$$(\rho _B^2)_{\mathrm{min}}=0.36$$
(49)
using the standard form factor function (47) and the inclusive rate calculated from the mass parameters in .
An upper bound on the slope parameter $`\rho _B`$ can be obtained using the spectator model bound discussed in Sec.II, which reads
$$\rho _B^22\rho _M^2\frac{1}{2}.$$
(50)
The mesonic slope parameter $`\rho _M^2`$ can be extracted from the exclusive semileptonic $`B`$ decays . The values are
$`(\rho _M^2)_1`$ $`=`$ $`0.66\pm 0.19\mathrm{from}\overline{B}Dl^{}\overline{\nu }`$ (51)
$`(\rho _M^2)_2`$ $`=`$ $`0.71\pm 0.11\mathrm{from}\overline{B}D^{}l^{}\overline{\nu }`$ (52)
with the world average values for $`V_{cb}`$ being $`|V_{cb}|=0.0394\pm 0.0050`$ and $`|V_{cb}|=0.0387\pm 0.0031`$, respectively. The weighted average of the two mesonic slope parameters are then $`\rho _M^2=0.70\pm 0.10`$. This translates into an upper bound for the baryonic slope parameter $`\rho _B^2`$ according to the spectator quark model bound (50). One has
$$(\rho _B^2)_{\mathrm{max}}=0.89\pm 0.19.$$
(53)
We mention that very likely the error on this bound will be considerable reduced in the near future with the new data expected from the bottom factories at SLAC and KEK. Combining both limits, (49) and (53), we obtain a prediction for the allowed values of the baryon slope parameter given by
$`0.36<\rho _B^2<0.89\pm 0.19.`$ (54)
According to these upper and lower bounds the first model (as remarked on before) and the last four models in Table I have to be excluded since they possess form factors which are too flat or too steep, respectively. The two QCD sum rule calculations as well as the simple quark model evaluation feature slope parameters that satisfy the bounds (54). We consider the two QCD sum rule calculations to be the most reliable of the three model calculations since they are the least model dependent. Our final prediction for the range of values of the exclusive/inclusive ratio will be based on the two slope parameter values $`\rho _B^2=0.85`$ and $`0.65`$, resulting from the analysis of the diagonal and nondiagonal sum rules in , respectively. This range also includes the sum rule result of . In determining our prediction for the range of $`R_E`$ we shall also allow for the smaller inclusive rate calculated by the method of . Thus our final prediction for the range of the exclusive/inclusive ratio in semileptonic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ decays is $`R_E=0.81÷0.92`$. This range is consistent with the range of values from the semiquantitave analysis performed in Sec.II. Our conclusion is that the exclusive/inclusive ratio of semileptonic $`\mathrm{\Lambda }_b`$-decays is considerable higher than in the corresponding bottom meson case.
## VI Missing final states
Besides the quasi-elastic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ contribution discussed before there are also $`\mathrm{\Lambda }_c^{}`$ resonant states and multiparticle final states contributing to the fully inclusive semileptonic $`\mathrm{\Lambda }_b`$ rate. Of course, if the quasi-elastic contribution dominates the total inclusive rate much more than by the 66% in the heavy meson case, there would not be much room left for the resonant and multiparticle final states. In the main body of this paper we have collected together theoretical evidence that the latter situation is very likely the case. One could turn this statement around in the following sense: if one would have theoretical reasons to believe that resonant and multiparticle final states are suppressed in inclusive semileptonic $`\mathrm{\Lambda }_bX_c`$ transitions then the quasi-elastic exclusive $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ contribution must dominate. As we shall see there are theoretical reasons to believe in such a suppression in as much as some of the transitions to orbitally excited $`\mathrm{\Lambda }_c^{}`$ charm baryon states involve spin-orbit coupling transitions which are believed to be suppressed.
The purpose of this section is to classify those final states in semileptonic $`\mathrm{\Lambda }_bX_c`$ transitions that form the complement of the quasi-elastic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ transition. We divide these into class $`A`$ contributions $`\mathrm{\Lambda }_b\mathrm{\Lambda }_cXl\nu `$, where the charm quark of the decay ends up in a charm $`\mathrm{\Lambda }_c`$ directly or indirectly, and class $`B`$ contributions, where the charm quark goes into a meson or a charm-strangeness baryon $`\mathrm{\Lambda }_bX_c(non\mathrm{\Lambda }_c)l\nu `$. Accordingly we define the two ratios
$$R_A=\frac{\mathrm{\Gamma }(\mathrm{\Lambda }_b\mathrm{\Lambda }_cXl\nu )}{\mathrm{\Gamma }(\mathrm{\Lambda }_bX_cl\nu )}.$$
(55)
and
$$R_B=\frac{\mathrm{\Gamma }(\mathrm{\Lambda }_bX_c(non\mathrm{\Lambda }_c)l\nu )}{\mathrm{\Gamma }(\mathrm{\Lambda }_bX_cl\nu )}.$$
(56)
Together with the exclusive/inclusive ratio defined before
$$R_E=\frac{\mathrm{\Gamma }(\mathrm{\Lambda }_b\mathrm{\Lambda }_cl\nu )}{\mathrm{\Gamma }(\mathrm{\Lambda }_bX_cl\nu )},$$
(57)
the three ratios must add up to one, i.e.
$$R_E+R_A+R_B=1.$$
(58)
Note that all three ratios are positive definite which makes the constraint (58) potentially quite powerful if $`R_E`$ is close to one as is indicated by our analysis of the quasielastic rate in the previous sections. As concerns the sizes of $`R_A`$ and $`R_B`$ one cannot even hope to provide semiquantitative answers at present. It is nevertheless useful to enumerate the final states belonging to the class $`A`$ and class $`B`$ transitions which we shall do in the following.
Class $`A`$ final states
Potentially prominent among the class $`A`$ final states are the transitions into the seven excited $`P`$-wave $`\mathrm{\Lambda }_c^{}`$-states. Taking the bottom meson case for comparison theoretical estimates show that the corresponding transitions into excited mesonic $`P`$-wave states make up approximately 10% of semileptonic B-decays . The $`\mathrm{\Lambda }_c^{}`$-states eventually decay down to the $`\mathrm{\Lambda }_c`$ ground state via (multiple) pion emission or, with a much smaller branching fraction, via photon emission. There are altogether seven such $`P`$-wave states which are grouped into the three HQS doublets $`\{\mathrm{\Lambda }_{cK1}\}`$, $`\{\mathrm{\Lambda }_{ck1}\}`$, $`\{\mathrm{\Lambda }_{ck2}\}`$, and the singlet $`\{\mathrm{\Lambda }_{ck0}\}`$. We use the terminology of such that the excited $`K`$\- and $`k`$-states are symmetric and antisymmetric under the exchange of the momenta of the light quarks. The five symmetric states $`\{\mathrm{\Lambda }_{ck0}\}`$, $`\{\mathrm{\Lambda }_{ck1}\}`$, and $`\{\mathrm{\Lambda }_{ck02}\}`$ are made from a heavy quark and a light spin-one diquark. $`\mathrm{\Lambda }_b`$ transitions into these five states involve spin-zero to spin-one light-side transitions which can be expected to be strongly suppressed since they involve spin-orbit interactions. In the spectator quark model, where one neglects spin-orbit interactions, transitions into these five states are forbidden . It would be interesting to experimentally confirm this suppression. One thus remains with the transitions into the HQS doublet $`\{\mathrm{\Lambda }_{cK1}\}`$ whose spin $`1/2^{}`$ and spin $`3/2^{}`$ members are very likely the recently discovered $`\mathrm{\Lambda }_c(2593)`$ and $`\mathrm{\Lambda }_c(2625)`$ states . $`\mathrm{\Lambda }_b`$ branching ratios into these states are not yet available. There could also be transitions into higher orbital $`\mathrm{\Lambda }_c^{}`$-states. These transitions are, however, expected to be suppressed because of angular momentum suppression factors. Besides, transitions into symmetric higher orbital $`\mathrm{\Lambda }_c^{}`$-states would again be suppressed due to spin-orbit coupling suppression. The suppression of transitions into the symmetric orbitally excited $`\mathrm{\Lambda }_c^{}`$-states could be the source of the possible depletion of class $`A`$ final states. For example, using spin counting only $`1/3`$ of the existing $`P`$-wave excitations can be reached in semileptonic $`\mathrm{\Lambda }_b`$ transitions if the spin-orbit coupling suppression is active.
Another source of class $`A`$ final states is accessible due to the creation of one or more additional $`(d\overline{d})`$\- or $`(u\overline{u})`$-quark pairs in the basic transition. The relevant transitions for $`(d\overline{d})`$ creation are (see Fig.3a)
$`\mathrm{\Lambda }_b^0`$ $``$ $`\mathrm{\Lambda }_c^+(\mathrm{\Sigma }_c^+)+X_M^0+l^{}+\overline{\nu }_l,`$ (59)
or, when exchanging the $`du`$ lines originating from the $`\mathrm{\Lambda }_b`$, one has
$`\mathrm{\Lambda }_b^0`$ $``$ $`\mathrm{\Sigma }_c^0+X_M^++l^{}+\overline{\nu }_l.`$ (60)
For $`(u\overline{u})`$ creation shown in Fig.3b one has
$`\mathrm{\Lambda }_b^0`$ $``$ $`\mathrm{\Sigma }_c^{++}+X_M^{}+l^{}+\overline{\nu }_l.`$ (61)
The exchange of the $`d,u`$ lines originating from the $`\mathrm{\Lambda }_b`$ brings one back to (59). $`X_M`$ stands for a charmless mesonic inclusive state. Excited charm baryon states such as $`\mathrm{\Lambda }_c^{}`$ and $`\mathrm{\Sigma }_c^{}`$ are not explicitly included in the listing (59 \- 61), but are implied. The $`\mathrm{\Sigma }_c^0`$, $`\mathrm{\Sigma }_c^+`$, and $`\mathrm{\Sigma }_c^{++}`$ appearing in (59-61) cascade down to the $`\mathrm{\Lambda }_c^+`$ state via pion emission making these processes class $`A`$ final states.
Class $`B`$ final states
There are two sources for class $`B`$ final states. First there is $`(s\overline{s})`$ quark pair creation where the strange quark ends up in a charm-strangeness baryon (Fig.4a) which decays weakly into noncharm states and therefore does not contribute to the class $`A`$ final states <sup>5</sup><sup>5</sup>5The weak decay $`\mathrm{\Xi }_c\mathrm{\Lambda }_c+\pi `$, though interesting, occurs only at the per mill level .. Second, the charm quark of the decay may end up in a charm meson accompanied by $`(u\overline{u})`$-, $`(d\overline{d})`$, and $`(s\overline{s})`$-quark pair creation as shown in Figs.4b-d. Let us list a few examples of such transitions. From $`(s\overline{s})`$ pair creation one has (Fig.4a)
$`\mathrm{\Lambda }_b^+`$ $``$ $`\mathrm{\Xi }_c^++X_{M_s}^0+l^{}+\overline{\nu }_l,`$ (62)
or, when exchanging the $`du`$ lines, one has
$`\mathrm{\Lambda }_b^+`$ $``$ $`\mathrm{\Xi }_c^0+X_{M_s}^++l^{}+\overline{\nu }_l,`$ (63)
$`X_{M_s}`$ now stands for a strangeness meson state. Then there are the transitions where the charm quark goes into a charm meson. These are
$$\mathrm{\Lambda }_b^+D^++X_B^0+l^{}+\overline{\nu }_l,$$
(64)
$$\mathrm{\Lambda }_b^+D^0+X_B^++l^{}+\overline{\nu }_l,$$
(65)
$$\mathrm{\Lambda }_b^+D_s^++X_{B_s}^0+l^{}+\overline{\nu }_l,$$
(66)
$`X_B`$ stands for a light baryon state and $`X_{B_s}`$ for a strangeness baryon state. Excitations of the charm meson and charm baryon states are again implied. We do not discuss $`(c\overline{c})`$ pair creation. The corresponding final states are barely accessible in semileptonic $`\mathrm{\Lambda }_b`$ decays for kinematical reasons and will have a spectacular signature anyhow.
## VII Summary and Conclusion
We have brought together various pieces of theoretical evidence that the exclusive/inclusive ratio $`R_E`$ in semileptonic $`\mathrm{\Lambda }_b`$-decays is larger than in semileptonic $`B`$ decays, where the exclusive/inclusive ratio amounts to 66%. We predict that the exclusive quasielastic semileptonic $`\mathrm{\Lambda }_b`$-decays make up between 81% and 92% of the total inclusive semileptonic $`\mathrm{\Lambda }_b`$ rate. At present there is no experimental information on either the exclusive or the inclusive branching ratio in semileptonic $`\mathrm{\Lambda }_b`$ decays. The problem is that present and planned experiments do not have access to reliable $`\mathrm{\Lambda }_b`$ tags which are necessary for a measurement of their branching fractions. Ideally one would run a $`e^+e^{}`$-machine right above $`\mathrm{\Lambda }_b\overline{\mathrm{\Lambda }}_b`$ threshold which would solve the tagging problem. However, such experiments are not planned in the foreseeable future. The above assertion about the dominance of the quasi-elastic mode in semileptonic $`\mathrm{\Lambda }_b`$ decays may take a long time to verify experimentally. It may nevertheless be used as a working hypothesis in the experimental analysis of semileptonic $`\mathrm{\Lambda }_b`$ decays in particular if further theoretical progress in the theoretical description of semileptonic $`\mathrm{\Lambda }_b`$ decays confirms the estimates made in this paper.
###### Acknowledgements.
This investigation was prompted by two questions of our experimental colleagues P. Roudeau and G. Sciolla who asked us about theoretical expectations for the size of the the $`\mathrm{\Lambda }_c+X`$ (class A) and $`(non\mathrm{\Lambda }_c)+X_c`$ (class B) contributions in inclusive semileptonic $`\mathrm{\Lambda }_b`$ decays. We did our best to try and provide answers to these questions in our paper, at least partially. We would like to thank N. Uraltsev for fruitful discussions. We also thank O. Yakovlev who participated in the early stages of this calculation. The work of B.M. was partially supported by the Ministry of Science and Technology of the Republic of Croatia under the contract Nr. 00980102.
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# Removal of the calcium underabundance in cool metal rich Galactic disk dwarfsBased on observations made at ESO, La Silla
## 1 Introduction
Calcium, being one of the so called $`\alpha `$-elements, is believed to be produced in SNII, exploding massive stars. Due to the early start of such events in the Galaxy as compared to the other stellar nucleosynthesis sites SNIa and AGB stars, old (metal poor) stars are usually rich in Ca compared to Fe as compared to the corresponding solar Ca/Fe ratio (e.g McWillliam et al., 1995; Chen et al., 2000). For metal rich stars the Ca/Fe trend seems to level out at the solar value. A considerable scatter was found for 48 metal rich disk dwarfs by Feltzing & Gustafsson (1998). As a mean \[Ca/Fe\]$`0`$, but for the cool dwarfs a significant underabundance was derived. To examine this, I observed 10 cool disk dwarfs with ESO CAT/CES in Nov-Dec 1995. The underabundance in Ca for K-dwarfs was suspected to be caused by overionization, as predicted by Drake (1991). A MULTI Carlsson (1986) NLTE investigation of this with a new Ca atom model was therefore performed by Thorén (2000).
## 2 Observations and analysis
The observational data and the LTE analysis are presented in detail in a separate paper with a wider scope Thorén and Feltzing (2000). The stars observed and the parameters used in the models are presented in Table 1. The LTE line profiles computed by the Uppsala synthetic spectrum code SPECTRUM agreed well with the observed spectra. With a new Ca model atom for the code MULTI it appeared that NLTE effects for Ca in the cool dwarfs in the sample actually were very small Thorén (2000). The recent MARCS atmospheres Asplund et al. (1997) have a higher, more realistic amount of UV metal line blocking than the models used by Drake (1991). This increases the opacity and decreases the non-local ionizing radiation field. Except in the line cores no strong NLTE effects could be seen, rather the Ca/Fe ratio appeared solar, as for the hotter stars in Feltzing & Gustafsson (1998). The suggested underabundance, if not real, had to have a different origin.
## 3 Atomic line data
The lines used in this analysis are presented in Table 2. They were selected from the VALD database Piskunov et al. (1995). Most oscillator strengths of the lines used in the analysis were modified, to fit our solar observations. VALD provides oscillator strengths typically correct to the order of magnitude. The van der Waals broadening parameters for lines used in the analysis were also changed. For Ca lines the van der Waals broadening width in terms of $`\delta \mathrm{\Gamma }_6`$ factors (being the corrections to the classical Unsöld value) calculated from lab measurements by Smith (1981) are used, except for the 6798 Å line for which no $`\delta \mathrm{\Gamma }_6`$ value was available. However, this line is very weak and is not affected significantly by pressure broadening. Its broadening was calculated according to Barklem & O’Mara (1997). For Fe i lines the damping was calculated according to Anstee & O’Mara (1995) and Barklem & O’Mara (1997) for the lines in Table 2 marked with asterisks. For the remaining Fe i lines the $`\delta \mathrm{\Gamma }_6`$ factor had to be used.
When the atomic line data for the project were examined, the explanation was highlighted. The Uppsala code EQWIDTH had been used for the analysis in Feltzing & Gustafsson (1998). This code uses a correction factor $`\delta \mathrm{\Gamma }_6`$ as input for atomic line broadening by H atoms. To get the broadening the factor is multiplied with the classical Unsöld broadening value. This factor is typically $`2`$. The Ca atomic line parameters used by Feltzing & Gustafsson (1998) and Thorén (2000) are both taken from Smith (1981). The Ca lines $`\delta \mathrm{\Gamma }_6`$ factors were, however, erroneously calculated in Feltzing & Gustafsson (1998), causing the line broadening to be much too large. In the cool dwarfs the overall line strength is larger than in the hotter ones. For lines with equivalent widths increasing beyond 50 mÅ this parameter soon becomes very important, because it brings the line out of the saturated region on the curve of growth faster. This means that a model line is too strong for any given abundance, forcing a reduction of the abundance in order to satisfy the measured equivalent width.
Since Feltzing & Gustafsson (1998) used equivalent widths, rather than synthetic spectra, the effect was not noticed until it showed up as an apparent underabundance for Ca in the cool dwarfs. Since such a NLTE effect had already been predicted Drake (1991), this was suggested to be the cause. Because of this it was decided to make the new K dwarf observations and a new MULTI analysis for Ca. The Ca NLTE properties of cool metal rich dwarfs will be discussed in Thorén (2000).
The abundances were obtained in the following way: First synthetic spectra were fitted with SPECTRUM to the observed spectra. The atmospheric model parameters used were those determined with photometric Strömgren calibrations for Pop I F-G dwarfs, presented by Olsen (1984). The Strömgren colors were obtained from the catalogues in Olsen (1994) and Olsen (1995).
The equivalent widths of the fitted synthetic lines to be used in the abundance analysis were exported from the fitted spectra. These equivalent widths were then used as measured ones from the observed spectra and used as input into the code EQWIDTH (which was used in the analysis of Feltzing & Gustafsson (1998)). The lines used in the analysis are presented in Table 2.
Figure 1 shows the difference in abundance pattern for the old and new calculated values of $`\mathrm{\Gamma }_6`$. Triangles represent abundances determined with the values in Feltzing & Gustafsson (1998), circles represent abundances with corrected $`\mathrm{\Gamma }_6`$ values. As seen in the Fig. 1 there remains an increasing trend in \[Ca/Fe\] with increasing temperature after the damping treatment correction. This effect is reduced to virtually zero when the photometric effective temperatures are adjusted, according to Thorén & Feltzing (2000), until Fe i lines of different excitation energy give similar iron abundances. Four of the stars required such changes. Their Ca abundances before ($``$) and after ($``$) temperature adjustments are shown in figure 1. The need for $`T_{\mathrm{eff}}`$ changes is illustrated for one of the objects in Fig. 2. HD32147 needed a positive temperature correction of 200 K, which also raised the Ca abundance by +0.20 dex.
## 4 Summary
Changing $`\mathrm{\Gamma }_6`$ to the ’true’ value reduces the underabundance by 0.3 dex for the coolest stars in this sample. This solves the major part of the underabundance problem.
The photometric effective temperatures adopted by Feltzing & Gustafsson (1998) lead in some cases to severe trends in the abundances derived from Fe i lines with different excitation energy. Modifications of these $`T_{\mathrm{eff}}`$ values dramatically decrease this abundance scatter and also modify the Ca abundances to values observed in hotter stars of similar metallicity.
The results in this work and in the forthcoming paper Thorén and Feltzing (2000) indicate that LTE abundance analysis is indeed useful even for cool metal rich dwarfs, increasing the number of objects available for studying the chemical evolution of the Galaxy.
###### Acknowledgements.
The author was supported by the Swedish Natural Science Research Council, NFR. Professor Bengt Gustafsson is thanked for suggesting the writing of this letter. Dr. Paul Barklem and Docent Bengt Edvardsson are thanked for valuable comments on the text.
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# PION PRODUCTION MODEL - CONNECTION BETWEEN DYNAMICS AND QUARK MODELS
## Acknowledgments
This work was supported U.S. DOE Nuclear Physics Division Contract No. W-31-109-ENG-38 and JSPS Grant-in-Aid for Scientific Research (C) 12640273.
## References
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# On the WKB Quantum Equivalence between Diverse 𝑝–brane Actions
## Abstract
We consider an action for a closed, bosonic, $`p`$–brane, where the brane tension is not an assigned parameter but rather it is induced by a maximal rank gauge $`p`$–form. This model is classically equivalent to the Nambu–Goto/Howe–Tucker model. We investigate how this classical equivalence can be implemented in the path integral framework. For this purpose we adopt a “first order” integration procedure over gauge $`p`$–forms and a “shortened” Fadeev–Popov procedure.
Diverse action functionals have been proposed to describe dynamics of a relativistic, bosonic, $`p`$–brane . The first brane action, proposed in 1975, was a generalization of the Nambu–Goto action for strings, i.e. the measure of the brane world history
$$S_{DNG}[Y]=m_{p+1}_\mathrm{\Sigma }d^{p+1}\sigma \sqrt{\gamma },\gamma det\left(_mY^\mu _nY_\mu \right),$$
(1)
where $`m_{p+1}`$ is the “$`p`$–tension”. We use “$`p`$” for the spatial dimensionality of the brane; thus, the coordinates $`\sigma ^m`$, $`m=0,1,\mathrm{}p`$, span the $`(p+1)`$–dimensional world manifold $`\mathrm{\Sigma }`$. The $`D`$ functions $`Y^\mu (\sigma )`$, $`\mu =0,1,\mathrm{}D`$, are the brane coordinates in the $`D`$–dimensional target spacetime. The special case $`p=2`$ and $`D=4`$ was already introduced in 1962 by Dirac in an attempt to resolve the electron–muon puzzle in terms of a relativistic membrane .
An alternative description, preserving world manifold reparametrization invariance, can be achieved by introducing an auxiliary world manifold metric $`g_{mn}(\sigma )`$ and a “cosmological term” , ,
$$S_{HTP}[Y,g]=\frac{m_{p+1}}{2}_\mathrm{\Sigma }d^{p+1}\sigma \sqrt{g}\left[g^{mn}_mY^\mu _nY_\mu (p1)\right],$$
(2)
where $`gdet(g_{mn})`$. In both functionals (1) and (2) the brane tension $`m_{p+1}`$ is a pre–assigned parameter.
The two actions (1) and (2) are classically equivalent as the “field equations” $`\delta S/\delta g^{mn}(\sigma )=0`$ require the auxiliary world metric to match the induced metric, i.e. $`g_{mn}=\gamma _{mn}=_mY^\mu _nY_\mu `$. Moreover they are also complementary: $`S_{DNG}`$ provides an “extrinsic” geometrical description in terms of the embedding functions $`Y^\mu (\sigma )`$ and the induced metric $`\gamma _{mn}`$, while $`S_{HTP}`$ assigns an “intrinsic” geometry to the world manifold $`\mathrm{\Sigma }`$ in terms of the metric $`g_{mn}`$ and the “cosmological constant” $`m_{p+1}`$; the $`Y^\mu (\sigma )`$ functions enter as a “multiplet of scalar fields” propagating on a curved $`(p+1)`$–dimensional manifold.
More recently new action functionals have been proposed where the brane tension, or world manifold cosmological constant, is not an a priori assigned parameter, but follows from the dynamics of the object itself and can attain both positive and vanishing values. Either Kaluza–Klein type mechanism and modified integration measure have been proposed as candidate dynamical processes to produce tension at the classical level.
The main purpose of this note is to investigate how the dynamical generation of the brane tension and the equivalence between diverse action functionals can be extended at the quantum level in the WKB approximation of a “sum over histories” approach. However, before considering the path–integral it is instrumental to review how classical dynamics leads to the action (1) as an effective, on–shell action.
The guiding principle to assign the $`p`$–brane tension the role of a dynamical variable is borrowed from modern cosmology, where the cosmological constant can be represented by a maximal rank gauge $`p`$–form . Thus, we introduce the following action functional
$`S[Y,g,A]`$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}d^{p+1}\sigma \sqrt{g}\left[{\displaystyle \frac{m_{p+1}}{2}}{\displaystyle \frac{(\gamma )}{(g)}}{\displaystyle \frac{1}{2(p+1)!}}F_{m_1\mathrm{}m_{p+1}}F^{m_1\mathrm{}m_{p+1}}\right]`$ (3)
$`=`$ $`{\displaystyle _\mathrm{\Sigma }}d^{p+1}\sigma \left[{\displaystyle \frac{m_{p+1}}{2}}{\displaystyle \frac{(\gamma )}{\sqrt{g}}}{\displaystyle \frac{\sqrt{g}}{2(p+1)!}}F_{m_1\mathrm{}m_{p+1}}F^{m_1\mathrm{}m_{p+1}}\right],`$ (4)
where the world manifold $`\mathrm{\Sigma }`$ has a space–like boundary $`\mathrm{\Sigma }`$ whose target space image will represent a closed, $`p`$–dimensional, relativistic object. Moreover, we introduce on the world manifold a maximal rank gauge field, $`A_{m_2\mathrm{}m_{p+1}}(\sigma )`$, with field strength $`F_{m_1\mathrm{}m_{p+1}}_{[m_1}A_{m_2\mathrm{}m_{p+1}]}(\sigma )`$. To preserve gauge invariance under $`\delta A_{m_2\mathrm{}m_{p+1}}=_{[m_2}\mathrm{\Lambda }_{m_3\mathrm{}m_{p+1}]}`$ in the presence of a boundary we must give up a current–potential interaction term and consider only a gravitational coupling $`A`$$`g`$. By a suitable rescaling of the brane coordinates the dimensional constant $`m_{p+1}`$ can be washed out, and the classical action (4) written without any dimensional scale. Our goal is to show that the Dirac–Nambu–Goto functional can be obtained as an effective action from (4) once the classical field equations for the $`p`$–form gauge potential are solved.
Varying the action (4) with respect to $`A`$ we get
$$\frac{\delta S[Y,g,A]}{\delta A_{m_2\mathrm{}m_{p+1}}(\sigma )}=0_m\left(\sqrt{g}F^{mm_2\mathrm{}m_{p+1}}\right)=0,$$
which, since $`A`$ is maximal on the $`p`$–brane, has the solution
$$F^{mm_2\mathrm{}m_{p+1}}=\mathrm{\Lambda }ϵ^{mm_2\mathrm{}m_{p+1}}=\mathrm{\Lambda }\frac{1}{\sqrt{g}}\delta ^{[mm_2\mathrm{}m_{p+1}]},$$
(5)
where $`\mathrm{\Lambda }`$ is an arbitrary integration constant. By inserting the solution (5) back into (4) we obtain
$$S[Y,g]=_\mathrm{\Sigma }d^{p+1}\sigma \left[\frac{m_{p+1}}{2}\frac{(\gamma )}{\sqrt{g}}+\frac{\mathrm{\Lambda }^2}{2}\sqrt{g}\right],$$
(6)
where, the world manifold cosmological constant $`\mathrm{\Lambda }^2`$ shows up as the on–shell value of the gauge field kinetic term.
The on-shell action (6) depends from the world metric only through the volume density $`\sqrt{g}`$. Hence, variations with respect to $`g_{mn}`$ reduce to variations with respect $`\sqrt{g}`$:
$$\frac{\delta S[Y,g]}{\delta g_{mn}(\sigma )}=0\frac{\delta S[Y,g]}{\delta \sqrt{g}}=0.$$
(7)
By inserting the solution (5) into (7) we get
$$m_{p+1}\frac{(\gamma )}{(g)}=\mathrm{\Lambda }^2\sqrt{g}=\frac{1}{\mathrm{\Lambda }}\sqrt{m_{p+1}}\sqrt{\gamma }$$
(8)
and
$$S=\mathrm{\Lambda }\sqrt{m_{p+1}}_\mathrm{\Sigma }d^{p+1}\sigma \sqrt{\gamma }\rho _p_\mathrm{\Sigma }d^{p+1}\sigma \sqrt{\gamma }.$$
(9)
After solving for the world metric in terms of the brane coordinates, the action (4) turns out to be equivalent to a Dirac–Nambu–Goto action with a dynamically induced brane tension given by $`\rho _p\mathrm{\Lambda }\sqrt{m_{p+1}}`$. Let us remark that $`\mathrm{\Lambda }`$ can take any value including zero. Accordingly, null branes, corresponding to the action
$$S_{null}[Y,g]=\frac{m_{p+1}}{2}_\mathrm{\Sigma }d^{p+1}\sigma \frac{(\gamma )}{\sqrt{g}},$$
(10)
are included in our description as well. This special case stresses how the parameter $`m_{p+1}`$ is not necessarily the brane tension, but only a dimensional constant needed to allow the various dynamical fields in the action to keep their canonical dimensions<sup>§</sup><sup>§</sup>§ If we keep $`m_{p+1}0`$ the brane coordinates have canonical dimension of length, while the world and the induced metric are dimensionless (in units $`\mathrm{}=1`$, $`c=1`$), i.e. $`[Y^\mu ]=M^1`$, $`[\gamma _{mn}]=[g_{mn}]=1`$. .
In the second part of this note we shall discuss the above equivalence at the quantum level. The basic quantity encoding the $`p`$–brane quantum dynamics is the boundary wave functional, or vacuum—one–brane amplitude
$$ZZ[\widehat{Y},\widehat{A},\widehat{g}]=^{\widehat{g}}[Dg_{mn}]^{\widehat{Y}}[DY^\mu ]^{\widehat{A}}[DA]\mathrm{exp}\left(iS[Y,g,A]\right),$$
(11)
where the sum is over all bulk fields configurations inducing “hatted” fields on the boundary of the brane. We are assuming that the brane world manifold has a single, $`p`$–dimensional boundary, parametrized as $`\sigma ^m=\sigma ^m(s^a)`$, $`a=1,\mathrm{}p`$, which is mapped into the physical brane $`\widehat{Y}^\mu (s)`$; $`\widehat{g}`$ and $`\widehat{A}`$ are the induced metric and gauge potential over $`\widehat{Y}^\mu (s)`$. The integration variables in $`Z`$ “live” in the brane bulk, while we let free the fields induced on the boundary, i.e. we do not assign an independent classical action to the hatted fields.
The first field to be integrated out is the gauge $`p`$–form $`A`$. The standard routine goes through a lengthy procedure of gauge fixing and Fadeev–Popov compensation to invert the classical kinetic operator and define an appropriate quantum propagator. On the other hand, one knows that a gauge $`p`$–form over a $`(p+1)`$–dimensional manifold has no dynamical degrees of freedom and can describe only a static interaction. In such a limiting case the Fadeev–Popov procedure leaves no propagating degree of freedom at the quantum level. To shorten the whole gauge fixing procedure of ghost terms with different rank , we shall provide an alternative “recipe” to kill all the apparent degrees of freedom. We write the path–integral in the first order version, where the gauge potential $`A`$ and field strength $`F`$ are introduced as independent integration variables and we integrate away the gauge part of $`A`$ after inserting gauge fixing Dirac delta’s and the corresponding ghost determinants in the functional measure. The remaining, gauge invariant part of $`A`$ enforces $`F`$ to be a classical solution of the field equations, which is a constant background field. No propagating degrees of freedom survive at the quantum level. A formal proof of the equivalence between second order and first order quantization procedures, in the general case of a $`p`$–form in $`p+1`$ dimensions, is beyond the purpose of this short note. Rather, we will briefly consider the simplest, non trivial case which is $`p=1`$ gauge form over a two-dimensional, flat manifold without boundary, and then translate the result to the case we are studying. The first order, gauge fixed and Fadeev–Popov compensated path integral is
$$Z_{p=1}=[DF][DA]\delta \left[^mA_m\right]\mathrm{\Delta }_{FP}\mathrm{exp}\left\{id^2\sigma \left[\frac{1}{4}F_{mn}F^{mn}\frac{1}{2}F^{mn}_{[m}A_{n]}\right]\right\}.$$
(12)
By splitting $`A_m`$ into the sum of a “transverse” vector $`A_m^T`$ and a “gauge part” $`_m\varphi `$, the integration measure $`[DA]`$ turns into $`[DA^T][D\varphi ]\left(det\mathrm{}\right)^{1/2}`$ and (12) reads
$`Z_{p=1}`$ $`=`$ $`{\displaystyle }[DF][DA^T][D\varphi ](det\mathrm{})^{1/2}\delta \left[\mathrm{}\varphi \right]\times `$ (14)
$`\times \mathrm{\Delta }_{FP}\mathrm{exp}\left\{i{\displaystyle d^2\sigma \left[\frac{1}{4}F_{mn}F^{mn}\frac{1}{2}F^{mn}_{[m}A_{n]}^T\right]}\right\}`$
$`=`$ $`{\displaystyle [DF][DA^T]\left(det\mathrm{}\right)^{1/2}\mathrm{exp}\left\{id^2\sigma \left[\frac{1}{4}F_{mn}F^{mn}+\frac{1}{2}A_n^T_mF^{mn}\right]\right\}},`$ (15)
where the the gauge part has been integrated away thanks to the Fadeev–Popov determinant $`\mathrm{\Delta }_{FP}=det\mathrm{}`$ and only the gauge invariant $`A^T`$ vector remains in the classical action. The extra Jacobian, coming from the change of the integration measure, will be cancelled in a while when integrating $`F`$. The pay off for relaxing the relationship between $`A`$ and $`F`$ and getting rid of the gauge part is that $`A^T`$ linearly enters the first order action, i.e. $`A^T`$ plays the role of Lagrange multiplier imposing $`F`$ to satisfy the classical field equations
$`Z_{p=1}`$ $`=`$ $`{\displaystyle [DF]\left(det\mathrm{}\right)^{1/2}\delta \left[_mF^{mn}\right]\mathrm{exp}\left\{id^2\sigma \left[\frac{1}{4}F_{mn}F^{mn}\right]\right\}}`$ (16)
$`=`$ $`{\displaystyle [DF]\delta \left[F^{mn}\mathrm{\Lambda }ϵ^{mn}\right]\mathrm{exp}\left\{id^2\sigma \left[\frac{1}{4}F_{mn}F^{mn}\right]\right\}}.`$ (17)
Equation (17) shows that the first order formulation of a limiting rank, abelian gauge theory, and the Fadeev–Popov prescription lead to a “trivial ” path integral for $`F`$. The Dirac–delta picks up the classical configurations of the world tensor $`F`$ and the whole path–integral “collapses” around the classical trajectory The same kind of “collapse” around the classical trajectory has been introduced in string theory to pick up the Eguchi “Area dynamics” . For a pedagogical introduction to this new path–integral manipulation see , where it has been applied to a simpler case, the non–relativistic point particle propagator.. Thus, $`F`$ is “frozen” to a constant value $`\mathrm{\Lambda }`$ and no degrees of freedom are left free to propagate. The same result can be obtained, with some additional work, for $`p>1`$ as well. Accordingly, we get
$`Z`$ $`=`$ $`\mathrm{exp}\{{\displaystyle \frac{1}{p!}}{\displaystyle _\mathrm{\Sigma }}dN_{k_1}\sqrt{\widehat{g}}\widehat{F}^{k_1\mathrm{}k_{p+1}}\widehat{A}_{k_2\mathrm{}k_{p+1}}\}\times `$ (20)
$`\times {\displaystyle }[DF](det\mathrm{})^{1/2}\delta \left[_{m_1}\left(\sqrt{g}F^{m_1m_2\mathrm{}m_{p+1}}\right)\right]\times `$
$`\times \mathrm{exp}\left({\displaystyle \frac{i}{2(p+1)!}}{\displaystyle _\mathrm{\Sigma }}d^{p+1}\sigma \sqrt{g}F_{m_1\mathrm{}m_{p+1}}^{\mathrm{\hspace{0.17em}2}}\right)`$
$`=`$ $`\mathrm{exp}\left\{{\displaystyle \frac{i\mathrm{\Lambda }}{p!}}{\displaystyle _\mathrm{\Sigma }}\widehat{A}_{k_1\mathrm{}k_{p+1}}𝑑s^{k_1}\mathrm{}ds^{k_p}\right\}\mathrm{exp}\left({\displaystyle \frac{i\mathrm{\Lambda }^2}{2}}{\displaystyle _\mathrm{\Sigma }}d^{p+1}\sigma \sqrt{g}\right).`$ (21)
The first term in (21) is a pure boundary quantity produced by a partial integration of the term $`F_[A_]`$. A similar term arises in string theory when boundary and bulk quantum dynamics are properly split .
After integrating out the gauge degrees of freedom the resulting path–integral reads
$`Z`$ $`=`$ $`{\displaystyle ^{\widehat{g}}}[Dg_{mn}]{\displaystyle ^{\widehat{Y}}}[DY^\mu ]\mathrm{exp}({\displaystyle \frac{i\mathrm{\Lambda }}{p!}}{\displaystyle _\mathrm{\Sigma }}\widehat{A}_{k_1\mathrm{}k_{p+1}}ds^{k_1}\mathrm{}ds^{k_p})\times `$ (23)
$`\times \mathrm{exp}\left(i{\displaystyle _\mathrm{\Sigma }}d^{p+1}\sigma \left[{\displaystyle \frac{m_{p+1}}{2}}{\displaystyle \frac{(\gamma )}{\sqrt{g}}}+{\displaystyle \frac{\mathrm{\Lambda }^2}{2}}\sqrt{g}\right]\right)`$
$``$ $`{\displaystyle ^{\widehat{g}}}[Dg_{mn}]{\displaystyle ^{\widehat{Y}}}[DY^\mu ]W_{\widehat{A}}[\mathrm{\Sigma }]\mathrm{exp}\left(i{\displaystyle _\mathrm{\Sigma }}d^{p+1}\sigma \left[{\displaystyle \frac{m_{p+1}}{2}}{\displaystyle \frac{(\gamma )}{\sqrt{g}}}+{\displaystyle \frac{\mathrm{\Lambda }^2}{2}}\sqrt{g}\right]\right).`$ (24)
We remark that this integration procedure is exact and leads to a bulk action plus a boundary correction represented by the generalized Wilson factor $`W_{\widehat{A}}[\mathrm{\Sigma }]`$.
We also notice that the world metric enters the path–integral only through the world volume density. Accordingly, we can “change” integration variable
$$^{\widehat{g}}[Dg_{mn}]^{\widehat{g}}[Dg_{mn}]^{\widehat{e}}[De]\delta \left[e(\sigma )\sqrt{g}\right]$$
(25)
and write (24) as
$$Z=^{\widehat{e}}[De]^{\widehat{Y}}[DY^\mu ]W_{\widehat{A}}[\mathrm{\Sigma }]\mathrm{exp}\left(i_\mathrm{\Sigma }d^{p+1}\sigma \left[\frac{m_{p+1}}{2}\frac{(\gamma )}{e(\sigma )}+\frac{\mathrm{\Lambda }^2}{2}e(\sigma )\right]\right).$$
(26)
The saddle point value for the auxiliary field $`e(\sigma )`$ is defined by:
$$\frac{\delta S}{\delta e(\sigma )}=0e_{cl.}(\sigma )=\frac{1}{\mathrm{\Lambda }}\sqrt{m_{p+1}}\sqrt{\gamma }.$$
(27)
By expanding $`Z`$ around the saddle point $`e_{cl.}(\sigma )`$ we obtain the Dirac–Nambu–Goto path–integral. Correspondingly, we get the following semi–classical equivalence relation
$`Z`$ $`=`$ $`{\displaystyle ^{\widehat{g}}}[Dg_{mn}]{\displaystyle ^{\widehat{Y}}}[DY^\mu ]{\displaystyle ^{\widehat{A}}}[DA]\times `$ (29)
$`\times \mathrm{exp}\left[im_{p+1}{\displaystyle _\mathrm{\Sigma }}d^{p+1}\sigma {\displaystyle \frac{(\gamma )}{2\sqrt{g}}}{\displaystyle \frac{i}{2(p+1)!}}{\displaystyle _\mathrm{\Sigma }}d^{p+1}\sigma \sqrt{g}F_{m_1\mathrm{}m_{p+1}}^{\mathrm{\hspace{0.17em}2}}(A)\right]`$
$``$ $`{\displaystyle [DY^\mu ]W_{\widehat{A}}[\mathrm{\Sigma }]\mathrm{exp}\left[i\rho _p_\mathrm{\Sigma }d^{p+1}\sigma \sqrt{\gamma }\right]}.`$ (30)
The extension of the relation (30) beyond the saddle point approximation is currently under investigation, and requires a proper treatment of the $`p`$–brane degrees of freedom at the quantum level. As is well known, bosonic branes viewed as non–linear $`\sigma `$–models are non–renormalizable perturbative quantum field theories when $`p>1`$. However, we can look at $`p`$–branes not as $`\sigma `$–models but as elements of a more fundamental theory, say $`M`$-Theory, which is in essence a non–perturbative theory. This approach is not new. For example the Einstein–Hilbert action in three and four dimensions is not perturbatively renormalizable; neverthless, three dimensional Einstein–Hilbert gravity can be reformulated as a Chern–Simon gauge theory which can be be exactly solved at the quantum level . In a similar way, four dimensional General Relativity can be written in terms of Ashtekar variables which provides an exact formulation of non-perturbative canonical quantum gravity . From the same point of view, we think that perturbation theory is not the ultimate way to approach the problem of brane quantization. Moreover, a supersymmetric membrane in $`D=11`$ spacetime dimensions is expected to be a finite quantum model , where both ultraviolet and infrared divergences are kept under control. For a general discussion of quantum super-membranes we refer to , and limit our considerations to the semi-classical level. Hopefully, a proper understanding of $`M`$-Theory will provide a background independent formulation of string/brane theory where the quantum path–integral will be well defined. In the meanwhile, we shall work in the WKB approximation where one can choose the action (4), in place of (1) or (2), as a starting point. The non–linearity and reparametrization invariance of the Nambu–Goto action make difficult, if not impossible, to implement the original Feynman construction of the path–integral as a sum of phase space trajectories. One is forced, almost unavoidably, to resort to standard perturbative approaches, e.g. normal modes expansion, or sigma–model effective field theory. Any perturbative approach captures some dynamical feature and misses all the other ones. The impressive results obtained in string theory through duality relations of several kind show that what is not accessible in a given perturbation scheme can be obtained through a different one. With this in mind, we hope that an action of the type (4), where the brane variables enter polynomially, and the tension is brought in by a generalized gauge principle, can be more appropriate to implement the Feynman’s original proposal, or, at least, to provide a different “perturbative” quantization scheme for the Nambu–Goto model itself. In such a, would be, “new regime” of the Nambu–Goto brane both massive and massless objects are present at once and correspond to different values of the world manifold strength $`F_{m_1\mathrm{}m_{p+1}}`$. It would be tempting to assign $`F`$ the role of “order parameter” and describe the dynamical generation of the brane tension as a sort of phase transition. Furthermore, it would be interesting to extend the action (4) in order to include negative tension branes as well. This kind of objects appear to play an important role in the realization of the brane world scenario , .
All these problems, are currently under investigation and eventual results will be reported in future publications.
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# Isospectral deformations of metrics on spheres
## Introduction
To what extent does the spectrum of the Laplacian on a Riemannian manifold determine the geometry of the manifold? The spectrum is known to determine certain global geometric properties such as the dimension, volume and total scalar curvature. Some Riemannian manifolds, such as the round spheres in low dimensions, are known to be uniquely determined by their spectra. Other Riemannian manifolds are known to be infinitesimally spectrally rigid; for example negatively curved metrics on closed manifolds cannot be continuously deformed through a family of isospectral, non-isometric metrics (see \[GuK\] for the two-dimensional case and \[CS\] for the general case). On the other hand, numerous examples of isospectral manifolds show that the spectrum does not always uniquely determine the geometry. The main result of this article is the following theorem (see also the somewhat stronger statements of Corollary 3.10 and Remark 3.11):
###### Theorem
(i) In every dimension $`n8`$, there exist continuous isospectral deformations of Riemannian metrics on the sphere $`S^n`$. The metrics can be chosen to be positively curved; in fact they can be taken to be arbitrarily close to the round metric.
(ii) In every dimension $`n9`$, there exist continuous isospectral deformations of Riemannian metrics on the ball in $`^n`$. The metrics can be chosen arbitrarily close to the flat metric.
The metrics in part (i) of the theorem are the ones induced on the boundary of the ball by the metrics in part (ii). In the case of the balls, we are using the word isospectral to mean that the metrics are both Dirichlet and Neumann isospectral; i.e., the Laplacians acting on functions with Dirichlet or, respectively, Neumann boundary conditions are isospectral.
The metrics on the sphere in (i) provide the first examples of isospectral deformations of positively curved metrics, in marked contrast with the infinitesimal rigidity of negatively curved metrics on closed manifolds. The metrics in the theorem are, to the author’s knowledge, the first examples of isospectral deformations of metrics on balls or spheres.
Z. I. Szabo \[Sz2\] recently gave a method for constructing pairs of isospectral metrics on balls and spheres. Both his technique, which involves a careful analysis of the function spaces, and his examples are distinct from ours.
Prior to 1992, all known isospectral manifolds were locally isometric; they differed only in their global geometry. See, for example, \[BGG\], \[Bu\], \[DG\], \[GWW\], \[GW1, 2\], \[Gt1, 2\], \[I\], \[M\], \[Su\], \[V\] or the expository articles \[Be\], \[Br\], \[G3\], or \[GGt\]. These examples reveal various global invariants which are not spectrally determined such as the diameter and the fundamental group.
In the past eight years, many examples of isospectral manifolds with different local geometry have appeared. Szabo (preprint 1992) used explicit computations to construct the first examples of isospectral manifolds with boundary having different local geometry. The later published version \[Sz1\] also includes closed manifolds, including a pair of isospectral manifolds one of which is homogeneous and the other not. In \[G1, 2\], the author constructed the first isospectral closed manifolds with different local geometry; the second article uses a method based on Riemannian submersions to prove the isospectrality of the metrics. This method was further developed in the series of papers \[GW3\] (giving isospectral deformations of manifolds with boundary), \[GGSWW\] (giving isospectral deformations of closed manifolds), \[Sc1\] (giving isospectral deformations of simply-connected manifolds), \[GSz\] (giving, for example, isospectral deformations of negatively curved manifolds with boundary), and \[Sc2\] (giving, for example, isospectral deformations of left-invariant Riemannian metrics on compact Lie groups). The method produces Riemannian manifolds on which a torus of dimension at least two acts freely by isometries. In order to obtain the deformations in Theorem 0.1, we generalize the technique, weakening the condition that the torus action be free. The beautiful habilitation thesis \[Sc2\] of Dorothee Schueth provided inspiration.
The paper is organized as follows. In $`\mathrm{\S }1`$, we explain the general technique for proving isospectrality by the use of Riemannian submersions. In $`\mathrm{\S }2`$, we first define the class of metrics to be considered on the sphere and ball and observe that each of the metrics admits an isometric torus action. We then show that the open submanifold given by the union of the principal orbits is foliated by Riemannian manifolds isometric to those studied in the earlier papers \[GW3\] and \[GGSWW\]; this fact is used later in $`\mathrm{\S }3`$ to prove that the metrics in the isospectral families are not isometric. In $`\mathrm{\S }3`$, we construct continuous families of isospectral, non-isometric metrics within the class discussed in $`\mathrm{\S }2`$. A reader who is interested only in the construction of the isospectral deformations and not the proofs can restrict attention to Notation 2.1, 3.1, and 3.2, Theorems 3.3 and 3.5, Proposition 3.9, Corollary 3.10 and Remark 3.11.
The author would like to thank Dorothee Schueth for helpful suggestions and to alert the reader to a follow-up article in preparation by Schueth in which she gives an elegant reformulation of the technique developed here and constructs pairs of isospectral metrics on the 6-sphere. The author would also like to thank the Universidad Nacional de Córdoba, and especially Roberto Miatello and Juan Pablo Rossetti, for their hospitality while this paper was being written.
## Section 1. Technique for constructing isospectral manifolds with different local geometry.
###### Definition 1.1 Background and Notation
Let $`T`$ be a torus, let $`\pi :MN`$ be a principal $`T`$-bundle and endow $`M`$ with a Riemannian metric so that the action of $`T`$ is by isometries. Give $`N`$ the induced Riemannian metric so that $`\pi `$ is a Riemannian submersion. For $`aM`$ and $`XT_a(M)`$, write $`X=X^v+X^h`$ where $`X^v`$ is vertical (i.e., tangent to the fiber at $`a`$) and $`X^h`$ is horizontal (i.e., orthogonal to the fiber). Let $`H_aT_a(M)`$ denote the mean curvature vector at $`a`$ of the fiber through $`a`$. Since $`T`$ acts by isometries, we have $`H_{z(a)}=z_{}H_a`$ for all $`aM`$ and $`zT`$. Hence $`H`$ is $`\pi `$-related to a vector field $`\overline{H}`$ on $`N`$. We will refer to $`\overline{H}`$ as the projected mean curvature vector field of the submersion.
Berard-Bergery and Bourguignon \[BB\] gave a decomposition of the Laplacian $`\mathrm{\Delta }_M`$ into vertical and horizontal components $`\mathrm{\Delta }_M=\mathrm{\Delta }^v+\mathrm{\Delta }^h`$. In the case of functions $`f`$ on $`M`$ which are constant on the fibers of the submersion, so $`f=\pi ^{}\overline{f}`$ for some function $`\overline{f}`$ on $`N`$, then $`\mathrm{\Delta }^v(f)=0`$ and
$$\mathrm{\Delta }_M(f)=\mathrm{\Delta }^h(f)=\pi ^{}(\mathrm{\Delta }_N(\overline{f})+\overline{H}(\overline{f})).$$
$`(1.1)`$
If $`N`$ has non-trivial boundary, then $`M=\pi ^1(N)`$. Since $`\pi :MN`$ is a Riemannian submersion, $`\pi ^{}:C^{\mathrm{}}(N)C^{\mathrm{}}(M)`$ maps functions on $`N`$ satisfying Neumann boundary conditions to functions on $`M`$ satisfying Neumann boundary conditions. Of course, the same is true for Dirichlet boundary conditions.
###### Remark Remark
In the situation of Notation 1.1, the space $`\pi ^{}(C^{\mathrm{}}(N))`$ of functions on $`M`$ which are constant on the fibers of the submersion is precisely the space of $`T`$-invariant functions. Since $`T`$ acts by isometries, it follows that $`\pi ^{}(C^{\mathrm{}}(N))`$ is invariant under $`\mathrm{\Delta }_M`$. Except in the case when the projected mean curvature vector field $`\overline{H}`$ is zero, the operator $`\mathrm{\Delta }_N+\overline{H}`$ on $`N`$ is not self-adjoint; however, it does have a discrete spectrum since it is similar to the restriction of $`\mathrm{\Delta }_M`$ to $`\pi ^{}(C^{\mathrm{}}(N))`$ via the (non-unitary) isomorphism $`\pi ^{}:C^{\mathrm{}}(N)\pi ^{}(C^{\mathrm{}}(N))`$.
###### 1.2. Theorem
Let $`T`$ be a torus. Suppose $`T`$ acts by isometries on two compact Riemannian manifolds $`M_1`$ and $`M_2`$ and that the action of $`T`$ on the principal orbits is free. Let $`M_i^{}`$ be the union of all principal orbits in $`M_i`$, so $`M_i^{}`$ is an open submanifold of $`M_i`$ and a principal $`T`$-bundle, $`i=1,2`$. For each subtorus $`K`$ of $`T`$ of codimension at most one, suppose that there exists a diffeomorphism $`\tau _K:M_1M_2`$ which intertwines the actions of $`T`$ and which induces an isometry $`\overline{\tau }_K`$ between the induced metrics on the quotient manifolds $`K\backslash M_1^{}`$ and $`K\backslash M_2^{}`$. Assume further that the isometry $`\overline{\tau }_K`$ satisfies $`\overline{\tau }_K(\overline{H}_K^{(1)})=\overline{H}_K^{(2)}`$, where $`\overline{H}_K^{(i)}`$ is the projected mean curvature vector field for the submersion $`M_i^{}K\backslash M_i^{}`$. (See Notation 1.1.) Then in the case that $`M_1`$ and $`M_2`$ are closed manifolds, they are isospectral. In case $`M_1`$ and $`M_2`$ have boundary, then they are Dirichlet isospectral; under the additional assumption that $`(M_i)M_i^{}`$ is dense in $`(M_i)`$ for $`i=1,2`$, then the manifolds are also Neumann isospectral.
This theorem was proven in \[GSz\] in the special case that the action of $`T`$ on $`M_1`$ and $`M_2`$ is free so that $`M_i^{}=M_i`$. In that case, the hypothesis that the isometry $`\overline{\tau }_K`$ arises from a diffeomorphism $`\tau _K`$ between $`M_1`$ and $`M_2`$ is unnecessary. Earlier versions, used in the papers discussed in the introduction, required that the orbits be totally geodesic.
###### Demonstration Proof
Let $`\mathrm{\Delta }_i`$ denote the Laplacian of $`M_i`$, and let $`\text{L}_{}^2(M_i)`$ denote the space of complex-valued square-integrable functions on $`M_i`$. The torus $`T`$ acts on $`\text{L}_{}^2(M_i)`$, $`i=1,2`$, and by a Fourier decomposition for this action, we have
$$\text{L}_{}^2(M_i)=\mathrm{\Sigma }_{\alpha \widehat{T}}_i^\alpha $$
where $`\widehat{T}`$ consists of all characters on $`T`$, i.e., all homomorphisms from the group $`T`$ to the unit complex numbers, and
$$_i^\alpha =\{f\text{L}_{}^2(M_i):zf=\alpha (z)f\text{ for all }zT\}.$$
The space of $`C^{\mathrm{}}`$ functions on $`M_i`$ decomposes into its intersection with each space $`_i^\alpha `$. Since the torus action on $`M_i`$ is by isometries, the Laplacian leaves each such subspace of smooth functions invariant. If $`M_i`$ has boundary, then the space of smooth functions satisfying Dirichlet, respectively Neumann, boundary conditions similarly decomposes into its intersections with the spaces $`_i^\alpha `$.
Define an equivalence relation on $`\widehat{T}`$ by $`\alpha \beta `$ if $`ker(\alpha )=ker(\beta )`$. Let $`[\alpha ]`$ denote the equivalence class of $`\alpha `$ and let $`[\widehat{T}]`$ denote the set of equivalence classes. Setting
$$_i^{[\alpha ]}=\mathrm{\Sigma }_{\beta [\alpha ]}_i^\beta ,$$
then
$$\text{L}_{}^2(M_i)=\mathrm{\Sigma }_{[\alpha ][\widehat{T}]}_i^{[\alpha ]}.$$
Define $`Q_{[\alpha ]}:_1^{[\alpha ]}_2^{[\alpha ]}`$ by
$$Q_{[\alpha ]}(f)=f\tau _{ker(\alpha )}^1$$
where $`\tau _{ker(\alpha )}`$ is the map whose existence is hypothesized in the theorem, and set
$$Q=_{[\alpha ][\widehat{T}]}Q_{[\alpha ]}.$$
We will show that $`Q`$ intertwines the Laplacians of $`M_1`$ and $`M_2`$.
Let $`\stackrel{~}{}_i^{[\alpha ]}`$ denote the subspace of $`\text{L}_{}^2(M_i^{})`$ obtained by restriction to $`M_i^{}`$ of the elements of $`_i^{[\alpha ]}`$.
For the trivial character $`\alpha =1`$, we have $`[1]=\{1\}`$, and the space $`\stackrel{~}{}_i^{[1]}`$ consists of those functions constant on the orbits of $`T`$. By equation (1.1), the projection $`\pi _i:M_i^{}T\backslash M_i^{}`$ intertwines the restriction of $`\mathrm{\Delta }_i`$ to $`\stackrel{~}{}_i^{[1]}`$ with the operator $`\overline{\mathrm{\Delta }}_i+\overline{H}_T^{(i)}`$ acting on $`\text{L}_{}^2(T\backslash M_i^{})`$, where $`\overline{\mathrm{\Delta }}_i`$ denotes the Laplacian of $`T\backslash M_i^{}`$. The isometry $`\overline{\tau }_T=\overline{\tau }_{ker(\alpha )}`$ gives a unitary isomorphism between $`\text{L}_{}^2(T\backslash M_1^{})`$ and $`\text{L}_{}^2(T\backslash M_2^{})`$ which intertwines $`\overline{\mathrm{\Delta }}_1+\overline{H}_T^{(1)}`$ with $`\overline{\mathrm{\Delta }}_2+\overline{H}_T^{(2)}`$. In case the manifolds have boundary, this unitary isomorphism carries eigenfunctions of the operator $`\overline{\mathrm{\Delta }}_1+\overline{H}_T^{(1)}`$ satisfying either Dirichlet or Neumann boundary conditions to eigenfunctions of $`\overline{\mathrm{\Delta }}_2+\overline{H}_T^{(2)}`$ satisfying the same conditions. This unitary isomorphism pulls back to the isomorphism $`\stackrel{~}{f}\stackrel{~}{f}\tau _T^1`$ from $`\stackrel{~}{}_1^{[1]}`$ to $`\stackrel{~}{}_2^{[1]}`$ satisfying $`\mathrm{\Delta }_2(\stackrel{~}{f}\tau _T^1)=\mathrm{\Delta }_1(\stackrel{~}{f})\tau _T^1`$. Since $`M_i^{}`$ is dense in $`M_i`$, we therefore have that $`\mathrm{\Delta }_2(f\tau _T^1)=\mathrm{\Delta }_1(f)\tau _T^1`$ for smooth functions $`f_1^1`$. Thus the map $`Q_{[1]}`$ defined above intertwines the Laplacians on $`_1`$ and $`_2`$. In the case of manifolds with boundary, Dirichlet boundary conditions are preserved since $`\tau _T`$ is a diffeomorphism. Under the additional hypothesis that $`(M_i)M_i^{}`$ is dense in $`(M_i)`$ for $`i=1,2`$, preservation of Neumann boundary conditions follows from the statement immediately following equation (1.1). Thus $`Q_{[1]}`$ carries eigenfunctions satisfying either boundary condition to eigenfunctions satisfying the same condition.
For non-trivial $`\alpha \widehat{T}`$, the kernel of $`\alpha `$ is a subtorus $`K`$ of $`T`$ of codimension one. The space of all functions on $`M_i^{}`$ constant on the fibers of the submersion $`\pi _i:M_i^{}K\backslash M_i^{}`$ is given by $`\stackrel{~}{}_i^{[\alpha ]}\stackrel{~}{}_i^{[1]}`$. Arguing as in the previous paragraph, we see that the restrictions of the Laplacians of $`M_1`$ and $`M_2`$ to the subspaces $`_1^{[\alpha ]}_1^{[1]}`$ and, respectively, $`_2^{[\alpha ]}_2^{[1]}`$ are intertwined by the map $`ff\tau _K^1`$. Moreover, since $`\tau _K`$ intertwines the actions of $`T`$, this intertwining map carries $`_1^{[\alpha ]}`$ to $`_2^{[\alpha ]}`$; it’s restriction to this subspace is precisely $`Q_{[\alpha ]}`$. Thus $`Q`$ intertwines the Laplacians of $`M_1`$ and $`M_2`$. Preservation of the appropriate boundary conditions follows exactly as in the previous paragraph.
## Section 2. Construction of the metrics on the ball and sphere
###### Definition 2.1 Notation
(i) Let $`T`$ denote the torus $`^k\backslash ^k`$. The Lie algebra $`𝔷`$ of translation invariant vector fields on $`T`$ is canonically identified with $`^k`$. Define a representation $`\rho :TSO(2k)`$ of $`T`$ on $`^{2k}`$ by diagonally embedding $`T`$ into $`SO(2k)`$ as the direct product of $`k`$ copies of the circle $`SO(2)`$. Thus $`\rho _{}:𝔷𝔰𝔬(2k)`$ is a representation of the Lie algebra $`𝔷`$. For $`Z𝔷`$, let $`Z^{}`$ denote the fundamental vector field on $`^{2k}`$ given by $`Z_u^{}=\rho _{}(Z)(u)`$ for $`u^{2k}`$. Let $`𝔷^{}=\{Z^{}:Z𝔷\}`$.
Given an alternating bilinear map $`B:^m\times ^m𝔷`$, define $`g_B`$ to be the unique Riemannian metric on $`^{m+2k}=^m\times ^{2k}`$ satisfying the following three conditions:
(a) the canonical projection $`(^m\times ^{2k},g_B)^m`$ is a Riemannian submersion where the base manifold has the standard Euclidean metric;
(b) the metric induced by $`g_B`$ on each fiber $`^{2k}`$ is the standard Euclidean metric;
(c) the horizontal space at each point $`(x,u)^m\times ^{2k}`$ is given by $`_{(x,u)}=\{y+\frac{1}{2}B(x,y)_u^{}:yT_x(^m)\}`$. (Here we are using the canonical identification of $`T_x(^m)`$ with $`^m`$ to define $`B(x,y)`$. The coefficient $`\frac{1}{2}`$ is included for later convenience.)
The action $`\rho `$ of the torus $`T`$ on $`^{2k}`$ gives rise to an action of $`T`$ by isometries on $`(^{m+2k},g_B)`$ preserving each fiber. The action is not free. However, $`T`$ does act freely on the principal orbits; these fill an open dense subset of $`^{m+2k}`$.
(ii) Let $`D`$, respectively $`S`$, denote the closed unit ball, respectively unit sphere, in $`^{m+2k}`$ relative to the Euclidean inner product. We continue to denote by $`g_B`$ both the restriction of $`g_B`$ to the bounded domain $`D`$ and the induced Riemannian metric on $`S`$. For $`zT`$, the associated isometry of $`(^{m+2k},g_B)`$ leaves both $`S`$ and $`D`$ invariant. Thus $`T`$ acts by isometries on $`(D,g_B)`$ and $`(S,g_B)`$ with the action on the principal orbits being free.
In $`\mathrm{\S }3`$ we will construct continuous families of anti-symmetric bilinear maps $`B:^m\times ^m𝔷`$ so that the associated Riemannian metrics $`g_B`$ on $`S`$ and on $`D`$ are pairwise isospectral but not isometric. The isospectrality proof will be an easy consequence of Theorem 1.2. To prepare for the proof in $`\mathrm{\S }3`$ that the deformations are non-trivial, we now give an alternate description of the metric $`g_B`$.
###### Definition 2.2 Notation and Remarks
(i) Denote elements of the torus $`T=^k\backslash ^k`$ by $`\overline{z}`$, with $`z^k`$. Given an alternating bilinear map $`B:^m\times ^m^k`$, endow $`^m\times T`$ with the structure of a two-step nilpotent Lie group by defining the group multiplication as
$$(x,\overline{z})(x^{},\overline{z}^{})=(x+x^{},\overline{z}+\overline{z}^{}+\frac{1}{2}\overline{B(x,x^{})}).$$
We will denote this Lie group by $`G_B`$. Letting $`𝔷=^k`$ be the Lie algebra of $`T`$, then the Lie algebra $`𝔤_B`$ of $`G_B`$ is given by $`𝔤_B=^m+𝔷`$ with Lie bracket $`B:^m\times ^m𝔷`$. In particular, $`𝔷`$ is central and contains the derived algebra.
(ii) The projection $`\pi :G_B^m`$ gives $`G_B`$ the structure of a principal $`T`$-bundle. Let
$$E_B=G_B\times _T^{2k}$$
be the bundle over $`^m`$ with fiber $`^{2k}`$ associated to this principal $`T`$-bundle via the action $`\rho `$ of $`T`$ on $`^{2k}`$ defined in 2.1. Elements of $`E_B`$ are equivalence classes $`[((x,\overline{z}),u)]`$, with $`(x,\overline{z})G_B`$ and $`u^{2k}`$, under the equivalence relation $`[((x,\overline{z+z^{}}),u)]=[((x,\overline{z}),\overline{z}^{}u)]`$, where $`\overline{z}^{}u`$ denotes $`\rho (\overline{z}^{})(u)`$. Since $`\pi :G_B^m`$ is a trivial bundle, $`E_B`$ is diffeomorphic to $`^{m+2k}`$.
The left-invariant vector fields on $`G_B`$ define vector fields on the product $`G_B\times ^{2k}`$ which are invariant under the diagonal action of $`T`$. Thus they induce vector fields on $`E_B`$. Define a Riemannian metric on $`E_B`$ so that the projection $`\stackrel{~}{\pi }:E_B^m`$ is a Riemannian submersion, the fibers have the standard Euclidean metric, and the horizontal distribution is given by the space of left-invariant vector fields on $`G_B`$ lying in the subspace $`^m`$ of $`𝔤_B`$, viewed as vector fields on $`E_B`$.
###### 2.3 Proposition
In the notation of 2.1 and 2.2, $`(^{m+2k},g_B)`$ is isometric to $`E_B`$.
###### Demonstration Proof
Define a bundle diffeomorphism $`\tau :E_B^{m+2k}`$ by $`\tau ([((x,\overline{z}),u)])=(x,\overline{z}u)`$. To see that $`\tau `$ is an isometry between the metric on $`E_B`$ defined in 2.2 and the metric $`g_B`$ on $`^{m+2k}`$, we need only show that $`\tau `$ carries the horizontal space of $`E_B`$ at $`[((x,\overline{z}),u)]`$ to that of $`(^{m+2k},g_B)`$ at $`(x,\overline{z}u)`$.
Each element $`u`$ of $`𝔤_B=^m+𝔷`$ defines two vector fields on $`G_B`$: a left-invariant vector field, which we temporarily denote by the corresponding upper case letter $`U`$, and a directional derivative $`D_u`$ given by ignoring the group structure on $`G_B`$ and viewing $`G_B`$ as $`^m\times T`$. For the central elements $`z𝔷`$ of $`𝔤_B`$, the two vector fields coincide. However, for $`y^m`$, the expression for multiplication in $`G_B`$ given in 2.2(i) shows that the vector fields $`Y`$ and $`D_y`$ are related by $`Y_{(x,\overline{z})}=D_{y+\frac{1}{2}B(x,y)}`$. Thus letting $`\stackrel{~}{Y}`$ denote the vector field on $`E_B`$ associated with the left-invariant vector field $`Y`$ and identifying $`y`$ with an element of $`T_x(^m)`$, we have $`\tau _{[((x,\overline{z}),u)]}(\stackrel{~}{Y})=y+\frac{1}{2}(B(x,y))_{\overline{z}u}^{}`$. Hence $`\tau `$ carries the horizontal space of $`E_B`$ at $`[((x,\overline{z}),u)]`$ to that of $`(^{m+2k},g_B)`$ at $`(x,\overline{z}u)`$, so $`\tau `$ is an isometry.
###### Definition 2.4 Notation
(i) Left-invariant Riemannian metrics on the Lie group $`G_B`$ defined in 2.2 correspond to inner products on its Lie algebra $`𝔤_B`$. Given an inner product $`h`$ on $`𝔷`$, extend $`h`$ to an inner product on $`𝔤_B`$ by taking the orthogonal direct sum with the standard inner product on $`^m`$. Denote by $`g_h`$ the associated left-invariant Riemannian metric on $`G_B`$. . The projection $`\pi :G_B^m`$ is a Riemannian submersion from the metric $`g_h`$ to the canonical metric on $`^m`$.
Recall that $`𝔷`$ may be canonically identified with $`^k`$. Given $`a=(a_1,\mathrm{},a_k)(^+)^k`$, define an inner product $`h_a`$ on $`^k`$ so that the standard basis $`\{Z_1,\mathrm{},Z_k\}`$ of $`^k`$ is orthogonal and so that $`Z_i_a=a_i`$, where $`_a`$ is the norm associated with $`h_a`$. We will write $`g_a`$ to mean $`g_{h_a}`$.
(ii) Let $`S_r`$, respectively $`U_r`$, denote the geodesic sphere, respectively geodesic ball, of radius $`r`$ in the Euclidean space $`^m`$, and let $`N_r(B,h)=\pi ^1(S_r)`$ and $`Q_r(B,h)=\pi ^1(U_r)`$ with the Riemannian metrics induced by $`g_h`$. Thus we have Riemannian submersions $`\pi :N_r(B,h)S_r`$ and $`\pi :Q_r(B,h)U_r`$. In case $`h`$ is the standard inner product on $`𝔷=^k`$, we will write $`N_r(B)`$ and $`Q_r(B)`$, suppressing the name of the inner product.
The manifolds $`N_r(B,h)`$ and $`Q_r(B,h)`$ were studied in \[GGSWW\] and \[GW3\], respectively.
###### 2.5 Proposition
We use Notation 2.4. For $`c^+`$, the map $`\mu :(G_B,g_{c^2h})(G_{cB},g_h)`$ given by $`\mu (x,\overline{z})=(x,c\overline{z})`$ is both an isometry and a Lie group isomorphism. Moreover, $`\mu `$ restricts to isometries $`N_r(B,c^2h)N_r(cB,h)`$ and $`Q_r(B,c^2h)Q_r(cB,h)`$.
We omit the elementary proof.
We next describe a foliation of the dense open subest $`S^{}`$ (respectively, $`D^{}`$) of the sphere (respectively, ball) on which $`T`$ acts freely.
###### 2.6 Proposition
Denote elements of $`^{2k}`$ by $`u=(u_1,\mathrm{},u_k)`$ with $`u_i^2`$ for each $`i`$. Let $`u^{2k}`$ and let $`a_i=u_i`$, $`i=1,\mathrm{},k`$, where $``$ is the Euclidean norm on $`^2`$. Then:
(i) Under the action $`\rho `$ of $`T`$ on $`^{2k}`$ defined in 2.1, the orbit of $`T`$ through $`u`$ is given by $`Tu=\{v:v_i=a_i,i=1,\mathrm{},k\}`$ .
(ii) If $`a_i0`$ for all $`i`$, then the $`T`$-saturated submanifold $`L(a):=^m\times (Tu)`$ of $`^{m+2k}`$ with the metric induced by $`g_B`$ is isometric to $`(G_B,g_a)`$ where $`a=(a_1,\mathrm{},a_k)`$. If $`u^2=\mathrm{\Sigma }_{i=1}^ka_i^2<1`$, then the intersection $`S(a):=L(a)S`$ with the sphere $`S`$ in $`^{m+2k}`$ is isometric to $`N_r(B,h_a)`$, where $`r=\sqrt{1u^2}`$, and the intersection $`D(a):=L(a)D`$ is isometric to $`Q_r(B,h_a)`$.
###### Demonstration Proof
(i) is elementary. For (ii), let $`\tau :E_B(^{m+2k},g_B)`$ be the isometry defined in the proof of Proposition 2.3. Then the inverse image of $`L(a)`$, given by $`\tau ^1(L(a))=\{[((x,\overline{z}),u)]:(x,\overline{z})G_B\}`$, with the Riemannian metric induced by the metric on $`E_B`$ is canonically isometric to $`(G_B,g_a)`$. The final statement of (ii) follows easily.
As $`a`$ varies, the manifolds $`S(a)`$ and $`D(a)`$ foliate $`S^{}`$ and $`D^{}`$, respectively. These foliations will play a key role in the proof of non-triviality of the isospectral deformations in $`\mathrm{\S }3`$. There we will show that any isometry between the metrics on the sphere (respectively, ball) must induce an isometry between the metrics on a suitable leaf. We will then appeal to \[GW3\] (respectively, \[GGSWW\]) to see that the metrics on the leaf are not isometric, thus obtaining a contradiction.
## Isospectral deformations
###### Definition 3.1 Notation and Remarks
Let $`𝔷`$ be a finite-dimensional vector space and $`B:^m\times ^m𝔷`$ an alternating bilinear map. Let $`,`$ be the standard inner product on $`^m`$ and $`h`$ an inner product on $`𝔷`$. We then obtain a linear map $`j:𝔷𝔰𝔬(m)`$ by the condition
$$h(B(x,y),z)=j(z)x,y$$
$`(3.1)`$
for all $`x,y^m`$ and all $`z𝔷`$. Conversely, given a linear map $`j:𝔷𝔰𝔬(m)`$ and an inner product $`h`$ on $`𝔷`$, then equation (3.1) defines an alternating bilinear map $`B:^m\times ^m𝔷`$.
We consider $`𝔷`$ as in 2.1, so that $`𝔷`$ is canonically identified with $`^k`$. Given $`j:𝔷𝔰𝔬(m)`$, we let $`B_j:^m\times ^m𝔷`$ be the bilinear map associated with $`j`$ as in equation (3.1), taking $`h`$ to be the standard inner product on $`𝔷=^k`$. We will use the abbreviated notation $`g_j`$ for the Riemannian metric $`g_{B_j}`$ on $`^{m+2k}`$ defined in Notation 2.1(i).
###### Definition 3.2. Definition
(i) Let $`𝔷`$ be a vector space. A pair $`j,j^{}`$ of linear maps from $`𝔷`$ to $`𝔰𝔬(m)`$ will be called isospectral, denoted $`jj^{}`$, if for each $`z𝔷`$, the eigenvalue spectra, with multiplicities, of $`j(z)`$ and $`j^{}(z)`$ coincide; i.e., for each $`z𝔷`$, there exists an orthogonal linear operator $`A_z`$ for which
$$A_zj(z)A_z^1=j^{}(z).$$
(ii) Let $`(𝔷,h)`$ be an inner product space. A pair $`j,j^{}`$ of linear maps from $`𝔷`$ to $`𝔰𝔬(m)`$ will be called equivalent, denoted $`jj^{}`$, if there exist orthogonal linear maps $`A`$ of $`^m`$ and $`C`$ of $`𝔷`$ such that
$$Aj(z)A^1=j^{}(C(z))$$
for all $`z𝔷`$.
###### 3.3 Theorem
Let $`𝔷=^k`$ and let $`j,j^{}`$ be isospectral linear maps from $`𝔷`$ to $`𝔰𝔬(m)`$ as in 3.2. In the notation of 2.1 and 3.1, the Riemannian metrics $`g_j`$ and $`g_j^{}`$ on the ball $`D`$ in $`^{m+2k}`$ are both Dirichlet and Neumann isospectral, and the metrics $`g_j`$ and $`g_j^{}`$ on the sphere $`S`$ in $`^{m+2k}`$ are isospectral.
###### 3.4 Lemma
Let $`j:𝔷𝔰𝔬(m)`$ be a linear map. We use the notation of 2.1 and 3.1; in particular, $`^{m+2k}`$ fibers over $`^m`$ with fibers invariant under the action of the torus $`T`$. Then relative to the metric $`g_j`$ on $`^{m+2k}`$ defined in 3.1, the fibers are totally geodesic Euclidean submanifolds. Consequently, if $`K`$ is any subtorus of $`T`$ and if $`(x,u)^{m+2k}`$, then the mean curvature of the orbit $`K(x,u)=(x,Ku)`$ in $`^{m+2k}`$ is completely determined by the mean curvature of the torus $`Ku`$ in $`^{2k}`$, independently of $`j`$.
We omit the straightforward proof of the lemma.
###### Demonstration Proof of Theorem 3.3
We will apply Theorem 1.2. Let $`(^{m+2k})^{}`$ be the union of the principal orbits for the action of $`T`$ on $`^{m+2k}`$ defined in 2.1. Then $`T`$ acts freely on $`(^{m+2k})^{}`$ and for any subtorus $`K`$ of $`T`$, the quotient $`K\backslash (^{m+2k})^{}`$ is diffeomorphic to $`^m\times (K\backslash (^{2k})^{})`$, where $`(^{2k})^{}`$ is the union of the principal orbits for the action $`\rho `$ of $`K`$ on $`^{2k}`$ in 2.1 (i.e., for the restriction to $`K`$ of the action $`\rho `$ of $`T`$). Let $`B=B_j`$ and $`B^{}=B_j^{}`$. For $`x,y^m`$, the fundamental vector field $`B(x,y)^{}`$ on $`^{2k}`$, defined as in 2.1, induces a vector field on $`K\backslash (^{2k})^{}`$ which we denote by $`\overline{B}(x,y)`$. We define $`\overline{B}^{}(x,y)`$ similarly.
The metric $`\overline{g}_j`$ on $`K\backslash (^{m+2k})^{}`$ induced by $`g_j`$ is determined by the following three properties:
(a) the projection $`K\backslash (^{m+2k})^{}^m`$ is a Riemannian submersion relative to the canonical metric on the base;
(b) the metric on the fibers $`K\backslash (^{2k})^{}`$ is the quotient metric induced by the canonical metric on the open subset $`(^{2k})^{}`$ of $`^{2k}`$; and
(c) the horizontal space at $`(x,\overline{u})K\backslash (^{m+2k})^{}`$, where $`x^m`$ and $`\overline{u}K\backslash (^{2k})^{}`$, is given by $`\{y+\frac{1}{2}\overline{B}(x,y)_{\overline{u}}:yT_x(^m)\}`$, where we use the canonical identification of $`T_x(^m)`$ with $`^m`$ in order to define $`\overline{B}(x,y)`$. The metric $`\overline{g}_j^{}`$ induced by $`g_j^{}`$ is defined analogously.
First consider the case $`K=T`$. In this case, the metrics $`\overline{g}_j`$ and $`\overline{g}_j^{}`$ on $`T\backslash (^{m+2k})^{}`$ are identical. Let $`M`$ denote either the ball $`D`$ or the sphere $`S`$. Taking $`\tau _T`$ to be the identity map on $`M`$, then $`\tau _T`$ satisfies the hypothesis of Theorem 1.2.
Next suppose that $`K`$ has codimension one in $`T`$. Let $`Z𝔷`$ be a non-zero vector orthogonal to the Lie algebra $`𝔨`$ of $`K`$. By the hypothesis that $`j`$ and $`j^{}`$ are isospectral maps, there exists an orthogonal transformation $`A=A_Z`$ of $`^m`$ such that $`j^{}(Z)=Aj(Z)A^1`$. Equivalently,
$$B^{}(Ax,Ay),Z=B(x,y),Z$$
for all $`x,y^m`$. Consequently,
$$\overline{B}^{}(Ax,Ay)=\overline{B}(x,y)$$
$`(3.2)`$
for all $`x,y^m`$.
From equation (3.2) and the defining properties (a)-(c) of the metrics $`\overline{g}_j`$ and $`\overline{g}_j^{}`$, we conclude that the restriction to $`(^{m+2k})^{}`$ of the diffeomorphism $`\tau _K:^{m+2k}^{m+2k}`$ given by $`\tau _K((x,u))=(A(x),u)`$ induces an isometry $`\overline{\tau }_K:(K\backslash (^{m+2k})^{},\overline{g}_j)(K\backslash (^{m+2k})^{},\overline{g}_j^{})`$. Since $`\tau _K`$ restricts to the identity on the fibers of the submersion onto $`^m`$, Lemma 3.4 shows that it intertwines the mean curvature vector fields for the metrics $`g_j`$ and $`g_j^{}`$. Now considering the restriction of $`\tau _K`$ to either the ball $`D`$ or the sphere $`S`$, we see that the hypothesis of Theorem 1.2 holds. Thus the metrics $`g_j`$ and $`g_j^{}`$ viewed either on $`D`$ or on $`S`$ are isospectral.
###### Theorem 3.5
Let $`M`$ denote either the ball $`D`$ or the sphere $`S`$ in $`^{m+2k}`$. Suppose that $`j:𝔷𝔰𝔬(m)`$ satisfies the property that there are only finitely many orthogonal maps of $`^m`$ which commute with all the transformations $`j(Z)`$, $`Z𝔷`$. In the notation of 2.1, 3.1 and 3.2, if $`j^{}:𝔷𝔰𝔬(m)`$ is any linear map for which the metrics $`g_j`$ and $`g_j^{}`$ on $`M`$ are isometric, then $`jj^{}`$.
The hypothesis holds for generic maps $`j`$.
###### Lemma 3.6
Let $`j,j^{}:𝔷𝔰𝔬(m)`$ be linear maps. In the notation of 2.4(ii):
(a) If $`Q_r(B_j)`$ is isometric to $`Q_r(B_j^{})`$, then $`jj^{}`$.
(b) If $`j`$ satisfies the genericity condition in Theorem 3.5 and if $`N_r(B_j)`$ is isometric to $`N_r(B_j^{})`$, then $`jj^{}`$.
Parts (a) and (b) of the lemma are proven in \[GW3\] and \[GGSWW\], respectively. The genericity condition on $`j`$ is stated differently in \[GGSWW\], Proposition 10. To clarify, let $`𝔤_j=𝔤_{B_j}`$ be the Lie algebra with bracket $`B_j`$ defined in 2.2. Any orthogonal linear map $`\alpha `$ of $`^m`$ which commutes with all the transformations $`j(Z)`$, $`Z𝔷`$, extends to an orthogonal automorphism of $`𝔤_j`$ by defining $`\alpha `$ to be the identity map on $`𝔷`$. (Here we are using the word “orthogonal” to mean that $`\alpha `$ preserves the standard inner product on $`𝔤_j`$.) Conversely, every orthogonal automorphism of $`𝔤_j`$ which restricts to the identity on $`𝔷`$ must be of this form. Thus the genericity condition in Theorem 3.5 is equivalent to the condition that there are only finitely many orthogonal automorphisms of $`𝔤_j`$ which restrict to the identity on $`𝔷`$. The latter condition is slightly weaker than the genericity condition used in \[GGSWW\] in that the word “orthogonal” has been inserted. However, a glance at the arguments in \[GGSWW\] show that only this weaker condition is actually used.
###### Lemma 3.7
Let $`M`$ denote $`D`$, respectively $`S`$, and let $`j,j^{}:𝔷𝔰𝔬(m)`$ be any linear maps. Suppose $`\tau :(M,g_j)(M,g_j^{})`$ is an isometry which carries $`T`$ orbits to $`T`$ orbits. Let $`a=(a_1,\mathrm{},a_k)(^+)^k`$ satisfy $`a<1`$ and $`a_1=\mathrm{}=a_k`$, and let $`c=\sqrt{a_i}`$. In the notation of Proposition 2.6, let $`M(a)`$ denote $`D(a)`$, respectively $`S(a)`$. Then $`\tau `$ leaves $`M(a)`$ invariant. Thus by Proposition 2.6(ii) and Proposition 2.5, the restriction of $`\tau `$ to $`M(a)`$ gives an isometry between $`Q_r(cB_j)`$ and $`Q_r(cB_j^{})`$, respectively between $`N_r(cB_j)`$ and $`N_r(cB_j^{})`$.
###### Demonstration Proof of Lemma 3.7
For $`(x,u)M`$ with $`u=(u_1,\mathrm{},u_k)`$, the $`T`$ orbit $`T(x,u)=(x,Tu)`$ is isometric to the direct product of circles of radii $`|u_1|,\mathrm{},|u_k|`$. Thus $`\tau `$ must carry the orbit $`T(x,u)`$ to an orbit $`T(y,v)`$ such that $`(|u_1|,\mathrm{}|u_k|)=(|v_1|,\mathrm{}|v_k|)`$ up to permutation. Since for $`a`$ as in the lemma, $`M(a)`$ is the union of all those $`T`$-orbits $`T(x,u)`$ for which $`|u_1|=\mathrm{}=|u_k|=c^2`$, it follows that $`\tau `$ leaves $`M(a)`$ invariant.
###### Lemma 3.8
Suppose that $`j`$ satisfies the hypothesis of Theorem 3.5. Then $`T`$ is a maximal torus in the full isometry group of $`(M,g_j)`$.
###### Demonstration Proof of Lemma 3.8
The proof of the Lemma will be based on the following analogous statements for the manifolds $`N_r(B_j)`$ and $`Q_r(B_j)`$: (These statements also use the genericity hypothesis on $`j`$.)
(i) $`T`$ is a maximal torus in the full isometry group of $`N_r(B_j)`$.
(ii) $`T`$ is a maximal torus in the full isometry group of $`Q_r(B_j)`$.
(i) was proven in \[GGSWW\] (see the proof of Proposition 10 there). (ii) can be seen either by a direct argument or by appealing to (i) as follows: Suppose $`\tau `$ is an isometry of $`Q_r(B_j)`$ which commutes with the action of $`T`$. Then $`\tau `$ restricts to an isometry of the boundary $`N_r(B_j)`$ also commuting with the action of $`T`$. An isometry $`\tau `$ of a connected manifold $`Q`$ is uniquely determined by its value and differential at a single point. Consequently, it is also determined up to two possibilities by its restriction to any submanifold $`N`$ of codimension one. Indeed, for $`pN`$, the restriction of $`\tau `$ to $`N`$ determines both $`\tau (p)`$ and the restriction of $`\tau _p`$ to $`T_p(N)`$. Since $`\tau _p`$ is an inner product space isometry, it is moreover determined up to sign on the orthogonal one dimensional subspace of $`T_p(Q)`$. (In the special case under consideration in which $`N`$ is the boundary of $`Q`$, $`\tau _p`$ is in fact uniquely determined by its restriction to $`N`$.) Consequently (ii) follows from (i).
We now prove the lemma. By Lemma 3.7, $`\tau `$ leaves $`M(a)`$ invariant and its restriction to $`M(a)`$ gives an isometry of $`Q_r(cB_j)`$ (if $`M=D`$) or of $`N_r(cB_j)`$ (if $`M=S`$). Moreover this isometry commutes with the action of $`T`$. By 3.1, $`N_r(cB_j)`$ is identical to $`N_r(B_{cj})`$. Note that $`cj`$ also satisfies the genericity hypothesis of Theorem 3.5. Thus by statement (i) or (ii) above, $`\tau _{|M(a)}`$, is determined up to countably many possibilities modulo composition with elements of $`T`$.
We now use an argument similar to the proof of (ii) above to show that $`\tau `$ is determined up to at most two possibilities by its restriction to $`M(a)`$. Fix a point $`p`$ of $`M(a)`$. The tangent space $`T_p(M(a))`$ has co-dimension $`k`$ in $`T_p(M)`$. First consider the case that $`M=D`$. Each of the fundamental vector fields $`Z^{}`$, $`Z𝔷`$, defined in 2.1 is $`\tau `$ invariant since $`\tau `$ commutes with $`T`$. Hence $`\tau `$ leaves $`_Z^{}Z^{}`$ invariant for each $`Z`$. As $`Z`$ ranges over $`𝔷`$, the vectors $`_Z^{}Z^{}(p)`$ span the orthogonal complement of $`T_p(M(a))`$ in $`T_p(M)`$. Thus $`\tau _p`$ is uniquely determined by its restriction to $`T_p(M(a))`$ in this case. Next in the case that $`M=S`$, we argue the same way, but now the vectors $`_Z^{}^SZ^{}(p)`$ span only a $`(k1)`$-dimensional subspace of $`T_p(S)`$ orthogonal to $`T_p(M(a))`$. Hence from the restriction of $`\tau `$ to $`M(a)`$, we know the restriction of $`\tau _p`$ to a subspace of co-dimension one in $`T_p(S)`$. As in the proof of statement (ii) above, $`\tau _p`$ is determined up to sign on the orthogonal one dimensional subspace. Thus $`\tau `$ is determined up to at most two choices by its restriction to $`M(a)`$. The lemma follows.
###### Demonstration Proof of Theorem 3.5
Suppose $`\tau :(M,g_j)(M,g_j^{})`$ is an isometry. Then it must conjugate the maximal torus $`T`$ in the isometry group of $`(M,g_j)`$ to that in the isometry group of $`(M,g_j^{})`$. We conclude that the maximal torus in the isometry group of $`(M,g_j^{})`$ has the same dimension as $`T`$, so $`T`$ itself is a maximal torus. Since any two maximal tori in the isometry group are conjugate, we may assume, after composing $`\tau `$ with an isometry of $`(M,g_j^{})`$, that $`\tau `$ conjugates the torus $`T`$ in the isometry group of $`(M,g_j)`$ to the torus $`T`$ in the isometry group of $`(M,g_j^{})`$. Thus the isometry $`\tau `$ carries $`T`$ orbits to $`T`$ orbits. By Lemma 3.7, it follows that $`N_r(cB_j)`$ (or $`Q_r(cB_j)`$) is isometric to $`N_r(cB_j^{})`$ (respectively, to $`Q_r(cB_j^{})`$). From \[GGSWW\] (respectively, \[GW3\]), we conclude that $`jj^{}`$.
###### 3.9 Proposition \cite{GW3}
Let $`𝔷`$ be an inner product space with $`dim(𝔷)=2`$, and let $`m`$ be any positive integer other than $`1,2,3,4`$, or $`6`$. Let $`W_m`$ be the real vector space consisting of all linear maps from $`𝔷`$ to $`𝔰𝔬(m)`$. Then there is a Zariski open subset $`𝒪_m`$ of $`W_m`$ (i.e., $`𝒪_m`$ is the complement of the zero locus of some non-zero polynomial function on $`W`$) such that each $`j𝒪_m`$ belongs to a $`d`$-parameter family of isospectral, inequivalent elements of $`W_m`$. Here $`d\frac{m(m1)}{2}[\frac{m}{2}]([\frac{m}{2}]+2)>1`$. In particular, $`d`$ is of order at least $`O(m^2)`$.
The families of isospectral, inequivalent $`j`$-maps may be parameterized so that $`j`$ depends smoothly on the parametrization. Although the proposition does not give any information when $`m=6`$, a specific example of a smooth family of isospectral, inequivalent $`j`$-maps was given in \[GW3\] when $`m=6`$.
###### 3.10 Corollary
There exist continuous families of isospectral, non-isometric Riemannian metrics on the ball of dimension $`n`$ and on the sphere of dimension $`n1`$ for each $`n9`$. The metrics depend smoothly on the parameter. Given $`ϵ>0`$, the metrics on the ball can be chosen so that the sectional curvature $`K`$ satisfies $`|K|<ϵ`$ and the metrics on the sphere can be chosen so that $`1ϵ<|K|<1+ϵ`$. In particular, there exist isospectral deformations of positively curved metrics.
###### Demonstration Proof
One of the defining properties of the Zariski set $`𝒪_m`$ in \[GW3\] is that the elements $`j`$ satisfy the hypothesis of Theorem 3.5. This condition also holds for the specific example in which $`m=6`$. Thus the first statement of the corollary follows immediately from Theorems 3.3 and 3.5. For the curvature statement, fix $`j`$ and consider the family of metrics $`g_{cj}`$, $`c^+`$, on $`^{m+2k}`$. (See Notation 2.1 and 3.1). As $`c0^+`$, this family of metrics converges to the Euclidean metric on $`^{m+2k}`$, and the induced family of metrics on $`S`$ converges to the round metric. Given a smooth family $`\{j_t\}`$ of isospectral, inequivalent maps, then $`\{cj_t\}`$ is also an isospectral family of inequivalent maps for any $`c^+`$. By taking $`c`$ small enough, we can obtain any desired curvature bounds for the metrics associated with $`cj_t`$ for all $`t`$ in a compact subset of the parameter space.
###### Remark 3.11 Remarks
We have actually shown that for each $`j`$ in the Zariski open set $`𝒪_m`$ in Proposition 3.9 and for each $`c>0`$, the Riemannian metric $`g_{cj}`$ on the ball, respectively sphere, in $`^{m+4}`$ lies in a continuous $`d`$-parameter family $`(cj)`$ of isospectral, non-isometric metrics, where $`d`$ is at least of order $`O(m^2)`$. Moreover, as $`c0^+`$, the metric $`g_{cj}`$ converges to the flat, respectively round, metric.
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# Real-time nonequilibrium dynamics in hot QED plasmas: dynamical renormalization group approach
## I Introduction
The study of nonequilibrium phenomena under extreme conditions play a fundamental role in the understanding of ultrarelativistic heavy ion collisions and early universe cosmology. Forthcoming relativistic heavy ion experiments at BNL Relativistic Heavy Ion Collider (RHIC) and CERN Large Hadron Collider (LHC) aim to search for a deconfined phase of quarks and gluons, the quark gluon plasma (QGP), which is predicted by lattice QCD simulations to emerge at a temperature scale $`T200\text{MeV}`$.
Recent results from CERN Super Proton Synchrotron (SPS) seem to confirm the main theoretical ideas that in the central region in ultrarelativistic heavy ion collision a deconfined plasma of quarks and gluons forms which expands and cools rapidly and eventually hadronizes.
Current estimates based on energy deposited in the central region for $`\sqrt{s}/\text{nucleon-pair}200\mathrm{GeV}`$ at RHIC suggest that the lifetime of a deconfined QGP is of order $`1050\text{fm}/c`$ . At such unprecedented short time scales, an important aspect is an assessment of thermalization time scales and the potential for non-equilibrium effects associated with the rapid expansion and finite lifetime of the plasma and their impact on experimental observables. Lattice QCD is simply unable to deal with these questions because simulations are restricted to thermodynamic equilibrium quantities and a field-theoretical nonequilibrium approach is needed for an accurate description of the formation and evolution of the QGP. An important and pioneering step in this direction was undertaken by Geiger who applied transport methods combined with perturbative QCD (pQCD) to obtain a quantitative picture of the evolution of partons in the early stages of formation and evolution of the plasma.
The consistent study of the evolution of partons in terms of pQCD cross sections that include screening corrections to avoid the infrared divergences associated with small angle scattering lead to the conclusion that quarks and gluons thermalize on time scales of a few $`\text{fm}/c`$ . The necessity of a deeper understanding of equilibrium and non-equilibrium aspects of the quark-gluon plasma motivated an intense study of the abelian and non-abelian plasmas in extreme environments. A major step towards a consistent description of non-perturbative aspects was taken by Braaten and Pisarski , who introduced a novel resummation method that re-organizes the perturbative expansion in terms of the degrees of freedom associated with collective modes , rather than bare particles. This program, called the hard thermal loop or HTL program is now at the heart of most treatments of equilibrium aspects of abelian and non-abelian plasmas .
Thermal field theory provides the tools to study the properties of plasmas in equilibrium , but the consistent study of nonequilibrium phenomena in real time requires the methods of nonequilibrium field theory (see also Refs. for further references). The study of the equilibrium and nonequilibrium properties of abelian and non-abelian plasmas as applied to the QGP has as ultimate goal a deeper understanding of the potential experimental signatures of the formation and evolution of the QGP in ultrarelativistic heavy ion collisions. Amongst these, photons and dileptons (electron and/or muon pairs) produced during the early stages of the QGP are considered as some of the most promising signals . Since photons and lepton pairs interact electromagnetically their mean free paths are longer than the estimated size of the QGP fireball $`1050\text{fm}`$ and unlike hadronic signals they do not undergo final state interactions. Therefore photons and dileptons produced during the early stages of QGP carry clean information from this phase.
The goals of this work. In this work we aim to provide a comprehensive study of several relevant aspects of the nonequilibrium dynamics of an abelian QED plasma in real time. Many features of full QCD are similar to those of the abelian (QED) theory, in particular leading contributions in the HTL limit can be straighfordarly generalized from QED to QCD, thus the leading results for relaxation and photon production from a QGP can be understood from the study of a QED plasma.
We utilize a gauge invariant formulation, available in abelian gauge theories that circumvents the possible ambiguities associated with gauge invariance . Thus gauge and fermionic mean fields and distribution functions are automatically gauge invariant. We implement and apply the method of the dynamical renormalization group, introduced recently to study non-equilibrium phenomena directly in real time to extract a consistent, non-perturbative description of real-time dynamics out of equilibrium.
In particular we focus on the following:
* The real time evolution of gauge mean fields in linear response in the HTL approximation. The goal here is to study directly in real time the relaxation of (coherent) gauge field configurations in the linearized approximation to leading order in the HTL program. Whereas a similar study has been carried out in scalar QED and confirmed numerically in the most relevant case of spinor QED has not yet been studied in detail.
* The quantum kinetic equation that describes the evolution of the distribution function of photons in the medium, again to leading order in the HTL approximation. This aspect is relevant to study photon production via off-shell effects directly in real time. As explained in detail, this quantum kinetic equation, obtained from a microscopic field theoretical approach based on the dynamical renormalization group displays novel off-shell effects that cannot be captured via the usual kinetic description that assumes completed collisions .
* The evolution in real time of fermionic mean fields features anomalous relaxation arising from the emission and absorption of magnetic photons (gluons) which are only dynamically screened by Landau damping . The fermion propagator was studied previously in real time in the Bloch-Nordsieck approximation which provides a resummation of the infrared divergences associated with soft photon (or gluon) bremsstrahlung in the medium . In this article we implement the dynamical renormalization group to study the evolution of fermionic mean fields providing an alternative to the Bloch-Nordsieck treatment.
* We obtain the quantum kinetic equation for the fermionic distribution function for hard fermions via the implementation of the dynamical renormalization group. There has recently been an important effort in trying to obtain the effective kinetic (Boltzmann) equations for hard charged (quasi)particles but the collision kernel in this equation features the logarithmic divergences associated with the emission and absorption of soft magnetic photons (or gluons) . The dynamical renormalization group leads to a quantum kinetic equation directly in real time bypassing the assumption of completed collisions and leads to a time-dependent collision kernel free of infrared divergences.
Summary of the main results. The main results of this study are summarized as follows.
* Relaxation of gauge mean fields. We studied the relaxation of a gauge mean field in linear response to leading order in the HTL approximation both for soft momentum $`keT`$ and for semihard momentum $`eTkT`$ under the assumption of weak electromagnetic coupling.
Soft momentum ($`keT`$): in this case the relaxation of the gauge mean field is dominated by the end-point contribution of the Landau damping cut. As a consequence, the soft gauge mean field relaxes with a power law long time tail of the form
$$𝐚_T(𝐤,t)\stackrel{kt1}{=}𝐚_T(𝐤,0)\left[\frac{k^2Z_T(k)}{\omega _T^2(k)}\mathrm{cos}[\omega _T(k)t]\frac{12}{e^2T^2}\frac{\mathrm{cos}kt}{t^2}\right],$$
where $`\omega _T(k)`$ is the transverse photon pole and $`Z_T(k)`$ is the corresponding residue. We note that in spite of the power law tail the gauge mean field relaxes towards the oscillatory mode determined by the transverse photon pole. This reveals that the soft collective excitation in a plasma is stable in the HTL approximation.
Ultrasoft momentum ($`keT`$): In the region of ultrasoft momentum the spectral density divided by the frequency features a sharp Breit-Wigner peak near zero frequency in the region of Landau damping, with width $`\mathrm{\Gamma }_k=12k^3/\pi e^2T^2`$. We find that the amplitude of a mean field of transverse photons prepared via a source that is adiabatically switched-on is given by
$$𝐚_T(𝐤,t)\stackrel{kt1,\mathrm{\Gamma }_kt1}{=}𝐚_T(𝐤,0)\left[\frac{k^2Z_T(k)}{\omega _T^2(k)}\mathrm{cos}[\omega _T(k)t]+e^{\mathrm{\Gamma }_kt}\right].$$
We emphasize that the exponential decay is a consequence of the sharp resonance near zero frequency in the Landau damping region of the spectral density and only arises if the mean field is prepared by an external source whose time Fourier transform has a simple pole at zero frequency, such is the case for an adiabatically prepared initial state. This result confirms those found numerically in Ref. .
Semi-hard momentum ($`eTkT`$): In this region both the HTL approximation and the perturbative expansion are formally valid. However the spectral density in the Landau damping region is sharply peaked near $`\omega =k`$ and the transverse photon pole approaches the edge of the Landau damping region from above. We find that although the perturbative expansion is in principle valid, the sharp spectral density near the edge of the continuum results in a breakdown of the perturbative expansion. The dynamical renormalization group provides a consistent resummation of the lowest order HTL perturbative contributions in real time, leading to the following relaxation of the mean field at intermediate asymptotic times:
$$𝐚_T(𝐤,t)\stackrel{kt1}{=}𝓐_T(𝐤,t)\left(\frac{t}{\tau _0}\right)^{\frac{e^2T^2}{12k^2}},$$
where $`\tau _01/k`$ and $`𝓐_T(𝐤,t)`$ is an oscillating function. The anomalous exponent is a consequence of an infrared enhancement arising from the sharp spectral density near the threshold of the Landau damping region for semihard momentum. The crossover to exponential relaxation due to collisional processes at higher orders is discussed.
* Quantum kinetic equation for the photon distribution function. Using the techniques of nonequilibrium field theory and the dynamical renormalization group, we obtain the quantum kinetic equation for the distribution function of semihard photons $`eTkT`$ to lowest order in the HTL approximation assuming that the fermions are thermalized. An important result is that the collision kernel is time-dependent and the dynamical renormalization group reveals that detailed balance emerges during microscopic time scales, i.e, much shorter than the relaxation scales. In the linearized approximation we find that the departure from equilibrium of the photon distribution function relaxes as:
$$\delta n_𝐤^\gamma (t)=\delta n_𝐤^\gamma (t_0)\left(\frac{tt_0}{\tau _0}\right)^{\frac{e^2T^2}{6k^2}}\text{for}k(tt_0)1,$$
where $`\tau _01/k`$, and $`t_0`$ is the initial time. Furthermore, this quantum kinetic equation allows us to study photon production by off-shell effects, which to leading order in the HTL approximation are determined by photon bremsstrahlung and is of order $`\alpha `$. Extrapolating the result from QED to thermalized QGP with two flavors of light quarks, we find that the total number of hard photons at time $`t`$ per invariant phase space volume to lowest order is
$$N(t)=\frac{5\alpha T^3}{18\pi ^2k^2}\left\{\mathrm{ln}\left[2k(tt_0)\right]+\gamma _E1\right\}\text{for}k(tt_0)>1.$$
with $`t_01\text{fm/c}`$ is the time scale at which the QGP plasma is thermalized.
We find that for a quark-gluon plasma at temperature $`T200\text{MeV}`$ and of lifetime $`10(tt_0)50\text{fm}/c`$, the hard ($`kT`$) photon production by off-shell bremsstrahlung is comparable to that from Compton scattering and pair annihilation of order $`\alpha \alpha _s`$ .
* Relaxation of fermion mean fields. We implement the dynamical renormalization group resummation to study the real-time relaxation of a fermion mean field for hard momentum. The emission and absorption of magnetic photons which are only dynamically screened by Landau damping introduce a logarithmic divergence in the spectral density near the fermion mass shell. The dynamical renormalization group resums these divergences in real time and leads to a relaxation of the fermion mean field for hard momentum given by:
$$\psi (𝐤,t)\stackrel{kt1}{=}e^{\alpha Tt\left[\mathrm{ln}(\omega _Pt)+0.12652\mathrm{}\right]}\times \text{oscillating phases},$$
with $`\omega _P`$ being the plasma frequency.
* Quantum kinetics for the fermion distribution function. We obtain a quantum kinetic equation for the distribution function of hard fermions using non-equilibrium field theory and the dynamical renormalization group resummation. Assuming thermalized photons we find that the collision kernel is infrared finite but time-dependent. In the linearized approximation the distribution function relaxes as:
$$\delta n_𝐤^f(t)\stackrel{kt1}{=}\delta n_𝐤^f(t_0)e^{2\alpha T(tt_0)\left[\mathrm{ln}(\omega _Pt)+0.12652\mathrm{}\right]},$$
where the anomalous relaxation exponent is twice that of the mean field.
The article is organized as follows. In Sec. II we review the main ingredients of nonequilibrium field theory, the initial value problem formulation for relaxation of nonequilibrium mean fields, and the dynamical renormalization group approach to quantum kinetics. In Sec. III we first study relaxation of the photon mean field in the hard thermal loop limit for soft $`keT`$ and semihard $`eTkT`$ photon momentum and obtain the quantum kinetic equation for the (hard and semihard) photon distribution function. In Sec. IV we study relaxation of fermionic mean fields for hard momentum and the quantum kinetic equation for the fermion distribution function. Our conclusions and some further questions are presented in Sec. V.
## II General Aspects
As stated in the introduction our goal is to provide a systematic study of non-equilibrium phenomena in a hot abelian plasma. In particular we focus on a detailed study of the real time relaxation of expectation values of fermions and gauge bosons, i.e, fermionic and gauge mean fields as well as the quantum kinetics for the evolution of the expectation value of the number operator associated with fermions and photons. A fundamental issue that must be addressed prior to setting up our study is that of gauge invariance. In the abelian case it is straightforward to reduce the Hilbert space to the gauge invariant subspace and to define gauge invariant charged fermionic operators. The description in terms of gauge invariant states and operators is best achieved within the canonical formulation which begins with the identification of the canonical field and conjugate momenta and the primary and secondary first class constraints associated with gauge invariance. The physical states are those annihilated by the constraints and physical operators commute with the (first class) constraints. This program has been implemented explicitly in the case of scalar quantum electrodynamics and the fermionic case can be treated in the same manner with few minor technical modifications.
The final result of this formulation is that the Hamiltonian acting on the gauge invariant states and written in terms of gauge invariant fields is exactly equivalent to that obtained in Coulomb gauge, which is the statement that Coulomb gauge describes the theory in terms of the physical degrees of freedom. Furthermore the instantaneous Coulomb interaction can be traded for a Lagrange multiplier leading to the following Lagrangian density
$``$ $`=`$ $`\overline{\mathrm{\Psi }}\left(i\overline{)}e\gamma _0A_0+e𝜸𝐀_Tm\right)\mathrm{\Psi }+{\displaystyle \frac{1}{2}}\left[\left(_\mu 𝐀_T\right)^2+\left(A_0\right)^2\right].`$
where $`𝐀_T`$ is the transverse component of the vector potential and $`A_0`$ is not to be interpreted as the time component of the gauge vector potential but is the Lagrange multiplier associated with the instantaneous Coulomb interaction. We emphasize that the fermionic charged fields $`\mathrm{\Psi }`$ as well as the Lagrange multiplier $`A_0`$ are gauge invariant fields (see Refs. ).
The fully renormalized equations of motion for the mean fields can be obtained by following the formulation presented in Refs. . However, the counterterms that eliminate the zero temperature ultraviolet divergences are temperature independent and play no rôle in the present context, thus we will neglect the zero temperature renormalization counterterms in our study.
Furthermore we consider a neutral plasma (i.e, with zero chemical potential for the charged fields) at a temperature $`Tm`$ with $`m`$ the (renormalized) mass of the fermions, hence in what follows we neglect the fermion mass unless otherwise stated.
### A Nonequilibrium Field Theory
The formulation of nonequilibrium quantum field theory in terms of the Schwinger-Keldysh or the closed-time-path (CTP) is standard . A path integral representation requires a contour in the complex time plane, running forward and then backwards in time corresponding to the unitary time evolution of an initially prepared density matrix, with an effective Lagrangian in terms of fields on the different branches of the path
$$_{\mathrm{noneq}}=[\mathrm{\Psi }^+,\overline{\mathrm{\Psi }}^+,A_T^+,A_0^+][\mathrm{\Psi }^{},\overline{\mathrm{\Psi }}^{},A_T^{},A_0^{}],$$
where the “$`+`$” (“$``$”) superscripts for the fields refer to fields defined in the forward (backward) time branches. This description leads to a straightforward diagrammatic expansion of the nonequilibrium Green’s functions in terms of real-time propagators but with modified Feynman rules (as compared to standard field theory) . The free fermion and photon propagators are given in detail in the Appendix.
Since we study two different type of situations: (i) the relaxation of expectation values in the linearized approximation (linear response) and (ii) the kinetic equation that describes the evolution of a non-equilibrium distribution function, we make a distinction between the real-time propagators used in each case.
* In the case of relaxation of expectation values in linear response, the distribution functions that enter in the fermion and photon propagators (125)-(135) are those with equilibrium distributions.
* When we study the kinetic equations for the expectation value of the fermion number and photon number operators, we assume that the initial density matrix is diagonal in the basis of the quasiparticles (see below) but with nonequilibrium initial distributions. Thus the propagators obtained with this initial density matrix have the free-field form, as given in (125)-(135) but the initial occupation numbers are not the equilibrium Fermi-Dirac or Bose-Einstein distribution.
### B Real-time relaxation of mean fields in linear response
The mean fields under consideration are the expectation value of either fermion or gauge field operators in the nonequilibrium state induced by the external source. Obviously in equilibrium and in the absence of external sources these must vanish, and we introduce time-dependent sources to induce an expectation value of these operators. Our strategy to study the relaxation of these mean fields as an initial value problem is to prepare these mean fields via the adiabatic switching-on of an external current. Once the current is switched off the expectation values must relax towards equilibrium and we study this real-time evolution. This formulation to study the real-time evolution of mean fields has the advantage that it leads to a straightforward and systematic implementation of the dynamical renormalization group method as explained in detail in Refs. .
The initial value problem formulation for studying the linear relaxation of the mean fields begins by introducing c-number external sources coupled to the quantum fields. Let $`\eta (x)`$ be the Grassmann-valued fermionic source and $`𝐉_T(x)`$ be the electromagnetic sources,<sup>*</sup><sup>*</sup>*In this article we will not discuss relaxation of the longitudinal gauge field $`A_0(x)`$ associated with the instantaneous Coulomb interaction, hence the corresponding external source is neglected. then the Lagrangian density becomes
$$+\overline{\mathrm{\Psi }}\eta +\overline{\eta }\mathrm{\Psi }+𝐉_T𝐀_T.$$
The presence of external sources will induce responses of the system. Let $`\mathrm{\Phi }(x)`$ be a generic quantum field (i.e., fermionic or bosonic) with $`\mathrm{\Phi }^\pm (x)`$ the fields defined on the forward and backward branches, respectively, and $`J(x)`$ the corresponding c-number external source (the same for both branches). The expectation value of $`\mathrm{\Phi }(x)`$ induced by $`J(x)`$ in a linear response analysis is given by
$`\varphi (𝐱,t)`$ $``$ $`\mathrm{\Phi }^\pm (𝐱,t)_J`$ (1)
$`=`$ $`{\displaystyle d^3x^{}_{\mathrm{}}^+\mathrm{}𝑑t^{}G_{\mathrm{ret}}(𝐱,t;𝐱^{},t^{})J(𝐱^{},t^{})},`$ (2)
with the retarded Green’s function
$`G_{\mathrm{ret}}(𝐱,t;𝐱^{},t^{})`$ $``$ $`i\left[\mathrm{\Phi }^+(𝐱,t)\overline{\mathrm{\Phi }}^+(𝐱^{},t^{})\mathrm{\Phi }^+(𝐱,t)\overline{\mathrm{\Phi }}^{}(𝐱^{},t^{})\right]`$ (3)
$`=`$ $`i[\mathrm{\Phi }(𝐱,t),\overline{\mathrm{\Phi }}(𝐱^{},t^{})]_{}\theta (tt^{}),`$ (4)
where expectation values are computed in the CTP formulation in the full interacting theory but with vanishing external source, $`\theta (tt^{})`$ is the step function and the subscripts $``$ refer to the commutator ($``$) for bosonic or anticommutator ($`+`$) for fermionic fields.
A practically useful initial value problem formulation for the real-time relaxation of mean fields is obtained by considering that the external source is adiabatically switched on in time at $`t=\mathrm{}`$ and suddenly switched off at $`t=t_0`$, i.e.,
$$J(𝐱,t)=J(𝐱)e^{ϵ(tt_0)}\theta (t_0t),ϵ0^+.$$
(5)
The adiabatic switching on of the external source induces a mean field that is dressed adiabatically by the interaction. Then for $`t>t_0`$, after the external current has been switched off, the mean field will relax towards equilibrium, and our aim is to study this relaxation directly in real time .
The equations of motion of the mean fields are obtained via the tadpole method and are automatically causal and retarded. The central idea of the tadpole method is to write the field into a c-numbered expectation value plus fluctuations around it, i.e., writing
$$\mathrm{\Phi }^\pm (𝐱,t)=\varphi (𝐱,t)+\chi ^\pm (𝐱,t)\text{with}\mathrm{\Phi }^\pm (𝐱,t)_J=\varphi (𝐱,t).$$
The equation of motion for $`\varphi (𝐱,t)`$ is obtained by requiring $`\chi ^\pm (𝐱,t)=0`$ to all orders in perturbation theory.
For the study of relaxation of the mean fields, we assume that at $`t=\mathrm{}`$ the system is in thermal equilibrium at a temperature $`T`$ and henceforth choose $`t_0=0`$ for convenience. The retarded and the equilibrium (i.e., time translational invariant) nature of $`G_{\mathrm{ret}}(𝐱,t;𝐱^{},t^{})`$ and the adiabatic switching on of $`J(𝐱,t)`$ entail that
$$\begin{array}{c}\varphi (𝐱,t=0)=\varphi _0(𝐱),\hfill \\ \dot{\varphi }(𝐱,t<0)=0,\hfill \end{array}$$
(6)
where $`\varphi _0(𝐱)`$ is determined by $`J(𝐱)`$ (or vice versa). It is worth pointing out that $`\dot{\varphi }(𝐱,t)`$ is not specified at $`t=0`$ although $`\dot{\varphi }(𝐱,t<0)=0`$. This is because the external source is switched off at $`t=0`$.There could be initial time singularities associated with our choice of the external source (see Ref. ). However the long time asymptotics is not sensitive to the initial time singularities , and we will not address this issue here as it is not relevant for the asymptotic dynamics.
### C Main ingredients for quantum kinetics
In Ref. a systematic treatment to obtain the kinetic equation that determines the real-time evolution of the distribution function (defined as the expectation value of the number operator) was established, leading to its interpretation as a dynamical renormalization group equation. Here we summarize the basic steps to obtain the relevant quantum kinetic equations from first principles and we refer the reader to Ref. for more details.
* Identify the proper degrees of freedom (quasiparticles) to be described by the kinetic equation, the corresponding microscopic time scale $`\tau _{\mathrm{micro}}`$ associated with their oscillation, and the number operator $`N(𝐤,t)`$ that counts these degrees of freedom with momentum $`𝐤`$ at time $`t`$.
* Use the Heisenberg equations of motion to derive a general rate equation for $`dn_𝐤(t)/dt`$ with $`n_𝐤(t)=N(𝐤,t)`$, where the expectation value is taken in the initial density matrix. Assuming that the initial density matrix is diagonal in the basis of free quasiparticles but with nonequilibrium distribution functions, we can expand the rate equation in perturbation theory using nonequilibrium Feynman rules and free real-time propagators in terms of nonequilibrium distribution functions. Since the resulting expression is a functional of the distribution function at initial time, the solution of the kinetic equation can be obtained by direct integration.
This solution is typically characterized by the emergence of secular terms, i.e, terms that grow in time in the intermediate asymptotic regime ($`\tau _{\mathrm{micro}}t\tau _{\mathrm{rel}}`$) where the scale $`\tau _{\mathrm{rel}}`$ is that at which perturbation theory breaks down .
* These secular terms growing in time become divergent if the perturbative solution is extrapolated to long times. However, perturbation theory is valid and dominated by the secular terms during the intermediate asymptotic time scales $`\tau _{\mathrm{micro}}t\tau _{\mathrm{rel}}`$. The dynamical renormalization group is used to resum the secular terms through a renormalization of the initial distribution functions. This introduces an arbitrary time scale $`\tau `$ which serves as a renormalization point. The $`\tau `$-independence of the renormalized solution leads to the dynamical renormalization equation . This dynamical renormalization equation describing the evolution of the quasiparticle distribution is recognized as the quantum kinetic equation.
## III Photons out of equilibrium
In this section we study non-equilibrium aspects of photon relaxation and production. We begin by analyzing the relaxation in real time of a photon condensate in linear response both in the case of soft $`keT`$ and semihard (or semisoft) $`eTkT`$ assuming the electromagnetic coupling to be small. We then continue with a study of the production of semihard photons via off-shell effects.
### A Relaxation of the gauge mean field
We begin with the relaxation of soft photons of momenta $`keT`$. As mentioned in the previous section, the equation of motion for the transverse photon mean field can be derived from the tadpole method by decomposing the full quantum fields into c-number expectation values and quantum fluctuations around the expectation values:
$`𝐀_T^\pm (𝐱,t)`$ $`=`$ $`𝐚_T(𝐱,t)+𝓐_T^\pm (𝐱,t),`$ (7)
$`𝐚_T(𝐱,t)`$ $`=`$ $`𝐀_T^\pm (𝐱,t).`$ (8)
With the external source $`𝐉_T(𝐱,t)`$ of the adiabatic form (5) chosen to vanish at $`t0`$, we find the linear equation of motion for the spatial Fourier transform of $`𝐚_T(𝐱,t)`$ for $`t0`$ to be given by
$$\left(\frac{^2}{t^2}+k^2\right)𝐚_T(𝐤,t)+_{\mathrm{}}^t𝑑t^{}\mathrm{\Pi }_T(𝐤,tt^{})𝐚_T(𝐤,t^{})=0,$$
(9)
where
$$𝐚_T(𝐤,t)d^3xe^{i𝐤𝐱}𝐚_T(𝐱,t).$$
Here $`\mathrm{\Pi }_T(𝐤,tt^{})`$ is the transverse part of the retarded photon self-energy and we have made explicit use of its retarded nature in writing Eq. (9), which evidently is causal and retarded.
Using the nonequilibrium Feynman rules and the real-time free fermion propagators given in the Appendix, we find the transverse polarization to be given by
$`\mathrm{\Pi }_T(𝐤,tt^{})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\omega \rho _T(\omega ,𝐤)\mathrm{sin}[\omega (tt^{})],`$ (10)
in terms of the spectral density to one-loop order
$`\rho _T(\omega ,𝐤)`$ $`=`$ $`e^2{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}\{[\delta (\omega pq)\delta (\omega +p+q)][1+(\widehat{𝐤}\widehat{𝐩})(\widehat{𝐤}\widehat{𝐪})][1n_F(p)n_F(q)]`$ (12)
$`+[\delta (\omega p+q)\delta (\omega +pq)][n_F(q)n_F(p)][1(\widehat{𝐤}\widehat{𝐩})(\widehat{𝐤}\widehat{𝐪})]\},`$
where $`𝐩=𝐤+𝐪`$.
It is now convenient to introduce an auxiliary quantity $`\pi _T(𝐤,tt^{})`$ as follows
$$\mathrm{\Pi }_T(𝐤,tt^{})=\frac{}{t^{}}\pi _T(𝐤,tt^{}).$$
and hence,
$$\pi _T(𝐤,tt^{})=_{\mathrm{}}^+\mathrm{}\frac{d\omega }{\omega }\rho _T(\omega ,𝐤)\mathrm{cos}[\omega (tt^{})].$$
Upon integration by parts the equation of motion takes a form that displays clearly the nature of the initial value problem:
$`\left({\displaystyle \frac{^2}{t^2}}+k^2\right)𝐚_T(𝐤,t)`$ $``$ $`{\displaystyle _0^t}𝑑t^{}\pi _T(𝐤,tt^{})\dot{𝐚}_T(𝐤,t^{})+\pi _T(𝐤,0)𝐚_T(𝐤,t)=0,`$ (13)
where a dot denotes derivative with respect to time and we used that $`\dot{𝐚}_T(𝐤,t)=0`$ for $`t<0`$. Eq. (13) together with the initial conditions specified at $`t=0`$ yields a well-defined initial value problem. As discussed above, the initial value $`𝐚_T(𝐤,t=0)=𝐚_T(𝐤,0)`$ is determined by $`𝐉_T(𝐤,t)`$, and we choose
$$\dot{𝐚}_T(𝐤,t=0)=0$$
(14)
consistently with the adiabatic switching-on condition (6).
Eq. (13) can be solved by the Laplace transform method. We define,
$`\stackrel{~}{𝐚}_T(s,𝐤)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑te^{st}𝐚_T(𝐤,t),`$ (15)
$`\stackrel{~}{\pi }_T(s,𝐤)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑te^{st}\pi _T(𝐤,t),`$ (16)
with Re $`s>0`$. In terms of which the Laplace transform of the equation of motion leads to
$$\left[s^2+k^2+\stackrel{~}{\mathrm{\Pi }}_T(s,𝐤)\right]\stackrel{~}{𝐚}_T(s,𝐤)=\left[s\stackrel{~}{\pi }_T(s,𝐤)\right]𝐚_T(𝐤,0),$$
(17)
where $`\stackrel{~}{\mathrm{\Pi }}_T(s,𝐤)`$ is the Laplace transform of $`\mathrm{\Pi }_T(𝐤,tt^{})`$:
$$\stackrel{~}{\mathrm{\Pi }}_T(s,𝐤)=\pi _T(𝐤,0)s\stackrel{~}{\pi }_T(s,𝐤).$$
The solution of Eq. (17) is given by
$$\stackrel{~}{𝐚}_T(s,𝐤)=\frac{1}{s}\left\{1\stackrel{~}{\mathrm{\Delta }}_T(s,𝐤)[k^2+\pi _T(𝐤,0)]\right\}𝐚_T(𝐤,0),$$
(18)
where the retarded transverse photon propagator $`\stackrel{~}{\mathrm{\Delta }}_T(s,𝐤)`$ is given by
$$\stackrel{~}{\mathrm{\Delta }}_T(s,𝐤)=[s^2+k^2+\stackrel{~}{\mathrm{\Pi }}_T(s,𝐤)]^1.$$
(19)
The real-time evolution of $`𝐚_T(𝐤,t)`$ is obtained by performing the inverse Laplace transform along the Bromwich contour which is to the right of all singularities of $`\stackrel{~}{𝐚}_T(s,𝐤)`$ in the complex $`s`$-plane.
We note that this result is rather different from that obtained in Ref. in the structure of the solution, in particular the prefactor $`1/s`$ in Eq. (18). This can be traced back to the different manner in which we have set up the initial value problem, as compared to the case studied in . The adiabatic source (5) (for $`t_0=0`$) has a Fourier transform that has a simple pole in the frequency plane, this translates into the $`1/s`$ factor in Eq. (18). This particular case was not contemplated in those studied in Ref. since the source (5) is not a regular function of the frequency. This difference will be seen to be at the heart of important aspects of relaxation of the mean field condensate in the soft momentum limit as discussed in detail below.
It is straightforward to see that the residue vanishes at $`s=0`$ in the solution (18). The singularities of $`\stackrel{~}{𝐚}_T(s,𝐤)`$ are those arising from the retarded transverse photon propagator $`\stackrel{~}{\mathrm{\Delta }}_T(s,𝐤)`$.
#### 1 Soft photons $`keT`$: real-time Landau damping
For soft photons of momenta $`keT`$, the leading ($`e^2T^2`$) contribution to the photon self-energy arises from loop momenta of order $`T`$ and is referred to as a hard thermal loop (HTL) . In this region of momenta, the contribution of the fermion loop is non-perturbative.
In the HTL approximation ($`skq`$), after some algebra we find that the HTL contribution to $`\stackrel{~}{\mathrm{\Pi }}_T(s,𝐤)`$ arises exclusively from the terms associated with $`\delta (\omega p\pm q)`$ in Eq. (12) and reads
$$\stackrel{~}{\mathrm{\Pi }}_T(s,𝐤)=\frac{e^2T^2}{12}\left[\frac{is}{k}\left(1+\frac{s^2}{k^2}\right)\mathrm{ln}\frac{is+k}{isk}2\frac{s^2}{k^2}\right].$$
(20)
The delta functions $`\delta (\omega p\pm q)`$ have support below the light cone (i.e., $`\omega ^2<k^2`$) and correspond to Landau damping processes in which the soft photon scatters a hard fermion in the plasma.
The analytic continuation of $`\stackrel{~}{\mathrm{\Pi }}_T(s,𝐤)`$ in the complex $`s`$-plane is defined by
$`\mathrm{\Pi }_T(\omega ,𝐤)`$ $``$ $`\stackrel{~}{\mathrm{\Pi }}_T(s=i\omega \pm 0^+,𝐤)=\text{Re}\mathrm{\Pi }_T(\omega ,𝐤)\pm i\text{Im}\mathrm{\Pi }_T(\omega ,𝐤).`$ (21)
From Eq. (10) it is straightforward to see that the real and imaginary parts are related by the usual dispersion relation. The analytical continuation of $`\stackrel{~}{\mathrm{\Delta }}_T(s,𝐤)`$ and $`\stackrel{~}{\pi }_T(s,𝐤)`$ can be defined analogously, and they are related to $`\mathrm{\Pi }_T(\omega ,𝐤)`$ by
$`\mathrm{\Delta }_T(\omega ,𝐤)`$ $`=`$ $`[\omega ^2k^2\mathrm{\Pi }_T(\omega ,𝐤)]^1,`$
$`\text{Re}\mathrm{\Pi }_T(\omega ,𝐤)`$ $`=`$ $`\pi _T(𝐤,0)\omega \text{Im}\pi _T(\omega ,𝐤),`$
$`\text{Im}\mathrm{\Pi }_T(\omega ,𝐤)`$ $`=`$ $`\omega \text{Re}\pi _T(\omega ,𝐤).`$
The analytical continuation of $`\stackrel{~}{\mathrm{\Pi }}_T(s,𝐤)`$ in the HTL limit thus reads
$`\mathrm{\Pi }_T(\omega ,𝐤)`$ $`=`$ $`{\displaystyle \frac{e^2T^2}{12}}\left[2{\displaystyle \frac{\omega ^2}{k^2}}+{\displaystyle \frac{\omega }{k}}\left(1{\displaystyle \frac{\omega ^2}{k^2}}\right)\mathrm{ln}\left|{\displaystyle \frac{\omega +k}{\omega k}}\right|\right]`$ (23)
$`i{\displaystyle \frac{\pi e^2T^2}{12}}{\displaystyle \frac{\omega }{k}}\left(1{\displaystyle \frac{\omega ^2}{k^2}}\right)\theta (k^2\omega ^2).`$
In the HTL limit it is straightforward to prove that $`\pi _T(𝐤,0)=0`$. We note that $`\text{Im}\mathrm{\Pi }_T(\omega ,𝐤)`$ only has support below the light cone and vanishes linearly as $`\omega k`$ from below. Since for soft momenta $`keT`$ the HTL photon self-energy is comparable in magnitude to or larger than the free inverse propagator, it has to be treated nonperturbatively.
Isolated poles of $`\mathrm{\Delta }_T(\omega ,𝐤)`$ in the complex $`\omega `$-plane correspond to collective excitations. The transverse photon pole $`\omega _T(k)`$ is real and determined by
$$\omega _T^2(k)k^2\text{Re}\mathrm{\Pi }_T[\omega _T(k),𝐤]=0.$$
For ultrasoft photons $`keT`$, the dispersion relation of the collective excitations is
$$\omega _T^2(k)=\omega _P^2+\frac{6}{5}k^2+𝒪\left(\frac{k^4}{e^2T^2}\right),$$
where $`\omega _P=eT/3`$ is the plasma frequency. Consequently, the retarded transverse photon propagator $`\stackrel{~}{\mathrm{\Delta }}_T(s,𝐤)`$ has two isolated poles at $`s=\pm i\omega _T(k)`$ and a branch cut from $`s=ik`$ to $`s=ik`$.
Having analyzed the analytic structure of $`\stackrel{~}{\mathrm{\Delta }}_T(s,𝐤)`$, we can now perform the integral along the Bromwich contour to obtain the real-time evolution of $`𝐚_T(𝐤,t)`$. Closing the contour in the half-plane $`\text{Re}s<0`$, we find $`𝐚_T(𝐤,t)`$ for $`t>0`$ to be given by
$$𝐚_T(𝐤,t)=𝐚_T^{\mathrm{𝑝𝑜𝑙𝑒}}(𝐤,t)+𝐚_T^{\mathrm{𝑐𝑢𝑡}}(𝐤,t),$$
with
$`𝐚_T^{\mathrm{𝑝𝑜𝑙𝑒}}(𝐤,t)`$ $`=`$ $`{\displaystyle \frac{k^2Z_T(k)}{\omega _T^2(k)}}\mathrm{cos}[\omega _T(k)t]𝐚_T(𝐤,0),`$ (24)
$`𝐚_T^{\mathrm{𝑐𝑢𝑡}}(𝐤,t)`$ $`=`$ $`k^2{\displaystyle _k^{+k}}{\displaystyle \frac{d\omega }{\omega }}\beta _T(\omega ,k)e^{i\omega t}𝐚_T(𝐤,0),`$ (25)
where
$`Z_T(k)`$ $`=`$ $`\left[1{\displaystyle \frac{\text{Re}\mathrm{\Pi }_T(\omega ,𝐤)}{\omega ^2}}\right]_{\omega =\omega _T(k)}^1,`$ (26)
$`\beta _T(\omega ,k)`$ $`=`$ $`{\displaystyle \frac{\frac{e^2T^2}{12}\frac{\omega }{k}\left(1\frac{\omega ^2}{k^2}\right)}{\left\{\omega ^2k^2\frac{e^2T^2}{12}\left[\frac{2\omega ^2}{k^2}+\frac{\omega }{k}\left(1\frac{\omega ^2}{k^2}\right)\mathrm{ln}\frac{k+\omega }{k\omega }\right]\right\}^2+\left[\frac{\pi e^2T^2}{12}\frac{\omega }{k}\left(1\frac{\omega ^2}{k^2}\right)\right]^2}}.`$ (28)
The solution evaluated at $`t=0`$ must match the initial condition, this requirement leads immediately to the sum rule
$$_{\mathrm{}}^+\mathrm{}\frac{dq_0}{q_0}\stackrel{~}{\rho }_T(q_0,q)=\frac{1}{q^2},$$
where $`\stackrel{~}{\rho }_T(q_0,q)`$ is the HTL-resummed spectral density for the transverse photon propagator:In our notation the spectral density for the self-energy is denoted by $`\rho `$, and the spectral density for the corresponding propagator is denoted by $`\stackrel{~}{\rho }`$.
$`\stackrel{~}{\rho }_T(q_0,q)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Delta }_T(q_0,q)`$ (29)
$`=`$ $`\text{sgn}(q_0)Z_T(q)\delta [q_0^2\omega _T^2(q)]+\beta _T(q_0,q)\theta (q^2q_0^2).`$ (30)
A noteworthy feature of the cut contribution (25) is the factor $`\omega `$ in the denominator. The presence of this factor can be traced back to the preparation of the initial state via the adiabatically switched-on external source, which is the switched off at $`t=0`$. The Fourier transform of this source is proportional to $`1/\omega `$ and results in the prefactor of $`\beta _T(\omega ,k)`$ in Eq. (25).
For $`keT`$, Eq. (25) cannot be evaluated in closed form but its long time asymptotics can be extracted by writing the integral along the cut as a contour integral in the complex $`\omega `$-plane. The integration contour $`𝒞`$ is chosen to run clockwise along the segment $`k<\omega <k`$ on the real axis, the line $`\omega =kiz`$ with $`0z<\mathrm{}`$, then around an arc at infinity and back to the real axis along the line $`\omega =kiz`$. After some algebra we obtain the following expression
$`𝐚_T^{\mathrm{𝑐𝑢𝑡}}(𝐤,t)`$ $`=`$ $`k^2\{e^{ikt}{\displaystyle _0^{\mathrm{}}}dze^{zt}{\displaystyle \frac{i\beta _T(kiz,k)}{kiz}}+\text{ c.c.}`$
$`\mathrm{\hspace{0.33em}2}\pi i{\displaystyle \underset{\genfrac{}{}{0pt}{}{\text{poles}}{\text{inside }𝒞}}{}}\text{Res}\left[{\displaystyle \frac{\beta _T(\omega ,k)}{\omega }}e^{i\omega t}\right]\}𝐚_T(𝐤,0).`$
Because of the exponential factor $`e^{zt}`$ in the integrals , the dominant contributions to the integral at long times ($`t1/k`$) arise from the end points of the branch cut with $`zk`$, and lead to a long time behavior characterized by a power law:
$$𝐚_T^{cut}(𝐤,t)\stackrel{kt1}{=}\frac{12}{e^2T^2}\frac{\mathrm{cos}kt}{t^2}𝐚_T(𝐤,0)\left[1+𝒪\left(\frac{1}{t}\right)\right],$$
(31)
which confirms the results of Ref. . The contribution from the poles inside the contour $`𝒞`$ is dominated at long times by the pole closest to the real axis.
For very soft momenta $`keT`$ the function $`\beta _T(\omega ,k)/\omega `$ is strongly peaked at $`\omega =0`$ as depicted in Fig. 1. For $`keT`$ and $`k\omega `$, Eq. (28) takes a Breit-Wigner form
$$\frac{1}{\omega }\beta _T(\omega ,k)\stackrel{\omega keT}{=}\frac{1}{\pi k^2}\frac{\mathrm{\Gamma }_k}{\omega ^2+\mathrm{\Gamma }_k^2}\text{with}\mathrm{\Gamma }_k=\frac{12k^3}{\pi e^2T^2}.$$
(32)
Using this narrow-width approximation to evaluate Eq. (25) yields,
$$𝐚_T^{\mathrm{𝑐𝑢𝑡}}(𝐤,t)\stackrel{kt1}{=}\left[e^{\mathrm{\Gamma }_kt}\frac{2\mathrm{\Gamma }_k}{\pi k^2t}\mathrm{sin}kt+𝒪\left(\frac{\mathrm{\Gamma }_k}{k^3t^2}\right)\right]𝐚_T(𝐤,0).$$
(33)
The end-point contribution which results in a power law as in Eq.(31) is very small compared with Eq. (33) for time $`t1/\mathrm{\Gamma }_k`$. This power law becomes dominant at time scales much longer than $`1/\mathrm{\Gamma }_k`$, because its amplitude at $`t1/\mathrm{\Gamma }_k`$ is of order $`(k/eT)^61`$ in the soft momentum limit.
Thus, for adiabatically prepared mean fields of soft momentum $`keT`$ we find the real-time behavior
$$𝐚_T(𝐤,t)\stackrel{kt1,\mathrm{\Gamma }_kt1}{=}𝐚_T(𝐤,0)\left\{\frac{k^2Z_T(k)}{\omega _T^2(k)}\mathrm{cos}[\omega _T(k)t]+e^{\mathrm{\Gamma }_kt}\right\}.$$
It is worth emphasizing that the exponential decay is not the same as the usual decay of an unstable particle (or even collisional broadening) for which the amplitude relaxes to zero at times longer than the relaxation time. In this case the asymptotic behavior of the amplitude is completely determined by the transverse photon pole $`\omega _T(k)`$, and to this order in the HTL approximation the collective excitation is stable. The exponential relaxation does not arise from a resonance at the position of the transverse photon pole but at zero frequency. Clearly we expect a true exponential damping of the collective excitation at higher order as a result of collisional broadening.
Our results thus confirm those found in hot scalar quantum electrodynamics wherein a numerical analysis of the relaxation of the mean field was carried out. In that reference the authors studied the real time evolution in terms of local Hamiltonian equations obtained by introducing a non-local auxiliary field. The initial conditions chosen there for this non-local auxiliary field correspond precisely to the choice of an adiabatically switched on current leading to an adiabatically prepared initial value problem.
We emphasize that the exponential relaxation (33) can only be probed by adiabatically switching on an external source \[see Eq. (5)\] whose temporal Fourier transform has a simple pole at $`\omega =0`$ which is the origin of the prefactor $`1/\omega `$ in Eq. (25). That is the adiabatic preparation of the initial state excites this pole in the Landau damping region. For external sources whose temporal Fourier transform is regular at $`\omega =0`$ no exponential relaxation arises, in agreement with the results of Ref. . Thus we find that the exponential relaxation is not a generic feature of the evolution of the mean field, but emerges only for particular (albeit physically motivated) initial conditions.
#### 2 Semi-hard photons $`eTkT`$: anomalous relaxation
Most of the studies of the photon self-energy in the hard thermal loop approximation focused on the soft external momentum region $`keT`$ (with $`e1`$ wherein the contribution of the hard thermal fermion loop must be treated non-nonperturbatively.
However there are important reasons that warrant a consideration of the region of semihard photons $`eTkT`$. From a phenomenological standpoint, hard and semihard photons produced in the QGP phase are important electromagnetic probes of the deconfined phase hence the study of all aspects of the space-time evolution, production and relaxation of photons to explore potential experimental signatures is warranted.
A more theoretical justification of the relevance of this region is that whereas the photon self-energy for this region of momenta is still dominated by the hard thermal fermion loop contribution (20), its contribution to the full photon propagator is now perturbative as compared with the free inverse propagator. The validity of perturbation theory in this regime licenses us to study the real-time evolution with the tools developed in refs. that provide a consistent implementation of the renormalization group to study real-time phenomena.
The dynamical renormalization group approach introduced in Refs. is particularly suited to study the real-time evolution in the case in which there are (infrared) threshold singularities in the spectral density. To understand the potential emergence of anomalous thresholds in the semihard and hard limit for the photon self-energy we note that at leading order in HTL the Laplace transform of the inverse photon propagator Eqs. (19)-(20) reads
$$\stackrel{~}{\mathrm{\Delta }}^1(s,𝐤)=(s^2+k^2)\left[1+\frac{e^2T^2}{12}\frac{is}{k^3}\mathrm{ln}\frac{is+k}{isk}\right]\frac{e^2T^2}{6}\frac{s^2}{k^2}.$$
The transverse photon poles are for $`s=\pm i[k+e^2T^2/12k+𝒪(T^4/k^3)]`$ which in the limit $`keT`$ approach $`s\pm ik`$. That is,
$$\omega _T(k)=k+\frac{e^2T^2}{12k}+𝒪\left(\frac{T^4}{k^3}\right).$$
(34)
In this limit the poles merge with the thresholds of the logarithmic branch cut from Landau damping and are no longer isolated from the continuum. The approach of the pole to the Landau damping threshold and the enhancement of the spectral density near threshold for $`keT`$ is clearly displayed in Fig. 2. While the perturbative expansion in the effective coupling $`e^2T^2/k^2`$ is warranted in the semihard case under consideration, the wave function renormalization constant evaluated as the residue at the pole is a non-analytic function of this coupling and given by
$$Z_T(k)=1+\frac{e^2T^2}{12k^2}\left[\mathrm{ln}\left(\frac{e^2T^2}{24k^2}\right)+3+𝒪\left(\frac{e^2T^2}{k^2}\mathrm{ln}\frac{eT}{k}\right)\right].$$
This is obviously a consequence of the logarithmic threshold singularities in the semihard regime. This threshold singularity is reminiscent of those studied in detail in Ref.. In this reference it was shown that the Bloch-Nordsieck resummation which is equivalent to a renormalization group improvement of the space-time Fourier transform of the self-energy, leads to a real-time evolution which is obtained by the implementation of the dynamical renormalization group in real time . Since the infrared threshold singularities in the semihard limit are akin to those studied in Ref. , we now implement the dynamical renormalization group program advocated in that reference to study the relaxation of the photon mean field in the semihard limit directly in real time.
Perturbation theory in terms of the effective coupling $`e^2T^2/k^2`$ is in principle reliable for semihard photons, hence we can try to solve Eq. (13) by perturbative expansion in powers of $`e^2`$ (the true dimensionless coupling is $`e^2T^2/k^2`$). Writing
$`𝐚_T(𝐤,t)`$ $`=`$ $`𝐚_T^{(0)}(𝐤,t)+e^2𝐚_T^{(1)}(𝐤,t)+𝒪(e^4),`$
$`\pi _T(𝐤,tt^{})`$ $`=`$ $`e^2\pi _T^{(1)}(𝐤,tt^{})+𝒪(e^4),`$
and expanding consistently in powers of $`e^2`$ we obtain a hierarchy of equations:
$`\left({\displaystyle \frac{^2}{t^2}}+k^2\right)𝐚_T^{(0)}(𝐤,t)`$ $`=`$ $`0,`$
$`\left({\displaystyle \frac{^2}{t^2}}+k^2\right)𝐚_T^{(1)}(𝐤,t)`$ $`=`$ $`{\displaystyle _0^t}𝑑t^{}\pi _T^{(1)}(𝐤,tt^{})\dot{𝐚}_T^{(0)}(𝐤,t^{}),`$
$`\mathrm{}`$ $`\mathrm{}`$
where we have explicitly used the fact that $`\pi _T(𝐤,0)=0`$ in the HTL limit. Starting from the solution to the zeroth-order equation
$$𝐚_T^{(0)}(𝐤,t)=\underset{\lambda =1}{\overset{2}{}}\left[A_\lambda (𝐤)e^{ikt}+A_\lambda ^{}(𝐤)e^{ikt}\right]𝓔_\lambda (𝐤),$$
(35)
with $`𝓔_\lambda (𝐤)`$ being the polarization vector, and the retarded Green’s function of the unperturbed problem:
$$G_{\mathrm{ret}}^{(0)}(k,tt^{})=\frac{\mathrm{sin}[k(tt^{})]}{k}\theta (tt^{}),$$
the solution to the first-order equation reads
$`e^2𝐚_T^{(1)}(𝐤,t)`$ $`=`$ $`{\displaystyle \frac{i}{4}}{\displaystyle \underset{\lambda =1}{\overset{2}{}}}A_\lambda (𝐤)𝓔_\lambda (𝐤){\displaystyle _{\mathrm{}}^+\mathrm{}}d\omega {\displaystyle \frac{\rho _T(\omega ,𝐤)}{\omega }}\{[{\displaystyle \frac{e^{ikt}}{\omega k}}(t{\displaystyle \frac{1e^{i(\omega k)t}}{i(\omega k)}})`$ (37)
$`+{\displaystyle \frac{e^{ikt}}{\omega +k}}({\displaystyle \frac{1e^{2ikt}}{2ik}}+{\displaystyle \frac{1e^{i(\omega k)t}}{i(\omega k)}})]+(\omega \omega )\}+\text{ c.c.},`$
where
$`\rho _T(\omega ,𝐤)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_T(\omega ,𝐤)={\displaystyle \frac{e^2T^2}{12}}{\displaystyle \frac{\omega }{k}}\left(1{\displaystyle \frac{\omega ^2}{k^2}}\right)\theta (k^2\omega ^2).`$
Potential secular terms (that grow in time) arise at long times from the regions in $`\omega `$ at which the denominators in Eq. (37) are resonant. These are extracted in the long time limit by using the formulae in an appendix of Ref. . Particularly relevant to the case under consideration are the following :
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dy}{y^2}}\left(1\mathrm{cos}yt\right)p(y)`$ $`\stackrel{t\mathrm{}}{=}`$ $`{\displaystyle \frac{\pi }{2}}tp(0)+p^{}(0)\left[\mathrm{ln}(\mu t)+\gamma _E\right]`$ (39)
$`+{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dy}{y^2}}\left[p(y)p(0)yp^{}(0)\theta (\mu y)\right]+𝒪\left(t^1\right),`$
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dy}{y}}\left(t{\displaystyle \frac{\mathrm{sin}yt}{y}}\right)p(y)`$ $`\stackrel{t\mathrm{}}{=}`$ $`tp(0)\left[\mathrm{ln}(\mu t)+\gamma _E1\right]`$ (41)
$`+t\mathrm{PV}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dy}{y}}\left[p(y)p(0)\theta (\mu y)\right]+𝒪\left(t^1\right),`$
where $`\gamma _E=0.5772157\mathrm{}`$ is the Euler-Mascheroni constant and as can be easily shown the dependence on $`\mu `$ in the above integrals cancels. For our analysis, the integration variable $`y=\omega \pm k`$ and $`p(y)=\rho _T(\omega ,𝐤)/\omega `$. The terms that grow linearly in time are recognized as those emerging in Fermi’s golden rule from elementary time-dependent perturbation theory in quantum mechanics. We note that $`p(0)=0`$ and $`p^{}(0)0`$, thus the first integral above gives a contribution to the real part of the mean field (i.e., the amplitude) that features a logarithmic secular term, whereas the second integral contributes to the imaginary part (i.e., the oscillating phase) with a linear secular term, which as will be recognized below determines a perturbative shift of the oscillation frequency.
Substituting $`\rho _T(\omega ,𝐤)`$ into Eq. (37), we obtain
$`e^2𝐚_T^{(1)}(𝐤,t)`$ $`=`$ $`{\displaystyle \underset{\lambda =1}{\overset{2}{}}}A_\lambda (𝐤)𝓔_\lambda (𝐤)e^{ikt}\left[{\displaystyle \frac{e^2T^2}{12k^2}}(\mathrm{ln}2kt+\gamma _E1)+i\delta _kt+𝒪\left(t^1\right)\right]+\text{c.c.},`$ (42)
where
$$\delta _k\frac{\text{ Re}\mathrm{\Pi }_T(\omega ,𝐤)}{2\omega }|_{\omega =k}=\frac{e^2T^2}{12k}.$$
(43)
Secular divergences are an ubiquitous feature in the perturbative solution of differential equations with oscillatory behaviour. A resummation method that improves the asymptotic solution of these differential equations was introduced in Ref. which is interpreted as a dynamical renormalization group. This method is very powerful and not only allows a consistent resummation of the secular terms in the perturbative expansion but provides also a consistent reduction of the dynamics to the slow degrees of freedom as is explicitly explained in the work of Chen et. al. and Kunihiro et. al. in. This method was recently generalized to the realm of quantum field theory as a dynamical renormalization group to resum the perturbative series for the real time evolution of non-equilibrium expectation values . This generalization of the work of reference for differential equations to non-equilibrium quantum field theory is a major step since in non-equilibrium systems the equations of motion for expectation values are non-local and as shown in detail in this work require a resummation of Feynman diagrams.
While the linear secular terms have a natural interpretation in terms of renormalization of the mass (the imaginary part) or a quasiparticle width (the real part) , the logarithmic secular term found above is akin to those found in Refs. that lead to anomalous relaxation. Furthermore the origin of these logarithmic secular terms is similar to the threshold infrared divergences and threshold enhancement of the spectral density due to the presence of nearby poles .
Thus following the method presented in Refs. we implement the dynamical renormalization group by introducing a (complex) amplitude renormalization factor in the following manner,
$`A_\lambda (𝐤)`$ $`=`$ $`𝒵_k(\tau )𝒜_\lambda (𝐤,\tau ),`$ (44)
where $`𝒵_k(\tau )=1+e^2z_k^{(1)}(\tau )+𝒪(e^4)`$ is a multiplicative renormalization constant, $`𝒜_\lambda (𝐤,\tau )`$ is the renormalized initial value, and $`\tau `$ is an arbitrary renormalization scale at which the secular divergences are cancelled . Choosing
$$e^2z_k^{(1)}(\tau )=\frac{e^2T^2}{12k^2}(\mathrm{ln}2k\tau +\gamma _E1)+i\delta _k\tau ,$$
we obtain to lowest order in $`e^2T^2/k^2`$ that the solution of the equation of motion is given by
$`𝐚_T(𝐤,t)`$ $`=`$ $`{\displaystyle \underset{\lambda =1}{\overset{2}{}}}𝒜_\lambda (𝐤,\tau )𝓔_\lambda (𝐤)e^{ikt}\left[1{\displaystyle \frac{e^2T^2}{12k^2}}\mathrm{ln}{\displaystyle \frac{t}{\tau }}i\delta _k(t\tau )\right]+\text{ c.c.}`$ (46)
$`+\text{nonsecular terms},`$
which remains bounded at large times $`t`$ provided that $`\tau `$ is chosen arbitrarily close to $`t`$. The solution does not depend on the renormalization scale $`\tau `$ and this independence leads to the dynamical renormalization group equation, which to this order is given by
$$\left[\frac{}{\tau }+\frac{e^2T^2}{12k^2\tau }+i\delta _k\right]𝒜_\lambda (𝐤,\tau )=0,$$
with the solution
$$𝒜_\lambda (𝐤,\tau )=𝒜_\lambda (𝐤,\tau _0)e^{i\delta _k(\tau \tau _0)}\left(\frac{\tau }{\tau _0}\right)^{\frac{e^2T^2}{12k^2}},$$
where $`\tau _0`$ is the time scale such that this intermediate asymptotic solution is valid and physically corresponds to a microscopic scale, i.e, $`\tau _01/k`$. Finally, setting $`\tau =t`$ in Eq. (46) we find that $`𝐚_T(𝐤,t)`$ evolves at intermediate asymptotic times $`t1/k`$ as
$$𝐚_T(𝐤,t)=\underset{\lambda =1}{\overset{2}{}}𝒜_\lambda (𝐤,\tau _0)𝓔_\lambda (𝐤)e^{i(k+\delta _k)(t\tau _0)}\left(\frac{t}{\tau _0}\right)^{\frac{e^2T^2}{12k^2}}+\mathrm{c}.\mathrm{c}..$$
(47)
From Eq. (43) we see that $`\delta _k`$ is consistent with the photon thermal mass $`m_\gamma ^2=e^2T^2/6`$ for $`k^2m_\gamma ^2`$.
As discussed in detail in Ref. the dynamical renormalization group solution (47) is also obtained via the Fourier transform of the renormalization group improved propagator in frequency-momentum space, hence the above solution corresponds to a renormalization group improved resummation of the self-energy .
This novel anomalous power law relaxation of the photon mean field will be confirmed below in our study of the kinetics of the photon distribution function in the linearized approximation.
This anomalous power law relaxation is obviously very slow in the semihard regime in which the HTL approximation and perturbation theory is valid. At higher orders we expect exponential relaxation due to collisional processes, which emerges from linear secular terms in a perturbative solution of the real-time equations of motion. The power law relaxation will then compete with the exponential relaxation and we expect a crossover time scale at which relaxation will change from a power law to an exponential. Clearly an assessment of this time scale requires a detailed calculation of higher order contributions which we expect to study in a forthcoming article.
A related crossover of behavior will be found below for the evolution of the photon distribution function and photon production.
#### 3 Secular terms from the Laplace transform for semihard photons
We show here how secular terms for semihard photons ($`eTkT`$) can be obtained from the Laplace transform representation of the photon mean field given by Eqs. (25)-(28) for large times $`kt1`$.
The integrand in Eq. (25) can be expanded as follows for semihard momentum,
$$\frac{1}{\omega }\beta _T(\omega ,k)\stackrel{eTkT}{=}\frac{e^2T^2}{12k^3}\frac{1}{\omega ^2k^2}\left[1+𝒪\left(\frac{e^2T^2}{k^2}\right)\right].$$
Inserting this expansion in Eq. (25) gives for $`𝐚_T^{\mathrm{𝑐𝑢𝑡}}(𝐤,t)`$,
$$𝐚_T^{\mathrm{𝑐𝑢𝑡}}(𝐤,t)=\frac{e^2T^2}{12k}\text{PV}_k^{+k}\frac{d\omega }{\omega ^2k^2}e^{i\omega t}𝐚_T(𝐤,0)\left[1+𝒪\left(\frac{e^2T^2}{k^2}\right)\right].$$
where PV stands for the principal value prescription. This integral yields for asymptotic times to first order in $`e^2T^2/k^2`$
$$𝐚_T^{\mathrm{𝑐𝑢𝑡}}(𝐤,t)\stackrel{kt1}{=}\frac{e^2T^2}{12k^2}\mathrm{cos}kt\left[\mathrm{ln}2kt+\gamma _E\right]𝐚_T(𝐤,0)+𝒪(\frac{1}{t},\frac{e^4T^4}{k^4}),$$
where we used ,
$$_0^{\mathrm{}}\frac{dy}{y}\left(1\mathrm{cos}yt\right)p(y)\stackrel{t\mathrm{}}{=}p(0)\left[\mathrm{ln}(\mu t)+\gamma \right]+_0^{\mathrm{}}\frac{dy}{y}\left[p(y)p(0)\theta (\mu y)\right]+𝒪\left(\frac{1}{t}\right).$$
The pole contribution (24) to first order in $`e^2T^2/k^2`$ reads
$$𝐚_T^{\mathrm{𝑝𝑜𝑙𝑒}}(𝐤,t)=𝐚_T(𝐤,0)\left[\mathrm{cos}kt\frac{e^2T^2}{12k^2}t\mathrm{sin}kt\right]+\text{nonsecular terms}+𝒪\left(\frac{e^4T^4}{k^4}\right).$$
Collecting pole and cut contributions we find
$`𝐚_T(𝐤,t)`$ $`=`$ $`𝐚_T(𝐤,0)\left\{\mathrm{cos}kt{\displaystyle \frac{e^2T^2}{12k^2}}\left[\left(\mathrm{ln}2kt+\gamma _E\right)\mathrm{cos}kt+t\mathrm{sin}kt\right]\right\}`$ (49)
$`+\text{nonsecular terms}+𝒪\left({\displaystyle \frac{e^4T^4}{k^4}}\right).`$
This result coincides with the perturbative solution (35), (42) and (43) after imposing the initial conditions (14) used for the Laplace transform to Eq. (35). That is,
$$A_\lambda (𝐤)=A_\lambda ^{}(𝐤).$$
In summary, the Laplace transform solution permits to compute the secular terms as well as the perturbative method used in sec. II.A2.
### B Quantum kinetics of the photon distribution function
As mentioned in Sec. II the first step towards a kinetic description is to identify the proper degrees of freedom (quasiparticles) and the corresponding microscopic time scale. For semihard photons of momenta $`eTkT`$, it is adequate to choose the free photons as the quasiparticles with the corresponding microscopic scale $`1/k`$. A kinetic description of the non-equilibrium evolution of the distribution function assumes a wide separation between the microscopic and the relaxation time scales. In the case under consideration the effective small coupling in the semihard momentum limit is $`e^2T^2/k^21`$. Furthermore assuming that $`e1`$ then in this regime both the HTL and the perturbative approximation are valid.
We begin by obtaining the photon number operator from the Heisenberg field operator and its conjugate momentum. Write
$`𝐀_T(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^{3/2}}𝐀_T(𝐤,t)e^{i𝐤𝐱}},`$
$`𝐏_T(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^{3/2}}𝐏_T(𝐤,t)e^{i𝐤𝐱}},`$
with
$`𝐀_T(𝐤,t)`$ $`=`$ $`{\displaystyle \underset{\lambda =1}{\overset{2}{}}}\sqrt{{\displaystyle \frac{1}{2k}}}\left[a_\lambda (𝐤,t)𝓔_\lambda (𝐤)+a_\lambda ^{}(𝐤,t)𝓔_\lambda (𝐤)\right],`$
$`𝐏_T(𝐤,t)`$ $`=`$ $`i{\displaystyle \underset{\lambda =1}{\overset{2}{}}}\sqrt{{\displaystyle \frac{k}{2}}}\left[a_\lambda ^{}(𝐤,t)𝓔_\lambda (𝐤)a_\lambda (𝐤,t)𝓔_\lambda (𝐤)\right],`$
where $`a_\lambda (𝐤,t)`$ \[$`a_\lambda ^{}(𝐤,t)`$\] is the annihilation (creation) operator that destroys (creates) a free photon of momentum $`𝐤`$ and polarization $`\lambda `$ at time $`t`$. The polarization-averaged number operator $`N_\gamma (𝐤,t)`$ that counts the semihard photons of momentum $`𝐤`$ is then defined by
$`N_\gamma (𝐤,t)`$ $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\lambda =1}{\overset{2}{}}}a_\lambda ^{}(𝐤,t)a_\lambda (𝐤,t)`$
$`=`$ $`{\displaystyle \frac{1}{4k}}\{𝐏_T(𝐤,t)𝐏_T(𝐤,t)+k^2𝐀_T(𝐤,t)𝐀_T(𝐤,t)`$
$`+ik[𝐀_T(𝐤,t)𝐏_T(𝐤,t)𝐏_T(𝐤,t)𝐀_T(𝐤,t)]\}.`$
The expectation value of this number operator is interpreted as the number of photons per polarization per unit phase space volume
$$n_𝐤^\gamma (t)=N_\gamma (𝐤,t)(2\pi )^3\frac{dN}{d^3xd^3k},$$
where $`N`$ is the total number of photons per polarization in the plasma. Taking the time derivative of $`N_\gamma (𝐤,t)`$ and using the Heisenberg equations of motion, we obtain the following expression for the time derivative of the expectation value
$$\dot{n}_𝐤^\gamma (t)=\underset{t^{}t}{lim}\frac{e}{4k}\left(\frac{}{t^{}}ik\right)\frac{d^3q}{(2\pi )^{3/2}}\overline{\psi }^{}(𝐩,t)𝜸𝐀_T^+(𝐤,t^{})\psi ^{}(𝐪,t)+\text{c.c.},$$
(50)
where the “$`+`$” (“$``$”) superscripts for the fields refer to fields defined in the forward (backward) time branch in the CTP formulation. We have separated the time arguments and the fields on different branches to be able to extract the time derivative from the expectation value.
The nonequilibrium expectation values on the right-hand side (RHS) of Eq. (50) can be computed perturbatively in powers of $`e`$ using the nonequilibrium Feynman rules and real-time propagators.
We assume the initial density matrix is diagonal in the basis of the photon occupation numbers with nonequilibrium initial photon populations given by $`n_𝐤^\gamma (t_0)`$ and there is no initial photon polarization asymmetry. Furthermore we assume that the fermions are in thermal equilibrium at a temperature $`T`$.
At $`𝒪(e)`$ the RHS of Eq. (50) vanishes identically. This is a consequence of our choice of initial density matrix diagonal in the basis of the photon number operator. To $`𝒪(e^2)`$, we obtain
$$\dot{n}_𝐤^\gamma (t)=[1+n_𝐤^\gamma (t_0)]\mathrm{\Gamma }_k^<(t)n_𝐤^\gamma (t_0)\mathrm{\Gamma }_k^>(t),$$
(51)
where the time-dependent rates are given by
$`\mathrm{\Gamma }_𝐤^{>\text{ }<}(t)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\omega _\gamma ^{>\text{ }<}(\omega ,k){\displaystyle \frac{\mathrm{sin}[(\omega k)(tt_0)]}{\omega k}},`$ (52)
$`_\gamma ^<(\omega ,k)`$ $`=`$ $`{\displaystyle \frac{e^2}{k}}{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}[[1(\widehat{𝐤}\widehat{𝐩})(\widehat{𝐤}\widehat{𝐪})]\{n_F(p)[1n_F(q)]\delta (\omega p+q)`$ (55)
$`+[1n_F(p)]n_F(q)\delta (\omega +pq)\}+[1+(\widehat{𝐤}\widehat{𝐩})(\widehat{𝐤}\widehat{𝐪})]`$
$`\times \{n_F(p)n_F(q)\delta (\omega pq)+[1n_F(p)][1n_F(q)]\delta (\omega +p+q)\}],`$
with $`𝐩=𝐤+𝐪`$ and $`_\gamma ^>(\omega ,k)`$ is obtained from $`_\gamma ^<(\omega ,k)`$ through the replacement $`n_F1n_F`$.
A comment here is in order. As explained above we are focusing on the leading HTL approximation, consequently, in obtaining $`_\gamma ^{>\text{ }<}(\omega ,k)`$ we use the free real-time fermion propagators which correspond to the hard part of the fermion loop momentum. In general there are contributions from the soft fermion loop momentum region which will require to use the HTL-resummed fermion propagators for consistency. We will postpone the detailed study of the contribution from soft loop momentum to a forthcoming article, focusing here on the comparison between the lowest order calculation in real time and the result available in the literature for the photon production rate that include HTL corrections.
Since the fermions are in thermal equilibrium, the Kubo-Martin-Schwinger (KMS) condition holds:
$$_\gamma ^>(\omega ,k)=e^{\beta \omega }_\gamma ^<(\omega ,k),$$
(56)
where $`\beta =1/T`$. It is apparent to see that $`_\gamma ^{<(>)}(\omega ,k)`$ has a physical interpretation in terms of the off-shell photon production (absorption) processes in the plasma. The first two terms in $`_\gamma ^<(\omega ,k)`$ describe the process that a fermion (or an anti-fermion) emits a photon, i.e, bremsstrahlung from the fermions in the medium, the third term describes annihilation of a fermion pair into a photon, and the fourth terms describes creation of a photon and a fermion pair out of the vacuum. The corresponding terms in $`_\gamma ^>(\omega ,k)`$ describe the inverse processes.
As argued, for semihard photons of momenta $`eTkT`$, the leading contribution of $`_\gamma ^{>\text{ }<}(\omega ,k)`$ arises from the hard loop momenta $`qk`$. A detailed analysis of $`_\gamma ^<(\omega ,k)`$ along the familiar lines in the hard thermal loop program shows that in the HTL approximation ($`qk`$)
$$_\gamma ^<(\omega ,k)|_{\mathrm{HTL}}=\frac{e^2T^3}{12k^2}\left(1\frac{\omega ^2}{k^2}\right)\theta (k^2\omega ^2).$$
(57)
Thus we recognize that in the HTL limit $`_\gamma ^{<(>)}(\omega ,k)`$ is completely determined to the off-shell Landau damping process in which a hard fermion in the plasma emits (absorbs) a semihard photon, i.e, bremsstrahlung from the fermions in the medium.
#### Dynamical Renormalization Group and the emergence of detailed balance
We now turn to the kinetics of semihard photons. To obtain a kinetic equation from Eq. (51), we implement the dynamical renormalization group resummation as introduced in Refs. . Direct integration with the initial condition yields
$`n_𝐤^\gamma (t)`$ $`=`$ $`n_𝐤^\gamma (t_0)+\left[1+n_𝐤^\gamma (t_0)\right]{\displaystyle _{t_0}^t}𝑑t^{}\mathrm{\Gamma }_k^<(t^{})n_𝐤^\gamma (t_0){\displaystyle _{t_0}^t}𝑑t^{}\mathrm{\Gamma }_k^>(t^{}).`$ (58)
The integrals that appear in the above expression:
$$_{t_0}^t𝑑t^{}\mathrm{\Gamma }_𝐤^{>\text{ }<}(t^{})=_{\mathrm{}}^+\mathrm{}𝑑\omega _\gamma ^{>\text{ }<}(\omega ,k)\frac{1\mathrm{cos}[(\omega k)(tt_0)]}{(\omega k)^2},$$
(59)
are dominated, in the long time limit, by the regions of $`\omega `$ for which the denominator is resonant, i.e, $`\omega k`$. The time dependence in the above integral along with the resonant denominator is the familiar form that leads to Fermi’s golden rule in elementary time dependent perturbation theory. In the long time limit $`tt_01/k`$ Fermi’s golden rule approximates the above integrals by
$$_{t_0}^t𝑑t^{}\mathrm{\Gamma }_𝐤^{>\text{ }<}(t^{})\frac{\pi }{2}(tt_0)_\gamma ^{>\text{ }<}(\omega =k,k)0.$$
Therefore the rate of photon production, i.e, the coefficient of the linear time dependence vanishes at this order because of the vanishing of the imaginary part of the photon self-energy on the photon mass shell. However the use of Fermi’s golden rule, which is the usual approach to extract (time-independent) rates, misses the important off-shell effects associated with finite lifetime processes. These can be understood explicitly by using the first of formulae (39). We find that for $`tt_01/k`$
$`{\displaystyle _{t_0}^t}𝑑t^{}\mathrm{\Gamma }_k^<(t^{})`$ $`\stackrel{k\left(tt_0\right)1}{=}`$ $`{\displaystyle \frac{e^2T^3}{6k^3}}\left\{\mathrm{ln}\left[2k(tt_0)\right]+\gamma _E1\right\}+𝒪\left({\displaystyle \frac{1}{k(tt_0)}}\right)`$ (60)
$`{\displaystyle _{t_0}^t}𝑑t^{}\mathrm{\Gamma }_k^>(t^{})`$ $`\stackrel{k\left(tt_0\right)1}{=}`$ $`e^{\beta k}{\displaystyle _{t_0}^t}𝑑t^{}\mathrm{\Gamma }_k^<(t^{}).`$ (61)
The second line of this equation displays the condition for detailed balance and holds for time scales $`tt_01/k`$. We emphasize that this condition is a consequence of the fact that the region $`\omega k`$ (i.e, the resonant denominator) dominates the long time behavior of the integrals in the time-dependent rates (59) and the KMS condition Eq. (56) holds under the assumption that the fermions are in thermal equilibrium.
Thus we obtain one of the important results of this article the time-dependent rates obey a detailed balance and that detailed balance emerges in the intermediate asymptotic regime $`tt_01/k`$ when the secular terms dominate the long time behavior. This is a noteworthy result because it clearly states that detailed balance emerges on microscopic time scales $`tt_0>1/k`$ at which the secular terms dominate the integrals but for which a perturbative expansion is still valid. Detailed balance then guarantees the existence of an asymptotic equilibrium solution which is reached on time scales of the order of the relaxation time.
The (logarithmic) secular terms in the time-dependent rates (61) can now be resummed using the dynamical renormalization group method by introducing a multiplicative renormalization of the distribution function
$`n_𝐤^\gamma (t_0)`$ $`=`$ $`𝒵_k(t_0,\tau )n_𝐤^\gamma (\tau ),`$
$`𝒵_k(t_0,\tau )`$ $`=`$ $`1+e^2z_k^{(1)}(t_0,\tau )+𝒪(e^4),`$
thus rewriting Eq. (58) consistently to $`𝒪(e^2)`$ as
$`n_𝐤^\gamma (t)`$ $`=`$ $`n_𝐤^\gamma (\tau )+e^2z_k^{(1)}(t_0,\tau )n_𝐤^\gamma (\tau )+\left[1+n_𝐤^\gamma (\tau )\right]{\displaystyle _{t_0}^t}𝑑t^{}\mathrm{\Gamma }_k^<(t^{})n_𝐤^\gamma (\tau ){\displaystyle _{t_0}^t}𝑑t^{}\mathrm{\Gamma }_k^>(t^{})+𝒪(e^4).`$ (62)
The renormalization coefficient $`z_k^{(1)}(t_0,\tau )`$ is chosen to cancel the secular divergence at a time scale $`t=\tau `$. Thus to $`𝒪(e^2)`$, the choice
$`e^2z_k^{(1)}(t,\tau )n_𝐤^\gamma (\tau )`$ $`=`$ $`\left[1+n_𝐤^\gamma (\tau )\right]{\displaystyle _{t_0}^\tau }𝑑t^{}\mathrm{\Gamma }_k^<(t^{})+n_𝐤^\gamma (\tau ){\displaystyle _{t_0}^\tau }𝑑t^{}\mathrm{\Gamma }_k^>(t^{}),`$
leads to an improved perturbative solution in terms of the “updated” occupation number $`n_𝐤^\gamma (\tau )`$,
$`n_𝐤^\gamma (t)`$ $`=`$ $`n_𝐤^\gamma (\tau )+\left[1+n_𝐤^\gamma (\tau )\right]{\displaystyle _\tau ^t}𝑑t^{}\mathrm{\Gamma }_k^<(t^{})n_𝐤^\gamma (\tau ){\displaystyle _\tau ^t}𝑑t^{}\mathrm{\Gamma }_k^>(t^{})+𝒪(e^4),`$ (63)
which is valid for large times $`tt_0`$ provided that $`\tau `$ is chosen arbitrarily close to $`t`$. A change in the time scale $`\tau `$ in the integrals is compensated by a change of the $`n_𝐤^\gamma (\tau )`$ in such a manner that $`n_𝐤^\gamma (t)`$ does not depend on the arbitrary scale $`\tau `$. This independence leads to the dynamical renormalization group equation which consistently to $`𝒪(e^2)`$ is given by
$$\frac{d}{d\tau }n_𝐤^\gamma (\tau )=[1+n_𝐤^\gamma (\tau )]\mathrm{\Gamma }_k^<(\tau )n_𝐤^\gamma (\tau )\mathrm{\Gamma }_k^>(\tau )+𝒪(e^4).$$
(64)
Choosing $`\tau `$ to coincide with $`t`$ in Eq. (64), we obtain the quantum kinetic equation to order $`e^2`$ given by
$$\dot{n}_𝐤^\gamma (t)=[1+n_𝐤^\gamma (t)]\mathrm{\Gamma }_k^<(t)n_𝐤^\gamma (t)\mathrm{\Gamma }_k^>(t).$$
(65)
For intermediate asymptotic times $`tt_01/k`$ at which the logarithmic secular terms dominate the integrals for the time-dependent rates and detailed balance emerges, we find
$`\mathrm{\Gamma }_k^<(t)`$ $`\stackrel{k\left(tt_0\right)1}{=}`$ $`{\displaystyle \frac{e^2T^3}{6k^3}}{\displaystyle \frac{1}{tt_0}}\left[1+𝒪\left({\displaystyle \frac{1}{k(tt_0)}}\right)\right],`$ (66)
$`\mathrm{\Gamma }_k^>(t)`$ $`\stackrel{k\left(tt_0\right)1}{=}`$ $`{\displaystyle \frac{e^2T^3}{6k^3}}{\displaystyle \frac{e^{\beta k}}{tt_0}}\left[1+𝒪\left({\displaystyle \frac{1}{k(tt_0)}}\right)\right],`$ (67)
where the detailed balance relation is explicitly displayed and the terms being neglected are oscillatory on time scales $`1/k`$ and fall off faster.
The detailed balance relation between the time-dependent rates (65) guarantees the existence of an equilibrium solution of the kinetic equation with the occupation number $`n_𝐤^\gamma (t=\mathrm{})=(e^{\beta k}1)^1`$.
The full solution of this kinetic equation is given by
$`n_𝐤^\gamma (t)`$ $`=`$ $`n_𝐤^\gamma (t_0)e^{_{t_0}^t𝑑t^{}\gamma _k(t^{})}+e^{_{t_0}^t𝑑t^{}\gamma _k(t^{})}{\displaystyle _{t_0}^t}𝑑t^{}\mathrm{\Gamma }_k^<(t^{})e^{_{t_0}^t^{}𝑑t^{\prime \prime }\gamma _k(t^{\prime \prime })},`$ (68)
$`\gamma _k(t)`$ $`=`$ $`\mathrm{\Gamma }_k^>(t)\mathrm{\Gamma }_k^<(t).`$ (69)
However, in order to understand the relaxation to equilibrium of the distribution function, we now focus on the linearized relaxation time approximation which describes the approach to equilibrium of the distribution function of one k-mode which has been displaced slightly off equilibrium while all other modes (and the fermions) are in equilibrium. Thus writing $`n_𝐤^\gamma (t_0)=n_B(k)+\delta n_𝐤^\gamma (t_0)`$ with $`n_B(k)`$ the Bose-Einstein distribution function, which as shown above is an equilibrium solution as a consequence of detailed balance (61), we obtain
$$\delta \dot{n}_𝐤^\gamma (t)=\delta n_𝐤^\gamma (t)(e^{\beta k}1)\mathrm{\Gamma }_k^<(t).$$
For semihard photons of momentum $`eTkT`$, we can simply replace $`e^{\beta k}1`$ by $`k/T`$ and upon integration we find
$$\delta n_𝐤^\gamma (t)=\delta n_𝐤^\gamma (t_0)\left(\frac{tt_0}{\tau _0}\right)^{\frac{e^2T^2}{6k^2}}\text{ for}k(tt_0)1,$$
(70)
where $`\tau _01/k`$.
A noteworthy feature of Eq. (70) is that the relaxation of the distribution function for semihard photons in the linearized approximation is governed by a power law with an anomalous exponent rather than by exponential damping. This result is similar to that found in scalar quantum electrodynamics in the Markovian approximation . Comparing Eq. (70) and Eq. (47), we clearly see that the anomalous exponent for the relaxation of the photon distribution function in the linearized approximation is twice that for the linear relaxation of the mean field (47).
This relationship is well known in the case of exponential relaxation but our analysis with the dynamical renormalization group reveals it to be a more robust feature, applying just as well to power law relaxation.
Our approach to derive kinetic equations in field theory is different from the one often used in the literature which involves a Wigner transform and assumptions about the separation of fast and slow variables . Our work differs in many important respects from these other formulations in that it reveals clearly the dynamics of off-shell effects associated with non-exponential relaxation. This aspects will acquire phenomenological relevance in our study of photon production by off-shell effects in a plasma with a finite lifetime in the next section.
### C Photon production
An important and phenomenologically relevant byproduct of the study of the kinetics of the photon distribution function is to give an assessment of photon production via off-shell effects by the lightest (up and down) quarks from the quark-gluon plasma in which quarks and gluons (but not photons) are in thermal equilibrium. To the order under consideration the only contribution to the photon self-energy is a fermion loop and although we have computed it assuming these fermions to be electrons, we can use this result to study photon production from thermalized quarks by simply accounting for the electromagnetic charges of the up and down quarks and for the number of colors. To this order $`𝒪(\alpha )`$ there are no gluon contributions to the photon self-energy and therefore to photon production and the rate for photon production obtained from the imaginary part of the photon self-energy evaluated on the photon mass shell vanishes. This can be seen directly from Eqs. (57) and (59) which show explicitly that the usual (time-independent) rate obtained from Fermi’s golden rule and determined by the imaginary part of the photon self-energy on the photon mass shell vanishes. The underlying assumption that motivates the use of Fermi’s golden rule is that the lifetime of the system is much larger than typical observation times and the limit of $`t\mathrm{}`$ can be taken replacing the oscillatory factors with resonant denominators in Eq. (59) by a delta function. This obviously extracts the leading time dependence but ignores corrections arising from the finite lifetime of the system under consideration. A quark-gluon plasma formed in an ultrarelativistic heavy ion collision is estimated to have a lifetime $`50\text{fm}/c`$ at RHIC energies thus photons are produced by a source of a finite lifetime and off-shell effects could lead to considerable photon production during the lifetime of the QGP perhaps comparable to those extracted from the usual Fermi’s golden rule results for infinite lifetime. The focus of this section is precisely the study of this possibility .
From the expression for the number of photons per phase space volume given by Eq. (68) we can obtain the number of photons produced at a given time $`t`$ per unit phase space from an initial (photon) vacuum state, by setting $`n_𝐤^\gamma (t_0)=0`$ and including the charge of the up and down quarks and the number of colors. Assuming that photons do not thermalize in the quark-gluon plasma, i.e, that their mean free path is much longer than the size of the plasma, we neglect the exponentials in the second term of Eq. (68) which are responsible for building the photon population up to the equilibrium value. To lowest order in the electromagnetic fine structure constant we find from Eq. (68) the total number of semihard photons at time $`t`$ per invariant phase space volume summed over the two polarizations given by
$$k_{t_0}^t𝑑t^{}\frac{dN_{\mathrm{tot}}(t^{})}{d^3xd^3kdt^{}}\frac{2k}{(2\pi )^3}n_𝐤^\gamma (t)\stackrel{k\left(tt_0\right)1}{=}\frac{5\alpha T^3}{18\pi ^2k^2}\left[\mathrm{ln}2k(tt_0)+\gamma _E1\right]+𝒪\left[(tt_0)^1\right].$$
(71)
This is a noteworthy result and one of the main points of this article, we find photon production to lowest (one-loop) order arising solely from off-shell effects: ($`qq\gamma `$ and $`\overline{q}\overline{q}\gamma `$) and annihilation of quarks ($`q\overline{q}\gamma `$). In the HTL limit, valid for semihard photons, the leading contribution arises from Landau damping, i.e, from quarks in the medium.
In the usual approach the photon production rate is obtained from the imaginary part of the photon self-energy on-shell and receives contributions only at order $`𝒪(\alpha \alpha _s)`$ and beyond. The lowest order contribution corresponds to a quark loop with a self-energy insertion from gluons as well as electromagnetic vertex correction by gluons. As a result the imaginary part of the photon self-energy on the mass shell arises from Compton scattering ($`qgq\gamma `$ and $`\overline{q}g\overline{q}\gamma `$) and pair-annihilation ($`q\overline{q}g\gamma `$. The lowest order $`𝒪(\alpha \alpha _s)`$ contribution displays an infrared divergence for massless quarks that is screened by the HTL resummation leading to the result quoted in Ref. for the hard photon production rate : <sup>§</sup><sup>§</sup>§We follow Kapusta et al. in adding 1 to the argument of the logarithm so as to make it agree with the numerical result for $`kT`$. For details, see discussion below Eq. (41) in Ref .
$$k\frac{dN_{\mathrm{tot}}(t)}{d^3xd^3kdt}|_{\text{on-shell}}=\frac{5}{18\pi ^2}\alpha \alpha _sT^2e^{k/T}\mathrm{ln}\left(\frac{2.912}{g^2}\frac{k}{T}+1\right),$$
where $`\alpha _s=g^2/4\pi `$ with $`g`$ the QCD coupling constant. The total number of photons produced per (invariant) phase space at time $`t`$ is then given by
$$k_{t_0}^t𝑑t^{}\frac{dN_{\mathrm{tot}}(t^{})}{d^3xd^3kdt^{}}|_{\text{on-shell}}=\frac{5}{18\pi ^2}\alpha \alpha _sT^2(tt_0)e^{k/T}\mathrm{ln}\left(\frac{2.912}{g^2}\frac{k}{T}+1\right),$$
(72)
which can obviously be interpreted as a linear secular term in the photon distribution function, which is the usual situation corresponding to exponential relaxation of the photon distribution function as well as for the photon mean field .
#### Comparison between on- and off-shell photon production
The important question to address is how do the on- and off-shell contributions to photon production compare? Obviously the answer to this question depends on the time scale involved, eventually at long time scales the linear growth of the photon number with time from the on-shell contribution will overwhelm the logarithmic growth from the off-shell, but at early times the off-shell will dominate. Thus a crossover regime is expected from the off-shell to the on-shell dominance to photon production. This crossover time scale will depend on the temperature and the energy of the photons as well as on the numerical value of the electromagnetic and strong couplings.
There are two important time scales that are of phenomenological relevance to photon production in a quark-gluon plasma. Firstly the assumption that quarks are in thermal equilibrium restricts the earliest time scale to be somewhat larger than about $`13\text{fm}/c`$ which is the estimate for quark thermalization . Secondly the lifetime of the QGP phase is estimated to be somewhere in the range $`1050\text{fm}/c`$. Thus the relevant time scale for comparison between the off- and on-shell contributions is $`13\text{fm}/ctt_01050\text{fm}/c`$.
Before we begin the comparison, we point out that Eqs. (71) and (72) have different regimes of validity in photon momenta. Whereas Eq. (71) is valid in the semihard regime $`eTkT`$ where the HTL approximation is reliable, Eq. (72) is valid in the hard regime $`kT`$ . We have studied numerically the off-shell contribution to photon production for hard momentum $`kT`$ directly in terms of the time-dependent rates given by Eqs. (51) and (55) without using the HTL approximation and compared it to the HTL approximation for $`1\text{fm}/ct100\text{fm}/c`$ (with $`t_0=0`$). The result is depicted in Fig. 3, which displays clearly a logarithmic time dependence in both cases for $`t>12\text{fm}/c`$ but with different slopes, a consequence of the different temperature and momentum dependence of the coefficients of the logarithmic time dependence. The HTL approximation overestimates the total number of photons by at most a factor two in the relevant regime $`1\text{fm}/ct100\text{fm}/c`$ for hard photons. The reliability of the HTL approximation (71) in the semihard momentum region is confirmed by Fig. 4, which compares the number of photons obtained from the numerical evaluation with the full time-dependent rates and the result from the HTL approximation in the weak coupling and semihard momentum region.
We can now compare the off-shell contribution to that of Refs. , which is a result of on-shell processes and valid in the hard momentum limit. For this we focus on the relevant scenario of a quark-gluon plasma with thermalized quarks at a temperature $`T200\text{MeV}`$ and of lifetime $`t1020\text{fm}/c`$. These are approximately the temperature and temporal scales expected to be reached at RHIC. The value of $`\alpha _s`$ at this temperature is not know with much certainty but expected to be $`1`$. Following we choose $`\alpha _s=0.4`$ which corresponds to $`g^2=5`$, we find that for hard photons of momenta $`kT`$ the ratio of photon produced by off-shell processes is comparable to that produced by on-shell processes.
Fig. 5 depicts the ratio of the number of hard ($`kT`$) photons produced by off-shell processes in the HTL approximation to that from on-shell processes in the regime of lifetimes for a quark-gluon plasma phase $`5\text{fm}/c<t<100\text{fm}/c`$ (here we set $`t_0=0`$). Accounting for the overestimate of the HTL approximation for hard photons, we then conclude that off-shell processes such as from quarks in the medium (Landau damping) are just as important as on-shell processes. This is another of the important results of this work and that cannot be obtained from the usual approach to photon production based on the computation of a time-independent rate.
Furthermore, an important point worth emphasizing is that whereas the on-shell calculation displays infrared divergences at lowest order, the calculation directly in real time is infrared finite because time acts as an infrared cutoff. Thus a calculation of photon production directly in real time does not require an HTL improvement of the fermion propagators to cutoff an infrared divergence. However, the region of soft loop momentum does require the HTL resummation of one of the fermion propagators and will be studied in a future article.
A thorough study of off-shell effects and their contribution to photon production in a wider range of temperatures and momenta including screening corrections to the internal quark lines and the study of potential observables will be the subject of a longer study which we think is worthwhile on its own and on which we expect to report soon.
## IV Hard fermions out of equilibrium
To provide a complete picture of relaxation and non-equilibrium aspects of a hot QED plasma in real time, we now focus on a detailed study of relaxation of fermionic mean fields (as induced by an adiabatically switched-on Grassmann source) as well as the quantum kinetics of the fermion distribution function. In this section we assume that the photons are in thermal equilibrium, and since we work to leading order in the HTL approximation we can translate the results vis à vis to the case of equilibrated gluons. In particular we seek to study the possibility of anomalous relaxation as a result of the emission and absorption of magnetic photons. In Refs. it was found that the relaxation of fermionic excitations is anomalous and not exponential as a result of the emission and absorption of magnetic photons that are only dynamically screened by Landau damping. The study of the real-time relaxation of the fermionic mean fields in these references was cast in terms of the Bloch-Nordsieck approximation which replaces the gamma matrices by the classical velocity of the fermion. In Ref. the relaxation of a charged scalar mean field as well as the quantum kinetics of the distribution function of charged scalars in scalar electrodynamics were studied using the dynamical renormalization group, both the charged scalar mean field and the distribution function of charged particles reveal anomalous non-exponential relaxation as a consequence of emission and absorption of soft magnetic photons. While electric photons (plasmons) are screened by a Debye mass which cuts off their infrared contribution, magnetic photons are only dynamically screened by Landau damping and their emission and absorption dominates the infrared behavior of the fermion propagator.
While the dynamical renormalization group has been implemented in scalar theories it has not yet been applied to fermionic theories. Thus the purpose of this section is twofold, (i) to implement the dynamical renormalization group to study the relaxation and kinetics of fermions with a detailed discussion of the technical differences with the bosonic case and (ii) to focus on the real-time manifestation of the infrared singularities associated with soft magnetic photons.
### A Relaxation of the fermionic mean field
The equation of motion for a fermionic mean field is obtained by following the strategy described in section II. We begin by writing the fermionic field as
$`\mathrm{\Psi }^\pm (𝐱,t)`$ $`=`$ $`\psi (𝐱,t)+\chi ^\pm (𝐱,t),`$
$`\psi (𝐱,t)`$ $`=`$ $`\mathrm{\Psi }^\pm (𝐱,t).`$
Then using the tadpole method , with an external Grassmann source that is adiabatically switched-on from $`t=\mathrm{}`$ and switched-off at $`t=0`$, we find the Dirac equation for the spatial Fourier transform of the fermion mean field for $`t>0`$ given by
$$\left(i\gamma _0\frac{}{t}𝜸𝐤\right)\psi (𝐤,t)_{\mathrm{}}^t𝑑t^{}\mathrm{\Sigma }(𝐤,tt^{})\psi (𝐤,t^{})=0,$$
(73)
where $`\mathrm{\Sigma }(𝐤,tt^{})`$ is the retarded fermion self-energy. A comment here is in order. To facilitate the study and maintain notational simplicity, in obtaining the equation we neglect the contribution from the instantaneous Coulomb interaction which is irrelevant to the relaxation of the mean field and only results in a perturbative frequency shift.
As noted above the relaxation of hard fermions is dominated by the soft photon contributions, thus in a perturbative expansion one needs to use the HTL-resummed photon propagators to account for the screening effects in the medium. To one-loop order but with the HTL-resummed photon propagators given in the Appendix, $`\mathrm{\Sigma }(𝐤,tt^{})`$ reads
$`\mathrm{\Sigma }(𝐤,tt^{})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\omega \left[i\gamma _0\rho _1(\omega ,𝐤)\mathrm{cos}[\omega (tt^{})]+𝜸\widehat{𝐤}\rho _2(\omega ,𝐤)\mathrm{sin}[\omega (tt^{})]\right],`$ (74)
where the spectral densities,
$`\rho _1(\omega ,𝐤)`$ $`=`$ $`e^2{\displaystyle \frac{d^3q}{(2\pi )^3}_{\mathrm{}}^+\mathrm{}𝑑q_0[1+n_B(q_0)n_F(p)]}`$ (76)
$`\times \left[\stackrel{~}{\rho }_T(q_0,q)+{\displaystyle \frac{1}{2}}\stackrel{~}{\rho }_L(q_0,q)\right][\delta (\omega pq_0)+\delta (\omega +p+q_0)],`$
$`\rho _2(\omega ,𝐤)`$ $`=`$ $`e^2{\displaystyle \frac{d^3q}{(2\pi )^3}_{\mathrm{}}^+\mathrm{}𝑑q_0[1+n_B(q_0)n_F(p)]}`$ (78)
$`\times \left[(\widehat{𝐤}\widehat{𝐪})(\widehat{𝐩}\widehat{𝐪})\stackrel{~}{\rho }_T(q_0,q){\displaystyle \frac{\widehat{𝐤}\widehat{𝐩}}{2}}\stackrel{~}{\rho }_L(q_0,q)\right][\delta (\omega pq_0)\delta (\omega +p+q_0)],`$
with $`𝐩=𝐤𝐪`$. Here $`\stackrel{~}{\rho }_T(q_0,q)`$ is the HTL-resummed spectral density for transverse photon propagator defined in Eq. (30) and $`\stackrel{~}{\rho }_L(q_0,q)`$ is the HTL-resummed spectral density for longitudinal photon propagator
$`\stackrel{~}{\rho }_L(q_0,q)`$ $`=`$ $`\text{sgn}(q_0)Z_L(q)\delta [q_0^2\omega _L^2(q)]+\beta _L(q_0,q)\theta (q^2q_0^2),`$ (79)
$`\beta _L(q_0,q)`$ $`=`$ $`{\displaystyle \frac{\frac{e^2T^2}{6}\frac{q_0}{q}}{\left[q^2+\frac{e^2T^2}{6}\left(2\frac{q_0}{q}\mathrm{ln}\frac{q+q_0}{qq_0}\right)\right]^2+\left[\frac{\pi e^2T^2}{6}\frac{q_0}{q}\right]^2}},`$ (80)
where $`\omega _L(q)`$ is the plasmon (longitudinal photon) pole and $`Z_L(q)`$ is the corresponding residue . It is worth pointing out that $`\rho _{1(2)}(\omega ,𝐤)`$ is an even (odd) function of $`\omega `$, a property that will be useful in the following analysis. Furthermore, to establish a connection with results in the literature, it proves convenient to introduce the Laplace transform of the retarded self-energy $`\stackrel{~}{\mathrm{\Sigma }}(s,𝐤)`$ just as in Eq. (16) and its analytic continuation $`\mathrm{\Sigma }(\omega ,𝐤)`$ as in Eq. (21) which is given by
$$\mathrm{\Sigma }(\omega ,𝐤)=_{\mathrm{}}^+\mathrm{}\frac{d\nu }{\nu \omega i0^+}\left[\gamma _0\rho _1(\nu ,𝐤)+𝜸\widehat{𝐤}\rho _2(\nu ,𝐤)\right].$$
Following the same strategy in the study of the photon mean field, we define $`\sigma (𝐤,tt^{})`$ as
$$\mathrm{\Sigma }(𝐤,tt^{})=\frac{}{t^{}}\sigma (𝐤,tt^{}),$$
and rewrite Eq. (73) as an initial value problem
$`\left(i\gamma _0{\displaystyle \frac{}{t}}𝜸𝐤\right)\psi (𝐤,t)`$ $`+`$ $`{\displaystyle _0^t}𝑑t^{}\sigma (𝐤,tt^{})\dot{\psi }(𝐤,t^{})\sigma (𝐤,0)\psi (𝐤,t)=0,`$ (81)
with the initial conditions $`\psi (𝐤,0)=\psi _0(𝐤)`$ and $`\dot{\psi }(𝐤,t<0)=0`$.
We are now ready to solve the equation of motion by perturbative expansion in powers of $`e^2`$ just as in the case of the gauge mean field. Let us begin by writing
$`\psi (𝐤,t)=\psi ^{(0)}(𝐤,t)+e^2\psi ^{(1)}(𝐤,t)+𝒪(e^4),`$ (82)
$`\sigma (𝐤,tt^{})=e^2\sigma ^{(1)}(𝐤,tt^{})+𝒪(e^4),`$ (83)
we obtain a hierarchy of equations:
$`\left(i\gamma _0{\displaystyle \frac{}{t}}𝜸𝐤\right)\psi ^{(0)}(𝐤,t)`$ $`=`$ $`0,`$
$`\left(i\gamma _0{\displaystyle \frac{}{t}}𝜸𝐤\right)\psi ^{(1)}(𝐤,t)`$ $`=`$ $`\sigma ^{(1)}(𝐤,0)\psi ^{(0)}(𝐤,t){\displaystyle _0^t}𝑑t^{}\sigma ^{(1)}(𝐤,tt^{})\dot{\psi }^{(0)}(𝐤,t^{}),`$
$`\mathrm{}`$ $`\mathrm{}`$
These equations can be solved iteratively by starting from the zeroth-order (free field) solution
$$\psi ^{(0)}(𝐤,t)=\underset{s}{}\left[A_s(𝐤)u_s(𝐤)e^{ikt}+B_s(𝐤)v_s(𝐤)e^{ikt}\right],$$
and the retarded Green’s function of the unperturbed problem
$$𝒢_{\mathrm{ret}}^{(0)}(𝐤,tt^{})=\left\{i\gamma _0\mathrm{cos}[k(tt^{})]+𝜸\widehat{𝐤}\mathrm{sin}[k(tt^{})]\right\}\theta (tt^{}).$$
Here $`u_s(𝐤)`$ and $`v_s(𝐤)`$ are free Dirac spinors that satisfy
$$(\gamma _0𝜸\widehat{𝐤})\left\{\begin{array}{c}u_s(𝐤)\hfill \\ v_s(𝐤)\hfill \end{array}\right\}=0.$$
The solution to the first-order equation is found to be given by
$$\psi ^{(1)}(𝐤,t)=\psi ^{(1,a)}(𝐤,t)+\psi ^{(1,b)}(𝐤,t),$$
where
$`e^2\psi ^{(1,a)}(𝐤,t)`$ $`=`$ $`i\gamma ^a(𝐤)t{\displaystyle \underset{s}{}}\left[A_s(𝐤)u_s(𝐤)e^{ikt}B_s(𝐤)v_s(𝐤)e^{ikt}\right],`$ (84)
$`e^2\psi ^{(1,b)}(𝐤,t)`$ $`=`$ $`{\displaystyle \frac{i}{\pi }}{\displaystyle \underset{s}{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{d\omega }{\omega k}}\gamma ^b(\omega ,𝐤)\{A_s(𝐤)u_s(𝐤)e^{ikt}[t{\displaystyle \frac{1e^{i(\omega k)t}}{i(\omega k)}}]`$ (86)
$`B_s(𝐤)v_s(𝐤)e^{ikt}[t+{\displaystyle \frac{1e^{i(\omega k)t}}{i(\omega k)}}]\},`$
with
$`\gamma ^a(𝐤)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{d\omega }{\omega }}\rho _2(\omega ,𝐤)={\displaystyle \frac{1}{4}}\text{Tr}\left[\text{Re}\mathrm{\Sigma }(0,𝐤)(\gamma _0𝜸\widehat{𝐤})\right],`$ (87)
$`\gamma ^b(\omega ,𝐤)`$ $`=`$ $`{\displaystyle \frac{\pi k}{\omega }}[\rho _1(\omega ,𝐤)\rho _2(\omega ,𝐤)]={\displaystyle \frac{k}{4\omega }}\text{Tr}\left[\text{ Im}\mathrm{\Sigma }(\omega ,𝐤)(\gamma _0𝜸\widehat{𝐤})\right].`$ (89)
In Eq. (84) secular terms are explicitly linear in time and are purely imaginary, whereas in Eq. (86) secular terms may arise at long times from the resonant denominators.
From the form of $`\gamma ^b(\omega ,𝐤)`$ in Eq. (89), we recognize that $`\gamma ^b(\omega ,𝐤)`$ evaluated at the fermion mass shell $`\omega =k`$ is the fermion damping rate computed in perturbation theory . It has been shown in the literature that due to the emission and absorption of soft quasi-static transverse photons which are only dynamically screened by Landau damping, the fermion damping rate exhibits infrared divergences near the mass shell in perturbation theory.
This becomes evident from the following analysis. For soft photons with $`q_0,qT`$, we can replace
$$1+n_B(q_0)n_F(p)T/q_0,pkq\mathrm{cos}\theta ,$$
where $`\mathrm{cos}\theta =\widehat{𝐤}\widehat{𝐪}`$, thus write
$`\gamma ^b(\omega ,𝐤)`$ $`=`$ $`{\displaystyle \frac{\pi e^2Tk}{\omega }}{\displaystyle \frac{d^3q}{(2\pi )^3}_q^{+q}\frac{dq_0}{q_0}\left[(1\mathrm{cos}^2\theta )\beta _T(q_0,q)+\beta _L(q_0,q)\right]}`$ (91)
$`\times \delta (\omega k+q\mathrm{cos}\theta q_0).`$
Here we have neglected the subleading pole contributions which corresponding to emission and absorption of on-shell photons . Recall that for very soft $`qeT`$ the function $`\beta _T(q_0,q)/q_0`$ is strongly peaked at $`q_0=0`$ \[see Eq. (32) and Fig. 1\], and as $`q0`$ it can be approximated by
$`{\displaystyle \frac{1}{q_0}}\beta _T(q_0q,q)`$ $``$ $`{\displaystyle \frac{\delta (q_0)}{q^2}}\text{as}q0.`$ (92)
The infrared divergences near the fermion mass shell become manifest after substituting Eq. (92) into Eq. (91). The physical origin of the behavior of the function $`\beta _T(q_0,q)/q_0`$ as $`q_0q0`$ is the absence of a magnetic mass.
In order to isolate the singular behavior of $`\gamma ^b(\omega ,𝐤)`$, we follow the steps in Ref. and write
$$\frac{1}{q_0}\beta _T(q_0,q)=\delta (q_0)\left(\frac{1}{q^2}\frac{1}{q^2+\omega _P^2}\right)+\frac{1}{q_0}\nu _T(q_0,q),$$
(93)
where $`\omega _P=eT/3`$ is the plasma frequency and $`\nu _T(q_0,q)`$ denotes the regular part of the transverse photon spectral density. Substituting Eq. (93) into Eq. (91), we can then separate $`\gamma ^b(\omega ,𝐤)`$ into an infrared singular part which is logarithmically divergent near the fermion mass shell, given by
$`\gamma _{\mathrm{sing}}^b(\omega ,𝐤)`$ $`=`$ $`{\displaystyle \frac{\pi e^2Tk}{\omega }}{\displaystyle \frac{d^3q}{(2\pi )^3}\left(\frac{1}{q^2}\frac{1}{q^2+\omega _P^2}\right)\delta (\omega k+q\mathrm{cos}\theta )},`$
and a contribution that remains finite and can be expanded near the fermion mass shell, given by
$`\gamma _{\mathrm{reg}}^b(\omega ,𝐤)`$ $`=`$ $`{\displaystyle \frac{\pi e^2Tk}{\omega }}{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}{\displaystyle _q^{+q}}{\displaystyle \frac{dq_0}{q_0}}[\beta _L(q_0,q)\mathrm{cos}^2\theta \beta _T(q_0,q)`$
$`+\nu _T(q_0,q)]\delta (\omega k+q\mathrm{cos}\theta q_0).`$
Using the delta function $`\delta (q\mathrm{cos}\theta q_0)`$ to perform the angular integration yields
$`\gamma ^b(\omega ,𝐤)`$ $`\stackrel{\omega k}{=}`$ $`\alpha T\left(\mathrm{ln}{\displaystyle \frac{\omega _P}{|\omega k|}}+\right)+𝒪(\omega k),`$ (94)
where
$``$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑qq{\displaystyle _q^{+q}}{\displaystyle \frac{dq_0}{q_0}}\left[\beta _L(q_0,q){\displaystyle \frac{q_0^2}{q^2}}\beta _T(q_0,q)+\nu _T(q_0,q)\right].`$ (95)
The above double integral has been computed analytically in Ref. with the result $`=\mathrm{ln}3/2`$,
$`\gamma ^b(\omega ,𝐤)`$ $`\stackrel{\omega k}{=}`$ $`\alpha T\left(\mathrm{ln}{\displaystyle \frac{|\omega k|}{\omega _P}}{\displaystyle \frac{\mathrm{ln}3}{2}}\right)+𝒪(\omega k).`$
We are now in position to find the secular terms in $`\psi ^{(1)}(𝐤,t)`$ that emerge in the intermediate asymptotic regime.
The imaginary part of the secular terms in $`e^2\psi ^{(1,a)}(𝐤,t)`$ and $`e^2\psi ^{(1,b)}(𝐤,t)`$ combine into a linear secular term given by
$$S_{I,k}(t)=i\alpha \delta _kt\frac{it}{4}\text{Tr}[\text{Re}\mathrm{\Sigma }(\omega ,𝐤)(\gamma _0𝜸\widehat{𝐤})]|_{\omega =k}.$$
for the positive (negative) energy spinors respectively, and no further secular terms arise from the higher order expansion around the fermion mass shell in Eq. (94). This purely imaginary linear secular term is thus identified with a perturbative shift of the oscillation frequency of the mean field and is determined by a dispersive integral of the spectral densities $`\rho _{1,2}(\omega ,𝐤)`$ which is rather difficult to obtain in closed form but a detailed analysis reveals that $`\delta _k`$ is finite.
The real secular terms are more involved. The contribution to $`\gamma ^b(\omega ,𝐤)`$ that is finite as $`\omega k`$ leads to a linear secular term, whereas for the logarithmically divergent contribution as $`\omega k`$ the following asymptotic result becomes useful:
$`{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{dy}{y^2}}(1\mathrm{cos}yt)\mathrm{ln}|y|\stackrel{t\mathrm{}}{=}\pi t(1\gamma _E\mathrm{ln}t)+𝒪\left(t^1\right),`$ (96)
where we have neglected terms that fall off at long times. Thus we find the real part of the secular terms to be given by
$$S_{R,k}(t)=\alpha Tt\left(\mathrm{ln}\omega _Pt+\gamma _E1+\frac{\mathrm{ln}3}{2}\right).$$
Gathering the above results, at large times $`t1/\omega _P`$ the perturbative solution reads
$`\psi (𝐤,t)`$ $`=`$ $`{\displaystyle \underset{s}{}}\left\{[1+S_k(t)]A_s(𝐤)u_s(𝐤)e^{ikt}+[1+S_k^{}(t)]B_s(𝐤)v_s(𝐤)e^{ikt}\right\}`$ (98)
$`+\text{nonsecular terms},`$
with
$$S_k(t)=\alpha t\left[\left(\mathrm{ln}\omega _Pt+\gamma _E1+\frac{\mathrm{ln}3}{2}\right)T+i\delta _k\right].$$
Obviously, this perturbative solution breaks down at a time scale $`\tau _{\mathrm{br}}[\alpha T\mathrm{ln}(1/\alpha )]^1`$. To obtain a uniformly valid solution for large times we now implement a resummation of the secular terms in the perturbative series via the dynamical renormalization group .
This method is implemented by defining the (complex) amplitude renormalization as follows
$$\begin{array}{c}A_s(𝐤)=𝒵_k(\tau )𝒜_s(𝐤,\tau ),\hfill \\ B_s(𝐤)=𝒵_k^{}(\tau )_s(𝐤,\tau ),\hfill \end{array}$$
where $`𝒵_k(\tau )=1+\alpha z_k^{(1)}(\tau )+𝒪(e^4)`$. The renormalization coefficient $`\alpha z_k^{(1)}(\tau )`$ is chosen to cancel the secular divergence at a time scale $`t=\tau `$, i.e, $`\alpha z_k^{(1)}(\tau )=S_k(\tau )`$, thus leading to
$`\psi (𝐤,t)`$ $`=`$ $`{\displaystyle \underset{s}{}}\{[1+S_k(t)S_k(\tau )]𝒜_s(𝐤,\tau )u_s(𝐤)e^{ikt}`$
$`+[1+S_k^{}(t)S_k^{}(\tau )]_s(𝐤,\tau )v_s(𝐤)e^{ikt}\},`$
which remains bounded at large times provided that $`\tau `$ is chosen arbitrarily close to $`t`$. The mean field does not depend on the arbitrary renormalization scale $`\tau `$ and this independence leads to the dynamical renormalization group equation, which to order $`𝒪(\alpha )`$ is given by
$`\left[{\displaystyle \frac{d}{d\tau }}{\displaystyle \frac{dS_k(\tau )}{d\tau }}\right]𝒜_s(𝐤,\tau )=0,`$ (99)
$`\left[{\displaystyle \frac{d}{d\tau }}{\displaystyle \frac{dS_k^{}(\tau )}{d\tau }}\right]_s(𝐤,\tau )=0,`$ (100)
with solutions
$`\left\{\begin{array}{c}𝒜_s(𝐤,\tau )\\ _s(𝐤,\tau )\end{array}\right\}`$ $`=`$ $`\left\{\begin{array}{c}𝒜_s(𝐤,\tau _0)e^{i\alpha \delta _k\tau }\\ _s(𝐤,\tau _0)e^{i\alpha \delta _k\tau }\end{array}\right\}e^{\alpha T\tau \left[\mathrm{ln}\omega _P\tau +0.12652\mathrm{}\right]},`$ (105)
where we have replaced $`\gamma _E1+\mathrm{ln}3/2=0.12652\mathrm{}`$, and $`\tau _01/\omega _P`$ is the time scale such that this intermediate asymptotic solution is valid. Hence choosing the renormalization point $`\tau `$ to coincide with the time $`t`$, we find the long time behavior (for $`\omega _Pt1`$) of the fermionic mean field for hard momentum is given by
$`\psi (𝐤,t)`$ $`=`$ $`{\displaystyle \underset{s}{}}\left[𝒜_s(𝐤,\tau _0)u_s(𝐤)e^{iE(k)t}+_s(𝐤,\tau _0)v_s(𝐤)e^{iE(k)t}\right]e^{\alpha Tt\left[\mathrm{ln}\omega _Pt+0.12652\mathrm{}\right]},`$ (106)
where $`E(k)`$ is the pole position shifted by one-loop corrections
$$E(k)=k+\frac{1}{4}\text{Tr}[\text{Re}\mathrm{\Sigma }(\omega ,𝐤)(\gamma _0𝜸\widehat{𝐤})]|_{\omega =k}.$$
Eq. (106) reveals a time scale for the relaxation of the fermionic mean field $`\tau _{\mathrm{rel}}[\alpha T\mathrm{ln}(1/\alpha )]^1`$, which coincides with the time scale at which the perturbative solution (98) breaks down . This highlights clearly the nonperturbative nature of relaxation phenomena.
This result coincides with that found in Refs. via the Bloch-Nordsieck approximation and in scalar quantum electrodynamics using the dynamical renormalization group.
The main purpose of this section was to introduce the dynamical renormalization group for fermionic theories. Furthermore, this is another important and relevant example of the reliability and consistency of this novel renormalization group applied to real-time nonequilibrium phenomena.
### B Quantum kinetics of the fermion distribution function
We now study the quantum kinetic equation for the distribution function of hard fermions. There has recently been an intense activity to obtain a Boltzmann equation for quasiparticles in gauge theories motivated in part by the necessity to obtain a consistent description for baryogenesis in non-abelian theories. Boltzmann equations with a diagrammatic interpretation were obtained in in which a collision-type kernel describes the scattering of hard quasiparticles. In these approaches this collision kernel reveals the infrared divergences associated with the emission and absorption of magnetic photons (or gluons) and must be cutoff by introducing a relaxation time scale $`\tau _{\mathrm{rel}}1/\alpha T`$ to leading logarithmic accuracy .
In a derivation of quantum kinetic equations for charged quasiparticles in scalar quantum electrodynamics using the dynamical renormalization group, it was understood that the origin of these infrared divergences is the implementation of Fermi’s golden rule that assumes completed collisions and takes the infinite time limit in the collision kernels. The dynamical renormalization group leads to quantum kinetic equations in real time in terms of time-dependent scattering kernels without any infrared ambiguity.
In this section we implement this program in spinor QED to derive the quantum kinetic equation for hard fermions. There are several important features of our study that must be emphasized: (i) as presented in detail in Sec. II, gauge invariance is automatically taken into account by working directly with gauge invariant operators, thus the operator that describes the number of fermionic quasiparticles is gauge invariant, (ii) a kinetic description relies on a separation between the microscopic and the relaxation time scales, this is warranted in a strict perturbative regime and applies to hard fermionic quasiparticles, (iii) the dynamical renormalization group leads to a quantum kinetic equation in real time without infrared divergences since time acts as an infrared cutoff.
The program begins by expanding the Heisenberg fermion field in terms of creation and annihilation operators as
$`\psi (𝐱,t)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^{3/2}}\psi (𝐤,t)e^{i𝐤𝐱}},`$
$`\psi ^{}(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^{3/2}}\psi ^{}(𝐤,t)e^{i𝐤𝐱}},`$
with
$`\psi (𝐤,t)`$ $`=`$ $`\sqrt{{\displaystyle \frac{m}{\omega _𝐤}}}{\displaystyle \underset{s}{}}\left[b_s(𝐤,t)u_s(𝐤)+d_s^{}(𝐤,t)v_s(𝐤)\right],`$
$`\psi ^{}(𝐤,t)`$ $`=`$ $`\sqrt{{\displaystyle \frac{m}{\omega _𝐤}}}{\displaystyle \underset{s}{}}\left[b_s^{}(𝐤,t)u_s^{}(𝐤)+d_s(𝐤,t)v_s^{}(𝐤)\right],`$
where $`\omega _𝐤=\sqrt{𝐤^2+m^2}`$ and $`b_s(𝐤,t)`$ \[$`b_s^{}(𝐤,t)`$\] is the annihilation (creation) operator that destroys (creates) a free fermion of momentum $`𝐤`$ and spin $`s`$ at time $`t`$. We have retained the fermion mass to avoid the subtleties associated with the normalization of massless spinors, the massless limit will be taken later.
The spin-averaged number operator for fermions with momentum $`𝐤`$ is then given by
$`N_f(𝐤,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s}{}}b_s^{}(𝐤,t)b_s(𝐤,t)={\displaystyle \frac{1}{4\omega _𝐤}}\psi ^{}(𝐤,t)(\overline{)}K+m)\gamma _0\psi (𝐤,t),`$ (107)
where $`K=(\omega _𝐤,𝐤)`$. Similarly, the spin-averaged number operator for anti-fermions with momentum $`𝐤`$ is then given by
$`\overline{N}_f(𝐤,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s}{}}d_s^{}(𝐤,t)d_s(𝐤,t)={\displaystyle \frac{1}{4\omega _𝐤}}\psi ^{}(𝐤,t)(\overline{)}Km)\gamma _0\psi (𝐤,t).`$ (108)
Taking time derivative of $`N_f(𝐤,t)`$ and using the Heisenberg equations of motion, we find
$`\dot{n}_𝐤^f(t)`$ $``$ $`\dot{N}_f(𝐤,t)`$
$`=`$ $`\underset{t^{}t}{lim}{\displaystyle \frac{ie}{4\omega _𝐤}}{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^{3/2}}}\{\overline{\psi }^{}(𝐩,t^{})[\gamma _0A_0^{}(𝐪,t^{})𝜸𝐀_T^{}(𝐪,t^{})]`$
$`\times (\overline{)}K+m)\gamma _0\psi ^+(𝐤,t)\}+\text{c.c.}.`$
As before here the “$`+`$” (“$``$”) superscripts for the fields refer to fields defined in the forward (backward) time branch in the CTP formulation. It is straightforward to check that the total fermion number (fermions minus antifermions) is conserved.
In the hard fermion limit $`kTm`$ we neglect the fermion mass and obtain to $`𝒪(e^2)`$
$`\dot{n}_𝐤^f(t)`$ $`=`$ $`e^2{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}{\displaystyle _{\mathrm{}}^+\mathrm{}}dq_0{\displaystyle _{t_0}^t}dt^{\prime \prime }\{𝒩_1(t_0)\mathrm{cos}[(kpq_0)(t^{\prime \prime }t_0)]`$ (112)
$`\times \left[\stackrel{~}{\rho }_L(q_0,q)𝒦_1^+(𝐤,𝐪)+2\stackrel{~}{\rho }_T(q_0,q)𝒦_2^+(𝐤,𝐪)\right]+𝒩_2(t_0)`$
$`\times \mathrm{cos}[(k+pq_0)(t^{\prime \prime }t_0)][\stackrel{~}{\rho }_L(q_0,q)𝒦_1^{}(𝐤,𝐪)+2\stackrel{~}{\rho }_T(q_0,q)𝒦_2^{}(𝐤,𝐪)]\},`$
where $`𝒩`$ and $`𝒦`$ denote respectively the following statistical and kinematic factors
$`𝒩_1(t)`$ $`=`$ $`[1n_𝐤^f(t)]n_𝐩^f(t)n_B(q_0)n_𝐤^f(t)[1n_𝐩^f(t)][1+n_B(q_0)],`$
$`𝒩_2(t)`$ $`=`$ $`[1n_𝐤^f(t)][1n_𝐩^f(t)]n_B(q_0)n_𝐤^f(t)n_𝐩^f(t)[1+n_B(q_0)],`$
$`𝒦_1^\pm (𝐤,𝐪)`$ $`=`$ $`1\pm \widehat{𝐤}\widehat{𝐩},𝒦_2^\pm (𝐤,𝐪)=1\widehat{𝐤}\widehat{𝐩}\pm {\displaystyle \frac{1(\widehat{𝐤}\widehat{𝐪})^2}{1\frac{q}{k}(\widehat{𝐤}\widehat{𝐪})}}.`$
Here $`\stackrel{~}{\rho }_T(q_0,q)`$ and $`\stackrel{~}{\rho }_L(q_0,q)`$ are the HTL-resummed spectral densities for the transverse and longitudinal photons given by Eqs.(30) and (80), respectively.
In the linearized relaxation time approximation we assume that $`n_𝐤^f(t_0)=n_F(k)+\delta n_𝐤^f(t_0)`$. Then upon integrating over $`t^{\prime \prime }`$, we obtain
$$\delta \dot{n}_𝐤^f(t)=\delta n_𝐤^f(t_0)_{\mathrm{}}^+\mathrm{}𝑑\omega _f(\omega ,𝐤)\frac{\mathrm{sin}[(\omega k)(tt_0)]}{\pi (\omega k)},$$
(113)
where
$`_f(\omega ,𝐤)`$ $`=`$ $`\pi e^2{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}{\displaystyle _{\mathrm{}}^+\mathrm{}}dq_0[1+n_B(q_0)n_F(p)]\{[\stackrel{~}{\rho }_L(q_0,q)𝒦_1^+(𝐤,𝐪)`$
$`+\mathrm{\hspace{0.33em}2}\stackrel{~}{\rho }_T(q_0,q)𝒦_2^+(𝐤,𝐪)]\delta (\omega pq_0)+[\stackrel{~}{\rho }_L(q_0,q)𝒦_1^{}(𝐤,𝐪)`$
$`+\mathrm{\hspace{0.33em}2}\stackrel{~}{\rho }_T(q_0,q)𝒦_2^{}(𝐤,𝐪)]\delta (\omega +p+q_0)\}.`$
Eq. (113) can be integrated directly to yield
$`\delta n_𝐤^f(t)`$ $`=`$ $`\delta n_𝐤^f(t_0)\left\{1{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\omega _f(\omega ,𝐤){\displaystyle \frac{1\mathrm{cos}[(\omega k)(tt_0)]}{\pi (\omega k)^2}}\right\}.`$ (114)
The time-dependent contribution above is now familiar from the previous discussions, potential secular terms will emerge at long times from the regions in which the resonant denominator vanishes. This is the region near the fermion mass shell $`\omega k`$, where $`_f(\omega ,𝐤)`$ is dominated by the regions of small $`q`$ and $`q_0`$, which physically corresponds to emission and absorption of soft photons. As before for soft photons with $`q,q_0T`$, we can replace
$`1+n_B(q_0)n_F(p)T/q_0,pkq\mathrm{cos}\theta ,𝒦_1^+2,𝒦_2^+1\mathrm{cos}^2\theta ,`$
thus write $`_f(\omega ,𝐤)`$ at $`\omega k`$ as
$`_f(\omega ,𝐤)`$ $`=`$ $`2\pi e^2T{\displaystyle \frac{d^3q}{(2\pi )^3}_q^q\frac{dq_0}{q_0}\left[(1\mathrm{cos}^2\theta )\beta _T(q_0,q)+\beta _L(q_0,q)\right]}`$ (116)
$`\times \delta (\omega k+q\mathrm{cos}\theta q_0).`$
Note that the double integral in Eq. (116) is exactly the same as that in Eq. (95). Thus $`_f(\omega ,𝐤)`$ features an infrared divergence near the fermion mass shell as shown in the previous subsection. Following the analysis carried out in the preceding subsection we obtain
$$_f(\omega ,𝐤)\stackrel{\omega k}{=}2\alpha T\left[\mathrm{ln}\frac{|\omega k|}{\omega _P}\frac{\mathrm{ln}3}{2}\right]+𝒪[(\omega k)^2].$$
(117)
Substituting Eq. (117) into Eq. (114), we find the number of fermions at intermediate asymptotic times $`tt_01/\omega _P`$ to be given by
$`\delta n_𝐤^f(t)`$ $`=`$ $`\delta n_𝐤^f(t_0)\{12\alpha T(tt_0)[\mathrm{ln}\omega _P(tt_0)+0.12652\mathrm{}]\}`$ (119)
$`+\text{nonsecular terms}.`$
As in the case of the fermion mean field relaxation \[cf. Eq. (98)\], the perturbative solution contains a secular term of the form $`t\mathrm{ln}t`$. Obviously, the secular term will invalidate the perturbative solution at time scales $`\tau _{\mathrm{rel}}[2\alpha T\mathrm{ln}(1/\alpha )]^1`$. In the intermediate asymptotic regime $`1/ktt_0\tau _{\mathrm{rel}}`$, the perturbative expansion can be improved by absorbing the contribution of the secular term at a time scale $`\tau `$ into a re-definition of the distribution function. Hence we apply the dynamical renormalization group method through a renormalization of the distribution function much in the same manner as the renormalization of the amplitude in the mean field discussed above,
$$\delta n_𝐤^f(t_0)=𝒵_k(\tau ,t_0)\delta n_𝐤^f(\tau ),$$
with
$$𝒵_k(\tau ,t_0)=1+\alpha z_k^{(1)}(\tau ,t_0)+𝒪(\alpha ^2).$$
The independence of the solution on the time scale $`\tau `$ leads to the dynamical renormalization group equation which to lowest order in $`\alpha `$ is given by
$$\left\{\frac{d}{d\tau }+2\alpha T[1+\mathrm{ln}\omega _P(\tau t_0)+0.12652\mathrm{}]\right\}\delta n_𝐤^f(\tau )=0.$$
(120)
which is obviously of the form of a quantum kinetic equation in the linearized relaxation time approximation, but with a time-dependent relaxation rate.
Solving Eq. (120) and choosing $`\tau `$ to coincide with $`t`$, we obtain the evolution of the fermion distribution function at asymptotic times $`tt_01/k`$ to be given by
$`\delta n_𝐤^f(t)`$ $`=`$ $`\delta n_𝐤^f(t_0)\mathrm{exp}\{2\alpha T(tt_0)[\mathrm{ln}\omega _P(tt_0)+\mathrm{\hspace{0.33em}0.12652}\mathrm{}]\}.`$ (121)
Comparing Eq. (121) with Eq. (106) we find that the anomalous exponent that describes the relaxation of the fermion distribution function in the linear approximation is twice that for the linear relaxation of the mean field. A similar relation is obtained between the damping rate for the single quasiparticle relaxation and the relaxation rate of the distribution function in the case of time-independent rates and true exponential relaxation . The dynamical renormalization group reveals this to be a generic feature even with time-dependent rates.
## V Conclusions and further questions
The goals of this article are the study of nonequilibrium effects in hot QED plasmas directly in real time. The focus is a systematic study of relaxation of mean fields as well as the distribution functions for photons and fermions. In particular the application of the dynamical renormalization group method to study anomalous relaxation as a consequence of the exchange of soft photons. To begin with, we have cast our study solely in terms of gauge invariant quantities, this can be done in the abelian theory in a straightforward manner, avoiding potential ambiguities associated with gauge invariance. The relaxation of photon mean fields revealed important features: for soft momentum $`keT`$ mean fields that are prepared by adiabatically switching-on an external source there is exponential relaxation towards the oscillatory behavior dominated by the transverse photon pole. The source that induces the mean field in this case has a Fourier transform that is singular at zero frequency and excites the resonance in the Landau damping region near zero frequency for the soft photon mean field. Sources that have a regular Fourier transform would not lead to the exponential relaxation. For semihard momentum $`eTkT`$ in principle both the HTL and the perturbative approximations are valid, however the spectral density for photons becomes sharply peaked at the edge of the Landau damping continuum consistent with the fact that the photon pole becomes perturbatively close to the Landau damping cut. This enhancement of the spectral density near the bare photon mass shell results in the breakdown of perturbation theory at large times $`kt1`$. The dynamical renormalization group provides a consistent resummation of the photon self-energy in real time and leads to anomalous relaxation of the mean field, given by Eq. (47). Clearly, higher order terms beyond the HTL approximation will include collisional contributions leading perhaps to exponential relaxation. We then expect a crossover between the anomalous power law obtained in lowest order and the exponential relaxation from collisional processes, the crossover time scale will depend on the details of the different contributions and requires a study beyond that presented here.
The dynamical renormalization group provides a consistent and systematic framework to obtain quantum kinetic equations directly in real time from the microscopic field theory . This method allows to extract information that is not available in the usual kinetic description in terms of time-independent collision kernels obtained under the assumption of completed collisions which only include on-shell processes. The dynamical renormalization group approach to quantum kinetics consistently includes off-shell processes and accounts for time-dependent collisional kernels. This is important and potentially phenomenologically relevant in the case of the quark-gluon plasma which has a finite lifetime.
We have implemented this approach to obtain the quantum kinetic equation for the distribution function of semihard photons in the case of thermalized fermions in lowest order in the HTL approximation. This equation features time-dependent rates and we established that detailed balance, a consequence of fermions being in thermal equilibrium, emerges on microscopic time scales. The linearization of the kinetic equation describes relaxation towards equilibrium with an anomalous exponent, twice as large as that of the photon mean field in the semihard case.
The kinetic equation for semihard photons allows us to study photon production from a thermalized quark-gluon plasma by simply replacing the fermions by two flavors of light quarks. To leading order in the HTL approximation photon production arises from off-shell bremsstrahlung ($`qq\gamma `$ and $`\overline{q}\overline{q}\gamma `$), and the total number of photons produced at time $`t`$ is given by Eq. (71). We study the hard region $`kT`$ numerically and find that for $`T200\text{MeV}`$ the number of photons produced by these off-shell processes form a thermalized QGP of lifetime $`10\text{fm}/ct50\text{fm}/c`$ is comparable to that by on-shell processes estimated in Refs. \[See Fig. 5\].
The relaxation of the fermion mean field for hard momentum is studied with the dynamical renormalization group, we find an anomalous exponential relaxation which confirms the results of Refs. where the Bloch-Nordsieck approximation was used.
We then obtain the quantum kinetic equation for the distribution function of hard fermions assuming that photons are in thermal equilibrium. The collisional kernel is time-dependent and infrared finite. The linearized kinetic equation describes approach to equilibrium with an anomalous exponential relaxation, which is twice that of the fermion mean field.
An important payoff of this approach to quantum kinetics is that it bypasses the assumption of completed collisions which lead to collisional kernels obtained by Fermi’s golden rule and only describe on-shell processes as in the usual kinetic approach, which in the case under consideration leads to infrared divergent collisional kernels .
Further questions. Perhaps the most phenomenologically pressing aspect that requires further and deeper study is the photon production by off-shell processes in a thermalized quark-gluon plasma. This is important in view of the fact that RHIC will begin physics runs very soon hence an assessment of potential experimental electromagnetic signatures is very relevant. The next step in the study of photon production must be to include hard thermal loop correction to the quark propagators, and to study the spectral distribution of the photons produced as a function of momentum, temperature and lifetime of the thermalized QGP. Furthermore, our calculations just as those in Refs. do not account for the expansion of the plasma, thus an important next step is to couple the quantum kinetic equation to the hydrodynamic equations of evolution of the QGP, this is required to obtain the photon distribution function integrated over the space-time volume of the QGP. Work on these aspects is currently underway.
Another important extension of this work is to obtain the relaxation of non-abelian gauge mean fields beyond the HTL approximation both in the soft and semihard regime. Similarly for the fermionic mean fields as well as the kinetics of the fermion distribution function. Of particular importance in this regard would be to include the nonequilibrium evolution of the distribution function of the gauge fields in the collision kernels for the fermionic distribution function. Obviously this will require a consistent analysis of the relevant time scales. These aspects bear on the physics of baryogenesis and we expect to report on some of our studies on this issues soon.
###### Acknowledgements.
D.B., H.J.d.V. and S.-Y.W. would like to thank U. Heinz, E. Mottola, D. Schiff, M. Simionato and L. Yaffe for useful discussions. They also thank the Institute for Nuclear Theory at the University of Washington for its hospitality during the early stages of this work. D.B. and S.-Y.W. thank the US NSF for partial support through grants PHY-9605186, CNRS-NSF-INT-9815064 and INT-9905954. S.-Y.W. thanks Andrew Mellon Foundation for partial support. D.-S.L. was supported by the ROC NSC through grant NSC89-2112-M-259-008-Y. LPTHE is UMR 7589 associated to CNRS.
## Real-Time Propagators
In this appendix we summarize the various CTP real-time propagators used in this article.
* The free fermion propagators (with zero chemical potential) are defined by
$`\mathrm{\Psi }^a(𝐱,t)\overline{\mathrm{\Psi }}^b(𝐱^{},t^{})=i{\displaystyle \frac{d^3k}{(2\pi )^3}S_𝐤^{ab}(t,t^{})e^{i𝐤(𝐱𝐱^{})}},`$ (122)
$`S_𝐤^{++}(t,t^{})=S_𝐤^>(t,t^{})\theta (tt^{})+S_𝐤^<(t,t^{})\theta (t^{}t),`$ (123)
$`S_𝐤^{}(t,t^{})=S_𝐤^>(t,t^{})\theta (t^{}t)+S_𝐤^<(t,t^{})\theta (tt^{}),`$ (124)
$`S_𝐤^\pm (t,t^{})=S_𝐤^{<\text{ }>}(t,t^{}),`$ (125)
where $`a,b=\pm `$ and the fermion Wightman functions read
$`S_𝐤^>(t,t^{})`$ $`=`$ $`{\displaystyle \frac{i}{2\omega _𝐤}}\{(\overline{)}K+m)[1n_F(k)]e^{i\omega _𝐤(tt^{})}`$ (127)
$`+\gamma _0(\overline{)}Km)\gamma _0n_F(k)e^{i\omega _𝐤(tt^{})}\},`$
$`S_𝐤^<(t,t^{})`$ $`=`$ $`{\displaystyle \frac{i}{2\omega _𝐤}}\{(\overline{)}K+m)n_F(k)e^{i\omega _𝐤(tt^{})}`$ (129)
$`+\gamma _0(\overline{)}Km)\gamma _0[1n_F(k)]e^{i\omega _𝐤(tt^{})}\},`$
with $`K=(\omega _𝐤,𝐤)`$, $`\omega _𝐤=\sqrt{𝐤^2+m^2}`$ and $`n_F(k)`$ the Fermi-Dirac distribution.
* The free (transverse) photon propagators are defined by
$`A_T^{i,a}(𝐱,t)A_T^{j,b}(𝐱^{},t^{})=i{\displaystyle \frac{d^3q}{(2\pi )^3}𝒢_{T,q}^{ab}(t,t^{})𝒫_T^{ij}(𝐪)e^{i𝐪(𝐱𝐱^{})}},`$ (130)
$`𝒢_{T,q}^{++}(t,t^{})=𝒢_{T,q}^>(t,t^{})\theta (tt^{})+𝒢_{T,q}^<(t,t^{})\theta (t^{}t),`$ (131)
$`𝒢_{T,q}^{}(t,t^{})=𝒢_{T,q}^>(t,t^{})\theta (t^{}t)+𝒢_{T,q}^<(t,t^{})\theta (tt^{}),`$ (132)
$`𝒢_{T,q}^\pm (t,t^{})=𝒢_{T,q}^{<\text{ }>}(t,t^{}),`$ (133)
where $`𝒫_T^{ij}(𝐪)=\delta ^{ij}\widehat{q}^i\widehat{q}^j`$ is the transverse projector and the photon Wightman functions read
$`𝒢_{T,q}^>(t,t^{})`$ $`=`$ $`{\displaystyle \frac{i}{2q}}\left[[1+n_B(q)]e^{iq(tt^{})}+n_B(q)e^{iq(tt^{})}\right],`$ (134)
$`𝒢_{T,q}^<(t,t^{})`$ $`=`$ $`{\displaystyle \frac{i}{2q}}\left[n_B(q)e^{iq(tt^{})}+[1+n_B(q)]e^{iq(tt^{})}\right],`$ (135)
with $`n_B(q)`$ the Bose-Einstein distribution.
* The HTL-resummed transverse photon propagators are obtained from the free ones (133) by using the following Wightman functions
$`𝒢_{T,q}^>(t,t^{})`$ $`=`$ $`i{\displaystyle 𝑑q_0\stackrel{~}{\rho }_T(q_0,q)[1+n_B(q_0)]e^{iq_0(tt^{})}},`$ (136)
$`𝒢_{T,q}^<(t,t^{})`$ $`=`$ $`i{\displaystyle 𝑑q_0\stackrel{~}{\rho }_T(q_0,q)n_B(q_0)e^{iq_0(tt^{})}},`$ (137)
where $`\stackrel{~}{\rho }_T`$ is the resumed spectral density.
* The HTL-resummed longitudinal photon propagators are given by
$`A_0^a(𝐱,t)A_0^b(𝐱^{},t^{})=i{\displaystyle \frac{d^3q}{(2\pi )^3}𝒢_{L,q}^{ab}(t,t^{})e^{i𝐪(𝐱𝐱^{})}},`$ (138)
$`𝒢_{L,q}^{++}(t,t^{})={\displaystyle \frac{1}{q^2}}\delta (tt^{})+𝒢_{L,q}^>(t,t^{})\theta (tt^{})+𝒢_{L,q}^<(t,t^{})\theta (t^{}t),`$ (139)
$`𝒢_{L,q}^{}(t,t^{})={\displaystyle \frac{1}{q^2}}\delta (tt^{})+𝒢_{L,q}^>(t,t^{})\theta (t^{}t)+𝒢_{L,q}^<(t,t^{})\theta (tt^{}),`$ (140)
$`𝒢_{L,q}^\pm (t,t^{})=𝒢_{L,q}^{<\text{ }>}(t,t^{}),`$ (141)
with the Wightman functions expressed in terms of the resummed spectral density $`\stackrel{~}{\rho }_L`$ as
$`𝒢_{L,q}^>(t,t^{})`$ $`=`$ $`i{\displaystyle 𝑑q_0\stackrel{~}{\rho }_L(q_0,q)[1+n_B(q_0)]e^{iq_0(tt^{})}},`$ (142)
$`𝒢_{L,q}^<(t,t^{})`$ $`=`$ $`i{\displaystyle 𝑑q_0\stackrel{~}{\rho }_L(q_0,q)n_B(q_0)e^{iq_0(tt^{})}}.`$ (143)
The free real-time propagators used in studying quantum kinetics are obtained from (125)-(135) with the equilibrium Fermi-Dirac and Bose-Einstein distributions replacing by the initial nonequilibrium ones. Finally, we note that the HTL-resummed photon propagators are only valid for photons in thermal equilibrium because in deriving their spectral representations (137) and (143) used has be made of the KMS condition .
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# Magnetic Field Decay in Neutron Stars. Analysis of General Relativistic Effects
## I Introduction
It is well known that a magnetic field in a plasma of finite conductivity is subject to diffusion and dissipation. Diffusion leads to a spreading of inhomogeneities while dissipation is due to the Ohmic decay of the currents producing the field. More concretely, a magnetic field $`\stackrel{}{B}(t,\stackrel{}{x})`$ in a plasma of uniform conductivity $`\sigma `$ evolves, in flat space-time, according to the following diffusion equation :
$$\frac{\stackrel{}{B}(t,\stackrel{}{x})}{t}=\frac{c^2}{4\pi \sigma }^2\stackrel{}{B}(t,\stackrel{}{x}).$$
(1)
Accordingly, if $`L`$ is a typical length scale of the field structure, then it will decay-diffuse in a characteristic time scale $`\tau _{\mathrm{Ohm}}`$ given by: $`\tau _{\mathrm{Ohm}}=\frac{4\pi \sigma L^2}{c^2}`$. Depending upon the prevailing conditions, the Ohmic decay time $`\tau _{\mathrm{Ohm}}`$ can range from seconds, in the case of a copper sphere of radius of a few centimeters , up to $`\tau _{\mathrm{Ohm}}=10^{10}`$ years or even much longer for astrophysical settings, as in the case of the sun or a neutron star .
The interactions of large scale cosmic magnetic fields with plasmas is a problem of great importance in astrophysics and cosmology. A particularly thorny issue nowadays concerns the origin and maintenance of cosmic magnetic fields. Although large scale fields have been observed , a satisfactory explanation of their origin is still lacking. Peebles considers the issue of the origin of the primordial magnetic field as one of the most important unsolved problems in cosmology. At the same time the gigantic field of the pulsars begs for an explanation . The general consensus of the astrophysical community is that such large scale fields have been generated via an episode of dynamo action , and then gradually suffer Ohmic decay due to the finite conductivity of the medium. It appears therefore that an understanding of the factors influencing the decay of large scale fields, combined with relevant observations, may offers important clues towards a better understanding of the initial scale involved as well as clues regarding its origin.
In neutron stars the decay of the magnetic field is an issue of out most importance by itself and accordingly there has been an intense effort by astrophysicists to understand the factors governing this decay. As far as we are aware all theoretical modeling of magnetic field decay in neutron stars utilized the familiar flat space time form of Maxwell’s equations (an exception to this rule constitutes the recent work of ref. ). Although the employment of such framework is a fruitful one and provides us with valuable informations, it altogether neglects the background curvature of the spacetime which for the case of neutron stars is not any longer weak. It would be worth to stress in that regard that curvature can modify considerably flat space-time solutions of Maxwell’s equations. For instance, the reader may compare the solution describing a dipole magnetic field on a $`Schwarzschild`$ background to that of a flat space time. The presence of the logarithmic term in the former (see eqs. 49 further below) is a sole consequence of the non vanishing curvature. This example suggests that the role of the spacetime curvature on the decay process of magnetic fields ought to be examined more thoroughly than hitherto. In that respect, we are aware only of the recent work of Sengupta, , where an investigation of general relativistic effects in the magnetic field decay of neutron stars have been attempted. However this work is restricted to the study of magnetic fields confined only to the outermost layers of a neutron star and furthermore it is assumed that those outermost layers (and thus also the magnetic field $`\stackrel{}{B}`$), are embedded on a $`Schwarzschild`$ background geometry. Thus strictly the framework of ref., deals exclusively with magnetic decay on a $`Schwarzschild`$ background. In addition to those approximations and according to the sentence following Eq. 15 of Sengupta’s second work, the author fails to include general relativistic effects on the outer boundary condition for matching the inner field with the outer vacuum dipolar field across the surface of the star. In contrast, in the present work, a broad framework dealing with general relativistic effects on the magnetic field decay on an arbitrary static geometry and with proper allowance of the correct general relativistic inner and outer boundary conditions is presented. Moreover, and in contrast to the approach of ref. , we formulate the entire problem avoiding the introduction of a vector potential and the associated ambiguities. Our analysis shows that general relativistic effects can influence the field decay, but the precise manner that this influence manifests itself depends upon the class of observers called in to describe the field decay. For the magnetic field of a non rotating neutron star it is natural to describe the field decay relative to the class observers that find themselves at rest relative to the star, ie the class of Killing observers. Relative to such observers, we find that relativistic effects are influencing the field decay via two major modes: the gravitational red shift as well as the intrinsic curved geometry of the spatial sections constituting the rest space of the Killing observers. Subsequent numerical analysis shows that the red shift factor is the dominant one in slowing down the field decay. Overall we find that the inclusion of relativistic effects make the decay time of the field larger than, but of the same order of magnitude, as in flat space-time. Nevertheless the preliminary study of the present paper utilizing a simple non rotating neutron star models suggests that general relativistic effects should be given further considerations. We explicitly illustrate the impact of relativistic effects upon the magnetic field decay, by examining the evolution of a magnetic field permeating a constant density neutron star, first in their presence and secondly without them. Although for both treatments we have obtained exponential decays, the decay time in the presence of relativistic effects, on the average, is enlarged by a factor that depends crucially upon the value of the compactness ratio $`ϵ=\frac{2GM}{c^2R}`$. Specifically for values of $`ϵ`$ in the domain $`(0.3,0.5)`$, characterizing realistic neutron star models, we find that the decay time is $`1.2`$ to $`1.3`$ larger than the corresponding flat decay time, while for higher values of $`ϵ`$, it can be larger. We may add parenthetically that the term ”average” increase in the decay time, is explained in details in section (IV) of the paper.
The present paper is organized as follows: In the following section, starting from Maxwell’s equations on a static spacetime we first derive the relevant induction equation taking into account the curved nature of the background spacetime geometry. It should be stressed however that the employment of a static geometry does not leave room for incorporating gravitomagnetic (Lense-Thirring) effects in the induction equation, as the latter would manifest themselves relative to non static backgrounds, but we do hope to present such analysis in a future work. In section III, we specialize the induction equation to a simple neutron star model and a detailed analysis of the content of the induction equation is presented. In the same section the sensitive issue of the boundary conditions accompanying the induction equation is also addressed. In the section IV, we discuss numerical solutions of the curved spacetime induction equation and an assessment of the relativistic factors influencing the field decay is discussed. Furthermore in the same section, a comparison of the field decay in curved and flat spacetime is also presented. In the concluding section, a brief discussion of the physical implications of our results to neutron stars physics is presented and possible extension of the present work is outlined. Finally we have included an Appendix where a few intermediate calculations leading to the main equations of section II are presented.
## II Induction equation on a static background geometry
Maxwell’s equations, in covariant form, are as follows :
$$^\alpha F_{\alpha \beta }=\frac{4\pi }{c}J_\beta $$
(3)
$$_{[\alpha }F_{\beta \gamma ]}=0$$
(4)
where $`F_{\alpha \beta }=F_{\beta \alpha }`$, $`J_\alpha `$ and $``$ are the coordinate components of the Maxwell tensor, the conserved four current and the derivative operator respectively. Given a solution $`F_{\alpha \beta }`$ of the above eqs., an observer with four velocity $`U^\alpha `$, $`U^\alpha U_\alpha =1`$, measures electric and magnetic fields $`(E,B)`$ with corresponding coordinate components given respectively by:
$$E_\alpha =F_{\alpha \beta }U^\beta ,B_\alpha =\frac{1}{2}ϵ_{\alpha \beta }^{}{}_{}{}^{\gamma \delta }F_{\gamma \delta }U^\beta $$
(5)
where $`ϵ_{\alpha \beta \gamma \delta }`$ stands for the four-dimensional Levi-Civita tensor density . We shall be concerned in the present paper with particular solutions of II where the current $`J`$ is described by the following relativistic extension of Ohm’s law, as it was first formulated by Weyl :
$$J^\alpha =\sigma g^{\alpha \beta }F_{\beta \gamma }V^\gamma $$
(6)
where in above equation, $`(V^\alpha ,\sigma )`$ stand for the four velocity of a conducting neutral plasma and its scalar electrical conductivity respectively. Although eqs. II to 6 are valid for any kind of background geometries and plasmas characterized by arbitrary four velocity and conductivity, hereafter we shall restrict our consideration to background geometries that are globally static. Staticity in turn allows us to select coordinates so that the spacetime geometry can be written in the form (see for instance discussion in ref. ):
$$ds^2=e^{2\mathrm{\Phi }}(dx^o)^2+\gamma _{ij}dx^idx^j$$
(7)
where $`x^o=ct`$, $`\gamma _{ij}`$ are functions of the the spatial coordinates $`x^i`$, $`(i=1,2,3)`$, and by $`\xi `$ denote the hypersurface orthogonal timelike Killing vector field obeying: $`\xi _\alpha \xi ^\alpha =e^{2\mathrm{\Phi }}`$. For the above form of the line element, Maxwell’s equations II and the current conservation law $`_\alpha J^\alpha =0`$ can be re-written in an equivalent form involving only the components $`(E^i,B^i)`$ of the electric and magnetic fields respectively, as well as the charge density $`c\rho =U_\mu J^\mu `$ and spatial current density $`J^i`$ as measured by the Killing observers . More precisely if by $`U^\mu `$ we denote their four velocity then eqs. II yield the following equivalent set (see Appendix for details, or ref. ):
$$D_iE^i=4\pi \rho ,D_iB^i=0$$
(9)
$$ϵ^{ijk}D_j(ZB_k)=\frac{4\pi }{c}ZJ^i+\frac{E^i}{x^o}$$
(10)
$$ϵ^{ijk}D_j(ZE_k)=\frac{B^i}{x^o}$$
(11)
$$U^\mu \frac{(c\rho )}{x^\mu }+D_iJ^i+J^iD_ilogZ=0$$
(12)
where in above $`D`$ stands for the covariant derivative operator associated with $`\gamma `$, $`ϵ^{ijk}`$ represent the (coordinate) components of the three dimensional totally antisymmetric Levi-Civita tensor density defined on the $`x^o=const`$ slices and $`Z=(\xi ^\alpha \xi _\alpha )^{\frac{1}{2}}=e^\mathrm{\Phi }`$ is the red shift factor which in the language of in the $`3+1`$ approach to spacetime or (and) electrodynamics, is also refered as the lapse function .
With Maxwell’s eqs. in the above form we can derive an induction equation by repeating the same steps leading to the derivation of its flat counterpart (see for example discussion in ). For a plasma at rest relative to the Killing observers, combined with Ohm’s law and the MHD approximation (i.e. neglecting the displacement current from the right hand side of 10), one obtains from 7-abc the following form of the generalized induction equation:
$$\frac{B^i}{x^o}+ϵ^{ijk}D_j\left[\frac{c}{4\pi \sigma }ϵ_{k}^{}{}_{}{}^{lm}D_l(ZB_m)\right]=0$$
(13)
This last equation describes the time evolution of a magnetic field configuration that finds itself in a conducting medium. In principle one could write down the explicit form of the dynamical evolution equation once a choice of background geometry has been made. However before we do so, we would like to make a further specialization of the eqs. 7 and 13, so that their interrelationship to the familiar flat space three plus one formalism of Maxwell’s eqs. is more transparent. Here, following the spirit of and particularly , we shall sacrifice the manifest three-covariance of eqs. 7 and 13 with respect to arbitrary coordinate transformations of the $`t=const`$ sections, for the benefits of practical usefulness. As was pointed out in ref., if one defines suitably the components of $`(E,B)`$ and under some weak constraints upon the spacetime geometry, then Maxwell’s equations can be recast in a more ”user friendly” form. This new form employs concepts familiar from the language of the three dimensional vector analysis expressed in orthogonal curvilinear coordinates and such approach to curved spacetime electrodynamics is particularly useful for astrophysical purposes. Having in mind further astrophysical applications of our results we shall recast eqs. 7-abc in such a form. Such form requires that the geometry of the spacetime permits the introduction of coordinates so that the spatial three element $`ds_{(3)}^2`$ of 7 could be recast in the following form:
$$ds_{(3)}^2=h_1^2(dx^1)^2+h_2^2(dx^2)^2+h_3^2(dx^3)^2$$
(14)
where the scale factors $`h_i=h_i(x^1,x^2,x^3)`$ are for the moment arbitrary functions of $`(x^1,x^2,x^3)`$. In the Appendix, (see also ) we show that for such geometries eqs. 7 can be written in the following form:
$$\stackrel{}{E}=4\pi \rho ,\stackrel{}{B}=0$$
(16)
$$\stackrel{}{}\times (Z\stackrel{}{B})=\frac{4\pi }{c}Z\stackrel{}{J}+\frac{1}{c}\frac{\stackrel{}{E}}{t}$$
(17)
$$\stackrel{}{}\times (Z\stackrel{}{E})=\frac{1}{c}\frac{\stackrel{}{B}}{t}$$
(18)
$$\stackrel{}{J}+\stackrel{}{J}(logZ)=0$$
(19)
where we have written the current conservation law for an electrically neutral plasma and in above equations the symbols $`(,\stackrel{}{}\times ,)`$ stand for the divergence, curl and gradient operators respectively, expressed entirely in terms of the scale factors $`h_i`$ (see Appendix for their explicit representation). We also remind the reader that all vector components in equations 14 are physical frame components taken with respect to the field of orthonormal frames $`e_i=\frac{1}{h_i}\frac{}{x^i},(i=1,2,3)`$, naturally singled out by the line element 14. Using now eqs. 14, or directly from eq. 13, upon eliminating the coordinate components of $`\stackrel{}{B}`$ in favor of its frame components, the induction equation 13 takes the following form:
$$\frac{1}{c}\frac{\stackrel{}{B}}{t}+\stackrel{}{}\times \left[\frac{c}{4\pi \sigma }\stackrel{}{}\times (Z\stackrel{}{B})\right]=0$$
(20)
Equations 14 and 20 are the main equations of this section. In the special case of a $`Schwarzschild`$ background, naturally they are reduced to those of ref., and in the case of a plasma of uniform conductivity the generalized induction eq. 20 reduces to eq. 1 in the limit of flat space.
## III Magnetic field decay interior to neutron stars
In the neutron star’s interiors the MHD approximation is well justified and we shall explore the content of the relativistic induction equation 20, by applying it to study the evolution of magnetic fields associated with neutron stars. Since the main purpose of the present work is to investigate the impact of the spacetime curvature upon the magnetic field decay, as a first preliminary step we shall adopt a rather simplified neutron star model. The chosen model primarily avoids technicalities that may obscure the issue at hand but at the same time it shows clearly the potential impact of the curvature on the magnetic field decay. Accordingly, and to avoid laborious numerical computations, we shall ignore the rotation of the neutron star and thus shall adopt as the background geometry a non-singular, static and spherically symmetric one. Hence, the scale factors of eq. 14 will be taken as:
$$h_r^2=\left(1\frac{2Gm(r)}{rc^2}\right)^1=\left(1\frac{2M(r)}{r}\right)^1,h_\theta ^2=r^2,h_\varphi ^2=r^2\mathrm{sin}^2\theta $$
(21)
while for the moment the lapse or red shift factor $`Z=Z(r)=e^{\mathrm{\Phi }(r)}`$ and the ”mass function” $`m=m(r)`$ are arbitrary functions of the radial coordinate.
We shall begin our analysis of the magnetic field decay by assuming that at some initial time $`t_o`$ an axially symmetric distribution of a magnetic field $`\stackrel{}{B}(t_o,r,\theta )`$ permeates the entire star. We are not concerned here upon the mechanism that brought such a field into existence but rather we are interested in its evolution. Its evolution is considerably affected by the electrical conductivity $`\sigma `$, but as a part of the adopted simplified picture and in order to emphasize the effects of space-time curvature we shall take $`\sigma `$ to be spherically symmetric and shall ignore any cooling effects that may influence its temporal evolution. For an axially symmetric field $`\stackrel{}{B}`$, it is convenient to decompose it into the so called poloidal $`\stackrel{}{B}_{(p)}`$ and toroidal part $`\stackrel{}{B}_{(t)}`$. In terms of the orthonormal basis vectors $`(e_r,e_\theta ,e_\varphi )`$ those parts are defined respectively by: $`\stackrel{}{B}_{(p)}=B^r\stackrel{}{e}_r+B^\theta \stackrel{}{e}_\theta `$ and $`\stackrel{}{B}_{(t)}=B^\varphi \stackrel{}{e}_\varphi `$ with $`(B^r,B^\theta ,B^\varphi )`$ arbitrary functions of $`(t,r,\theta )`$ respectively. One can then easily conclude from the induction eq. 20 that, as long as the scalar conductivity is spherically symmetric, the toroidal and poloidal parts of $`\stackrel{}{B}`$, evolve independently of each other . Such decoupling is rather convenient since it implies that if the initial distribution of the magnetic field is purely poloidal then it will not develop a toroidal component in the course of its evolution and vice versa. For simplicity, in the present paper we shall examine the effects of the spacetime curvature only on the evolution of a purely poloidal field $`\stackrel{}{B}_{(p)}=B^r\stackrel{}{e}_r+B^\theta \stackrel{}{e}_\theta `$. For such field $`\stackrel{}{B}`$, it follows from 17 that the current $`\stackrel{}{J}`$ is along the $`\stackrel{}{e}_\varphi `$ direction, and thus the current conservation eq. 19 is identically satisfied. Furthermore via Ohm’s law, and use of 17 (with the displacement current ignored), it follows that the electric field $`\stackrel{}{E}`$ is a purely toroidal and axisymmetric field and Gauss law $`\stackrel{}{E}=0`$ is satisfied as well. Consequently, from the system of eqs. 14, we are left to satisfy the constraint $`\stackrel{}{B}=0`$, solutions of which will be evolved by the induction eq. 20.
Taking into account the poloidal and axisymmetric nature of $`\stackrel{}{B}`$ as well as the formula of the div operator $``$, listed in the Appendix, in view of the scale factors of 21, one easily finds that $`B=0`$ implies:
$$\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{1}{r}\frac{(r^2B^r)}{r}+\frac{1}{\mathrm{sin}\theta }\frac{(B^\theta \mathrm{sin}\theta )}{\theta }=0$$
(22)
We shall look for separable solutions of the above equations in the form:
$$B^r=F(t,r)\mathrm{\Theta }_1(\theta ),B^\theta =G(t,r)\mathrm{\Theta }_2(\theta )$$
(23)
with the functions $`F,G,\mathrm{\Theta }_1,\mathrm{\Theta }_2`$ to be determined. Substituting the above representations of $`(B^r,B^\theta )`$ in 22 and separating variables one gets the following equivalent system:
$$\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{1}{r}\frac{(r^2F)}{r}\lambda G=0$$
(25)
$$\frac{1}{\mathrm{sin}\theta }\frac{(\mathrm{sin}\theta \mathrm{\Theta }_2)}{\theta }+\lambda \mathrm{\Theta }_1=0$$
(26)
where $`\lambda `$ stands for a separation constant. The second equation can be solved in terms of the Legendre polynomials by taking $`\lambda =l(l+1),l=0,1,2\mathrm{}`$, and:
$$\mathrm{\Theta }_2=\mathrm{sin}\theta \frac{dP_l(y)}{dy},\mathrm{\Theta }_1=P_l(y),y=\mathrm{cos}\theta $$
(28)
On the other hand, for such $`\lambda `$, eq. 25 is satisfied provided, for $`l0`$, one chooses $`G(r,t)`$ in the following form;
$$G(t,r)=\frac{1}{l(l+1)}\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{1}{r}\frac{(r^2F)}{r}$$
(29)
We shall disregard the $`l=0`$ mode since, as it is clear form above, it corresponds to a monopole field $`B`$. With the exclusion of monopole fields, the components of an arbitrary axisymmetric poloidal field can be written as a superposition of ”l-poles” in the form:
$$B^r(t,r,\theta )=\underset{l=1}{\overset{\mathrm{}}{}}F_l(t,r)P_l(y)$$
(31)
$$B^\theta (t,r,\theta )=\underset{l=1}{\overset{\mathrm{}}{}}\frac{1}{l(l+1)}\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{1}{r}\frac{(r^2F)}{r}\mathrm{sin}\theta \frac{dP_l(y)}{dy}$$
(32)
To simplify algebra, and on physical grounds, we shall restrict our considerations to the detailed analysis of only the $`l=1`$ mode. Such mode corresponds to a dipole field and such configuration is expected to be present and dominant within neutron stars. For $`l=1`$, eqs. 29 yield:
$$B^r(t,r,\theta )=F(t,r)\mathrm{cos}\theta ,B^\theta (t,r,\theta )=\frac{1}{2r}\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{(r^2F)}{r}\mathrm{sin}\theta $$
(33)
where for notational simplicity we write here after $`F`$ instead of $`F_1`$. On the other hand, for any poloidal axisymmetric field with components $`(B^r,B^\theta )`$, the induction equation 20 on the background geometry of 14, yields the following two non-trivial evolution equations:
$$\frac{B^r}{x^o}+\frac{1}{h_\theta h_\varphi }\frac{}{\theta }\left[\frac{cA}{4\pi \sigma }\right]=0$$
(35)
$$\frac{B^\theta }{x^o}\frac{1}{h_rh_\varphi }\frac{}{r}\left[\frac{cA}{4\pi \sigma }\right]=0$$
(36)
where:
$$A=\frac{h_\varphi }{h_rh_\theta }\left[\frac{}{r}\left(h_\theta B^\theta Z\right)\frac{}{\theta }\left(h_rB^rZ\right)\right]$$
(37)
When one now inserts in eq. 35 the explicit forms of the components of $`(B^r,B^\theta )`$ corresponding to a dipole field in the form shown in eq. 33, as well as the scale factors of eq. 21, then gets:
$$\frac{4\pi \sigma }{c}\frac{F}{x^o}=\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{1}{r^2}\frac{}{r}\left[Z\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{(r^2F)}{r}\right]\frac{2ZF}{r^2}$$
(38)
In arriving at the above equation we have taken explicitly into account the spherically symmetric nature of the scalar conductivity $`\sigma `$. We may point out that for non spherical $`\sigma `$, the right hand side of 38 contains gradients of $`\sigma `$ along the meridian directions but for our simple neutron star model a spherical conductivity is rather adequate. On the other hand identical manipulations of (3.8b) leads to:
$$\frac{}{r}\left[r^2\frac{F}{x^o}\frac{c}{4\pi \sigma }\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{}{r}\left[Z\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{(r^2F)}{r}\right]+2\frac{c}{4\pi \sigma }ZF\right]=0$$
(40)
from which we infer that
$$r^2\frac{F}{x^o}\frac{c}{4\pi \sigma }\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{}{r}\left[Z\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{(r^2F)}{r}\right]+2\frac{c}{4\pi \sigma }ZF=g(\theta ,\varphi ,t)$$
(41)
where $`g(\theta ,\varphi ,t)`$ an integration ”constant”. A comparison then between (3.9) and (3.10b) shows that necessary $`g=0`$. If we further define: $`\widehat{F}=r^2F`$, then one finds either from 38 or 41 that $`\widehat{F}`$ satisfies:
$$\frac{4\pi \sigma }{c}\frac{\widehat{F}}{x^o}=\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{}{r}\left[Z\left(1\frac{2M}{r}\right)^{\frac{1}{2}}\frac{\widehat{F}}{r}\right]2\frac{Z}{r^2}\widehat{F}$$
(42)
The above equation essentially describes the evolution of the the dipole field components as they are measured by the Killing observers. Linearity of Maxwell’s and induction equation implies that 42 specifies a unique solution up to an arbitrary rescaling. This rescaling freedom will be fixed later on by a suitable matching of the interior dipole field to a corresponding asymptotically vanishing exterior dipole one.
Taking $`Z=1,M=0`$ in eq. 42 one recovers the standard equation describing the evolution of a dipole axisymmetric poloidal field in flat space namely:
$$\frac{4\pi \sigma }{c^2}\frac{S}{t}=\frac{^2S}{r^2}\frac{2S}{r^2}$$
(43)
where in order to avoid confusion we have denoted the analogue of $`\widehat{F}(t,r)`$ for the flat space time case by $`S(t,r)`$. The latter function often in the astrophysics literature is referred to as the Stoke’s function . In order to get some insights into the significance of the various terms appearing in (3.11), the general relativistic counterpart of 43, we shall rewrite the former in an equivalent form so that a clean comparison between the two could be afforded. Eliminating the (areal) radial coordinate $`r`$ in favor of the physical proper radius $`l(r)`$ of the $`r=constant`$ spheres, via $`dl=dr(1\frac{2M}{r})^{\frac{1}{2}}`$, Eq. 42 takes then the following form:
$$\frac{4\pi \sigma }{c^2}\frac{\widehat{F}}{t}=\frac{}{l}(Z\frac{\widehat{F}}{l})\frac{2Z}{r(l)^2}\widehat{F}$$
(45)
A comparison then to eq. 43 shows that relative to the Killing observers, general relativistic effects can influence B-decay in three ways. Namely via the presence of the red shift factor $`Z`$, its gradient and as well as via the intrinsically curved nature of the rest space of the Killing observers, ie the $`t=const`$ hyperfaces. The latter manifest itself in 45 via the term $`r(l)`$, a term which in general satisfies $`r(l)l`$, implying that that the rest spaces of the Killing observers are intrinsically curved. From the above mentioned three factors the spatial gradient of $`Z`$ makes negligible contribution to the field decay and this has been verified numerically. Neglecting this gradient then equation 45 takes the following form:
$$\frac{4\pi \sigma }{c^2Z(l)}\frac{\widehat{F}}{t}=\frac{^2\widehat{F}}{l^2}\frac{2\widehat{F}}{r(l)^2}$$
(46)
In this form a clean comparison to equation (3.12) can be afforded. The right hand sides of the two equations involve physical spatial gradients and the differ only by terms of the order $`O(\frac{2Gm}{c^2R})`$. On the other hand their left hand sides as they stand cannot be compared. If however one reasonably replaces $`Z(l)`$ by some averaged value $`<Z>`$, then the left hand side of (3.13) involves also physical temporal gradients. In that event one gets a first flavor of the magnitude of the general relativistic effects. They modify the corresponding flat spacetime results by terms of order unity. Of course such conclusions has to be also documented at the solution level as well, and as we shall further ahead this indeed is the case.
Having thus identified the manner by which relativistic gravity effects the magnetic field decay, our assignment is now to access the relative importance of each of the above two factors. In the following section we shall do so by resorting to numerical computations. However, before we pass to that issue let us first record the suitable boundary conditions to be imposed upon the corresponding $`S(t,r),\widehat{F}(t,r)`$ in order to describe sensible physics. The required conditions for both, ie the flat and the general relativistic case, are drawn by demanding that the interior $`\stackrel{}{B}`$ ought to be a non singular field at all times and at all spatial points, and in addition it ought to join smoothly across the surface of the star to an exterior asymptotically vanishing dipole field. For the flat space case, taking $`M=0`$ and $`F(t,r)=\frac{2S}{r^2}`$ in eq. 33 one gets: $`\stackrel{}{B}=\frac{2S}{r^2}\mathrm{cos}\theta \stackrel{}{e}_r\frac{1}{r}\frac{S}{r}\mathrm{sin}\theta \stackrel{}{e}_\theta `$ from which we infer that the magnitude $`\stackrel{}{B}^2`$ of the interior magnetic field is given by: $`|\stackrel{}{B}|^2=\frac{4S^2}{r^4}\mathrm{cos}^2\theta +\frac{1}{r^2}(\frac{S}{r})^2\mathrm{sin}^2\theta `$. Accordingly a regular field at the star’s center requires $`lim\frac{S(t,r)}{r^2}`$ to be finite as the center of the star is approached. On the other hand an asymptotically vanishing dipole magnetic field in flat space due to a magnetic moment $`\mu `$, is described by : $`\stackrel{}{B}=\frac{2\mu \mathrm{cos}\theta }{r^3}\stackrel{}{e}_r+\frac{\mu \mathrm{sin}\theta }{r^3}\stackrel{}{e}_\theta `$. It follows then from the above expressions that a $`C^0`$ matching of the interior magnetic field to slow varying exterior dipole field , requires that across the star’s surface, ie at radius $`R`$, $`S(t,r)`$ should satisfy:
$$R\frac{S}{r}|_R=S$$
(47)
The relatively simple nature of (3.12) as well the simple form of the boundary-regularity conditions outlined above, permit us to construct exact closed form solutions. In fact, it is not difficult to verify that for the case of a star of a uniform conductivity $`\sigma `$, a sequence of exact solutions of eq. (3.12) obeying the above described conditions is given by :
$$S_n(t,x)=\left[\frac{\mathrm{sin}(n\pi x)}{n^2\pi ^2x}\frac{\mathrm{cos}(n\pi x)}{n\pi }\right]e^{\frac{t}{\tau _n}}=\psi _n(x)e^{\frac{t}{\tau _n}}$$
(48)
where $`\tau _n=\frac{4\sigma R^2}{\pi c^2n^2}=\frac{1}{n^2}\frac{\tau _{\mathrm{Ohm}}}{\pi ^2}`$, $`x=\frac{r}{R}`$ and $`n`$ takes the values $`(1,2,3\mathrm{}.)`$.
The above sequence of exact solutions offers a clear picture regarding the behavior of a magnetic field in a conducting medium that finds itself in a flat space time. Constructing for instance the field $`\stackrel{}{B}_1(t,r,\theta )`$ corresponding to $`S_1(t,r)`$, one immediately sees that an observer at fixed $`(r,\theta )`$ finds that the magnitude of $`\stackrel{}{B}_1(t,r,\theta )`$ decays exponentially with a characteristic e-folding time given by $`\tau _1=\frac{4\sigma R^2}{\pi c^2}`$. On the other hand expanding an arbitrary initial field configuration $`\stackrel{}{B}(t_0,r,\theta )`$ in terms of the eigenfunctions $`(\psi _n,n=1,2..)`$, one can easily see the spatial diffusion of the initial distribution. For a plasma characterized by an arbitrary $`\sigma `$, although the decay and diffusive nature of the initial $`\stackrel{}{B}`$ field remain intact, it is rather difficult to estimate analytically the characteristic decay time as well as to find out whether the decay will be channeled into an exponential phase. It is sufficient, however, to stress that as long as we are in flat spacetime the decay process is controlled by the conducting properties of the background medium and, of course, the length scale of the initial field distribution.
Let us now turn the discussion to the formulation of the appropriate conditions to be imposed on solutions of eq. 42. Since as already indicated in the introduction section, a dipole field on a $`Schwarzschild`$ is modified considerably from its flat form and, as eq. 42 shows, relativistic effects modify the local behavior of the relativistic Stoke’s function $`\widehat{F}(t,r)`$, one expects modification of the boundary conditions as well. As far as the behavior of $`\widehat{F}(t,r)`$ at the star’s center is concerned, by arguing in the same manner as in the flat space case, a non singular dipole field requires $`\widehat{F}(t,r)`$ to satisfy identical conditions at the star’s center as its flat counterpart, namely $`lim\frac{\widehat{F}(t,r)}{r^2}`$ should be finite as the center of the star is approached (after all, the principle of equivalence holds). Although this is the case at the star’s center the corresponding boundary conditions at the star’s surface are markedly different. Recalling that the frame component of the vector potential $`A=A_\mu e^\mu =A_\varphi e^\varphi `$ describing a magnetic dipole on a $`Schwarzschild`$ background is described by :
$$A_\varphi =\frac{3\mu \mathrm{sin}\theta }{4M^2}\left[x\mathrm{ln}\left(1\frac{1}{x}\right)+\frac{1}{2x}+1\right]$$
(49)
where $`x=\frac{r}{2M(R)}`$, $`Rr<\mathrm{}`$, and $`\mu `$ is the dipole moment. From 49, one then obtains the corresponding field $`\stackrel{}{B}=B^r\stackrel{}{e}_r+B^\theta \stackrel{}{e}_\theta `$ where the physical components $`(B^r`$, $`B^\theta )`$ as measured by the Killing observers are given by:
$$B^r=\frac{2\mu \mathrm{cos}\theta }{r^3}\left[3x^3\mathrm{ln}(1x^1)+3x^2+\frac{3}{2}x\right]$$
(51)
$$B^\theta =\frac{\mu \mathrm{sin}\theta }{r^3}\left[6x^3(1x^1)^{\frac{1}{2}}\mathrm{ln}(1x^1)+6x^2\frac{1\frac{1}{2x}}{(1x^1)^{\frac{1}{2}}}\right]$$
(52)
A comparison of (3.17, ab)) with the corresponding eq. 33 and a $`C^0`$ matching between the two along the star’s surface, requires that the exterior dipole magnetic moment $`\mu `$ should be identified with the generally slow time varying part of the function $`\widehat{F}(t,r)`$ . Moreover the gradient $`\widehat{F}(t,r)`$ along the radial direction should obey:
$$R\frac{\widehat{F}(t,r)}{r}|_R=G(y)\widehat{F}(t,R)$$
(54)
with:
$$G(y)=y\frac{2y\mathrm{ln}(1y^1)+\frac{2y1}{y1}}{y^2\mathrm{ln}(1y^1)+y+\frac{1}{2}}$$
(55)
and $`y=\frac{R}{2M(R)}`$.
Thus the behavior of the interior dipole field in the presence of curvature is described by $`\widehat{F}(t,r)`$ satisfying the differential eq. 42, subject to boundedness of $`\frac{\widehat{F}(t,r)}{r^2}`$ as the star’s center is approached, and additionally obeying 54 at its surface. Before we turn our discussion to the construction of solutions of eq. 42 subject to the above discussed conditions, we would like to write an explicit formula for the time evolution of $`\widehat{F}(t,r)`$ under the assumption that geometry of the spacetime corresponds to a static spherically symmetric star, solution of Einstein’s equations. We may recall that in the derivation of eqs. (3.11) and (3.18,ab) we assumed an arbitrary non-singular, static, spherically symmetric background geometry with the only constraint that it joins smoothly to an exterior $`Schwarzschild`$ field across the surface of the star. Nowhere in the derivation we needed the explicit form of the $`M(r)=\frac{Gm(r)}{c^2}`$ nor the form of $`Z=Z(r)`$. Hereafter we shall become more explicit and shall take the background interior geometry to be a non singular solution of the coupled Einstein-perfect fluid system. As is well known, and under the assumption that the background Maxwell field makes negligible contribution to the structure of the star , static spherically symmetric, perfect fluid solutions of Einstein’s equations imply satisfaction of the following differential equations between the metric functions $`\mathrm{\Phi }(r)`$, $`M(r)`$, the hydrostatic pressure $`P(r)`$, and mass density $`\rho (r)`$ (see for instance ):
$$\frac{d\mathrm{\Phi }}{dr}=\frac{M(r)+4\pi r^3P(r)}{r^2\left(1\frac{2M(r)}{r}\right)}$$
(57)
$$\frac{dM(r)}{dr}=4\pi r^2\rho $$
(58)
$$\frac{dP(r)}{dr}=(\rho +P)\frac{M(r)+4\pi r^3P}{r^2\left(1\frac{2M(r)}{r}\right)}$$
(59)
Making use of those equations, and restoring the fundamental units, we obtain from 42 the following equation to be satisfied by the relativistic Stokes function:
$$\frac{4\pi \sigma }{c^2}e^{\mathrm{\Phi }(r)}\frac{F}{t}=\left(1\frac{2Gm(r)}{c^2r}\right)\frac{^2F}{r^2}+\frac{1}{r^2}\left[\frac{2Gm(r)}{c^2}+\frac{4\pi G}{c^2}r^3\left(\frac{P(r)}{c^2}\rho (r)\right)\right]\frac{F}{r}\frac{2}{r^2}F$$
(60)
where for typographical convenience we shall write here after $`F(r,t)`$ instead of $`\widehat{F}(r,t)`$. The above equation via 33, describes the evolution of any axisymmetric, dipole, poloidal field $`\stackrel{}{B}`$ that finds itself interior to a spherical perfect fluid star. In the above form it includes all three relativistic factors influencing the field decay. Since the distributions of $`(m(r),P(r),\rho (r))`$ are related via Einstein’s equations directly to the spacetime curvature, eq. 60 shows implicitly that the influence of spacetime curvature on the decay of the magnetic field is a real effect and cannot be removed via coordinate transformations. In principle one could insert in 60 the appropriate distributions of $`m(r),P(r),\rho (r)`$ resulting from integrating the Oppenheimer-Tolman-Volkov equation, specify $`\sigma =\sigma (r,t)`$, and construct the history of the $`\stackrel{}{B}`$-decay. We shall report elsewhere our findings of this rather laborious numerical integration . For the purpose of the present paper we shall integrate 60 for a rather simple system introduced and discussed in the following section. The goal of this section is to show that the general relativistic eq. 60 (or its approximate forms corresponding to eq. (3.13b) under the assumption of a uniform conductivity admits decay modes analogous to those of the flat space case with one important difference: The corresponding e-fold decaying times are longer in the relativistic case. We interpret this amplification of the e-fold decay times as resulting from the non vanishing space time curvature. Unfortunately however it is not easy to construct analytically the exact decaying modes of the full relativistic system 60 or its approximate versions (3.13a), and thus we shall resort to numerical computations. The emphasis in those computations is the probing of the dependence of the corresponding e-folding times upon the value of the red shift factor or (and) upon the strength of the curvature of the spatial sections.
## IV Magnetic field decay in a constant density star. Explicit results.
We shall consider in this section the decay of a magnetic field in a neutron star of constant density. The assumption of a constant density star, although not a very reliable approximation of a real neutron star, offers the advantage that the Einsteins equations can be solved analytically (see for instance ), and thus provides us with closed form expressions for the coefficients of the induction equation, eq. 60, and the boundary condition, eq. 52. In particular, in Box 23.2 of the ref. the distribution of the various hydrodynamical and geometrical variables are plotted as functions of the areal coordinate $`r`$. As we have already discussed, the exact decaying modes of the flat space-time induction equation for a uniform conductivity are explicitly known (given by eq. 48), while the corresponding decaying modes of the full curved space-time equation 60 are presently unknown. We shall therefore resort to numerical techniques in an attempt to get insights into behavior of the space of solutions of eq. 60.
Viewed as an initial-boundary value problem, (3.20) is a diffusive initial value problem for which the standard numerical technique is the Crank-Nicholson implicit integration scheme (see for example ref. for a description). We checked our numerical code by evolving the fundamental mode of the flat space-time case (ie take n=1 in eq. 48) and compare the numerical solution with the analytical one: we obtained an accuracy better than 1% until times up to $`10\tau _1`$, with $`\tau _1=\frac{4\sigma R^2}{\pi c^2}=\frac{\tau _{\mathrm{Ohm}}}{\pi ^2}`$ the corresponding decay time of the $`n=1`$ flat fundamental mode.
Before we turn our discussion of the numerical results it is helpful to view the time evolution of a chosen initial distribution from a complementary point of view. On general grounds $`F(t=0,r)`$ as well as its time evolution can be (formally) expanded in a series of the following form:
$$F(t,x)=\mathrm{\Sigma }a_ne^{\frac{c^2\lambda _nt}{4\pi \sigma R^2}}g_n(x)$$
(61)
where the summation is extended over all eigenmodes $`g_n(x)`$ of the corresponding (singular) Sturm-Liouville eigenvalue problem arising from eq. 60 and the associated boundary-regularity conditions:
$$Lg_n+\lambda _ne^\mathrm{\Phi }g_n=0$$
(63)
,
$$L=\left(1\frac{2Gm(x)}{c^2x}\right)\frac{^2}{x^2}+\frac{1}{x^2}\left[\frac{2Gm(x)}{c^2}+\frac{4\pi G}{c^2}x^3\left(\frac{P(x)}{c^2}\rho \right)\right]\frac{}{x}\frac{2}{x^2}$$
(64)
where in the present case the coefficients in $`L`$ are determined by the geometrical and hydrodynamical variables of the constant density star solution, $`x=\frac{r}{R}`$ and $`R`$ is the areal radius of the star. It follows now from (4.1) that if the eigenvalues are positive and well spaced, then after $`t>>\frac{t_{Ohm}}{\lambda _1}`$ where $`\lambda _1`$ is the lowest eigenvalue of the above system, then the evolution of the distribution will channel into an exponentially decreasing phase with the dominant contribution in the sum (4.1) coming from the ”first” term. Our subsequent described numerical computations exhibits such feature and this property allows to construct numerically the lowest eigenvalue of the above system .
In the following numerical calculations we have taken the areal radius to be $`R=10`$ km, a constant uniform conductivity $`\sigma =10^{25}`$ s<sup>-1</sup> typical of neutron star values (which implies $`\tau _{\mathrm{Ohm}}=\mathrm{4.44\hspace{0.33em}10}^9`$ yrs), and we consider various neutron star masses characterized by different values of the dimensionless compactness ratio: $`ϵ=\frac{2GM}{Rc^2}`$ = 0, 0.3, 0.4, 0.5, 0.6, 0.740, 0.810, 0.865, and 0.889. The first one corresponds to a flat background space time, the last four are those values used in the numerical plots of ref. , while current realistic neutron star models are characterized by $`ϵ`$ in the range 0.3 to 0.5 . For each value of $`ϵ`$, at first we have solved numerically the full relativistic induction equation 60, by taking the initial $`F(t=0,r)`$ to be equal to the Stoke’s function $`S(t=0,r)`$ of the corresponding first fundamental decay modes of 48 (ie take $`n=1`$ and $`t=0`$ in 48). After performing a long time-integration of eq. 60 subject to the conditions cited earlier on, we find that the evolution of $`F(t,x)`$ channels into an exponentially decaying mode which means, according to 61, that the evolution of the initial distribution eventually is described by the first non vanishing term in the series expansion 61. This behavior of $`F(t,x)`$ allows us to determine only the lowest eigenvalue $`\lambda _1`$ of (4.2) from our numerical outputs. Besides the explicit determination of $`\lambda _1`$, our numerical treatment allows us to construct the magnetic field as well. In Fig. 1, we plot as a function of coordinate time $`t`$, the magnetic field as perceived by a Killing observer located at the star’s pole for the various values of the compactness ratio. Fig.(1) shows that once curvature effects are incorporated and upon ignoring the initial transit time during which the field is in a superposition of various curved eigenmodes, the field follows an exponential decay law (as would have done in the absence of gravity) but now the corresponding e-folding time is longer than the corresponding flat spacetime case. Thus even though we have started with identically prepared systems their evolution is distinct, a distinction traced in the influence of relativistic effects. It should be stressed however that the content of Fig.1 does not by itself provide us with a clear overall picture of the field decay. It rather provides us with a characteristic physical decay time as perceived by a Killing observer situated at the surface of the star and this decay time should not be extrapolated as being the physical decay time over the entire star . In fact, each Killing observer located at some $`r`$, will compute a physical e-fold decay time $`\tau (r)`$ given by: $`\tau (r)=Z(r)(\lambda _1)^1=\frac{Z(r)\tau _{Ohm}}{\beta \pi ^2}`$ and obviously this value changes across the star. (In this formula, we have parametrized $`\lambda _1`$ so that $`\beta =1`$ corresponds to flat spacetime). Because of this spatial dependence of $`\tau (r)`$, in order to get a better insights into the dynamics of the decay, in Fig.2, we have plotted $`\lambda _1=\beta \frac{\pi ^2}{\tau _{\mathrm{Ohm}}}`$, as a function of the compactness ratio $`ϵ`$. In the same figure, for comparison purposes, we have plotted the value of the red shift factor at the stars center ($`Z_o=e^{\mathrm{\Phi }(o)}`$) and surface ($`Z_s=e^{\mathrm{\Phi }(s)}`$) respectively. Thus it follows from Fig.(2), that the relativistic corrections to the lowest eigenvalue $`\lambda _1`$, are bounded from above by $`Z_s`$ while from bellow (almost) by $`Z_0`$. It is more instructive however and complements the content of Fig.(2), a plot showing the physical decay times as measured by Killing observers located at the center and at the surface of the star respectively. Fig.(3) stands for such plot, and its content shows that the physical decay time can vary considerably across the star. In fact there are regions around the star’s center, where the physical decay time is shorter than the corresponding flat space case and this effect is more pronounced as the compactness ratio increases. In contrast to what occurs in the vicinity of the star’s center, in the crust region the physical decay is always larger than the corresponding flat case. Because of this behavior, largely due to gravitational time dilation effect, we assign an overall physical decay time, by averaging the physical e-fold decay at the center and the surface of the star respectively. This amounts to assigning an overral a red shift factor $`Z`$ for the entire star equal roughly to its value at the middle of the star. With this type of averaging, Figs.(2,3) shows that for small values of the compactness ratio the overall physical decay time is almost identical to the flat space time case. However, as the compactness ratio increases, the relativistic effects become more apparent. For the case of neutrons stars with range in the realistic domain ie $`ϵ`$ in the range $`(0.3,0.5)`$, and via the averaging procedure outlined above, the overall physical e-fold decay time is $`(1.21.3)`$ larger than the corresponding flat case. Although the content of Fig. 1, Fig. 2 and Fig. 3 show the impact of relativistic effects upon the field decay, by themselves they do not offer a clear insight as which (if any) of the two factors, ie red shift or spatial curvature, are responsible for the dominant contribution in the field decay. In order to access their relative importance we solve numerically eq. 45 in two extreme cases and show the numerical outputs in Fig.(4). First eq. 45 is solved under the assumption $`r(l)=l`$ and in this approximation the relativistic effects on the decay are solely due to the red shift factor $`Z(l)`$. In fig.(4) the numerical outputs are indicated by the label: ”curved time”. In the opposite extreme, we adopt $`Z(l)=1`$ in eq. 45 and take $`r(l)`$ as given by the metric corresponding to a constant density star. Thus in this approximation, the only relativistic effect influencing the decay is due to the spatial curvature. The resulting numerical outputs in Fig.(4) are marked by the label ”curved space”. It follows then clearly from the content of Fig.(4) that for a constant density neutron star, the dominant effect in the field decay is due to the red-shift factor $`Z`$, since the corresponding eigenvalues indicated by ”curved-time” graph are much closer to the corresponding exact eigenvalue indicated by ”curved space-time” in Fig.(4). Moreover, the dominance of the red-shift factor holds through for all values of the compactness ratio $`ϵ`$, and increases as $`ϵ`$ increases.
From the analysis presented so far it is clear that the more compact the star is, the longer is the e-folding time. As a consequence one expects that models of pulsars with soft equation of state to maintain a strong magnetic field for longer period of time than the corresponding models with a stiff equation of state. In turn such slow $`B`$ field decay implies additional source of heating ie Joule heating and such additional heating, may explain the relatively high temperature observed in old neutron stars. However based on the present analysis, it is rather premature to draw definite conclusions. For instance cooling effects leading to the temporal variation of the conductivity as well as the the detailed structure of the star and its rotation has to be taken into account. Such study is currently under way and we expect to report in a future communication.
## V Conclusions
The behavior of the surface magnetic fields of neutron stars is a complicated and controversial issue. Many processes are believed to influence its magnitude and its subsequent evolution. Trapping for instance, of the field in the superconducting core is one possibility. The expulsion of the field out of this region is a delicate matter involving many different branches of physics . Another possibility that in principle influences enormously the magnetic properties of neutron stars is related to the accretion processes immediately after the core collapse . Accretion and particularly hypercritical accretion, can submerge the field of the new born neutron star beneath a layer of accreting matter thus in principle producing a delayed switched on mechanism for the pulsar activity . Furthermore according to recent work the neutron star may never turn on as a pulsar if the accretion is hyperctical. Besides the above mechanisms influencing the evolution of neutron star’s magnetic fields, many more have been introduced and discussed at length in the current literature. In this work, we have present a limited framework taking into account the effects of space time curvature on the field decay. For the simple neutron star models with a corresponding compactness ratio in the range $`(0.3,0.5)`$, considered in the present work, we have seen an overall increase in the decay time, $`(1.2`$ to $`1.3)`$ times larger than the flat spacetime value. Although the present work is preliminary and to assess the new effect more work is needed , it point towards to the direction that in a strongly gravitating system, effects due to space time curvature should not be neglected.
###### Acknowledgements.
This work was supported by a binational grant DFG (grant #444 - MEX - 1131410) - Conacyt (grant #E130.443), Conacyt (grant #2127P - E9507), UNAM - DGAPA (grant #IN105495) and Coordinación Científica - UMSNH. D.P. and T.Z. are thankful to the Astrophysikalisches Institut Potsdam for its kind hospitality and U.G. to the Instituto de Astronomía of UNAM.
## A $`(3+1)`$ Form of Maxwell’s Eqs. on Static spacetimes
In this Appendix we shall sketch a derivation of eqs.(2.8) starting from the covariant form of Maxwell’s eqs. II. The derivation makes use of the existence of the hypersurface orthogonal timelike Killing field and although all the following computations can be done in a covariant fashion for brevity we work explicitly in the coordinate gauge of 7. We shall also present formulas required for the derivation of equations in section II.
Starting from the temporal component of the inhomogeneous Maxwell eq. 3, combined with the line element 7 and taking into account the fact that $`E^j=F^{jo}e^\mathrm{\Phi }`$ one immediately obtains:
$$D_iE^i=\frac{4\pi }{c}J^oe^\mathrm{\Phi }=4\pi \rho $$
(A1)
where we have defined the charge density $`\rho `$ measured by a Killing observer by: $`c\rho =J^\mu U_\mu `$ and we denote by $`D`$ the covariant derivative operator associated with the Riemannian metric of the $`t=const`$ spaces. Due to the fact that the Maxwell tensor $`F_{\alpha \beta }`$ admits the following easy verifiable decomposition: $`F_{\alpha \beta }=U_\alpha E_\beta U_\beta E_\alpha +ϵ_{\alpha \beta \gamma \delta }U^\gamma B^\delta `$ one gets the following expression for its spatial part $`F_{ij}`$: $`F_{ij}=ϵ_{oijl}B^lU^o=ϵ_{ijl}B^l`$. Passing now to the spatial components of 3 one gets
$$\frac{E^i}{x^o}ϵ^{ijl}D_j(ZB_l)=\frac{4\pi }{c}ZJ^i$$
(A2)
where we have introduced the red shift factor $`Z`$ via $`Z=(\xi ^a\xi _a)^{\frac{1}{2}}=e^\varphi `$ instead of $`e^\mathrm{\Phi }`$. On the other hand the second pair of Maxwell’s equations 4 can be written equivalently as:
$$\frac{F_{\mu \nu }}{x^\lambda }+\frac{F_{\nu \lambda }}{x^\mu }+\frac{F_{\lambda \mu }}{x^\nu }=0$$
(A3)
Taking now all the indices to be spatial, and eliminating $`F_{ij}`$ one gets: $`D_iB^i=0`$ The other information encoded in the second pair of Maxwell eqs. can be revealed by considering the following arrangement of the spacetime indices: $`(\mu ,\nu ,\lambda =m,n,x^o)`$. For such arrangement one obtains:
$$\frac{(ϵ_{mnl}B^l)}{x^o}+\frac{(U_oE_n)}{x^m}+\frac{(U_oE_m)}{x^n}=0$$
(A4)
from which one easily obtains:
$$\frac{B^l}{x^o}+ϵ^{lmn}D_m(ZE_n)=0$$
(A5)
The current conservation eq. $`_\mu J^\mu =0`$ after a trivial rearrangement yields:
$$U^\mu \frac{c\rho }{x^\mu }+D_iJ^i+J^i\frac{logZ}{x^i}=0$$
(A6)
To pass into the equivalent set (14-abcd) and eq. 20 involving physical orthonormal components, we project all tensors involved onto the natural set of orthonormal vectors $`(e_i)`$ and one forms $`(e^i),(i=1,2,3)`$ respectively, associated with the line element 14. Thus for instance the electric field $`E`$ can be written as: $`E=E^i\frac{}{x^i}=E^{\widehat{i}}e_i`$ where $`E^{\widehat{i}}=h_iE^i`$ expresses the relationship between coordinate and frame components of the vector field $`E`$ and it is understood that no summation is involved over the repeated indices. With the help of the orthonormal components one for instance may rewrite eq. A1 in terms of orthonormal components. Writing $`\gamma ^{\frac{1}{2}}=h_1h_2h_3`$ and eliminating the coordinate components in terms of the frame components of $`E`$, equation A1 takes the following form:
$$\frac{1}{h_1h_2h_3}[\frac{}{x^1}(h_2h_3E^{\widehat{1}})+\frac{}{x^2}(h_1h_3E^{\widehat{2}})+\frac{}{x^3}(h_1h_2E^{\widehat{3}})]=\stackrel{}{E}=4\pi \rho $$
(A7)
where $``$ stands for the familiar divergence operator expressed in arbitrary orthogonal curvilinear coordinates defined by the line element 14. Similarly $`D_iB^i=0`$ can be written as $`\stackrel{}{B}=0`$. As far as the other set of Maxwell’s eqs. are concerned, one can proceed in a similar manner. For instance starting from A2, one first multiplies the corresponding equation by the scale factor $`h_i`$ thus leading into:
$$\frac{(E^{\widehat{i}})}{x^o}\frac{ϵ^{\widehat{i}\widehat{j}\widehat{k}}h_i}{h_1h_2h_3}\frac{(h_kB_{\widehat{k}}Z)}{x^j}=\frac{4\pi }{c}J^{\widehat{i}}Z$$
(A8)
Recalling that the orthonormal components of the $`\mathrm{𝑐𝑢𝑟𝑙}`$ operator of an arbitrary three dimensional differentiable vector field $`A`$ are given by :
$$(\times A)^{\widehat{i}}=\frac{ϵ^{\widehat{i}\widehat{j}\widehat{k}}h_i}{h_1h_2h_3}\frac{(h_kA_{\widehat{k}})}{x^j}$$
(A9)
one is lead immediately into eq. 17 used in the text. Note also that the action of the gradient operator $``$ acting on scalars is defined via:
$$f=\frac{1}{h_1}\frac{f}{x^1}+\frac{1}{h_2}\frac{f}{x^2}+\frac{1}{h_3}\frac{f}{x^3}$$
(A10)
Also in deriving eqs. (2.8-abc) of the main text we have used the following properties of the unit basis vectors $`(e_i)`$: $`e_{\widehat{i}}e_{\widehat{j}}=\delta _{\widehat{i}\widehat{j}}`$, $`e_{\widehat{i}}\times e_{\widehat{j}}=ϵ_{\widehat{i}\widehat{j}\widehat{k}}e_{\widehat{k}}`$ and the normalization $`ϵ_{\widehat{r}\widehat{\theta }\widehat{\varphi }}=1`$. We may also indicate that for typographical convenience the caret-symbol over frame components of the various tensors has been dropped. In particularly all vector, tensor components appearing anywhere in the main text after eq. 14, are frame components.
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# Are lepton flavor mixings in the democratic mass matrix stable against quantum corrections?
(May 8, 2000)
## Abstract
We investigate whether the lepton flavor mixing angles in the so-called democratic type of mass matrix are stable against quantum corrections or not in the minimal supersymmetric standard model with dimension five operator which induces neutrino mass matrix. By taking simple breaking patterns of $`S_3{}_{L}{}^{}\times S_3_R`$ or $`O(3)_L\times O(3)_R`$ flavor symmetries and the scale where democratic textures are induced as $`O(10^{13})`$ GeV, we find that the stability of the lepton flavor mixing angles in the democratic type of mass matrix against quantum corrections depends on the solar neutrino solutions. The maximal flavor mixing of the vacuum oscillation solution is spoiled by quantum corrections in the experimental allowed region of $`\mathrm{tan}\beta `$. The large angle MSW solution is spoiled by quantum corrections in the region of $`\mathrm{tan}\beta >10`$. The condition of $`\mathrm{tan}\beta 10`$ is needed in order to obtain the suitable mass squared difference of the small angle MSW solution. These strong constraints must be regarded for the model building of the democratic type of mass matrix.
hep-ph/0005064
DPNU-00-19
KEK-TH-691
PACS:14.60.Pq, 12.15.Ff
Recent neutrino oscillation experiments suggest the strong evidences of tiny neutrino masses and lepton flavor mixings . Studies of the lepton flavor mixing matrix, which is so-called Maki-Nakagawa-Sakata (MNS) matrix, will give us important cues of the physics beyond the standard model. One important study is finding the suitable texture of quark and lepton mass matrices in order to search the flavor symmetry existing behind. The democratic type of mass matrix is one of the most interesting candidate of the texture of quark and lepton mass matrices, since it can naturally explain the reason why only masses of third generation particles are large comparing to those of other generations. This type of mass matrix can be derived by flavor symmetries of $`S_3{}_{L}{}^{}\times S_3_R`$ or $`O(3)_L\times O(3)_R`$. As for the neutrino sector, it has been said that the democratic type of mass matrix can induce the suitable solutions of the atmospheric and the solar neutrino problems . However, are lepton flavor mixing angles in the democratic type of mass matrix stable against quantum corrections?
In this paper, we investigate whether the lepton flavor mixing angles in the democratic type of mass matrix are stable against quantum corrections or not in the minimal supersymmetric standard model with the dimension five operator which induces the neutrino Majorana mass matrix. The superpotential of the lepton-Higgs interactions is given by
$$𝒲=y_{ij}^\mathrm{e}(H_dL_i)E_j\frac{1}{2}\kappa _{ij}(H_uL_i)(H_uL_j),$$
(1)
where $`\kappa `$ is the coefficient of the dimension five operator, and the indices $`i,j`$ $`(=13)`$ stand for the generation number. $`L_i`$ and $`E_i`$ are chiral super-fields of $`i`$-th generation lepton doublet and right-handed charged lepton, respectively. $`H_u`$ ($`H_d`$) is the Higgs doublet which gives Dirac masses to the up- (down-) type fermions. We will show the renormalization group equation (RGE) analyses of the lepton flavor mixing angles .
We take the simple breaking patterns of $`S_3{}_{L}{}^{}\times S_3_R`$ or $`O(3)_L\times O(3)_R`$ symmetries, and the scale where democratic textures are induced as $`O(10^{13})`$ GeV. Under the above conditions, we find that the stability of the lepton flavor mixing angles in the democratic type of mass matrix against quantum corrections depends on the solar neutrino solutions. The maximal flavor mixing of the vacuum oscillation (VO) solution is spoiled by quantum corrections in the experimental allowed region of $`\mathrm{tan}\beta `$. The value of $`\mathrm{tan}\beta `$ is the ratio between the vacuum expectation values (VEVs) of the Higgs particles. The large angle MSW (MSW-L) solution is spoiled by quantum corrections in the region of $`\mathrm{tan}\beta >10`$. On the other hand, the condition of $`\mathrm{tan}\beta 10`$ is needed in order to obtain the suitable mass squared difference of the small angle MSW solution (MSW-S). These strong constraints must be regarded for the model building of the democratic type of mass matrix.
At first, we discuss the democratic mass matrix, which is based on $`S_{3L}\times S_{3R}`$ or $`O(3)_L\times O(3)_R`$ flavor symmetries. In the democratic type of mass matrix, the charged lepton mass matrix is given by
$$M_l=\frac{c_l}{3}\left(\begin{array}{ccc}1& 1& 1\\ 1& 1& 1\\ 1& 1& 1\end{array}\right)+M_l^{(c)},$$
(2)
where $`M_l^{(c)}`$ includes flavor symmetry breaking masses, which must be introduced to obtain the suitable values of $`m_e`$ and $`m_\mu `$. The matrix $`M_l`$ is diagonalized by the unitary matrix $`V_l=FL`$ from the side of left-handed fields, where
$$F=\left(\begin{array}{ccc}1/\sqrt{2}& 1/\sqrt{6}& 1/\sqrt{3}\\ 1/\sqrt{2}& 1/\sqrt{6}& 1/\sqrt{3}\\ 0& 2/\sqrt{6}& 1/\sqrt{3}\end{array}\right).$$
(3)
Since we do not know the definite structure of $`M_l^{(c)}`$, we can not determine the explicit form of the unitary matrix $`L`$. Here we assume that off diagonal elements of $`L`$ are small as $`L_{ij}1`$ $`(ij)`$ from the analogy of the quark sector . Thus, we obtain the relation of $`V_lF`$, which is used in the following discussions<sup>*</sup><sup>*</sup>* We will revive $`L_{ij}`$ in the MSW-S solution later. .
The neutrino mass matrix To obtain this mass matrix, there is an alternative way using $`S_3`$ symmetry instead of $`S_{3L}\times S_{3R}`$ or $`O(3)_L\times O(3)_R`$. is given by
$$M_\nu =c_\nu \left\{\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)+r\left(\begin{array}{ccc}1& 1& 1\\ 1& 1& 1\\ 1& 1& 1\end{array}\right)\right\}+M_\nu ^{(b)},$$
(4)
since the neutrinos are Majorana particles. In Eq.(4), $`c_\nu `$ and $`r`$ can be taken as real and non-negative parameters, and we neglect $`CP`$ phases, for simplicity. The mass matrix $`M_\nu ^{(b)}`$ breaks flavor symmetries, which must be introduced in order to obtain the suitable mass squared differences and mixing angles of neutrinos. There are following two simple breaking patterns according to the solar neutrino solutions.
(i): The simplest example of $`M_\nu ^{(b)}`$ for the MSW-L and the VO solutions is to introduce two real and non-negative parameters $`ϵ`$ and $`\delta `$ $`(1\delta ϵ)`$ in (2,2) and (3,3) elements in $`M_\nu ^{(b)}`$, where the neutrino mass matrix $`M_\nu ^{(i)}`$ is given by
$$M_\nu ^{(i)}=c_\nu \left(\begin{array}{ccc}1+r& r& r\\ r& 1+r+ϵ& r\\ r& r& 1+r+\delta \end{array}\right).$$
(5)
(ii): The simplest example of $`M_\nu ^{(b)}`$ for the MSW-S solution is to introduce two real and non-negative parameters $`ϵ`$ and $`\delta `$ $`(1\delta ϵ)`$ in (1,2), (2,1), and (3,3) elements in $`M_\nu ^{(b)}`$, where the neutrino mass matrix $`M_\nu ^{(ii)}`$ is given by
$$M_\nu ^{(ii)}=c_\nu \left(\begin{array}{ccc}1+r& r+ϵ& r\\ r+ϵ& 1+r& r\\ r& r& 1+r+\delta \end{array}\right).$$
(6)
When $`rϵ`$, $`\delta `$, the unitary matrix $`U_\nu `$, which diagonalizes $`M_\nu `$, becomes $`F`$ in both cases of (i) and (ii). In this case the MNS matrix approaches to the unit matrix as
$$V_{MNS}=V_l^{}U_\nu L^{}F^{}FL^{},$$
(7)
which does not have any large mixing angles. Therefore the magnitude of $`r`$ must be smaller than $`ϵ`$, $`\delta `$ in order to obtain the large flavor mixing of the atmospheric neutrino solution We do not consider the case of $`r=2/3`$ which gives degenerate neutrinos, where one of mass eigenvalues has a opposite sign from others. It is because the simple symmetry breaking patterns such as (i) and (ii) can not induce the large lepton flavor mixing as shown above. In the case of $`rϵ,\delta `$, the flavor symmetry breaking textures must be complicated in order to solve both the solar and the atmospheric neutrino problems. . Thus, in the democratic type of mass matrix, three neutrinos are degenerate with the same signs, where the RGE effects cannot be negligible as shown in case (c4) in Ref.. This is the reason why we need RGE analyses for the democratic type of mass matrix.
Under the condition of $`rϵ,\delta `$, both simple breaking patterns of (i) and (ii) induce the large mixing angle of $`\mathrm{sin}^22\theta _{23}8/9`$ which is suitable for the atmospheric neutrino solution, and negligibly small mixing between the first and the third generations as $`\mathrm{sin}^22\theta _{13}0`$ which is consistent with the CHOOZ experiment . Case (i) induces the maximal mixing between the first and the second generations for the solar neutrino solution as $`\mathrm{sin}^22\theta _{12}1`$. On the other hand, case (ii) induces the small mixing between the first and the second generations, since the maximal mixing angles induced from both $`M_l`$ and $`M_\nu `$ are canceled with each other.
Now let us estimate quantum corrections of the MNS matrix in the democratic type of mass matrix. We take the diagonal base of the charged lepton mass matrix at the high energy scale $`m_h`$, where the democratic textures are induced, for the RGE analysis. In this base the neutrino mass matrix in Eqs.(5) and (6) are written by $`V_l^TM_\nu (m_h)V_lF^TM_\nu (m_h)F`$, where we use the approximation of $`V_lF`$. Then, the neutrino mass matrix at $`m_Z`$ scale is given by
$$M_\nu \left(m_Z\right)=\frac{M_\nu \left(m_Z\right)_{33}}{M_\nu \left(m_h\right)_{33}}R_GF^TM_\nu (m_h)FR_G,$$
(8)
where the matrix $`R_G`$ shows the renormalization effects, which is defined as
$$R_G\left(\begin{array}{ccc}1+\eta & 0& 0\\ 0& 1+\eta & 0\\ 0& 0& 1\end{array}\right).$$
(9)
The small parameter $`\eta `$ is given by
$`\eta `$ $``$ $`1\mathrm{exp}\left({\displaystyle \frac{1}{16\pi ^2}}{\displaystyle _{\mathrm{ln}\left(m_Z\right)}^{\mathrm{ln}\left(m_h\right)}}y_\tau ^2𝑑t\right),`$ (10)
$``$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle \frac{m_\tau ^2}{v^2}}\left(1+\mathrm{tan}^2\beta \right)\mathrm{ln}\left({\displaystyle \frac{m_h}{m_Z}}\right),`$
where $`y_\tau `$ is the Yukawa coupling of $`\tau `$ and $`v^2H_u^2+H_d^2`$. We neglect the Yukawa couplings of $`e`$ and $`\mu `$ in Eq.(9), since those contributions to the RGE are negligibly small comparing to that of $`\tau `$. Therefore the first and the second generations receive the same RGE corrections as in Eq.(9). Now let us check whether the mixing angles receive significant changes by quantum corrections or not in both cases of (i) and (ii).
In the base of charged lepton democratic mass matrix, $`M_\nu (m_Z)`$ in case (i) is written as
$`M_\nu ^{(i)}(m_Z)`$ $`=`$ $`FR_GF^TM_\nu ^{(i)}(m_h)FR_GF^T,`$ (11)
$``$ $`\overline{c_\nu }\left(\begin{array}{ccc}1+\overline{r}& \overline{r}& \overline{r}\\ \overline{r}& 1+\overline{r}+ϵ& \overline{r}\\ \overline{r}& \overline{r}& 1+\overline{r}+\delta \end{array}\right)+2\eta \overline{c_\nu }\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right),`$ (18)
where
$$\overline{r}r\frac{2}{3}\eta ,\overline{c_\nu }\frac{M_\nu \left(m_Z\right)_{33}}{M_\nu \left(m_h\right)_{33}}c_\nu .$$
(19)
Here we neglect the small parameters of order $`ϵ^2`$, $`ϵ\eta `$, and $`ϵ\delta `$. Equation (18) means that the MNS matrix at $`m_Z`$ scale is obtained only by using $`\overline{r}`$ instead of $`r`$ in Eq.(5). This had been already shown in Ref.. Therefore all we have to do for the RGE analyses of the mixing angles is to trace the change of $`\overline{r}`$. Here we remind that the magnitude of $`\eta `$ is completely determined by the value of $`\mathrm{tan}\beta `$ and the scale of $`m_h`$, then $`\overline{r}`$ is determined by Eq.(19).
Now let us show how the MNS matrix changes as the change of $`\overline{r}`$.
(i-a): $`1\delta ϵ\left|\overline{r}\right|`$
Neglecting the second order of small parameters of $`\delta ,\overline{r}`$, and $`ϵ`$, mass eigenvalues of $`M_\nu ^{(i)}(m_Z)`$ are give by
$$\overline{c_\nu }(1+\overline{r}+2\eta ),\overline{c_\nu }(1+\overline{r}+ϵ+2\eta ),\overline{c_\nu }(1+\overline{r}+\delta +2\eta ).$$
(20)
Then the unitary matrix $`U_\nu `$ becomes
$$U_\nu \left(\begin{array}{ccc}1& \frac{\overline{r}}{ϵ}& \frac{\overline{r}}{\delta }\\ \frac{\overline{r}}{ϵ}& 1& \frac{\overline{r}}{\delta }\\ \frac{\overline{r}}{\delta }& \frac{\overline{r}}{\delta }& 1\end{array}\right),$$
(21)
which induces the MNS matrix $`V_{MNS}`$ as
$$V_{MNS}F^TU_\nu =\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\left(1+\frac{\overline{r}}{ϵ}\right)& \frac{1}{\sqrt{2}}\left(1\frac{\overline{r}}{ϵ}\right)& 0\\ \frac{1}{\sqrt{6}}\left(1+2\frac{\overline{r}}{\delta }\frac{\overline{r}}{ϵ}\right)& \frac{1}{\sqrt{6}}\left(1+2\frac{\overline{r}}{\delta }+\frac{\overline{r}}{ϵ}\right)& \sqrt{\frac{2}{3}}\left(1\frac{\overline{r}}{\delta }\right)\\ \frac{1}{\sqrt{3}}\left(1\frac{\overline{r}}{\delta }\frac{\overline{r}}{ϵ}\right)& \frac{1}{\sqrt{3}}\left(1\frac{\overline{r}}{\delta }+\frac{\overline{r}}{ϵ}\right)& \frac{1}{\sqrt{3}}\left(1+2\frac{\overline{r}}{\delta }\right)\end{array}\right).$$
(22)
Thus the mixing angles are given by
$$\mathrm{sin}^22\theta _{12}14(\frac{\overline{r}}{ϵ})^2,\mathrm{sin}^22\theta _{13}0,\mathrm{sin}^22\theta _{23}\frac{8}{9}\left(1+2\frac{\overline{r}}{\delta }9(\frac{\overline{r}}{\delta })^2\right).$$
(23)
This means all mixing angles are not changed by quantum corrections in the region of $`1\delta ϵ|\overline{r}|`$. The value of $`\mathrm{sin}^22\theta _{13}`$ is negligible, since $`\mathrm{sin}^22\theta _{13}O(L_{21}^2)1`$, where $`L_{21}`$ is the (21) component of L. Equation (20) shows that $`\mathrm{\Delta }m_{12}^22\overline{c_\nu }^2ϵ`$ and $`\mathrm{\Delta }m_{23}^22\overline{c_\nu }^2\delta `$. In order for the symmetry breaking parameter $`\delta `$ to be smaller than symmetric terms of order one, it must be that $`\delta O(0.1)`$. On the other hand, neutrinoless $`\beta \beta `$-decay experiments suggest $`\overline{c_\nu }O(0.1)`$ eV. Then, $`\overline{c_\nu }=O(0.1)`$ eV and $`\delta =O(0.1)`$ are obtained from the experimental results of $`\mathrm{\Delta }m_{\mathrm{ATM}}^210^3`$ eV<sup>2</sup>. This means that $`ϵ=O(10^3)`$ for the MSW-L solution, and $`ϵ=O(10^8)`$ for the VO solution. Then, the region of $`ϵ\overline{r}`$ corresponds to $`\mathrm{tan}\beta <10`$<sup>§</sup><sup>§</sup>§ The higher the scale of $`m_h`$ becomes, the smaller the value of $`\mathrm{tan}\beta `$ must be in order for the maximal flavor mixing not to be destroyed by quantum corrections. for the MSW-L solution and $`\mathrm{tan}\beta 1`$ for the VO solution. Since the region of $`\mathrm{tan}\beta 1`$ is excluded by the Higgs search experiments, we can conclude the maximal mixing of the VO solution in the democratic type of mass matrix, discussed in Refs., is completely spoiled by quantum corrections The maximal mixing of the VO solution is not spoiled by quantum corrections (even in $`\mathrm{tan}\beta =3`$), if $`m_hO(1)`$ TeV. However, such a low energy scale of $`m_h`$ is not suitable from the view point of model building. . For the MSW-L solution, discussed in Refs., the sufficient condition of $`\mathrm{tan}\beta <10`$ must be satisfied.
(i-b): $`1\delta \left|\overline{r}\right|ϵ`$
Neglecting the second order of small parameters of $`\delta ,\overline{r}`$, and $`ϵ`$, mass eigenvalues of $`M_\nu ^{(i)}(m_Z)`$ are give by
$$\overline{c_\nu }(1+\frac{1}{2}ϵ+2\eta ),\overline{c_\nu }(1+2\overline{r}+\frac{1}{2}ϵ+2\eta ),\overline{c_\nu }(1+\overline{r}+\delta +2\eta ).$$
(24)
The unitary matrix $`U_\nu `$ becomes
$$U_\nu \left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\left(1+\frac{1}{4}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{2}}\left(1\frac{1}{4}\frac{ϵ}{\overline{r}}\right)& \frac{\overline{r}}{\delta }\\ \frac{1}{\sqrt{2}}\left(1\frac{1}{4}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{2}}\left(1+\frac{1}{4}\frac{ϵ}{\overline{r}}\right)& \frac{\overline{r}}{\delta }\\ \frac{1}{2\sqrt{2}}\frac{ϵ}{\delta }& \sqrt{2}\frac{\overline{r}}{\delta }& 1\end{array}\right),$$
(25)
which induces the MNS matrix as
$$V_{MNS}F^TU_\nu =\left(\begin{array}{ccc}1& \frac{1}{4}\frac{ϵ}{\overline{r}}& 0\\ \frac{1}{2\sqrt{3}}\frac{\overline{r}}{\delta }& \frac{1}{\sqrt{3}}\left(1+\frac{1}{2}\frac{\overline{r}}{\delta }\right)& \sqrt{\frac{2}{3}}\left(1\frac{\overline{r}}{\delta }\right)\\ \frac{1}{2\sqrt{6}}\left(\frac{ϵ}{\delta }\frac{ϵ}{\overline{r}}\right)& \sqrt{\frac{2}{3}}\left(1\frac{\overline{r}}{\delta }\right)& \frac{1}{\sqrt{3}}\left(1+2\frac{\overline{r}}{\delta }\right)\end{array}\right).$$
(26)
Then the mixing angles are given by
$$\mathrm{sin}^22\theta _{12}\frac{1}{4}\left(\frac{ϵ}{\overline{r}}\right)^2,\mathrm{sin}^22\theta _{13}0,\mathrm{sin}^22\theta _{23}\frac{8}{9}\left(1+2\frac{\overline{r}}{\delta }9(\frac{\overline{r}}{\delta })^2\right).$$
(27)
The value of $`\mathrm{sin}^22\theta _{13}`$ is negligible, since $`\mathrm{sin}^22\theta _{13}O(L_{21}^2)1`$. This means that maximal mixings of all solar solutions in the democratic type of mass matrix of Eq.(18) are spoiled by quantum corrections in the region of $`1\delta |\overline{r}|ϵ`$, although the mixings between the first and the third generations, and between the second and the third generations are stable against quantum corrections If $`L_{21}O(1)`$, the mixing angle of $`\mathrm{sin}2\theta _{12}`$ can be stable agaist quantum corrections. However, it is difficult to obtain $`L_{21}O(1)`$ from $`M_l^{(c)}`$ .
(i-c): $`1\left|\overline{r}\right|\delta ϵ`$
Neglecting the second order of small parameters of $`\delta ,\overline{r}`$, and $`ϵ`$, mass eigenvalues of $`M_\nu ^{(i)}(m_Z)`$ are give by
$$\overline{c_\nu }(1+\frac{1}{2}ϵ+2\eta ),\overline{c_\nu }(1+\frac{2}{3}\delta +\frac{1}{6}ϵ+2\eta ),\overline{c_\nu }(1+3\overline{r}+\frac{1}{3}\delta +\frac{1}{3}ϵ+2\eta ).$$
(28)
In this case $`U_\nu `$ becomes
$$U_\nu \left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\left(1+\frac{1}{4}\frac{ϵ}{\delta }+\frac{1}{6}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{6}}\left(1\frac{3}{4}\frac{ϵ}{\delta }+\frac{2}{9}\frac{\delta }{\overline{r}}\frac{1}{9}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{3}}\left(1\frac{1}{9}\frac{\delta }{\overline{r}}\frac{1}{9}\frac{ϵ}{\overline{r}}\right)\\ \frac{1}{\sqrt{2}}\left(1\frac{1}{4}\frac{ϵ}{\delta }\frac{1}{6}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{6}}\left(1+\frac{3}{4}\frac{ϵ}{\delta }+\frac{2}{9}\frac{\delta }{\overline{r}}\frac{1}{9}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{3}}\left(1\frac{1}{9}\frac{\delta }{\overline{r}}+\frac{2}{9}\frac{ϵ}{\overline{r}}\right)\\ \frac{1}{2\sqrt{2}}\left(\frac{ϵ}{\delta }\frac{1}{3}\frac{ϵ}{\overline{r}}\right)& \sqrt{\frac{2}{3}}\left(1\frac{1}{9}\frac{\delta }{\overline{r}}+\frac{1}{18}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{3}}\left(1+\frac{2}{9}\frac{\delta }{\overline{r}}\frac{1}{9}\frac{ϵ}{\overline{r}}\right)\end{array}\right).$$
(29)
Thus, the MNS matrix is given by
$$V_{MNS}F^TU_\nu =\left(\begin{array}{ccc}1& \frac{\sqrt{3}}{4}\frac{ϵ}{\delta }& \frac{1}{3\sqrt{6}}\frac{ϵ}{\overline{r}}\\ \frac{\sqrt{3}}{4}\frac{ϵ}{\delta }& 1& \frac{\sqrt{2}}{9}\left(\frac{\delta }{\overline{r}}\frac{1}{2}\frac{ϵ}{\overline{r}}\right)\\ \frac{1}{2\sqrt{6}}\frac{ϵ}{\overline{r}}& \frac{\sqrt{2}}{9}\left(\frac{\delta }{\overline{r}}\frac{1}{2}\frac{ϵ}{\overline{r}}\right)& 1\end{array}\right),$$
(30)
which gives the mixing angles as
$$\mathrm{sin}^22\theta _{12}\frac{3}{4}\left(\frac{ϵ}{\delta }\right)^2,\mathrm{sin}^22\theta _{13}\frac{1}{54}\left(\frac{ϵ}{\overline{r}}\right)^2,\mathrm{sin}^22\theta _{23}\frac{8}{81}\left(\frac{\delta }{\overline{r}}+\frac{ϵ}{4\overline{r}}\right)^2.$$
(31)
This means that large mixing angles in both the solar and the atmospheric neutrino solutions are spoiled by quantum corrections in the region of $`1|\overline{r}|\delta ϵ`$. It is because the condition of $`|\overline{r}|\delta ,ϵ`$ induces $`U_\nu F`$, which is just the case of Eq.(7).
The conclusion in case (i) are that (1): the maximal mixing of the VO solution is destroyed by quantum corrections, and (2): the sufficient condition of $`\mathrm{tan}\beta <10`$ must be satisfied for the MSW-L solution.
Next, let us show the case (ii) in the base of charged lepton democratic mass matrix, where $`M_\nu ^{(ii)}(m_Z)`$ is written by
$$M_\nu ^{(ii)}(m_Z)\overline{c_\nu }\left(\begin{array}{ccc}1+\overline{r}& \overline{r}+ϵ& \overline{r}\\ \overline{r}+ϵ& 1+\overline{r}& \overline{r}\\ \overline{r}& \overline{r}& 1+\overline{r}+\delta \end{array}\right)+2\eta \overline{c_\nu }\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right).$$
(32)
Let us show how the MNS matrix changes according to the change of $`\overline{r}`$ as in case (i).
(ii-a): $`1\delta ϵ,|\overline{r}|`$
Neglecting the second order of small parameters of $`\delta ,\overline{r}`$, and $`ϵ`$, the mass eigenvalues of $`M_\nu ^{(ii)}(m_Z)`$ are give by
$$\overline{c_\nu }(1ϵ+2\eta ),\overline{c_\nu }(1+2\overline{r}+ϵ+2\eta ),\overline{c_\nu }(1+\overline{r}+\delta +2\eta ).$$
(33)
In this case $`U_\nu `$ becomes
$$U_\nu \left(\begin{array}{ccc}\frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}& \frac{\overline{r}}{\delta }\\ \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}& \frac{\overline{r}}{\delta }\\ 0& \sqrt{2}\frac{\overline{r}}{\delta }& 1\end{array}\right),$$
(34)
which induces the MNS matrix as
$`V_{MNS}=L^{}F^TU_\nu `$ (35)
$`\left(\begin{array}{ccc}1& {\displaystyle \frac{1}{\sqrt{3}}}L_{21}\left(1+2{\displaystyle \frac{\overline{r}}{\delta }}\right)& \sqrt{{\displaystyle \frac{2}{3}}}L_{21}\left(1{\displaystyle \frac{\overline{r}}{\delta }}\right)\\ L_{12}& {\displaystyle \frac{1}{\sqrt{3}}}\left(1+2{\displaystyle \frac{\overline{r}}{\delta }}\right)+{\displaystyle \frac{1}{\sqrt{3}}}L_{31}\left(1+{\displaystyle \frac{\overline{r}}{\delta }}\right)& \sqrt{{\displaystyle \frac{2}{3}}}\left(1{\displaystyle \frac{\overline{r}}{\delta }}\right)+{\displaystyle \frac{1}{\sqrt{3}}}L_{32}\left(1+2{\displaystyle \frac{\overline{r}}{\delta }}\right)\\ L_{13}& \sqrt{{\displaystyle \frac{2}{3}}}\left(1{\displaystyle \frac{\overline{r}}{\delta }}\right)+{\displaystyle \frac{1}{\sqrt{3}}}L_{23}\left(1+2{\displaystyle \frac{\overline{r}}{\delta }}\right)& {\displaystyle \frac{1}{\sqrt{3}}}\left(1+2{\displaystyle \frac{\overline{r}}{\delta }}\right){\displaystyle \frac{3}{2}}L_{23}\left(1{\displaystyle \frac{\overline{r}}{\delta }}\right)\end{array}\right),`$ (39)
where we revive the small elements of $`L_{ij}`$ $`(ij)`$. This shows that the mixing angles are given by
$$\mathrm{sin}^22\theta _{12}\frac{4}{3}L_{21}^2\left(1+2\frac{\overline{r}}{\delta }\right)^2,\mathrm{sin}^22\theta _{13}\frac{8}{3}L_{21}^2\left(1\frac{\overline{r}}{\delta }\right)^2,\mathrm{sin}^22\theta _{23}\frac{8}{9}\left(1+2\frac{\overline{r}}{\delta }9(\frac{\overline{r}}{\delta })^2\right).$$
(40)
This means that all flavor mixings are not spoiled by quantum corrections in the region of $`1\delta |\overline{r}|,ϵ`$. Equation (33) suggests that $`\mathrm{\Delta }m_{12}^24\overline{c_\nu }^2(\overline{r}+ϵ)`$ and $`\mathrm{\Delta }m_{23}^22\overline{c_\nu }^2\delta `$. Thus, when $`|\overline{r}|ϵ`$, quantum correction is the origin of mass squared difference for the solar neutrino solution. Where we must tune the value of $`\mathrm{tan}\beta `$ in order to obtain the suitable mass squared difference. The case of $`\mathrm{tan}\beta 10`$ induces $`\mathrm{\Delta }m_{12}^210^5`$ eV<sup>2</sup> at $`m_Z`$. Therefore the condition of $`\mathrm{tan}\beta 10`$ must be satisfied in order to obtain the suitable magnitude of mass squared difference, for the MSW-S solution discussed in Ref.. As for the mixings between the first and the third generations, and between the second and the third generations, Eq.(40) shows that they are stable against quantum corrections.
(ii-b): $`1|\overline{r}|\delta ϵ`$
Neglecting the second order of small parameters of $`\delta ,\overline{r}`$, and $`ϵ`$, mass eigenvalues of $`M_\nu ^{(ii)}(m_Z)`$ are give by
$$\overline{c_\nu }(1ϵ+2\eta ),\overline{c_\nu }(1+\frac{2}{3}\delta +\frac{1}{3}ϵ+2\eta ),\overline{c_\nu }(1+3\overline{r}+\frac{1}{3}\delta +\frac{2}{3}ϵ+2\eta ).$$
(41)
In this case $`U_\nu `$ becomes
$$U_\nu \left(\begin{array}{ccc}\frac{1}{\sqrt{2}}& \frac{1}{\sqrt{6}}\left(1+\frac{2}{9}\frac{\delta }{\overline{r}}\frac{2}{9}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{3}}\left(1\frac{1}{9}\frac{\delta }{\overline{r}}\frac{2}{9}\frac{ϵ}{\overline{r}}\right)\\ \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{6}}\left(1+\frac{2}{9}\frac{\delta }{\overline{r}}\frac{2}{9}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{3}}\left(1\frac{1}{9}\frac{\delta }{\overline{r}}\frac{2}{9}\frac{ϵ}{\overline{r}}\right)\\ 0& \sqrt{\frac{2}{3}}\left(1\frac{1}{9}\frac{\delta }{\overline{r}}+\frac{1}{9}\frac{ϵ}{\overline{r}}\right)& \frac{1}{\sqrt{3}}\left(1+\frac{2}{9}\frac{\delta }{\overline{r}}+\frac{4}{9}\frac{ϵ}{\overline{r}}\right)\end{array}\right),$$
(42)
which induces the MNS matrix as
$$V_{MNS}=L^{}F^TU_\nu =\left(\begin{array}{ccc}1& L_{21}& \frac{\sqrt{2}}{9}L_{21}\left(\frac{\delta }{\overline{r}}+\frac{ϵ}{\overline{r}}\right)\\ L_{12}& 1+L_{32}\left(\frac{\delta }{\overline{r}}\frac{ϵ}{\overline{r}}\right)& \frac{\sqrt{2}}{9}\left(\frac{\delta }{\overline{r}}2\frac{ϵ}{\overline{r}}\right)+L_{32}\\ L_{13}& \frac{\sqrt{2}}{9}\left(\frac{\delta }{\overline{r}}\frac{ϵ}{\overline{r}}\right)+L_{23}& 1\frac{\sqrt{2}}{9}\left(\frac{\delta }{\overline{r}}\frac{ϵ}{\overline{r}}\right)\end{array}\right).$$
(43)
This suggests that the mixing angles are given by
$$\mathrm{sin}^22\theta _{12}4L_{21}^2,\mathrm{sin}^22\theta _{13}\frac{2}{81}L_{21}^2\left(\frac{\delta }{\overline{r}}+\frac{ϵ}{\overline{r}}\right)^2,\mathrm{sin}^22\theta _{23}\frac{8}{81}\left(\frac{\delta }{\overline{r}}+2\frac{ϵ}{\overline{r}}\right)^2,$$
(44)
which means that the large mixing of the atmospheric neutrino solution is destroyed in the region of $`1|\overline{r}|\delta ϵ`$. It is because the condition of $`|\overline{r}|\delta ,ϵ`$ induces $`U_\nu F`$, which is just the case of Eq.(7).
The conclusion in case (ii) is that sufficient condition of $`\mathrm{tan}\beta 10`$ must be satisfied for the MSW-S solution.
The democratic type of mass matrix texture is one of the most interesting candidate of quark and lepton mass matrices, which has been said to be able to induce the suitable solutions of the atmospheric and the solar neutrino problems. In this paper, we investigate whether the lepton flavor mixing angles in the democratic type of mass matrix are stable against quantum corrections or not in the minimal supersymmetric standard model with the dimension five operator which induces the neutrino Majorana mass matrix. We take the simple breaking patterns of $`S_3{}_{L}{}^{}\times S_3_R`$ or $`O(3)_L\times O(3)_R`$ symmetries, and the scale where democratic textures are induced as $`O(10^{13})`$ GeV. Under the above conditions, we find that the stability of mixing angles in the democratic type of mass matrix against quantum corrections depends on the solar neutrino solutions. The maximal mixing of the VO solution is spoiled by quantum corrections in the experimentally allowed region of $`\mathrm{tan}\beta `$. The MSW-L solution is spoiled by quantum corrections in the region of $`\mathrm{tan}\beta >10`$. On the other hand, the condition of $`\mathrm{tan}\beta 10`$ is needed in order to obtain the suitable mass squared difference of the MSW-S solution. These strong constraints must be regarded for the model building of the democratic type of mass matrix. If we take $`m_h`$ as the GUT scale of $`O(10^{16})`$ GeV, the constraints for the stability of the mixing angles against quantum corrections become more severe, that is, the MSW-L solution needs $`\mathrm{tan}\beta <8`$, and the MSW-S solution needs $`\mathrm{tan}\beta 8`$.
We would like to thank T. Yanagida and M. Bando for the suggestion of this work. The work of NO is supported by the JSPS Research Fellowship for Young Scientists, No.2996.
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# The Fate of Cosmic String Zero Modes
## 1 Introduction
Although much of the evolution of the Universe is well understood, there are still many cosmological phenomena for which a completely satisfactory explanation has yet to be found. Topological defects, such as cosmic strings, could provide mechanisms for structure formation, CMB anisotropy, and high energy cosmic rays . Such defects form in many realistic particle physics theories, including those involving supersymmetry.
In the past it has been difficult to evaluate the usefulness of such ideas due to a lack of data. This is now changing, and predictions of CMB anisotropies from simple cosmic string models have been made . While these predictions show poor agreement with the observations, they do not take into account the full physics of string models. Indeed, recent analysis which includes the effect of particle production as an energy loss mechanism from the string network shows much improved agreement with data .
One significant possibility is that the strings carry conserved currents . These currents will alter the evolution of a string network, which could lead to better agreement with observation. Indeed, an analytic analysis showed that a much denser string network results for electromagnetically coupled strings . However, one significant implication of conserved currents is that they can stabilise loops of string. If persistent, these stabilised loops or ‘vortons’ , can easily dominate the energy density of the Universe, placing stringent constraints on the parameters of the model. An analysis has been made of the implications of this for particle physics models that predict current carrying strings .
Fermions are a natural choice for the charge carriers of such currents. Fermion zero modes exist in a wide class of cosmic string models. The fermions can be excited and move along the string, resulting in a current. Consequently, fermion conductivity occurs naturally in many supersymmetric and grand unified theories , such as SO(10). In particular, in SUSY models with a D term, there is a single chiral zero mode, either a left or right mover. This zero mode survives supersymmetry breaking . This class of theory is potentially important since it arises in many superstring models.
One criticism of fermion currents is that, unlike scalar boson currents, they are not topologically stable. It is possible that they could decay, either directly into particles off the string, or through interactions with the surrounding plasma. If the decay rate is too high, currents will not last long enough to have any significant effect. On the other hand if the decay rate is too low the Universe could become vorton dominated. We address this issue and show that these processes can remove current carriers close to the phase transition, but not otherwise . We then consider vorton constraints in chiral theories, showing this leads to stringent constraints on the underlying particle physics theory.
## 2 $`D`$-term Supersymmetric Cosmic Strings
Consider a supersymmetric theory with a $`\text{U(1)}I`$ phase transition. The simplest way to achieve this symmetry breaking is to use one charged chiral superfield and a non-zero Fayet-Iliopoulos term. Expanding the Lagrangian in terms of component fields gives (in Wess-Zumino gauge)
$`=|(D_\mu \varphi |^2{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }+|F|^2+{\displaystyle \frac{1}{2}}D^2+D(g|\varphi |^2+\xi )`$
$`i\psi \sigma ^\mu D_\mu ^{}\overline{\psi }i\lambda \sigma ^\mu _\mu \overline{\lambda }+ig\sqrt{2}\varphi ^{}\psi \lambda +(\text{c.c.}).`$ (1)
The two fermion fields $`\psi `$ and $`\lambda `$ are the Higgsino and gaugino respectively. Eliminating the two auxillary fields $`D`$ and $`F`$ will give potential terms. If $`\xi =g\eta ^2<0`$ the resulting $`D`$-term will give rise to spontaneous symmetry breaking, resulting in an expectation value of $`\eta `$ for $`\varphi `$. In order to avoid gauge anomalies, the model must contain other charged superfields. We will assume these have zero expectation values.
As well as the $`\varphi =`$ constant solution there also exist string solutions obtained from the ansatz
$`\varphi `$ $`=`$ $`\eta e^{in\theta }f(r)`$ (2)
$`A_\mu `$ $`=`$ $`n{\displaystyle \frac{a(r)}{gr}}\delta _\mu ^\theta `$ (3)
$`D`$ $`=`$ $`g\eta ^2(1f(r)^2),`$ (4)
where$`f(r)`$ and $`a(r)`$ are the usual Neilsen-Olesen profile functions.
Now consider the fermionic sector of the theory. Performing a SUSY transformation with 2-spinor parameter $`ϵ`$ gives
$`\delta \lambda _1`$ $`=`$ $`2ig\eta ^2(1f^2)ϵ_1`$
$`\delta \psi _1`$ $`=`$ $`2\sqrt{2}i\eta {\displaystyle \frac{n}{r}}(1a)fe^{i(n1)\theta }ϵ_1^{}.`$ (5)
$`\delta \psi _2`$ and $`\delta \lambda _2`$ are both zero. The expression (5) is a zero energy solution of the fermion field equations. It is trivial to add $`z`$ and $`t`$ dependence to it, and the resulting solution is a null current moving along the string.
Since the string solution is not invariant under transformations with $`ϵ_10`$, supersymmetry has been broken inside the string core. However it is only partially broken there since the string is still invaraint under transformations parametrised by $`ϵ_2`$.
It is also possible to get a phase transition using an $`F`$-term. This requires a non-trivial superpotential and at least three scalar fields. The corresponding string solutions have twice as many zero modes as the $`D`$-term theory, and move in both directions along the string. In this case supersymmetry is completely broken . In both theories the zero modes and SUSY breaking are confined to the string core.
When soft supersymmetry-breaking terms are included in the Lagrangian, oppositely moving zero modes in the $`F`$-term theory mix to form massive states . This cannot happen in the $`D`$-term theory since the zero modes all move in the same direction.
Thus, a generic property of cosmic strings in SUSY theories is that supersymmetry is broken in the string core and the resulting strings have fermion zero modes. As a consequence, cosmic strings arising in SUSY theories are automatically current-carrying and can give rise to vortons.
## 3 Decay Rates of Charge Carriers
Vortons can only be stable if the currents they carry are stable, thus the stabilty of the charge carriers is a crucial consideration. In the absence of other particles, currents carried by zero modes on isolated, straight strings are stable on grounds of energy and momentum conservation. However, in a realistic setting there are many processes which can depopulate zero modes. Strings are not isolated and in the early Universe they and their bound states will interact with the hot plasma. SO(10) is considered as a specific example. In this case a heavy neutrino zero mode may scatter from a light plasma particle to produce a light fermion-antifermion pair via an intermediate electroweak Higgs boson. A massive neutrino bound mode may also simply decay into light fermions. Further, bound states on different strings, or different parts of a single curved string, can scatter from one another by exchanging a Higgs particle.
The decay rates of both massive and massless modes have been calculated at tree level . The main difference between scattering in the string background and that in a trivial background is the lack of translation invariance transverse to the string. This enters the calculation via the bound mode wavefunctions which are localised around the string. Integration over the initial vertex position does not yield the standard momentum conserving $`\delta `$-function, instead it produces an approximately gaussian function that permits non-conservation of transverse momentum on the scale of the off-string fermion mass.
The violation of momentum conservation in the string background opens up significant areas of phase space that are forbidden in the trivial background. Of particular importance is the possibility of resonant scattering. For example the lifetime of a bound mode of mass $`m_\mathrm{B}`$ decaying into a light fermion and an electroweak Higgs is,
$$\tau (|g_\nu |^2m_\mathrm{B})^1,$$
(6)
where $`g_\nu `$ is the Yukawa coupling in the neutrino’s electroweak mass term. Massive modes with the appropriate couplings to the electroweak sector thus have a very short lifetime. The above result neglects the possibility that incoming and outgoing fermionic wavefunctions may be amplified near the string and should be taken as an overestimate.
The fate of massless bound states is also complicated by momentum non-conservation. Of particular interest are the bound states that stabilise vortons. Contraction of the vorton will ensure that the states at the Fermi surface will have GUT scale momenta. The lifetime of these high momentum states is critical; if they decay the vorton will contract and promote low momentum states to high momentum. Energy and $`z`$-momentum conservation prevent massless states from decaying spontaneously. However it is possible for them to decay by interaction with plasma particles or other zero modes. If the string Higgs mass, $`m_\mathrm{s}`$ exceeds the off-string fermion mass, $`m`$, it is also possible for high energy massless currents on a curved string to decay by tunneling to free heavy neutrinos . However the rate will not be significant unless $`m_\mathrm{s}m`$.
If we consider a massless bound state at the Fermi surface scattering from a typical plasma particle, we have a centre of mass energy of order $`\sqrt{m_\mathrm{s}T}`$. In the SO(10) case this is well above the mass of the electroweak Higgs intermediate particle and transverse momentum non-conservation again allows for resonant scattering. Including amplification of the incoming plasma particle wavefunction, the lifetime of these high momentum zero modes is found to be,
$$\tau \frac{\stackrel{~}{m}^2}{|g_\nu |^2T^3}\left(\frac{T}{\stackrel{~}{m}}\right)^{2Q},$$
(7)
where $`\stackrel{~}{m}=m_\mathrm{s}m/2`$ and $`Q`$ is the charge of the plasma particle under the string gauge field.
In the radiation dominated era the time is given by $`t=\alpha T^2`$, where $`\alpha m_{\mathrm{Pl}}/10`$. In the case of SO(10), $`Q=3/10`$, and the probability of a zero mode state scattering after some time $`t_i`$ is small if,
$$t_i>O\left([\stackrel{~}{m}^7|g_\nu |^{10}\alpha ^6]\right)>O\left(\left[\frac{m_{\mathrm{Pl}}}{10\stackrel{~}{m}}\right]^6|g_\nu |^{10}\stackrel{~}{m}^1\right).$$
(8)
As the lifetime varies only slightly faster than $`T^2`$, this result for $`t_i`$ is very sensitive to the Yukawa coupling. For $`\stackrel{~}{m}10^{15}`$GeV, if $`|g_\nu |=1`$, zero mode states populated after $`t_i10^{15}t_{\mathrm{GUT}}`$ will be stable, while if $`|g_\nu |<0.03`$, this scattering is never significant. In the SO(10) model $`g_\nu `$ is also the Yukawa coupling for the corresponding quarks, thus there is an epoch when $`\nu _\tau ^c`$ zero modes will scatter from the string, but $`\nu _e^c`$ and $`\nu _\mu ^c`$ zero modes will never scatter by this process. Thus the interaction with plasma particles can not significantly remove zero modes from the string. Note that it is also possible to create currents using the above interactions in reverse. Hence, if thermal equilibrium is reached the number density of zero modes will be of order $`T`$.
Within the SO(10) model there is also the possibility of mediating these processes by GUT mass Higgs fields with zero VEV. In this case the Yukawa coupling need not be small, but the centre of mass energy of the interaction is only of order the intermediate particle mass for $`TT_{\mathrm{GUT}}`$. Thus below the GUT temperature the reaction rates for these processes are rapidly suppressed by powers of $`T/T_{\mathrm{GUT}}`$.
None of these plasma scattering processes can remove $`\nu _e^c`$ and $`\nu _\mu ^c`$ zero modes, and are only significant for $`\nu _\tau ^c`$ immediately after the phase transition. Thus, they are unable to prevent the vorton density from dominating the energy density of the Universe.
The plasma scattering processes considered above failed to remove zero modes due to the decreasing plasma density at late times. A distinct category of process is the scattering of a zero mode on one string by a second zero mode on another string. This is particularly relevant for vortons as they form small loops with a typical radius only one or two orders of magnitude larger than the string width. We thus have zero modes moving in opposite directions on opposite sides of the vorton. For simplicity the decay rate can be calculated by considering two straight, anti-parallel strings with spacing $`2R`$.
If the intermediate particle is an electroweak Higgs boson, $`R`$ is much smaller than the electroweak length scale and there is no exponential range suppression. Resonant scattering is not possible and the cross-section is found to be,
$$\sigma \frac{|g_\nu |^4}{(m_{\mathrm{GUT}}R)^4}.$$
(9)
This cross-section is dimensionless as the scattering is effectively in one spatial dimension.
Conversely, for a GUT mass intermediate particle, $`m_{\mathrm{GUT}}R10100`$ and the reaction rate displays exponential range suppression,
$$\sigma \frac{|g_\nu |^4}{(m_{\mathrm{GUT}}R)^3}e^{4m_{\mathrm{GUT}}R}.$$
(10)
While the string density may be high at formation, it drops rapidly and the exponential suppression in (10) makes such processes irrelevant in all physical situations.
Taken at face value, if electroweak particles mediate current–current scattering on different segments of string then (9) gives a short lifetime for charge carriers on a vorton. However, for a circular loop the angular momentum of the fermions must be conserved and, combined with energy conservation, this would prevent massless modes on the string scattering into massive modes. However, the fermionic spectrum has not been calculated for a circular loop, so there is no reason to expect the zero modes to remain massless. Thus angular momentum conservation can only be considered in a consistent, rotationally invariant calculation. The calculation above works consistently with straight strings, thus while (9) may not be directly applicable to vorton decay, it is relevant for interactions on non-circular loops and scattering of currents on a string network.
## 4 Chiral Strings and Vortons
We have seen that plasma interactions donot significantly remove zero modes from the string. This is particularly true in SUSY $`D`$-term theories (see section 2) where we would expect the zero mode to be isolated from the electroweak sector. As a consequence, stable vortons would be expected to form. In this section we constrain the underlying particle physics theory which gives rise to chiral zero modes on the string. We assume that the strings are formed at a phase transition ocurring at temperature $`T_\mathrm{x}`$ and become current-carrying at a scale $`T_\mathrm{s}`$. The string loop is characterised by two currents, the topologically conserved phase current $`N`$ and the dynamically conserved particle number current $`Z`$. In the chiral case these are exactly identical.
A non conducting string loop must ultimately decay by radiative and frictional drag processes until it disappears completely. However, a conducting string loop could reach a state in which the energy attains a minimum for given non zero values of $`N`$ and $`Z`$. This is the vorton state. It should be emphasised that the existence of such vorton states does not require that the carrier field be electromagnetically coupled. Indeed, in the case of the $`D`$-term zero modes, it is not.
The physical properties of a vorton state are determined by the quantum numbers, $`N`$ and $`Z`$. However, these are not arbitrary. Indeed, to avoid the fate of the usual loops, the quantum numbers on a conducting loop must be large compared with unity. This in turn implies a minimum length for the loop. In our work we calculate the number density of protovorton loops, subject to the loops being greater than a minimum length. We then estimate the vorton density and constrain the underlying theory by the requirement that the universe is not vorton dominated. We apply two constraints. Firstly we take a conservative assumption that the vortons only live a few minutes and constrain the theory by requiring that the universe is radiation dominated at nucleosynthesis; we then take the more realistic assumptions that chiral vorton loops are stable and constrain the theory by requiring that the vorton density is less than the closure density today. The full details are presented in ref. ; here we sumarise the main results.
For chiral vortons the vorton energy, $`E_\mathrm{v}`$, is given by
$$E_\mathrm{v}NT_\mathrm{x}.$$
(11)
In the friction dominated era the number density of vortons can be estimated by considering the damping length scale and the resulting correlation length below which microstructure is damped. This automatically satisfies the minimum length criteria mentioned above. For a loop of length L, the conserved quantum number is then
$$|Z|=NLT_\mathrm{s}=\left(\frac{m_{\mathrm{Pl}}}{T_\mathrm{s}}\right)^{1/2}\frac{T_\mathrm{x}}{T_\mathrm{s}}.$$
(12)
The resulting vorton density was found to be
$$\frac{\rho _\mathrm{v}}{T^3}\frac{T_\mathrm{s}^3}{m_{\mathrm{Pl}}T_\mathrm{x}},$$
(13)
where factors of order unity have been dropped.
Requiring then that this vorton density be less than the radiation density at nucleosynthesis gives a constraint for strings which become current-carrying at formation of
$$T_\mathrm{x}10^8\text{GeV}.$$
(14)
This is the condition that must be satisfied by the formation temperature of cosmic strings that become current-carrying immediately, subject to the rather conservative assumption that the resulting vortons last for at least a few minutes.
If the zero modes condense on the string at a separate phase transition then the constraint takes us outside the friction dominated regime. The calculation is much more involved and we refer the reader to ref. for details. However, the resulting constraint is that GUT scale strings becoming current-carrying at a temperature above $`10^9\text{GeV}`$ are inconsistent with data.
Being less conservative we can make the assumption that the vortons are absolutely stable and survive to the present time. Requiring that the vorton density is less than the closure density gives a stronger constraint. In this case we find
$$T_\mathrm{x}=T_\mathrm{s}10^5\text{GeV}.$$
(15)
For strings which become current-carrying at a later transition, the details are again more complicated. However, for GUT scale strings we find that strings becoming current-carrying at a temperature above $`10^5\text{GeV}`$ are inconsistent with data.
It is amusing to point out that if strings formed, or became current-carrying, just below this temperature, then their vortons would contribute substantially to the dark matter of the universe.
## 5 Discussion
We have shown that cosmic strings arising in SUSY theories are generically current-carrying. For D-term theories there is a chiral fermion zero mode which survives SUSY breaking. Consequently, the current persists.
We have investigated ways in which the fermion bound modes could be destabilised. Zero modes are the most resilient: various scattering mechanisms were investigated, but, because of the low density of the surrounding plasma, they were unable to scatter all zero modes off the string. Zero modes on neighbouring string segments can also scatter off each other. This may be relevant for vortons.
Strings with a chiral zero mode lead to the most stringent constraints on the underlying particle physics theory, because in this case the fermion travels in a single direction, resulting in the current being maximal rather than random. We have shown that this leads to very strong constraints indeed.
However, if the theory somehow manages to evade the vorton constraints then the resulting cosmology could be very different. For example, the scaling solution of this type of cosmic string theory is completely unknown. Whilst the cosmology of strings resulting from the abelian Higgs model has been investigated, that resulting from realistic GUT theories has not. As a consequence. it is premature to rule out cosmic strings as being incompatible with microwave background data. This is currently under investigation .
## Acknowledgments
This work was supported in part by PPARC and an ESF network. We wish to thank our collaborators Brandon Carter and Mark Trodden, and also Rachel Jeannerot and the other organisers of COSMO99.
## References
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# Photon-Photon Interaction in a Photon Gas
## Abstract
Using the effective Lagrangian for the low energy photon-photon interaction the lowest order photon self energy at finite temperature and in non-equilibrium is calculated within the real time formalism. The Debye mass, the dispersion relation, the dielectric tensor, and the velocity of light following from the photon self energy are discussed. As an application we consider the interaction of photons with the cosmic microwave background radiation.
Photon-photon scattering has been considered already a long time ago . To lowest order QED perturbation theory it is caused by the so-called box diagram, which contains an electron loop. For low energy photons, i.e. for center of mass energies below the threshold for $`e^+`$-$`e^{}`$ pair creation, an effective Lagrangian for the photon-photon interaction has been derived by integrating out the electrons,
$$_I=a(F_{\mu \nu }F^{\mu \nu })^2+bF_{\mu \nu }F^{\nu \rho }F_{\rho \sigma }F^{\sigma \mu },$$
(1)
where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$ is the electromagnetic field strength tensor and
$$a=\frac{5\alpha ^2}{180m_e^4},b=\frac{7\alpha ^2}{90m_e^4}$$
(2)
with the fine structure constant $`\alpha 1/137`$ and the electron mass $`m_e`$. This effective Lagrangian containing an effective 4-photon interaction describes the deviation from the classical Maxwell theory by quantum effects. The 4-photon vertex in momentum space following from this Lagrangian has been derived only recently . The coupling constant corresponding to this vertex is of the order $`\alpha ^2/m_e^4`$.
The lowest order photon self energy $`\mathrm{\Pi }_{\mu \nu }`$, which is given by the tadpole diagram of Fig.1, vanishes at zero temperature after dimensional regularization . However, at finite temperature tadpole diagrams lead to a finite result . The in-medium photon self energy determines the Debye screening, the photon dispersion relations, and the dielectric functions of the system. Some of the results presented here are already discussed in the literature . Here we want to treat the photon self energy and its consequences in a systematic and comprehensive way starting from the real time formalism, which also allows an extension to non-equilibrium situations. As an application we consider the influence of the cosmic microwave background radiation on low energy photons. Although medium effects of the photon gas are expected to be very small due to the extremely weak photon-photon coupling at low energies, these effects might be interesting after all since the properties of the photons are experimentally very well known. For example there are very restrictive upper limits for the photon mass , namely $`m_\gamma <2\times 10^{16}`$ eV measured in the laboratory and $`m_\gamma <10^{27}`$ eV using arguments about the galactic magnetic field. These upper limits follow from searching for violations of the Maxwell theory, in particular of the Coulomb law . Hence these limits should be compared to the Debye mass caused by the background of thermal particles. It should be noted that the Debye mass, which follows from the photon self energy, does not violate gauge symmetry of QED.
Mass effects due to charged particles of the thermal background, of which the electron is the lightest, are suppressed exponentially by a factor $`\mathrm{exp}(m_e/T)`$, where $`T=2.7`$ K is the temperature of the background . This argument, however, does not hold in the case of the photon-photon interaction according to (1), since the electrons, which have been integrated out, come from vacuum polarization. Therefore it appears to be worthwhile to reconsider the effective photon mass due to the cosmic background radiation. Further interesting quantities following from the photon self energy are the dispersion relations of photons in a thermal photon gas and its dielectric function, related to the index of refraction and the velocity of light.
Using the notation $`P=(p_0,𝐩)`$ and $`p=|𝐩|`$ the retarded photon self energy according to Fig.1 reads
$$\mathrm{\Pi }_{\mu \nu }(P)=\frac{1}{2}\frac{d^4Q}{(2\pi )^4}D^{\rho \sigma }(Q)\mathrm{\Gamma }_{\rho \mu \sigma \nu }(Q,P,Q,P).$$
(3)
where $`\mathrm{\Gamma }_{\rho \mu \sigma \nu }`$ is the effective 4-photon vertex and the factor $`1/2`$ a symmetry factor associated with the tadpole diagram. Adopting the real time formalism in the Keldysh representation the photon propagator in a general covariant gauge with gauge parameter $`\xi `$ reads
$$D^{\rho \sigma }(Q)=\left(g^{\rho \sigma }\xi \frac{Q^\rho Q^\sigma }{Q^2}\right)\frac{1}{2}\left[D_R(Q)+D_A(Q)+D_F(Q)\right].$$
(4)
The retarded (R), advanced (A), and symmetric (F) propagators are given by
$`D_{R,A}(Q)`$ $`=`$ $`{\displaystyle \frac{1}{Q^2\pm i\mathrm{sgn}(q_0)\epsilon }},`$ (5)
$`D_F(Q)`$ $`=`$ $`2\pi i[1+2n_B(Q,x)]\delta (Q^2)`$ (6)
whith the non-equilibrium photon distribution $`n_B`$ depending on the momentum and the space-time coordinate. In equilibrium it reads $`n_B^{\mathrm{eq}}=1/[\mathrm{exp}(|q_0|/T)1]`$. In the following we restrict ourselves to isotropic momentum distributionsThe anisotropic case has been discussed in Ref. for QED and QCD., i.e. $`n_B=n_B(q_0,q,x)`$. Then there are only two independent components of $`\mathrm{\Pi }_{\mu \nu }`$, which depend on $`p_0`$ and $`p`$ . For these components we choose
$`\mathrm{\Pi }_L(p_0,p)`$ $`=`$ $`\mathrm{\Pi }_{00}(P),`$ (7)
$`\mathrm{\Pi }_T(p_0,p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\delta _{ij}{\displaystyle \frac{p_ip_j}{p^2}}\right)\mathrm{\Pi }_{ij}(P).`$ (8)
It should be noted that the longitudinal component is sometimes defined differently .
Since the zero temperature contributions vanish we find
$`\mathrm{\Pi }_L(p_0,p)`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \frac{d^4Q}{(2\pi )^3}n_B(q_0,q,x)\delta (Q^2)\mathrm{\Gamma }_{}^{\rho }{}_{0\rho 0}{}^{}(Q,P,Q,P)},`$ (9)
$`\mathrm{\Pi }_T(p_0,p)`$ $`=`$ $`{\displaystyle \frac{i}{4}}\left(\delta _{ij}{\displaystyle \frac{p_ip_j}{p^2}}\right){\displaystyle \frac{d^4Q}{(2\pi )^3}n_B(q_0,q,x)\delta (Q^2)\mathrm{\Gamma }_{}^{\rho }{}_{i\rho j}{}^{}(Q,P,Q,P)}.`$ (10)
Adopting the expression for the 4-photon vertex given in Ref., where some factors of 2 and the sign are corrected, we obtain
$`\mathrm{\Gamma }_{}^{\rho }{}_{0\rho 0}{}^{}(Q,P,Q,P)|_{Q^2=0}=16i(4a+3b)[(𝐩𝐪)^2p^2q^2],`$ (11)
$`(\delta _{ij}{\displaystyle \frac{p_ip_j}{p^2}})\mathrm{\Gamma }_{}^{\rho }{}_{i\rho j}{}^{}(Q,P,Q,P)|_{Q^2=0}=16i(4a+3b)(p_0^2+p^2)q^2[1+{\displaystyle \frac{(𝐩𝐪)^2}{p^2q^2}}].`$ (12)
In (10) the gauge fixing parameter $`\xi `$ does not appear since $`Q^\rho Q^\sigma \mathrm{\Gamma }_{\rho \mu \sigma \nu }(Q,P,Q,P)=0`$ as can be shown explicitly. Hence the lowest order photon self energy is gauge invariant.
Inserting (12) into (10) we end up with the final result
$`\mathrm{\Pi }_L(p_0,p)`$ $`=`$ $`\gamma p^2,`$ (13)
$`\mathrm{\Pi }_T(p_0,p)`$ $`=`$ $`\gamma (p_0^2+p^2),`$ (14)
where
$$\gamma =\frac{44}{135\pi ^2}\frac{\alpha ^2}{m_e^4}_0^{\mathrm{}}𝑑qq^3n_B(q,x).$$
(15)
In equilibrium (15) reduces to $`\gamma =(44\pi ^2/2025)\alpha ^2(T/m_e)^4`$. Then (14) agrees apart from the sign for $`\mathrm{\Pi }_T`$ with Ref..
Now we want to discuss the physical consequences of our result. First we consider Debye screening in the photon gas. It has been argued that there is no Debye mass due to the photon-photon interaction since the corresponding vertex vanishes if one of the external legs has zero momentum . However, this argument is based on an inconsistent definition for the Debye mass
$$m_D^2=\mathrm{\Pi }_L(p_0=0,p0),$$
(16)
where $`p=|𝐩|`$. This definition can lead to gauge dependent results and is not renormalization-group invariant . Instead of (16) the Debye mass should be determined self consistently from the pole of the longitudinal photon propagator, i.e. from
$$m_D^2\mathrm{\Pi }_L(p_0=0,p^2=m_D^2)=0.$$
(17)
Only if $`\mathrm{\Pi }_L(p_0=0,p)`$ does not depend on $`p`$, the definition (16) agrees with (17). Although this is not the case here, the Debye mass following from (14) and (17) vanishes also. This means that the photon-photon interaction in a photon gas does not lead to a screening of the Coulomb potential. Of course, there is also no static magnetic screening.
As another point we mention that the tadpole self energy of Fig.1 has no imaginary part, i.e. neither real nor virtual photons are damped to this order. In QED damping arises from the box diagram only above the threshold for electron-positron pair production. This effect, however, is not included in the effective theory (1) for low energy photons. Damping will be present in the effective theory at the two-loop level (sunset diagram) corresponding to photon-photon scattering.
Next we discuss the photon dispersion relations in the photon gas, which follow from the pole of the resummed photon propagator. In Coulomb gauge the resummed propagator is given by
$`D_L^1(p_0,p)`$ $`=`$ $`p^2\mathrm{\Pi }_L(p_0,p)=(1+\gamma )p^2,`$ (18)
$`D_T^1(p_0,p)`$ $`=`$ $`p_0^2p^2\mathrm{\Pi }_T(p_0,p)=(1+\gamma )p_0^2(1\gamma )p^2.`$ (19)
Whereas there is no dispersion relation for longitudinal photons, i.e., there are no plasmons in the photon gas, the dispersion relation of the transverse (physical) photons is modified compared to the vacuum. It is given by
$$\omega (p)=\sqrt{\frac{1\gamma }{1+\gamma }}p(1\gamma )p,$$
(20)
i.e., the dispersion is located below the light cone $`\omega <p`$ and the plasma frequency $`\omega _{pl}=\omega (p=0)`$ vanishes. The phase velocity $`v_p=\omega /p1\gamma `$ is identical to the group velocity $`v_g=\omega /p`$ and smaller than the speed of light in the vacuum. The result agrees with , where it has been derived in QED. Note also that the phase as well as the group velocity are independent of the momentum, i.e., there is no dispersion.
Finally we turn to the dielectric tensor. In an isotropic medium there are only two independent components of the dielectric tensor, for which we choose the longitudinal and the transverse dielectric functions, related to the photon self energy via
$`ϵ_L(p_0,p)`$ $`=`$ $`1{\displaystyle \frac{\mathrm{\Pi }_L(p_0,p)}{p^2}}=1+\gamma ,`$ (21)
$`ϵ_T(p_0,p)`$ $`=`$ $`1{\displaystyle \frac{\mathrm{\Pi }_T(p_0,p)}{p_0^2}}=1+\gamma {\displaystyle \frac{p_0^2+p^2}{p_0^2}}.`$ (22)
Since there is no direction preferred for $`𝐩=0`$ , the longitudinal and transverse dielectric functions coincide in this limit, $`ϵ_L(p_0,p=0)=ϵ_T(p_0,p=0)=1+\gamma `$.
The electric permittivity and the magnetic permeability given by
$`ϵ`$ $`=`$ $`ϵ_L=1+\gamma ,`$ (23)
$`{\displaystyle \frac{1}{\mu }}`$ $`=`$ $`1+{\displaystyle \frac{\mathrm{\Pi }_Tp_0^2\mathrm{\Pi }_L/p^2}{p^2}}=1\gamma `$ (24)
are independent of $`p_0`$ and $`p`$. The phase velocity following from , related to the index of refraction $`n`$,
$$v_p=\frac{1}{n}=\frac{1}{\sqrt{\mu ϵ}}1\gamma $$
(25)
agrees with the result found from the dispersion relation (20).
Owing to the cosmic microwave background the velocity of light in the Universe is reduced compared to the vacuum. Actually it increases continuously with time as the temperature drops. Today at a temperature of 2.7 K it is given by (25) with $`\gamma =4.7\times 10^{43}`$. In the early Universe, when the radiation decoupled from matter at a temperature of about 3000 K we had $`\gamma =6.5\times 10^{31}`$. Although this is probably not a measurable effect, the speed of light is not a constant in our Universe.
Summarizing, we have calculated the photon self energy in an isotropic, non-equilibrium photon gas using the real time formalism. For this purpose we considered the effective Lagrangian for photon-photon interaction and calculated the photon self energy to lowest order perturbation theory using an effective 4-photon vertex in momentum space. As physical consequences of this self energy we showed the absence of Debye and static magnetic screening in the photon gas. Also there are no longitudinal collective modes (plasmons). However, the transverse collective modes exhibit a modified dispersion with a vanishing plasma frequency and a smaller slope compared to the vacuum modes. This results in a reduced, dispersion free velocity of light, which increases during the evolution of the Universe as the temperature of the cosmic microwave background drops. We also determined the dielectric tensor, the electric permittivity and the magnetic permeability of the photon gas.
ACKNOWLEDGMENTS
The author would like to thank U. Mosel and A. Rebhan for stimulating and helpful discussions.
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# Quintessence models in Supergravity
## I Introduction
Several cosmological observations appear to suggest that the present Universe is dominated by an unknown form of energy with negative pressure, acting as an effective cosmological constant $`\mathrm{\Lambda }`$. CMB and cluster distribution data, when combined with the measurements of the redshift–luminosity relation using high redshift type Ia supernovae, single out a flat cosmological model with $`\mathrm{\Omega }_{\mathrm{matter}}1/3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }2/3`$ . The contribution to the energy density which has led to the observed acceleration is named “quintessence” and could well be some more general component than a bare cosmological term $`\mathrm{\Lambda }`$ . In particular, a scalar field $`Q`$ rolling down a potential $`V(Q)`$ can act as an effective dynamical cosmological constant with equation of state
$$w_Q=\frac{p_Q}{\rho _Q}=\frac{\dot{Q}^2/2V(Q)}{\dot{Q}^2/2+V(Q)}.$$
(1)
Dynamical models are appealing because they could potentially solve the smallness and coincidence problems which arise when postulating that there is today a cosmological constant energy density contribution, $`\rho _\mathrm{\Lambda }=\mathrm{\Lambda }/8\pi G`$, of the same order of magnitude as the critical energy density $`\rho _c10^{47}\mathrm{GeV}^4`$. The advantage of dealing with a dynamical scalar field is that there exist attractor solutions in a cosmological setting for some classes of potentials . If this is the case, the scalar field will join the attractor solution before the present epoch for a very wide range of initial conditions, thus relieving the fine–tuning issue.
For example, the potential $`V(Q)=M^{4+\alpha }Q^\alpha `$, with $`\alpha >0`$, leads to a ‘tracking’ attractor solution for which the scalar energy density scales as $`\rho _Q/\rho _Ba^{3(w_Bw_Q)}`$, where $`\rho _B`$ and $`w_B`$ are the background energy density and equation of state (the latter being $`=0,1/3`$ during matter and radiation domination respectively), and $`a`$ is the scale factor of the Universe. The equation of state of the attractor is given by $`w_Q=(w_B\alpha 2)/(\alpha +2)`$ and is always negative during matter domination . If $`Q`$ has already reached the attractor solution today, it must be of order of the Planck mass $`M_{\mathrm{Pl}}`$. This follows from the fact that for the tracking solutions $`V^{\prime \prime }(Q)H^2`$ and moreover we require $`V\rho _c`$ today. The scale $`M`$ in the potential is fixed by requiring that $`V(QM_{\mathrm{Pl}})\rho _c`$ and depends on the exponent $`\alpha `$. This class of potentials is then a suitable candidate for quintessence and has been shown to be compatible with observations .
If instead an exponential potential is considered, $`V(Q)=M^4\mathrm{exp}(\lambda Q)`$, two attractors are found , depending on the value of $`\lambda `$. If $`\lambda >3(w_B+1)`$, the late time attractor is a ‘scaling’ one, characterized by $`w_Q=w_B`$ and $`\mathrm{\Omega }_Q=3(w_B+1)/\lambda ^2`$. Note that in this case the ratio of the scalar to critical energy density is independent of the mass scale $`M`$ in the potential. Unfortunately the equation of state is non negative and so cannot drive the Universe to an accelerated expansion. If $`\lambda <3(w_B+1)`$, instead, then the late time attractor is the scalar field dominated solution with $`\mathrm{\Omega }_Q=1`$ and $`w_Q=1+\lambda ^2/3`$. Now the equation of state is negative, but clearly cannot be the solution for all times, as there is the tight constraint on the allowed contribution of the scalar field at nucleosynthesis of $`\mathrm{\Omega }_Q(1\mathrm{M}\mathrm{e}\mathrm{V})<0.13`$. Besides, one must allow time for structure formation before the Universe starts accelerating. A scenario like this requires that the scalar field starts with roughly the same energy density as today, and so provides no improvement on solving the fine tuning problem. The exponential potentials, then, cannot be used individually for modeling quintessence, but interesting modifications have been proposed in order to fit the data .
### A Particle physics models
While dealing with what is thought to be one of the fundamental components of the present universe, it is mandatory to ask if we can find any firmer motivation for introducing a cosmological scalar field. We actually know that scalar fields arise quite naturally in unified theories and are a necessity in supersymmetric theories where they play the role of ‘superpartners’ of the Standard model fermion fields. Scalar fields are also required by another crucial cosmological mechanism, inflation. Indeed, it has been recently shown that the same unique field could possibly dominate both the inflationary and quintessence phases of our universe. The issue of finding deeper roots for the quintessence scalar into ‘realistic’ particle physics models is then a pressing one.
It is well known that inverse power law scalar potentials arise in supersymmetry (SUSY) gauge theories due to non–perturbative effects . They have been studied as candidates for quintessence and shown to provide a viable phenomenology both in the one–scalar and multi–scalar cases. Exponential potentials are also common in high energy physics, since they arise as a consequence of Kaluza–Klein type compactifications of string theory. The attempt of building a particle physics motivated model for quintessence is then a well-posed problem, although many points remain to be clarified. For example, exponential-like models as those proposed in still await a firmer high energy motivation. At the same time, inverse power law models based on SUSY QCD need to be revisited taking into account the possibility that the scalar potential receive corrections originating from quantum or supergravity effects.
This last issue was recently addressed by Brax and Martin . In particular they showed that inverse power law models which lead to quintessence are stable against quantum corrections, both in the supersymmetric and non–supersymmetric cases. Furthermore, discussing curvature and Kählerian corrections in the globally supersymmetric case, they found that while curvature corrections are negligible, Kählerian ones might be important in the early evolution of $`Q`$.
The most interesting question, though, concerns supergravity (SUGRA) corrections . Since today the vacuum expectation value (vev) of the scalar field is of order Planck mass, these corrections are unavoidable and may lead to problems for these power law quintessence models. In particular, if the scalar potential is dramatically altered for field values corresponding to the most recent epoch, then we might have a completely different phenomenology from the attractor case just described. In ref. , a possible way of avoiding the dangerous SUGRA corrections was outlined, which involved the restrictive condition $`W=0`$. The aim of this paper is to further develop this line of attack by proposing an alternative way of tackling SUGRA corrections. As an example, we will study the cosmology of a specific SUGRA–inspired quintessence model and show that it is compatible with observations for a wide range of initial conditions.
## II Supergravity models
### A Global supersymmetry
Globally supersymmetric QCD theories with $`N_c`$ colors and $`N_f<N_c`$ flavors may give an explicit realization of a model of quintessence with an inverse power law scalar potential . The matter content of the theory is given by the chiral superfields $`Q_i`$ and $`\overline{Q}_i`$ ($`i=1\mathrm{}N_f`$) transforming according to the $`N_c`$ and $`\overline{N}_c`$ representations of $`SU(N_c)`$, respectively.
Supersymmetry and anomaly-free global symmetries constrain the superpotential to take the form
$$W=(N_cN_f)\left(\frac{\mathrm{\Lambda }^{(3N_cN_f)}}{\mathrm{det}T}\right)^{\frac{1}{N_cN_f}}$$
(2)
where the gauge-invariant matrix superfield $`T_{ij}=Q_i\overline{Q}_j`$ appears. $`\mathrm{\Lambda }`$ is the only mass scale of the theory. It is the supersymmetric analogue of $`\mathrm{\Lambda }_{QCD}`$, the renormalization group invariant scale at which the gauge coupling of $`SU(N_c)`$ becomes non-perturbative.
If the Kähler metric is flat, the scalar potential is given by<sup>*</sup><sup>*</sup>*In the following, the same symbols will be used for the superfields $`Q_i`$, $`\overline{Q}_i`$, and their scalar components.
$`V(Q_i,\overline{Q}_i)`$ $`=`$ $`V_F+V_D`$ (3)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N_f}{}}}\left(|F_i|^2+|F_{\overline{i}}|^2\right)+{\displaystyle \frac{1}{2}}D^aD^a`$ (4)
where $`F_i=W/Q_i`$, $`F_{\overline{i}}=W/\overline{Q}_i`$, and
$$D^a=Q_i^{}t^aQ_i\overline{Q}_it^a\overline{Q}_i^{},$$
(5)
with the $`t^a`$’s being the generators of the gauge group. Of interest to us are the dynamics of the field expectation values which take place along directions in field space in which the above D-term vanishes, i.e. the perturbatively flat directions $`Q_{i\alpha }=\overline{Q}_{i\alpha }^{}`$, where $`\alpha =1\mathrm{}N_c`$ is the gauge index. At the non-perturbative level these directions acquire a non vanishing potential from the F-terms in (4), which are zero to any order in perturbation theory. Gauge and flavor rotations can be used to diagonalize the $`Q_{i\alpha }`$ and put them in the form
$`Q_{i\alpha }=\overline{Q}_{i\alpha }^{}=\{\begin{array}{cc}Q_i\delta _{i\alpha }1\alpha N_f\hfill & \\ 0N_f\alpha N_c\hfill & \end{array}.`$
Along these directions, if the expectation values of all the $`N_f`$ scalars are taken to be equal, $`Q_i=Q,i=1,\mathrm{},N_f`$, the cosmological evolution of the scalar vev $`Q`$ is given by
$`\ddot{Q}`$ $`=`$ $`3H\dot{Q}+\beta {\displaystyle \frac{\mathrm{\Lambda }^{4+2\beta }}{Q^{2\beta +1}}},`$ (6)
$`\beta `$ $`=`$ $`{\displaystyle \frac{N_c+N_f}{N_cN_f}},`$ (7)
thus reproducing exactly the case of a single scalar field $`Q`$ in a potential $`V=\frac{\mathrm{\Lambda }^{4+2\beta }}{2}Q^{2\beta }`$. The ‘tracker’ solution is characterized by an equation of state
$$w_Q=\frac{N_c+N_f}{2N_c}w_B\frac{N_cN_f}{2N_c}$$
(8)
as a function of the parameters of the theory, with the scalar energy density growing with respect to the matter one as
$$\rho _Q/\rho _\mathrm{m}=a^{3(N_cN_f)/2N_c}.$$
(9)
Requiring that the scalar $`Q`$ both has reached the tracker today and is starting to dominate the energy density, we obtain that at the present epoch $`QM_{\mathrm{Pl}}`$ and that the mass scale in the potential must satisfy $`\mathrm{\Lambda }^{4+\beta }\rho _cM_{\mathrm{Pl}}^\beta `$, introducing a degree of fine tuning in the problem.
### B Supergravity corrections
The above discussion is valid as far as the global SUSY limit can be taken as a good approximation. This is correct for most of the cosmological evolution of the scalar $`Q`$. However, for the present epoch we are not allowed to neglect SUGRA corrections, since we enter the regime for which they might become important. In this case, the F-term in the scalar potential in general is
$`V(Q)`$ $`=`$ $`F^2e^{\kappa ^2K}3\kappa ^2|W|^2`$ (10)
$`=`$ $`e^{\kappa ^2K}[(W_i+\kappa ^2WK_i)K^{j^{}i}(W_j+\kappa ^2WK_j)^{}`$ (11)
$``$ $`3\kappa ^2|W|^2]`$ (12)
where the subscript $`i`$ indicates the derivative with respect to the $`i`$-th field, and $`\kappa ^2=8\pi G=8\pi M_{\mathrm{Pl}}^2`$.
Brax and Martin discuss the case of a theory with superpotential $`W=\mathrm{\Lambda }^{3+\alpha }Q^\alpha `$ and a flat Kähler potential, $`K=QQ^{}`$. It is straightforward to compute the resulting scalar potential:
$$V(Q)=e^{\frac{\kappa ^2}{2}Q^2}\frac{\mathrm{\Lambda }^{4+\beta }}{Q^\beta }\left(\frac{(\beta 2)^2}{4}(\beta +1)\frac{\kappa ^2}{2}Q^2+\frac{\kappa ^4}{4}Q^4\right)$$
(13)
where $`\beta =2\alpha +2`$. The main effect of the supergravity corrections is that the scalar potential can now become negative due to the presence of the second term. This is a serious drawback for the model, which becomes ill defined for $`QM_{\mathrm{Pl}}`$. They go on to propose a possible solution by imposing the condition that the expectation value of the superpotential vanishes, $`W=0`$. We then see from equation (12) that the negative contribution to the scalar potential disappears, and it takes the form
$$V(Q)=\frac{\mathrm{\Lambda }^{4+\alpha }}{Q^\alpha }e^{\frac{\kappa ^2}{2}Q^2}.$$
(14)
The condition $`W=0`$ is possible to realize, for example, in a model in which we allow matter fields to be present in addition to the quintessence scalar field . Then, if at least one of the gradients of the superpotential with respect to the matter fields is non–zero, the scalar potential will always be positive.
This is not the only possibility, though. There are two obvious problems with the potential (13): one, as already stated, is the negative term in the general expression (12) and the second is the choice of the Kähler metric which makes this term grow with the field’s vev, relative to the other terms in the potential. As mentioned above, setting $`W=0`$ is a tight restriction, so we will address the issue by relaxing this constraint but allow for more general forms of the Kähler metric.
Such an approach was recently adopted in as a method of obtaining a minimum for the dilaton field in string theory. It had the advantage of relying on only one gaugino condensate and provided an alternative approach to the phenomenology associated with ‘racetrack’ models . In this scenario, the Kähler potential acquires string inspired non-perturbative corrections. A further nice feature of these models is that it is possible to have a minimum with zero or small positive cosmological constant , and moreover it is possible to establish that the dilaton can be stabilized in such a minimum in a cosmological setting .
In general, for different choices of the Kähler metric, the negative term in (12) does not always lead to the disaster of a negative minimum in the scalar potential. For a general Kähler, we do not know a priori the shape of the potential or the location of the minimum. In fact, in what follows we will show through explicit examples that the scalar potential might always remain positive through a suitable choice of the Kähler metric. Moreover, with this approach there is no need to introduce additional fields in the model.
Let us now go on to study SUGRA corrections to inverse power law quintessence models by choosing more general Kähler potentials. Consider, for example, a theory with superpotential $`W=\mathrm{\Lambda }^{3+\alpha }\stackrel{~}{Q}^\alpha `$ and a Kähler $`K=\mathrm{ln}(\kappa \stackrel{~}{Q}+\kappa \stackrel{~}{Q}^{})/\kappa ^2`$, the type of term which is present at tree level for both the dilaton and moduli fields in string theory. In this case, the resulting scalar potential, expressed in terms of the canonically normalized field $`Q=(\mathrm{ln}\kappa \stackrel{~}{Q})/\sqrt{2}\kappa `$, is
$$V(Q)=M^4e^{\sqrt{2}\beta \kappa Q}$$
(15)
where $`M^4=\mathrm{\Lambda }^{5+\beta }\kappa ^{1+\beta }(\beta ^23)/2`$ with $`\beta =2\alpha +1>\sqrt{3}`$ to allow for positivity of the potential. This corresponds to the ‘scaling’ solution discussed in the introduction and so cannot lead to a negative equation of state for the field in a matter dominated regime.
Another example follows as a natural extension of the one just described and leads to potentials with more than one exponential. For a superpotential of the form $`W=\mathrm{\Lambda }^{3+\alpha }\stackrel{~}{Q}^\alpha +\mathrm{\Lambda }^{3+\beta }\stackrel{~}{Q}^\beta `$ and the same Kähler metric as above, then in terms of the same canonically normalized field $`Q=(\mathrm{ln}\kappa \stackrel{~}{Q})/\sqrt{2}\kappa `$, the scalar potential becomes
$`V(Q)`$ $`=`$ $`(M_1)^4e^{\sqrt{2}\gamma \kappa Q}+(M_2)^4e^{\sqrt{2}\delta \kappa Q}`$ (16)
$`+`$ $`(M_3)^4e^{\frac{\gamma +\delta }{\sqrt{2}}\kappa Q},`$ (17)
where $`\gamma =2\alpha +1`$, $`\delta =2\beta +1`$ and
$`(M_1)^4`$ $`=`$ $`\mathrm{\Lambda }^{5+\gamma }\kappa ^{1+\gamma }(\gamma ^23)/2,`$ (18)
$`(M_2)^4`$ $`=`$ $`\mathrm{\Lambda }^{5+\delta }\kappa ^{1+\delta }(\delta ^23)/2,`$ (19)
$`(M_3)^4`$ $`=`$ $`\mathrm{\Lambda }^{5+\frac{\gamma +\delta }{2}}\kappa ^{1+\frac{\gamma +\delta }{2}}(\gamma \delta 3).`$ (20)
At first sight this appears to be of the form required in in that it involves multiple exponential terms. However, closer analysis indicates that the slopes of the exponentials are not adequate to satisfy the bounds arising from nucleosynthesis constraints, whilst also providing a positive cosmological constant type contribution today.
As we mentioned earlier, it is possible to have more general Kähler potentials, and with that in mind we now consider the original model $`W=\mathrm{\Lambda }^{3+\alpha }\stackrel{~}{Q}^\alpha `$, but with a Kähler potential which depends on a parameter $`\gamma `$
$$K=\frac{1}{\kappa ^2}\left[\mathrm{ln}(\kappa \stackrel{~}{Q}+\kappa \stackrel{~}{Q}^{})\right]^\gamma ,\gamma >1.$$
(21)
In this case, the second derivative of the Kähler is
$`K_{\stackrel{~}{Q}\stackrel{~}{Q}^{}}`$ $`=`$ $`{\displaystyle \frac{\gamma (\gamma 1)}{\kappa ^2(\stackrel{~}{Q}+\stackrel{~}{Q}^{})^2}}[\mathrm{ln}(\kappa \stackrel{~}{Q}+\kappa \stackrel{~}{Q}^{})]^{\gamma 2}`$ (22)
$`\times `$ $`\left(1{\displaystyle \frac{\mathrm{ln}(\kappa \stackrel{~}{Q}+\kappa \stackrel{~}{Q}^{})}{\gamma 1}}\right)`$ (23)
and the canonically normalized field $`Q`$ can be obtained as a function of $`\stackrel{~}{Q}`$ by integrating the following expression
$$dQ=\sqrt{2K_{\stackrel{~}{Q}\stackrel{~}{Q}^{}}}d\stackrel{~}{Q}.$$
(24)
In order to avoid the singularity at $`\stackrel{~}{Q}+\stackrel{~}{Q}^{}=1/\kappa ^2`$, when $`\mathrm{ln}(\kappa \stackrel{~}{Q}+\kappa \stackrel{~}{Q}^{})`$ passes through zero (see equation (23)), the only possible choice is $`\gamma =2`$. We then obtain:
$$K_{\stackrel{~}{Q}\stackrel{~}{Q}^{}}=\frac{2[1\mathrm{ln}(\kappa \stackrel{~}{Q}+\kappa \stackrel{~}{Q}^{})]}{\kappa ^2(\stackrel{~}{Q}+\stackrel{~}{Q}^{})^2}$$
(25)
and as a consequence
$$Q=\frac{2}{3\kappa }[1\mathrm{ln}(2\kappa \stackrel{~}{Q})]^{3/2}.$$
(26)
Implying that the theory is well defined for
$`\mathrm{}<\mathrm{ln}(2\kappa \stackrel{~}{Q})<1`$
which corresponds to $`0<\stackrel{~}{Q}<e/2\kappa `$.
The scalar potential in terms of the canonically normalized field $`Q`$ reads
$`V`$ $`=`$ $`M^4\left[2x^2+(4\alpha 7)x+2(\alpha 1)^2\right]{\displaystyle \frac{1}{x}}`$ (27)
$`\times \mathrm{exp}[(1x)^22\alpha (1x)]`$ (28)
where for notational convenience we have defined the quantities
$$x\left(\frac{3}{2}\kappa Q\right)^{2/3}=1\mathrm{ln}(2\kappa \stackrel{~}{Q}),$$
(29)
$$M^4=\mathrm{\Lambda }^{6+2\alpha }\kappa ^{2+2\alpha }2^{2\alpha }.$$
(30)
Note that the canonically normalized field $`Q`$ has a range $`\mathrm{}<Q<0`$.
We can see from equations (28)–(29) that the scalar potential behaves like an exponential for $`|Q|1`$ and like an inverse power law for $`|Q|1`$, and thus develops a minimum at a finite value $`Q_\mathrm{m}`$. Note that the potential is always positive for any $`\alpha >1.25`$. Thus, we have found that in this case the supergravity corrections induce a finite minimum in the potential but do not spoil its positivity. Note also that the field’s value in the minimum is exactly in the region where we expect the supergravity corrections to become relevant. For example, with $`\alpha =5`$ we obtain $`Q_\mathrm{m}0.02`$ (in $`8\pi G=1`$ units), which corresponds to $`\stackrel{~}{Q}1.2`$. Imposing that the minimum of the potential equals the critical energy density today, we can also estimate the mass scale $`\mathrm{\Lambda }`$, depending on $`\alpha `$. In the case $`\alpha =5`$ we have that $`V(Q=Q_\mathrm{m})10^{47}\mathrm{GeV}^4`$ which corresponds to $`\mathrm{\Lambda }610^{10}\mathrm{GeV}`$.
### C Supersymmetry breaking
If supersymmetry is to be realized in nature, it must be broken at a mass scale $`M_S`$ such that $`M_S^2F(10^{10}\mathrm{GeV})^2`$ or $`M_S^2F(10^4\mathrm{GeV})^2`$ (for gravity and gauge mediated cases respectively), in order to lift the supersymmetric scalar particle masses above $`10^2\mathrm{GeV}`$.
This then requires the superpotential in (12) to be $`WF\kappa ^1m_{3/2}\kappa ^2`$ in order to cancel the F-term contribution and consequently to give a negligible total vacuum energy density ($`m_{3/2}`$ is the gravitino mass).
From the discussion in the last sectionWe thank K. Choi and D. Lyth for discussions on this point., it is clear that the dynamical cosmological constant provided by the quintessence potential cannot be the dominant source of SUSY breaking in the Universe as $`W\mathrm{\Lambda }^{3+\alpha }\kappa ^\alpha (10^3\mathrm{eV})^2\kappa ^1F\kappa ^1`$. Therefore, we need some additional source of SUSY breaking. If we consider now the superpotential,
$$W=\mathrm{\Lambda }^{3+\alpha }Q^\alpha +m_{3/2}\kappa ^2,$$
(31)
then one gains a correction to the scalar potential in (28) of,
$$\delta V\mathrm{\Lambda }^{3+\alpha }m_{3/2}\kappa ^\alpha +m_{3/2}^2\kappa ^2,$$
(32)
for $`Q\mathrm{M}_{\mathrm{Pl}}`$ today .
The first term can in principle be controlled for sufficiently large $`\alpha `$, however, the constant second term unavoidably leads to a disruption of the quintessence potential. This is a very serious problem of all supergravity models in quintessence In order to avoid the SUGRA corrections problem, Choi proposes a Goldstone-type quintessence model inspired in heterotic M-Theory ..
The situation gets worse, since, as pointed out in , it appears that even if we imagine that the amount of SUSY breaking that we observe in the universe today comes from a hidden sector other than the quintessence one, there will still be gravitational couplings between the two sectors that rekindle the original problem.
However, some recent proposals point in a slightly different direction for solving this problem. The basic idea is that the traditional approach to SUSY breaking might not be the best way to explain the world we live in. For example, the mass difference between the superpartners could arise in a 4D world with unbroken SUSY through some higher dimensional effects . If this is the case, we would not need to break supersymmetry, and the quintessence potential would be preserved. Another possibility is that the relation between the SUSY breaking scale $`M_S`$ and the cosmological constant $`F^2`$ is not what is usually considered. If $`M_S=\kappa ^1[F^2\kappa ^4]^\beta `$ and $`\beta =1/8`$, instead of the usual $`1/4`$, then the observed cosmological constant would provide just the right amount of SUSY breaking. In this case we wouldn’t have any dangerous F-term of order $`\kappa ^2M_S^2`$ which spoils the quintessence potential.
## III Cosmology of the model
We can now study in further detail the model we presented in the previous section. It can be easily checked that for $`Q0`$ the potential (28) goes as $`VQ^{2/3}`$, while for $`Q\mathrm{}`$ it is $`VQ^{2/3}e^{Q^{4/3}}`$. This behavior is independent of the parameter $`\alpha `$ in our theory, which plays no role in the asymptotic form of the potential. For all the values of $`\alpha `$ we then find the same qualitative behavior for the scalar $`Q`$. For a very wide range of the initial conditions, indeed, we obtain scalar field dominance and negative equation of state at the present epoch.
In Figure 1 and Figure 2 we plot the evolution of the scalar energy density and equation of state for $`\alpha =5`$.
If the scalar field is rolling down the potential towards the minimum from the side $`Q_{\mathrm{in}}/Q_\mathrm{m}1`$, then it will exhibit a ‘tracking’ behavior as in the general case $`VQ^\beta `$, with $`\beta =2/3`$. This is characterized by an equation of state $`w_Q=(w_B\beta 2)/(\beta +2)=2/3,3/4`$ during radiation and matter domination respectively. When the field $`Q`$ approaches the minimum, it will depart from the attractor solution and enter a regime of damped oscillations. The non–zero vacuum energy will rapidly take over the kinetic term and the equation of state be driven towards $`w_Q=1`$ (see Figure 2).
If, instead, $`Q`$ rolls down from the side $`Q_{\mathrm{in}}/Q_\mathrm{m}1`$ the exponential in (28) is then important. As it will be shown below, this leads to an attractor with $`w_Qw_B`$ up to a scale factor dependent logarithmic correction. Eventually, the field will settle down in the minimum with $`w_Q=1`$, after a stage of small oscillations about the minimum as before. The earlier requirement made on $`\alpha `$, namely $`\alpha >1.25`$ is also sufficient to respect the nucleosynthesis bound of $`\mathrm{\Omega }_Q(1\mathrm{M}\mathrm{e}\mathrm{V})<0.13`$.
In the second reference of a general formula for the scalar equation of state in the case of ‘tracking’ potentials was derived:
$$w_Q\frac{w_B2(\mathrm{\Gamma }1)}{1+2(\mathrm{\Gamma }1)},$$
(33)
where $`\mathrm{\Gamma }`$ is defined as $`\mathrm{\Gamma }V^{\prime \prime }V/(V^{})^2`$, with a prime denoting the derivative with respect to the scalar $`Q`$.
In deriving eq. (33), it is assumed that $`\mathrm{\Gamma }`$ is nearly constant. A tracker solution with $`w_Q<w_B`$ is reached if $`\mathrm{\Gamma }>1`$. For a potential of the form
$$V=M^4e^{\alpha (\kappa Q)^\beta },$$
(34)
we obtain
$$\mathrm{\Gamma }=1+\frac{\beta 1}{\alpha \beta }\frac{1}{(\kappa Q)^\beta },$$
(35)
which is larger than 1 and slowly varying as required. Actually the dependence of $`\mathrm{\Gamma }`$ on the scale factor $`a(t)`$ through $`Q(a)`$ will give the logarithmic variation of the equation of state mentioned earlier.
By substituting the expression (35) for $`\mathrm{\Gamma }`$ in eq. (33), we obtain
$$w_B=w_Q+(w_Q+1)\frac{2(\beta 1)}{\beta }\frac{1}{(\kappa Q)^\beta },$$
(36)
with
$$(\kappa Q)^\beta =\mathrm{ln}\left(\frac{V_0}{M^4}\right)3(w_Q+1)\mathrm{ln}\left(\frac{a}{a_0}\right),$$
(37)
where $`V_0`$ is the value of the scalar field potential at some time $`a=a_0`$, when the attractor has already been reached, and we have assumed the following dependence of the potential on the scale factor
$$V(a)=V_0\left(\frac{a}{a_0}\right)^{3(w_Q+1)}$$
(38)
which is exact if $`w_Q`$ is constant and can be taken as a good approximation when $`w_Q`$ is slowly varying.
Note that, when $`\beta =1`$, then $`V`$ reduces to a simple exponential and from (36) we recover the ‘scaling’ behaviour with $`w_Q=w_B`$.
## IV Conclusions
In this paper we have studied supergravity corrections to quintessence models with a superpotential $`W=\mathrm{\Lambda }^{\alpha +3}Q^\alpha `$. The motivation for this is simple. Scalar field models of quintessence typically require that the expectation value of the field today is of order the Planck mass, if we want it to explain the observed acceleration of the Universe. This suggests that we should be considering models in the context of supergravity.
We have proposed a new line of attack to a serious problem that these models share, i.e. the negativity of the minimum in the resulting scalar potential. Allowing for nonlinear modifications to the Kähler potential this problem can effectively be cured. Such modifications are, as far as we are aware, perfectly acceptable and lead to some interesting features.
In particular we discussed in detail a model with Kähler potential $`K=[\mathrm{ln}(\kappa Q+\kappa Q^{})]^2/\kappa ^2`$ and superpotential $`W=\mathrm{\Lambda }^{\alpha +3}Q^\alpha `$ . In this case, the minimum of the resulting scalar potential is always positive for any $`\alpha >1.25`$ and is located close to the point where the SUGRA corrections start to be important, i.e. at $`QM_{\mathrm{Pl}}`$. We found that the resulting scalar potential yields two possible attractor solutions, both of which lead the scalar field towards the minimum without dominating the Universe dynamics before today. After reaching the minimum, the scalar field will mimic a cosmological constant with $`w_Q1`$.
We also noted that the dynamical cosmological constant described above cannot be the dominant source of SUSY breaking in the universe, since the F-term is not large enough. This means that we need some additional source of SUSY breaking, which in turn implies we must ensure that the SUSY breaking sector does not interact with the quintessence sector, in order not to spoil the behaviour described above. Unfortunately, this particular problem still remains to be resolved , without invoking some extra symmetry principle in the action . A possible, but still speculative resolution has been proposed by Witten , who has suggested that we may actually live in a vacuum with unbroken supersymmetry. As an alternative, Banks has proposed amechanism through which the measured cosmological constant might well be able to provide the observed amount of supersymmetry breaking.
In closing we note that an effective equation of state $`w_Q1`$ for the present time, as the one we find, is favored by the available data . Unfortunately this makes it even harder to observationally distinguish the scalar field contribution from the pure cosmological constant case.
###### Acknowledgements.
We would like to thank Bobby Acharya, David Bailin, Tiago Barreiro, Philippe Brax, Beatriz de Carlos, Kiwoon Choi, Sacha Davidson, David Lyth, Cindy Ng, Massimo Pietroni and Fernando Quevedo for useful comments and discussions. EJC is supported by PPARC, NJN by the Fundação para a Ciência e a Tecnologia (Portugal) and FR by the Ministero dell’Università e della Ricerca Scientifica (Italy).
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# References
HZPP-0004
May. 5, 2000
Thermal Equilibrium and Non-uniform Longitudinal Flow
in Relativistic Heavy Ion Collisions<sup>1</sup><sup>1</sup>1Work supported in part by the NSFC under project 19775018.
FENG Shengqin, LIU Feng and LIU Lianshou
Institute of Particle Physics, Huazhong Normal University, Wuhan, 430079, China
ABSTRACT
A model with non-uniform flow in the longitudinal direction is proposed for the relativistic heavy-ion collisions and compared with the 14.6 A GeV/$`c`$ Si-Al and 10.8 A GeV/$`c`$ Au-Au collision data. The stronger influence of transparency on the distribution of heavier produced particles and the larger stopping in the heavier collision system are accounted for by using a new geometrical parameterization picture. The central dips in the proton and deuteron rapidity distributions for Si-Al collision are reproduced.
PACS number(s): 25.75.-q
I. Introduction
The experimental finding that colliding nuclei are not transparent but undergo a violent reaction in central collisions represents one of the major motivations for the study of ultra-relativistic heavy ion collisions at the CERN/SPS, BNL/AGS and also at the future BNL/RHIC and CERN/LHC. Of central importance is the ability of understanding to what extent the nuclear matter has been compressed and heated.
The use of thermal models to interpret data on particle distribution from nuclear collisions is motivated by the hope and expectation that in collision between sufficiently large nuclei at sufficiently high energies a state of exicted nuclear matter close to local thermodynamic equilibrium can be formed, allowing us to study the thermodynamics of QCD and the possible phase transition from a hadronic gas to a QGP. Since we do not yet reliably know how big the collision system and how large the beam energy have to be for this to occur (if at all), thermal models should be considered as an importent phenomenological tool to test for such a behavior.
The study of collective flow in high energy nuclear collisions has attracted increasing attentions from both experimental<sup></sup> and theoretical<sup></sup> point of view. The rich physics of longitudinal and transverse flows is due to their sensitivity to the system evolution at early time. The expansion and cooling of the heated and highly compressed matter could lead to a considerable collectivity in the final state. Due to the high pressure, particles might be boosted in the transverse and longitudinal directions. The collective expansion of the system created during a heavy-ion collision implies space-momentum correlation in particle distributions at freeze-out.
The experimental data on the rapidity distributions of produced particles in 14.6 A GeV/$`c`$ Si-Al collisions has been ultilized to study the collective expansion using a cylindrical-symmetric longitudinal flow model (CSFM)<sup></sup>. In this model, the distribution of final-state particles comes from the superposition of a number of fire-balls, distributed uniformly within the kinematical limit along the longitudinal (rapidity $`y`$) axis. The model-results fit well the experimental distribution of pion, but are too narrow in the case of heavier particles, such as proton and deuteron. In particular, the central dip, which can be clearly seen in the distribution of proton, is not reproduced.
More recently, E877 Collaboration<sup></sup> has published their data for 10.8 A GeV/$`c`$ Au-Au collisions, which provides a good chance to compare the stopping power in collision systems of different sizes. A possible central peak of the rapidity distribution of proton at around mid-rapidity, which was obtained through extrapolating the experimental data to mid-rapidity using RQMD model<sup></sup>, has been taken as an evidence for the increasing of stopping power, but the reliability of this extrapolation is model-dependent.
Stopping and transparency describe the same physical aspect in relativistic heavy ion collision from two opposite sides. In order to settle a reliable model for relativistic heavy ion collision, this aspect must be taken into account carefully. The assumption of uniformly distributed fire-balls in the CSFM model is a crude one. It does not account for the memory of the fire-balls on the motion of the incident nuclei. In the present paper we propose a non-uniform longitudinal flow model (NUFM) using a new geometrical parameterization picture<sup></sup> to describe the nuclear stopping/transparency in a more proper way. The stronger influence of transparency on the distribution of heavier produced particles (proton and deuteron) as well as the larger stopping in the heavier collision system (Au-Au) are described by using this picture. The central dips in the proton and deuteron rapidity distributions for Si-Al collisions are reproduced, and at the same time the central peak in the pion rapidity distributions is maintained.
In section II the non-uniform longitudinal flow model (NUFM) with a geometrical parameterization picture is formulated. The results of the model are given and compared with the experimental data in section III. A short summary and conclusion is given in section IV. In order to avoid the complexity associated with the production of strange particles and concentrate on the expansion of the system, we will discuss in this paper only normal non-strange particles —— pions, protons and deuterons.
II. Non-uniform longitudinal flow
Firstly, let us briefly recall the fireball scenario of relativistic heavy ion collisions.
Since the temperature at freeze-out exceeds 100 MeV, the Boltzmann approximation is used. Transformed into rapidity $`y`$ and transverse momentum $`p_t`$ this implies<sup></sup>:
$$E\frac{\mathrm{d}^3n}{\mathrm{d}^3p}E\mathrm{e}^{(E/T)}=m_t\mathrm{cosh}y\mathrm{e}^{(m_t\mathrm{cosh}y/T)}$$
(1)
Here $`m_t=\sqrt{m^2+p_t^2}`$ is the transverse mass, $`m`$ is the mass of the produced particles at freeze-out.
The rapidity is defined as $`y=\mathrm{tanh}^1(p_l/E)`$, where $`p_l`$ is the longitudinal momentum of the produced particle. Substituting into Eqn.(1) and integrating over $`m_t`$, we get the rapidity distribution of the isotropic thermal source,
$$\frac{\mathrm{d}n_{\mathrm{iso}}}{\mathrm{d}y}\frac{m^2T}{(2\pi )^2}(1+2\xi _0+2\xi _0^2)\mathrm{e}^{(1/\xi _0)}.$$
(2)
where $`\xi _0=T/m\mathrm{cosh}y`$.
Eqn’s.1 and 2 give the isotropic thermal distribution. As mentioned in Ref., for pions with $`mT`$, this distribution is close to that of massless particles, i.e. proportional to $`1/\mathrm{cosh}^2y`$. For heavier particles isotropic emission implies a strongly narrow distribution, in contradiction to the experimental findings, see e.g. the dashed lines in Fig.6.
The measured momentum distribution of the final-state particles is certainly aniso-tropic. It is privileged in the direction of the incident nuclei. This is because the produced hadrons still carry their parent’s kinematic information, making the longitudinal direction more populated than the transverse ones. The simplest way<sup></sup> to account for this anisotropy is to add up the contributions from a set of fire-balls with centers located uniformly in the rapidity region \[$`y_0,y_0`$\], as sketched schematically in Fig.1. The corresponding rapidity distribution is obtained through changing the $`\xi _0`$ in Eqn.(2) into $`\xi =T/m\mathrm{cosh}(yy^{})`$ and integrating over $`y^{}`$ from $`y_0`$ to $`y_0`$:
$$\frac{\mathrm{d}n_{_{\mathrm{CSFM}}}}{\mathrm{d}y}_{y_0}^{y_0}dy^{}\frac{m^2T}{(2\pi )^2}(1+2\xi +2\xi ^2)\mathrm{e}^{(1/\xi )},$$
(3)
where $`\xi =T/m\mathrm{cosh}(yy^{})`$. Equivalently, we can also use the angular variable $`\mathrm{\Theta }`$ defined by $`\mathrm{\Theta }=2\mathrm{tan}^1\mathrm{exp}(y^{})`$, and change the integration variable in Eqn.(3) to $`\mathrm{\Theta }`$,
$`{\displaystyle \frac{\mathrm{d}n_{_{\mathrm{CSFM}}}}{\mathrm{d}y}}`$ $`=`$ $`eKm^2T{\displaystyle _{\mathrm{\Theta }_{\mathrm{min}}}^{\mathrm{\Theta }_{\mathrm{max}}}}{\displaystyle \frac{\mathrm{d}\mathrm{\Theta }}{\mathrm{sin}\mathrm{\Theta }}}\left(1+{\displaystyle \frac{2T}{m\mathrm{cosh}(yy^{})}}+{\displaystyle \frac{2T^2}{m^2\mathrm{cosh}^2(yy^{})}}\right)`$ (4)
$`\times \mathrm{exp}(m\mathrm{cosh}(yy^{})/T),`$
cf. the solid circle and lines in Fig.2.
This simple approach fits the rapidity distribution of pions well but fails to reproduce the central dip in heavier produced particles, which is clearly seen in the experimental distributions of protons and deuterons.
Note that in this CSFM model the distribution of fire-balls in the longitudinal direction of phase space is uniform which does not account for the interaction of the incident nuclei when they pass through each other properly. This is a crude approximation. A more reasonable assumption is that the fire-balls keep some memory on the motion of the incident nuclei, and therefore the distribution of fire-balls, instead of being uniform in the longitudinal direction, is more concentrated in the direction of motion of the incident nuclei, i.e. more dense at large absolute value of rapidity. A parametrization for such a distribution can be obtained by using an ellipse like picture on emission angle distribution, as shown in Fig.3. In this scenario, the emission angle is
$$\theta =\mathrm{tan}^1(e\mathrm{tan}\mathrm{\Theta }),$$
(5)
where the parameter $`e(0e1`$) represents the ellipticity of the introduced ellipse which describes the non-uniformity of fire-ball distribution in the longitudinal direction, as sketched in Fig.4.
Subsituting Eqn.(5) together with $`y_\mathrm{e}=\mathrm{ln}\mathrm{tan}(\theta /2)`$ into Eqn.(3), a distribution function $`Q(\theta )`$ for the emission angle of non-uniform flow is introduced in the integration
$`{\displaystyle \frac{\mathrm{d}n_{_{\mathrm{NUFM}}}}{\mathrm{d}y}}`$ $`=`$ $`eKm^2T{\displaystyle _{\theta _{\mathrm{min}}}^{\theta _{\mathrm{max}}}}{\displaystyle \frac{Q(\theta )\mathrm{d}\theta }{\mathrm{sin}\theta }}\left(1+{\displaystyle \frac{2T}{m\mathrm{cosh}(yy_\mathrm{e})}}+{\displaystyle \frac{2T^2}{m^2\mathrm{cosh}^2(yy_\mathrm{e})}}\right)`$ (6)
$`\times \mathrm{exp}(m\mathrm{cosh}(yy_\mathrm{e})/T),`$
$$y_\mathrm{e}=\mathrm{ln}\mathrm{tan}(\theta /2),Q(\theta )=\frac{1}{\sqrt{e^2+\mathrm{tan}^2\theta }|\mathrm{cos}\theta |}.$$
(7)
Here $`\theta _{\mathrm{min}}=2\mathrm{tan}^1(e^{y_{\mathrm{e0}}})`$, $`\theta _{\mathrm{max}}=2\mathrm{tan}^1(e^{y_{\mathrm{e0}}})`$. $`y_{\mathrm{e0}}`$ is the rapidity limit which confines the rapidity interval of longitudinal flow.
Changing the integration variable in Eqn.(6) back to $`y_\mathrm{e}`$, the rapidity distribution can be rewritten as follows:
$`{\displaystyle \frac{\mathrm{d}n_{_{\mathrm{NUFM}}}}{\mathrm{d}y}}`$ $`=`$ $`eKm^2T{\displaystyle _{y_{\mathrm{e0}}}^{y_{\mathrm{e0}}}}\left(1+{\displaystyle \frac{2T}{m\mathrm{cosh}(yy_\mathrm{e})}}+{\displaystyle \frac{2T^2}{m^2\mathrm{cosh}^2(yy_\mathrm{e})}}\right)`$ (8)
$`\times \mathrm{exp}(m\mathrm{cosh}(yy_\mathrm{e})/T)\rho (y_\mathrm{e})\mathrm{d}y_\mathrm{e}.`$
$$\rho (y_\mathrm{e})=\sqrt{\frac{1+\mathrm{sinh}^2(y_\mathrm{e})}{1+\mathrm{e}^2\mathrm{sinh}^2(y_\mathrm{e})}}.$$
(9)
Here $`\rho (y_\mathrm{e})`$ is the distribution function of non-uniform flow in the longitudinal direction. It can be seen from fig.5, that the larger is the parameter $`e`$, the flatter is the distribution $`\rho (y_\mathrm{e})`$ and the more uniform is the longitudinal-flow distribution. When $`e1`$, the longitudinal-flow distribution is completely uniform ($`\rho (y_\mathrm{e})1`$), returning back to the CSFM model.
III. Comparison with experiments
We now proceed to compare the model-results with the AGS data.
The rapidity distributions of pion, proton and deuteron for 14.6 A GeV/$`c`$ Si-Al collisions<sup></sup> are given in Fig.6 ($`a,b`$ and $`c)`$. The dashed, dotted and solid lines (band) correspond to the results from isotropical thermal model, uniform longitudinal flow model (CSFM) and the non-uniform longitudinal flow model (NUFM) respectively. The rapidity limit $`y_{\mathrm{e0}}`$ and the ellipticity $`e`$ used in the calculation are listed in Table I and illustrated in Fig.7. The rapidity limit $`y_0`$ used in the CSFM model of Ref. is also listed in Table I for comparison. The parameter $`T`$ is chosen to be 0.12 GeV following Ref..
Note that the distribution of the light particle (pion) is insensitive to the ellipticity $`e`$. Changing $`e`$ from 0.35 to 0.7 results in a narrow band shown in Fig.6($`a`$). Since pions are produced through the interaction of colliding nuclei, they have less memory on the motion of the incident nuclei before interaction. Therefore, physically the value of ellipticity $`e`$ for pion should be bigger than that for proton and deuteron, but the exact value cannot be fixed through fitting the model results to the experimental data.
Table I The value of model-parameters
| | Si-Al Collisions | | | Au-Au Collisions | |
| --- | --- | --- | --- | --- | --- |
| Parameter | $`\pi `$ | p | d | $`\pi `$ | p |
| $`e`$ | 0.35 – 0.7 | 0.52 | 0.56 | 0.35 – 0.7 | 0.58 |
| $`y_{\mathrm{e0}}`$ | 1.35 | 1.35 | 1.35 | 1.05 | 1.05 |
| $`y_0`$ | 1.15 | 1.15 | 1.15 | | |
It can be seen from the figures that the NUFM model reproduces the central dip of the rapidity distribution of heavier particles (proton and deuteron) in agreement with the experimental findings, while for light particles (pions) there is a peak instead of dip at central rapidity. Note that the appearance or disappearance of central dip is insensitive to the rapidity limit $`y_{\mathrm{e0}}`$ but depends strongly on the magnitude of the ellipticity $`e`$ and the mass $`m`$ of the produced particles. For the heavier particles (proton and deuteron) a central dip appears for $`e<0.8`$, but for light particles (pions) there is no dip even when $`e`$ is as small as 0.35, cf. Table I and Fig. 6.
On the other hand, the width of the rapidity distributions is mainly controlled by the amplitude $`y_{\mathrm{e0}}`$ of the longitudinal flow. A single value (1.35) of $`y_{\mathrm{e0}}`$ can account for the wide distribution of heavier particles (protons and deuterons) and at the same time fit the pion-distribution well. The difference in the width of the $`\mathrm{d}N/\mathrm{d}y`$ distribution originates mainly from the difference in the mass of the particles<sup></sup>.
In Fig.8 are shown the rapidity distributions of pions and protons for Au+Au collisions at 10.8 A GeV/$`c`$<sup></sup>. The dotted and solid lines (band) correspond to the results of CSFM and NUFM models (with parameters listed in Table I) respectively. In the calculation the CSFM results we have also used the NUFM model with the same rapidity limit $`y_{\mathrm{e0}}`$ listed in Table I but with ellipticity $`e=1`$. The histogram is the result from the RQMD model.
It can be seen from Fig.8 that in the NUFM model there is a shallow dip in the central rapidity of the distribution of proton, instead of a central peak as predicted by the CSFM model. However, the presently available experimental data are restricted to the large rapidity. The peak at central rapidity is the extrapolation of data using RQMD and is model dependent. It is interesting to see whether the prediction of a central dip (plateau) or a central peak will be observed in future experiments.
Comparing the parameter values for Si-Al (smaller colliding nuclei) and Au-Au (larger colliding nuclei) collisions listed in Table I, it can be seen that the rapidity limit $`y_{\mathrm{e0}}`$ is smaller and the ellipticity $`e`$ for proton is bigger for the larger colliding nuclei than for the smaller ones. Both of these two show that the hadronic system formed from the larger colliding nuclei is less spread out in rapidity, i.e. there is stronger nuclear stopping in the collision of larger nuclei.
IV. Summary and Conclusions
In high energy heavy-ion collisions, due to the transparency of nucleus the participants will not lose the historical memory and the produced hadrons will carry some of their parent’s memory of motion, leading to the unequivalence in longitudinal and transverse directions. So it is reasonable to assume that the flow of produced particles is privileged in the longitudinal direction. This picture has been used by lots of models<sup></sup>. Here we would mention two thermal and hydrodynamic models, one is the boost-invariant longitudinal expansion model postulated by Bjorken <sup></sup> which can explain such an anisotropy already at the level of particle production in hadron-hadron collisions. This model has been formulated for asymptotically high energies, where the rapidity distribution of produced particles establishes a plateau at midrapidity. The second model is the CSFM postulated first by Schnedermann, Sollfrank and Heinz<sup></sup> which accounts for the anisotropy of longitudinal and transverse direction by adding the contribution from a set of fire-balls with centers located uniformly in the rapidity region \[-$`y_0`$,$`y_0`$\] in the longitudinal direction, sketched schematically in Fig’s.1 and 2 as dashed circles. It can account for the wider rapidity distribution when comparing to the prediction of the pure thermal isotropic model but fails to reproduce the central dip in the proton and deuteron rapidity distributions.
In this paper, we argue that the transparency/stopping of relativistic heavy ion collisions should be taken into account more carefully. It will not only lead to the anisotropy in longitudinal-transverse directions, but also render the fire-balls (especially for those of proton and deuteron) concentrate more in the direction of motion of the incident nuclei. A non-uniform longitudinal flow model is proposed, which assumes that the centers of fire-balls are distributed non-uniformly in the longitudinal phase space. A parameter $`e`$ is introduced through a geometrical parametrization which can express the non-uniformity of flow in the longitudinal direction, i.e. the centers of fire-balls of proton and deuteron prefer to accumulate in the two extreme rapidity regions ($`y_\mathrm{e}\pm y_{\mathrm{e0}}`$) in the c.m.s. frame of relativistic heavy-ion collisions, and accordingly the distribution is diluted in the central rapidity region ($`y_\mathrm{e}0`$).
It is found that the depth of the central dip depends on the magnitude of the parameter $`e`$ and the mass of produced particles, i.e. the non-uniformity of longitudinal flow which is described by the parameter $`e`$ determines the depth of the central dip for heavier particles. On the other hand, the central peak in the pion distribution turns out to be insensitive to the value of $`e`$ and can be well reproduced from this model simultaneously with the central dip in the proton and deuteron rapidity distributions.
Through comparing the feature of collision systems of different size, it is found that the maximum flow velocities are smaller for the heavier collision system than for the lighter ones, which suggests, together with the larger $`e`$, a larger stopping in the bigger collision system.
Figure captions
Fig.1 Schematic sketch of the distribution of fire-balls in the uniform flow model (CSFM).
Fig.2 Schematic sketch of the emission angle $`\mathrm{\Theta }`$ (solid circle and lines) and the corresponding distribution of fire-balls (dashed circles) in the uniform flow model (CSFM).
Fig.3 Schematic sketch of the emission angle $`\theta `$ in the non-uniform flow model (NUFM).
Fig.4 Schematic sketch of the distribution of fire-balls in the non-uniform flow model (NUFM).
Fig.5 The distributions $`\rho (y_\mathrm{e})`$ of the center of fire-balls as a function of $`y_\mathrm{e}`$. When $`e1`$, the non-uniform distribution turns to the uniform distribution $`\rho (y_\mathrm{e})1`$.
Fig.6 Rapidity distributions for central 14.6 A GeV/$`c`$ Si+Al collisions. Open circles are the experimental data for Si+Al collisions taken from Ref.’s . Dashed, dotted and solid lines are the distributions from the isotropical thermal model, cylindrical-symmetric longitudinal flow model (CSFM) and non-uniform longitudinal flow model (NUFM) respectively. Fig’s. ($`a`$), ($`b`$) and ($`c`$) are the pion, proton and deuteron distributions respectively. The ellipticity $`e`$ for pion is within a range \[0.35, 0.7\] and therefore the NUFM results are present as a band. The temperature $`T=0.12`$ GeV.
Fig.7 The fire-ball distribution functions $`\rho (y_\mathrm{e})`$ verus rapidity $`y_\mathrm{e}`$ in the non-uniform flow of different particles for the collisions of Si+Al at 14.6 A GeV/$`c`$ ($`a`$) and Au+Au at 10.8 A GeV/$`c`$ ($`b`$) respectively. For pion the distribution functions $`\rho (y_\mathrm{e})`$ are only ploted for the two boundaries of the used region $`0.35e0.7`$.
Fig.8 Rapidity distributions for pions ($`a`$) and protons ($`b`$) in central Au+Au collisions at 10.8 A GeV/$`c`$. Full circles represent the experimental data taken from Ref.. Open circles are the reflection of the data. The solid line is our calculation using the NUFM model. The ellipticity $`e`$ for pion is within a range \[0.35, 0.7\] and therefore the results for pion are present as a band. The histogram shows the results from RQMD and the dotted lines are the results from the prediction of CSFM model. The temperature $`T=0.14`$ GeV.
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# Early Starbursts and Magnetic Field Generation in Galaxy Clusters
## 1 Introduction
Rich clusters of galaxies are the largest gravitationally bound structures in the Universe and should confine a representative fraction of its mass. Therefore the study of their dynamical properties and radiation content should allow, amongst other things, interesting cosmological conclusions on the relative amounts of visible and dark baryonic matter, and of nonbaryonic matter (e.g. White & Fabian (1995); Turner (1999)).
Another basic characteristic, due to energetic particle confinement, is the ratio of nonthermal to thermal energy in these objects. To a significant extent that ratio should be pre-determined during the epoch of early starburst activity and thus preserve the energetic history of cluster formation. The necessary confinement of the nonthermal particle components is intimately related to the existence of strong and chaotic magnetic fields in the Intracluster Medium (ICM), and we shall propose physical mechanisms for their early generation as well as for their present fluctuations.
In principle, detailed ab initio simulations of the dynamics of cluster formation under the dominant gravitational influence of the dark matter component (see e.g. Kauffmann et al. (1999)) should establish the overall cosmological framework for the present considerations. We rather start in a complementary way with the discussion of a simplified model of cluster formation and of chemical enrichment of the Intracluster gas. It has the advantage that it directly allows a discussion of the physical processes of nonthermal particle production and confinement. The main part of the paper concerns a proposal of cluster magnetic field generation in terms of galactic winds due to early starbursts and their amplification effect on magnetic fields drawn out from the progenitors of today’s cluster galaxies into Intracluster space. It is argued that due to these dynamical processes there is no need for the operation of a dissipative turbulent dynamo in the ICM. The ongoing cluster accretion naturally leads to a strong fluctuating part of the Intracluster magnetic fields. A detailed discussion of the nonthermal radiation from galaxy clusters will be given in a separate paper (Atoyan & Völk (1999)).
## 2 Rich Clusters
We shall be concerned here with rich clusters, i.e. conglomerates with typically more than 100 member galaxies. They have typical radii $`R_{\mathrm{cl}}`$ few Mpc and baryonic masses $`M_{\mathrm{cl}}10^{14}\mathrm{to}\mathrm{\hspace{0.17em}10}^{15}\mathrm{M}_{}`$. Many such clusters are rather evolved and contain predominantly early type S0 and E-galaxies, at least in their inner parts. Examples for bright and relatively nearby clusters of this type are the Perseus and the Coma clusters with distances $`d100\mathrm{Mpc}`$. The Perseus cluster is the brightest cluster in soft X-rays. The large X-ray luminosity is due to the very hot ($`T10^7\mathrm{to}\mathrm{\hspace{0.17em}10}^8`$K), massive ($`M_{\mathrm{gas}}\mathrm{few}\times M_{\mathrm{gal}}`$), and metal-rich ($`[\mathrm{Fe}]_{\mathrm{cl}}0.35[\mathrm{Fe}]_{}`$) ICM gas (e.g. Böhringer (1996)). As a consequence the gas pressures are extremely high, with $`nT`$ ranging from $`10^3\mathrm{to}\mathrm{\hspace{0.17em}10}^5\mathrm{Kcm}^3`$.
### 2.1 Cluster Formation
The metallicity of the ICM gas, for instance in terms of the fractional ICM iron mass, is correlated with the total optical luminosity in the E and S0 galaxies of rich clusters (Arnaud et al. (1992)). The correlation supports the view that early starbursts due to galaxy-galaxy interactions of gas-rich progenitors have produced a large number of core collapse Supernovae due to massive stars (for simplicity referred to here as SNe). They should have heated the originally present interstellar gas and generated violent galactic winds which removed the Interstellar Medium, leaving gas-poor E and S0 galaxies behind. This mass loss should have led to the observed strong chemical enrichment of the ICM gas. We also conjecture that the ionizing radiation, the winds, and the large-scale shocks from these early galaxy mergers - together with the hard radiation from AGNs - strongly heated the remaining primordial ICM gas, and thus prevented further galaxy formation. A quantitative discussion of the dynamical prerequisites for galactic winds and of the total number of SNe in clusters is given by Völk et al. ((paper1, 1996, hereafter referred to as Paper I)).
The total number of SNe since galaxy formation in the cluster, roughly a Hubble time $`T_\mathrm{H}1.5\times 10^{10}\mathrm{yr}`$ ago, is then given by
$$N_{\mathrm{SN}}=_{T_\mathrm{H}}^0𝑑t\times \nu _{\mathrm{SN}}(t)=\frac{0.35[Fe]_{}\times M_{\mathrm{cl}}}{\delta M_{\mathrm{Fe}}},$$
where $`\delta M_{\mathrm{Fe}}`$ is the amount of iron produced per event. In such starbursts we dominantly expect core collapse SNe from massive progenitor stars to occur, with $`\delta M_{\mathrm{Fe}}0.1M_{}`$ on average. For the Perseus cluster this implies $`N_{\mathrm{SN}}^{\mathrm{Perseus}}3\times 10^{12}`$. The corresponding total energy input into the interstellar medium is $`N_{\mathrm{SN}}E_{\mathrm{SN}}3\times 10^{63}E_{51}\mathrm{erg}`$, where $`E_{51}=(E_{\mathrm{SN}}/10^{51}\mathrm{erg})`$ is the average hydrodynamic energy release per SN in units of $`10^{51}\mathrm{ergs}`$.
Assuming the early starbursts to occur at a typical redshift of $`z2`$ due to the merging of gas-rich progenitors in an overdense protocluster environment (Steinmetz (1993)), with a duration of $`T_{\mathrm{SB}}<\mathrm{\hspace{0.17em}10}^9\mathrm{yr}`$, we obtain
$$\frac{(N_{\mathrm{SN}}^{\mathrm{Perseus}}/N_{\mathrm{gal}}^{\mathrm{Perseus}})}{T_{\mathrm{SB}}}>\mathrm{\hspace{0.17em}100}\times \nu _{\mathrm{SN}}^{\mathrm{Milky}\mathrm{Way}},$$
where $`\nu _{\mathrm{SN}}^{\mathrm{Milky}\mathrm{Way}}`$ is taken as 1/(30 yr), and $`N_{\mathrm{gal}}^{\mathrm{Perseus}}500`$ denotes the number of galaxies in the Perseus cluster. As an example we can compare to the archetypical contemporary starburst galaxy $`M82`$. It has a current SN rate $`\nu _{\mathrm{SN}}^{\mathrm{M82}}10\times \nu _{\mathrm{SN}}^{\mathrm{Milky}\mathrm{Way}}`$, a wind velocity $`v_{\mathrm{wind}}2300\mathrm{km}/\mathrm{sec}`$, and a mass-loss rate of $`\dot{M}0.8M_{}/\mathrm{yr}`$ (Breitschwerdt (1994)). The starburst nucleus of M82 is characterized by the following values for the interstellar gas temperature $`T`$, gas density $`n`$, and thermal gas pressure $`p`$ at the base of the wind: $`T_{\mathrm{base}}10^8\mathrm{K}`$, $`n_{\mathrm{base}}0.3\mathrm{cm}^3`$, and $`p_{\mathrm{base}}/k_\mathrm{B}10^7\mathrm{K}\mathrm{cm}^3`$ (Schaaf et al. (1989)). Since the thermal ICM gas pressure in the Perseus cluster is $`p_{\mathrm{cl}}^{\mathrm{Perseus}}/k_\mathrm{B}10^4\mathrm{K}\mathrm{cm}^3`$, it is clear that an object like M82 could readily drive a wind even against the present-day ICM pressure. At the galaxy formation epoch the ICM pressure should have been much smaller than this value.
In an expanding galactic wind flow the SN-heated gas will cool adiabatically to quite small temperatures. However it will be reheated in the termination shock, where the ram pressure of the wind adjusts to the ICM pressure. Much beyond this point the ejected galactic gas is expected to slowly exchange energy and some metal-rich material with the unprocessed ICM gas.
### 2.2 Nonthermal particle production
Cluster formation also implies the production of a strong nonthermal component of relativistic particles. They will be accelerated during the early phase - and presumably also in later accretion events that have shock waves associated with them. In the Supernova Remnants (SNRs) the main acceleration should occur at the outer shock, with a very high efficiency of $`10\mathrm{percent}`$ (Drury et al. (1989); Berezhko & Völk (2000)), the rest of $`E_{\mathrm{SN}}`$ going into the thermal gas ($`10\mathrm{to}\mathrm{\hspace{0.17em}20}\mathrm{percent}`$) (Dorfi (1993)), and radiation ($`>\mathrm{\hspace{0.17em}70}`$ percent). However, since the particles are ultimately removed from the galaxies in a strong galactic wind (see Paper I and section 3. below), they will cool adiabatically like the thermal gas and transfer their energy to the kinetic energy of the wind flow. As a consequence the original SNR-accelerated particles should constitute a negligible fraction of the present-day nonthermal particle content of the cluster. But this is not the end of the story. At distances $`100\mathrm{kpc}`$ from the galaxies, fresh particle acceleration will occur at the strong galactic wind termination shocks (Fig. 1). We estimate the overall acceleration efficiency in these shocks to be again of the order of 10 percent<sup>1</sup><sup>1</sup>1Note that this acceleration efficiency might as well be lower, if the termination shocks were essentially perpendicular.. Over the early phase of galaxy formation, using our estimate for the total number of SNe, this should result in a total gas internal energy in the postshock region $`E_{\mathrm{gas}}^{\mathrm{GW}}\mathrm{few}\times 10^{62}`$ erg and a nonthermal energy $`E_{\mathrm{CR}}^{\mathrm{GW}}\mathrm{few}\times 10^{61}`$ erg for a system like the Perseus cluster, ultimately driven by star formation and the subsequent SN explosions. Since the confinement time in the cluster (see next subsection) exceeds the cluster lifetime, the energy spectrum of these particles in the ICM is basically the same as the source spectrum at the termination shocks. In Cosmic Ray (CR) parlance these are cosmological CRs.
Since the galaxies are distributed across the cluster quasi-uniformly, this should originally also be true for the nonthermal particle population. The ensuing gravitational contraction/accretion of the cluster will subsequently energize CRs and thermal gas at least adiabatically or, more likely, shock accelerate/irreversibly heat both components so that finally the total energy $`E_{CR}`$ of energetic particles should reach at least the adiabatic value $`3\times 10^{62}\mathrm{erg}1/30E_{\mathrm{gas}}`$ in the cluster; $`E_{\mathrm{gas}}`$ now denotes the total present internal energy of the ICM gas (Paper I).
It is worthwhile to compare the expected nonthermal energy with the thermal energy content of the cluster galaxies. Assuming the stars internally to be in virial equilibrium and, for purposes of estimate, all of them to have a solar mass and radius, then $`E_{\mathrm{th}}^{\mathrm{star}}(3/10)GM_{}^2/R_{}10^{48}\mathrm{erg}`$. For a total mass of about $`10^{14}M_{}`$ contained in the galaxies of the Perseus cluster this gives a total thermal energy in stars $`10^{62}\mathrm{erg}`$, and thus $`E_{\mathrm{CR}}>_{\mathrm{gal}}_{\mathrm{stars}}E_{\mathrm{th}}^{\mathrm{star}}`$. This means that the nonthermal ICM energy should be at least as large as the total thermal energy content of all the stars in all the galaxies together.
It has been argued more recently that, apart from star formation and overall gravitational contraction, also individual giant radio galaxies should have injected large and in fact comparable amounts of nonthermal particles during the life time of a cluster (Enßlin et al. (1997); Berezinsky et al. (1997)). This is no doubt an important additional possibility. A weakness of this argument consists in the fact that per se it is predicated on statistical knowledge about the luminosity function for active galaxies in clusters in general, and not on direct observations of the individual cluster to which it is applied.
### 2.3 Particle Confinement
Energetic particle confinement has been discussed in a number of papers in recent years (Paper I, Berezinsky et al. 1997, Colafrancesco & Blasi 1998). We review here this important point, also in the light of the magnetic field structure to be discussed in the next section.
The large-scale magnetic field in the ICM gas may be quite chaotic and not well connected over distances exceeding typical intergalactic distances. Thus energetic particles may not readily escape from the cluster due to such topological characteristics. However, already a consideration of pure pitch angle diffusion along straight magnetic field lines with superposed turbulent fluctuations gives important insights into the confinement properties of galaxy clusters. Standard quasilinear theory yields a spatial diffusion coefficient $`\kappa _{}`$ along the large scale field $`B`$ due to a power spectrum $`P(k)`$ of magnetic field fluctuations with wavelength $`\lambda =2\pi /k`$ in the following form:
$$\kappa _{}=(1/3)cr_\mathrm{g}(p)\frac{B^2}{_k^{\mathrm{}}dk^{}P(k^{})},$$
where $`r_\mathrm{g}(p)`$ is the gyro radius of a particle with momentum $`p`$, $`kr_\mathrm{g}(p)1`$, and $`k`$ denotes the wavenumber of the field fluctuations. Let us assume a relative fluctuation field strength of order unity at the inter-galaxy distance $`1/k_0`$, i.e. a totally turbulent field $`P(k_0)\times k_0B^2`$ on this scale, and a power law form of $`P(k)=P(k_0)(k/k_0)^n`$. Then the diffusion time across the cluster $`T_{\mathrm{esc}}R_{\mathrm{cl}}^2/\kappa _{}>T_H`$ for $`(cp)_{\mathrm{protons}}<\mathrm{\hspace{0.17em}10}^{17}\mathrm{eV}\mathrm{and}<\mathrm{\hspace{0.17em}10}^{15}\mathrm{eV}`$, for $`n=3/2`$ and $`n=5/3`$, respectively (Paper I). Also $`t_{\mathrm{loss}}^{\mathrm{protons}}T_H`$ for nuclear collisions in the ICM gas. Therefore - except at subrelativistic energies with their prevailing Coulomb losses - up to these energies CR hadrons should accumulate in the cluster since the galaxy formation epoch, and that is what we called cosmological CRs before. The situation is different for relativistic electrons, which suffer radiative losses. Electrons, and the nonthermal radiation from galaxy clusters are discussed extensively in the paper by Atoyan & Völk (2000).
### 2.4 Gamma-rays from cosmological CRs
The accumulated cosmological CR protons and nuclei will produce high energy $`\gamma `$ rays from inelastic pp-collisions on the Intracluster gas which, in particular, lead to $`\pi ^0`$-production and subsequent decay. Fig. 2 shows the $`\pi ^0`$-decay energy fluxes expected from the hadronic CRs of the Coma cluster with a differential energy spectrum $`E^{\alpha _{\mathrm{CR}}}\mathrm{exp}(E/E_c)`$, with $`\alpha _{\mathrm{CR}}=2.1`$ and an upper cutoff energy $`E_c=200\mathrm{TeV}`$. The solid and the dashed curve assume $`E_{\mathrm{CR}}=3\times 10^{62}\mathrm{erg}`$ and $`E_{\mathrm{CR}}=3\times 10^{61}`$ erg, respectively, in an ICM with a gas density of $`n=10^3\mathrm{cm}^3`$. Also the EGRET upper limit by Sreekumar et al. (sre96 (1996)) is shown. The observed size of the radio emission produced by high-energy electrons in Coma is about half a degree. The brightness of the hadronic TeV emission should be significantly higher in the central region of the cluster, since the ICM gas density strongly increases towards the center. Detection of an extended and weak ($``$0.1 Crab) TeV flux by current instruments is problematic, but such fluxes are quite accessible for future imaging atmospheric Cherenkov telescope arrays like H.E.S.S., VERITAS, and CANGAROO III (see Fig. 2).
## 3 Intracluster Magnetic Fields
The magnetic field strengths in the ICM of rich clusters, which may be as large as $`B1\mu \mathrm{G}`$ as inferred by Faraday rotation measurements (e.g. Kronberg (1994)), are not easily explained by a contemporary dissipative mechanism because present-day turbulent dynamo effects in such a large-scale system should be extremely slow. This problem is compounded by the extraordinary smallness of the expected intergalactic seed fields. Therefore we suggest a field configuration that is due to the early formation history of galaxy clusters as discussed above; it should be preserved in its essential features to this day (Völk & Atoyan (1999)). This field should even be still in a state of development at the present epoch. The model derives from the violent early galactic winds which accompany the starbursts responsible for the predominance of the early type galaxies in rich clusters.
In a general form the ejection of galactic fields has been discussed by Kronberg (e.g. Kronberg (1994). For the generation of the general Intergalactic magnetic field, related arguments have been advanced independently by Kronberg et al. (kro99 (1999)), assuming very early formation of dwarf galaxies with wind outflows at redshifts $`10`$.
We assume first of all that the gas-rich progenitors whose mergers supposedly constitute the building blocks for the E and S0 galaxies, can be specifically pictured as protospirals that had already generated galactic magnetic fields of $`\mu `$G strength. This should indeed be possible within a time of $`10^8`$ yr or less, i.e. a time scale of the order of a rotation period of our Galaxy or even shorter. The explanation derives from fast turbulent dynamo action that invokes boyancy effects due to CRs that inflate the magnetic flux tubes together with magnetic reconnection over spatial scales of order $`100\mathrm{p}\mathrm{c}`$ (Parker (1992)). In starburst galaxies also systematic dynamo effects might play an important role (section 3.1).
In a second stage, i.e. during the galaxy mergers, the resulting supersonic galactic winds will extend these fields from the interacting galaxies to almost intergalactic distances. In the final and by far longest stage that lasts until now, the fields should be recompressed by the contraction of the cluster to its present size<sup>2</sup><sup>2</sup>2For a different view, emphasizing internal shear flows, see Dolag et al. (dol99 (1998).). In addition, the continuing accretion of subclusters and individual galaxies constantly perturbs this field, keeping the fluctuations around this large scale field at a high level.
The ICM fields do not reconnect on the intergalactic scale in a Hubble time. Consequently there is no need for a continuous regeneration of these fields since their formation. However, this also implies that a topologically connected overall ICM field will on average not be formed either and that the ICM field is chaotic on a scale smaller or equal to the present intergalactic distance.
In detail we draw on arguments which we have used in the past for the field configuration in a galactic wind from our own Galaxy (Zirakashvili et al. (1996); Ptuskin et al. (1997)), see also Fig. 3. They are based on estimates of the relative amount of field line reconnection vs. the extension of Galactic field lines by a wind to ”infinity” (Breitschwerdt et al. (1993)). The basic result is that the rates of reconnection - and thus of the formation of ”Parker bubbles” leaving the Galaxy by their boyancy and allowing the generation of the disk magnetic field - and the rates of extension of this field into the Galactic Halo by the pressure forces of the wind are roughly equal. Thus both effects occur with about equal probability. For the cluster galaxies this means that magnetic energy can be generated on the large scale of the wind at the expense of the thermal and nonthermal enthalpies produced in the starburst. The geometry of the field should roughly correspond to straight field lines out to meridional distances $`s`$ of the order of the starburst (SB) radius, $`R_{\mathrm{gal}}^{\mathrm{SB}}1`$ kpc in the protogalactic disk, and spherically diverging field lines beyond that. The slow rotation of the system should then lead to an azimuthal field component, decreasing with the wind distance $`1/s`$, which dominates at large distances over the radial component. However, in contrast to the familiar situation in the Solar Wind equatorial plane, the axis of rotation is rather parallel than perpendicular to the flow at the base of the wind, and thus the dominance of the azimuthal field component is by no means as drastic as in the case of a stellar wind (Fig. 3). The wind becomes supersonic at about the same critical distance, $`s_{crit}R_{\mathrm{gal}}^{\mathrm{SB}}`$. Far beyond this critical point the mass velocity $`u(s)`$ becomes constant and the density falls off $`s^2`$.
### 3.1 Mean magnetic field
Choosing a present average baryon density $`n_\mathrm{b}(z=0)=3\times 10^7\mathrm{cm}^3`$, i.e. assuming most of the baryonic matter to be in the form of Intergalactic gas, the mean density at the formation stage of the early type galaxies (at redshifts $`z2`$) was then $`n_\mathrm{b}(z=2)=27n_\mathrm{b}(z=0)n_{\mathrm{cl}}(z=2)`$. With a present ICM density $`n_{\mathrm{cl}}>10^4\mathrm{cm}^3`$, we have $`n_{\mathrm{cl}}(z=0)/n_{\mathrm{cl}}(z=2)>12`$ due to gravitational compression of the ICM gas. With a dominant thermal gas pressure the corresponding adiabatic pressure increase $`p_{\mathrm{cl}}n^{5/3}`$ amounts to $`p_{\mathrm{cl}}(z=0)/p_{\mathrm{cl}}(z=2)>63`$; for the following we shall use $`p_{\mathrm{cl}}(z=0)/p_{\mathrm{cl}}(z=2)=10^2`$. The wind termination shock distance $`s_{\mathrm{sh}}`$ is then given by $`\rho (r_{\mathrm{sh}})u^2p_{\mathrm{cl}}(z=2)`$, where $`\rho (r_{\mathrm{sh}})`$ is the wind mass density upstream of the shock.
To estimate the wind characteristics we assume a quasi-steady state and a strong starburst for which we may disregard gravity, magnetic forces and CR pressure gradients in the overall energy balance equation (e.g. Zirakashvili et al. (1996)). It reads in this case:
$$\left[\frac{u^2}{2}+\frac{\gamma }{\gamma 1}\frac{p}{\rho }\right]_s\frac{\gamma }{\gamma 1}\frac{p_{\mathrm{base}}}{\rho _{\mathrm{base}}},$$
where in addition $`\rho u^2(s)/2\gamma p/(\gamma 1)`$ is assumed at the base $`s=s_{\mathrm{base}}`$ of the wind. The critical point $`s_{\mathrm{crit}}`$ is given by $`\rho _{\mathrm{crit}}u_{\mathrm{crit}}^2=\gamma p_{\mathrm{crit}}`$. Approximately $`u_{\mathrm{crit}}u_{\mathrm{}}/2`$, where $`u_{\mathrm{}}`$ is the asymptotic wind speed. Beyond the critical point at $`s_{\mathrm{crit}}R_{\mathrm{gal}}^{\mathrm{SB}}`$ the wind achieves spherical symmetry so that from mass conservation
$$\left(\frac{s_{\mathrm{sh}}}{R_{\mathrm{gal}}^{\mathrm{SB}}}\right)^2\frac{\rho _{\mathrm{crit}}u_{\mathrm{crit}}}{\rho _{\mathrm{sh}}u_{\mathrm{sh}}}$$
If the ICM pressure is small compared to the base pressure then we can assume that the wind is highly supersonic at the shock distance which implies $`u_{\mathrm{sh}}u_{\mathrm{}}`$. This should be true for M82-like objects - although they should have a considerably larger scale $`R_{\mathrm{gal}}^{\mathrm{SB}}1`$ kpc - which have $`p_{\mathrm{gas}}/k_\mathrm{B}10^7\mathrm{Kcm}^3`$. In this case we can assume the wind pressure to evolve adiabatically to lowest order, so that
$$\frac{p_{\mathrm{crit}}}{p_{\mathrm{base}}}\left(\frac{\rho _{\mathrm{crit}}}{\rho _{\mathrm{base}}}\right)^{5/3}$$
As a result, taking into account that in a flow dominated by the thermal gas the adiabatic index $`\gamma =5/3`$
$$\frac{s_{\mathrm{sh}}}{R_{\mathrm{gal}}^{\mathrm{SB}}}\left[\frac{4\gamma }{\gamma +1}\left(\frac{1}{2(\gamma 1)}\right)^{\frac{1}{\gamma 1}}\right]^{1/2}\left(\frac{p_{\mathrm{base}}}{p_{\mathrm{cl}}}\right)^{1/2}1.27\left(\frac{p_{\mathrm{base}}}{p_{\mathrm{cl}}}\right)^{1/2}400$$
for the above starburst parameters.
Therefore
$$\frac{s_{\mathrm{sh}}}{d_{\mathrm{gal}}^{\mathrm{field}}(0)/(1+z)}\frac{400\mathrm{kpc}}{(2\mathrm{Mpc}/3)}0.6$$
The Wind Bubble containing the shock-heated wind gas will have a radius still exceeding $`s_{\mathrm{sh}}`$. Beyond the Bubble, part of the external gas will be shocked by the rapidly expanding Bubble gas. This shock heated gas will exchange energy with the cold ambient gas by heat conduction and instabilities. To some extent such exchanges will also take place with the Bubble gas.
Even though the volume of at least initially unmagnetized ICM gas may be large enough so that the ICM gas mass exceeds the mass associated with galaxies by a factor of a few - as observed - it may therefore be that some Wind Bubbles touch. Thus we should consider whether the magnetic fields of the bubbles can reconnect with each other to produce field lines that pervade the entire cluster. Let us assume that the fastest rate of reconnection proceeds with a speed between 1 and 10 percent of the Alfv$`\stackrel{´}{\mathrm{e}}`$n speed $`v_\mathrm{A}`$, to be specific, with $`v_\mathrm{A}/50`$ (Parker (1992)). Then the present-day reconnection time across a termination shock scale is
$$t_{\mathrm{rec}}\frac{s_{\mathrm{sh}}}{(v_\mathrm{A}/50)}9\times 10^{10}\mathrm{yr}\left\{\left(\frac{s_{\mathrm{sh}}}{400\mathrm{k}\mathrm{p}\mathrm{c}}\right)\left(\frac{B}{10^6\mathrm{G}}\right)^1\left(\frac{n}{10^4}\right)^{1/2}\right\},$$
greater than a Hubble time, even for a present-day ICM magnetic field as high as $`10^6`$ G. Therefore, many bubble fields might not yet be reconnected and much of the field structure could well remain topologically disconnected until today.
The field strength $`B_{\mathrm{cl}}(z=2)`$ in the early Wind Bubbles should be of the order of
$$B_{\mathrm{cl}}(z=2)4B_{\mathrm{gal}}r_{\mathrm{gal}}^{\mathrm{SB}}/r_{\mathrm{sh}}10^2B_{\mathrm{gal}}10^8\mathrm{G}$$
or somewhat larger, if the field in the bubble increases in the decelerating postshock flow. The subsequent and still ongoing overall cluster contraction/accretion compresses the field to lowest order isotropically $`l^2`$, with the scale factor
$$l[n_{\mathrm{cl}}(0)/n_{\mathrm{bar}}(0)]^{1/3}/(1+z),$$
where $`n_{\mathrm{cl}}(0)(10^3\mathrm{to}10^4)\mathrm{cm}^3`$ and $`z=2`$. For a present mean baryon number density $`n_{\mathrm{bar}}(0)3\times 10^7\mathrm{cm}^3`$, we obtain $`l(2.3\mathrm{to}5)`$.
From these estimates, we obtain for the present-day ICM field: $`B_{\mathrm{cl}}(z=0)/B_{\mathrm{cl}}(z=2)5\mathrm{to}25`$. Therefore the present-day ICM magnetic field should have a mean strength of the order of $`10^7`$G, from $`1\mu \mathrm{G}`$ ”primordial” seed fields, and should be randomly directed on an intergalactic scale. Although smaller by about one order of magnitude than estimated from Faraday rotation measurements, such fields need not necessarily be unrealistic, considering that observations might emphasize regions of high magnetic fields. Indeed the simplest inverse Compton interpretation of the EUV excess and the excess X-ray flux in the Coma cluster requires such low field strengths (Fusco-Femiano et al. (1999)).
On the other hand, the increase of the field in the galactic Wind Bubbles beyond their postshock value might be more than a factor of unity as assumed above.
An additional possibility is that the ”initial” fields for such starburst galaxies might be an order of magnitude stronger than assumed. Apart from a fast dynamo whose strength could be directly proportional to the star formation rate (Parker (1992)), we could invoke a systematic field amplification through the commencing galactic outflow that ”combs” the field outward. It might occur on very short time scales of about $`10^7`$ yr. Empirically this field amplification is suggested by the statistical time-independence of the radio synchrotron to far infrared emission ratio in starburst galaxies (Lisenfeld et al. (1996)) which can hardly be understood otherwise than through a field that increase almost simultaneously in strength with the star formation rate.
A final argument is that the large scale shear deformations induced by the later accretion of large subclusters may amplify the field even further. Thus we cannot exclude $`\mu `$G fields; the force balance certainly allows them.
### 3.2 Field evolution and structure
The Wind Bubbles and the associated magnetic fields should at some stage decouple from the galaxies they emanated from, simply by magnetic reconnection which is fastest near the galaxies: the Alfv$`\stackrel{´}{\mathrm{e}}`$n velocity is approximately independent of the meridional distance s in the wind, whereas the distance between oppositely directed field lines that emanate from the galaxy is obviously the smaller the smaller s is. Thus, after the termination of the starburst, the magnetized bubbles and the stellar component of the remaining early type galaxies should acquire independent identities and their dynamics should decouple.
The topologically disconnected structure of the mean magnetic field in the cluster has also some bearing on the evolution of magnetic field strengths during the development of cooling flows towards the cluster center (Fabian (1994)). Instead of building up global magnetic pressure gradients and tension forces such a subsonic and sub-Alfv$`\stackrel{´}{\mathrm{e}}`$nic flow allows optimal internal segregation of high and low field regions on galaxy - galaxy seperation scales. Thus we should expect that the field strength towards the cluster center increases by less than by isotropic compression, $`Bn^{2/3}`$, even though at the compressed spatial scales reconnection will be more effective. The two-thirds law should hold only for the overall cluster gas compression discussed earlier.
A question of direct importance for the interpretation of Faraday rotation measurements is the scale of reversal changes of these fields. Field reversals are a natural consequence of the suggested field structure, given the field pattern in the interstellar media of the starburst galaxies. If the Milky Way can serve as a guide, then this pattern is determined by the 100 pc scale of the Parker instability (Parker 1966). From the radius $`r_{\mathrm{gal}}^{\mathrm{SB}}`$ to distances $`r_{\mathrm{sh}}`$ this scale projects like the ratio $`r_{\mathrm{sh}}/r_{\mathrm{gal}}^{\mathrm{SB}}`$. Using the above estimates this implies a field reversal scale $`40`$ kpc in clusters, rather well in line with observational estimates which indicate reversal scales of $`10`$ to $`100`$ kpc (e.g. Kronberg (1994); clarke99).
### 3.3 Magnetic field fluctuations
From the foregoing arguments there is hardly any need for a contemporary ”turbulent ICM dynamo”. However, the ongoing accretion will perturb this mean intracluster field randomly in space and time, maintaining a turbulent magnetic fluctuation field that develops smaller and smaller spatial scales. Due to the topology of the mean field, the largest turbulent scale is given by the distance between galaxies. The accretion will probably span the range from single field galaxies falling into the cluster to accreting massive subclusters. As long as this accretion process remains important, one has to expect that the cluster rings with it. At the largest scale the relative magnetic fluctuation level should approach unity (see subsection 2.3).
The authors thank F. A. Aharonian, E. N. Parker, and R. Rosner for illuminating discussions. The work of AMA was supported through the Verbundforschung Astronomie/Astrophysik of the German BMBF under the grant No. 05-2HD66A(7).
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# Multifractal Behaviour of n-Simplex Lattice
## 1 Introduction
In recent years considerable attention has been devoted to studying the properties of disordered systems with the hope of understanding percolative phenomena. Key to several such approaches has been the concept of randomness and also of frustration \[1-12\]. However, many of the patterns we encounter in nature are not random but self-similar and scale invariant \[13-14\]. For instance, the complicated and scale invariant structures that occur when a solid mixture evolves via an aggregation process . To understand such systems the concept of fractals has been found to be very useful. Fractals are scale invariant objects that may be considered as intermediate lattices between regular and random (disordered) lattices \[13-14,16-17\]. Such a fractal lattice describes a class of random systems where the consequence of the loss of translational invariance of a lattice can be studied in detail. Additionally, resulting from their dialational symmetry, statistical, mechanical and transport problems are solvable; hence the attraction of the model in such studies .
In this paper we consider a particular class of fractal known as the n-simplex lattices, to model various properties of inhomogeneous materials \[16-17\]. The lattice is defined recursively. The map of the zero-order truncated n-simplex lattice is a complete set of (n+1) points. The map of the (r+1)th order n-simplex lattice is obtained by replacing each of the lattice points of the rth order map by the entire rth order map. Each of the resulting n points is connected to one of the lines connecting the original rth order vertices. The fractal and spectral dimensions of this lattice are given by:
$$\overline{d}=\frac{\mathrm{ln}(n)}{\mathrm{ln}2}$$
(1)
and
$$\stackrel{~}{d}=\frac{2\mathrm{ln}(n)}{\mathrm{ln}(n+2)}.$$
(2)
The lattices with $`n3`$ are of particular interest as they provide a family of fractals in which $`\overline{d}`$ varies with $`n`$ leaving $`\stackrel{~}{d}`$ almost constant.
In order to understand how resistance scales with the size of the system, in a homogeneous system, we study the distribution of currents in a network modeled by an n-simplex lattice. We consider each bond, where bond refers to a line joining two lattice points, of the zero-order network as a unit resistor offering resistance $`𝐑`$. A unit current enters the network at one of the external nodes and leaves through another, the rest of nodes being left open. It is known that the distribution of currents in such a network is found to be multifractal in the sense that different moments of the distribution scale with different exponents.
For resistance scaling analysis two methods may be adopted, either (a) to obtain the distribution of current over the entire network and measure the energy dissipated in the system or (b) to simplify the network and obtain a closed-form solution. This second method has been used rigorously by us for resistance scaling and the results obtained by this method match those from the current distribution method. The moments of the current in a n-simplex are:
$$S_r^a(I_1,I_2,\mathrm{}..I_n)=\underset{p}{}|I_p|^a,$$
(3)
where $`I_p`$ is the current in the pth bond and $`p`$ goes from $`1`$ to $`n(n1)/2`$; $`S_r^a`$ is the cumulant for an arbitrary exponent $`a`$. The currents flowing in at the external nodes of a $`n`$-simplex are represented by $`I_1`$, $`I_2`$,…………$`I_n`$ respectively (see figure 1) with the condition $`I_1+I_2+\mathrm{}\mathrm{}..I_n=0`$. A scaling factor independent of $`I_1`$, $`I_2`$ …….. $`I_n`$ can then be defined as:
$$\lambda (a)=\frac{S_{r+1}^a(I_1,I_2,\mathrm{}\mathrm{}..I_n)}{S_r^a(I_1,I_2,\mathrm{}\mathrm{}\mathrm{}I_n)}.$$
(4)
Note that $`\lambda (a)`$ is related to the fractal scaling exponent $`D(a)`$. For a fractal with a resistance scaling parameter of 2, the rth generation length scales as $`L_r=2^r`$. Using the definitions $`S_r^a(I_1,I_2,\mathrm{}\mathrm{}\mathrm{}I_n)=L_r^{D(a)}`$ and $`S_r^a(I_1,I_2,\mathrm{}\mathrm{}\mathrm{}I_n)\lambda ^r(a)`$, we get:
$$D(a)=\frac{\mathrm{ln}\lambda (a)}{\mathrm{ln}2}.$$
(5)
The case $`a=0`$ determines the fractal dimension of the simplex because $`\lambda (0)`$ is simply the ratio of number of bonds in successive order of the n-simplex lattice; $`a=2`$ measures the heat loss in the network and gives resistance scaling. It has been shown \[19-20\] that:
$$R(L)L^{\beta _l}(L1)$$
(6)
where $`\beta _l`$ is an exponent controlling the transport properties.
In disordered material, the elastic scattering of the carriers at impurities leads to the random conductance or resistance fluctuation. The fluctuation arises from the interference of the scattered waves, and they are random. The magnitude of the resistance noise spectrum (flicker noise 1/f) depends on a new exponent, $`b`$, pertaining to the fractal lattice. This exponent (corresponding to $`a=4`$) is a member of infinite number of exponents required to characterize the fractal lattice . The exact reason as to why this fluctuation occurs is unknown though it is believed that it appears in response to changes in many extrinsic parameters such as the carrier density, the applied measuring current, external electric fields and external magnetic fields. The spectrum of resistance fluctuation is given by
$$S_R(w)=e^{iwt}<R(t)R(0)>𝑑t.$$
(7)
The exponent $`b`$ associated with the scaling behaviour of normalized noise is given by :
$$\rho _R=\mathrm{\S }_R/R^2L^b(L>>1)$$
(8)
As long as each bond resistance fluctuates independently with the same spectrum, the explicit frequency dependence can be discarded. The upper and lower bounds of $`b`$ \[19-20\] are given by:
$$\beta _L<b<\overline{d}.$$
(9)
The paper is organised as follows: In Section 2 we derive a closed-form solution to calculate $`\beta _L`$ for any value of n. In Section 3 we use the current cumulant method to calculate the noise exponent for n-simplex. The paper ends with a brief discussion on the bounds proposed and comparison with our results with experimental data.
## 2 Calculation of $`\beta _L`$ associated with resistance scaling
In this section we propose a simple method of calculating $`\beta _l`$ for any n-simplex. Consider a fixed current $`I_1`$ entering at one of the external nodes of the lattice, and leaving from another, all the remaining external nodes being left open ($`I_2=I_3=\mathrm{}I_n=0`$). We calculate the equivalent resistance between these two external nodes and establish a recursion relation between rth and (r+1)th order lattices and use the Real Space Renormalization Group Technique to find the exponents \[18,21-22\].
¿From the symmetry properties of the simplex, it is apparent that all (n-2) external nodes apart from those through which current enters and leaves are at equipotential. Redrawing just those bonds through which currents flow, we have (n-1) parallel paths for current to flow. Of these paths, one offers unit resistance and each of the others offer twice the unit resistance (since they include two resistances in series). Hence the equivalent resistance is given by:
$$\frac{1}{R_E}=\frac{1}{R}+[\frac{1}{2R}+\frac{1}{2R}+\mathrm{}(n2)\mathrm{terms}]$$
(10)
where R is the unit resistance and the square bracket contains exactly (n-2) identical terms. This directly leads to:
$$R_E=\frac{2R}{n}.$$
(11)
Now if we consider a star of n-branches, each offering a resistance of $`R/n`$, the effective resistance between any two external nodes through which current flows will be $`2R/n`$ as they are in series and all other nodes being left open. It is then straight forward to show using these transformation for $`n`$-simplex lattice that the following scaling holds good:
$$\lambda (2)\frac{R(2L)}{R(L)}=\frac{R_{r+1}}{R_r}=\frac{n+2}{n}.$$
(12)
¿From equation (11) we know that for any n-simplex the equivalent resistance of first order is:
$$R_{E1}=\frac{2R}{n},$$
(13)
combined with equation (12) gives the equivalent resistance of rth order as:
$$R_{Er}=\frac{2(n+2)^{r1}}{n^r}.$$
(14)
Thus we see how, by merely knowing the simplex one can calculate the equivalent resistance of any iteration. No long winded applications of Kirchoff’s Laws are required to obtain the resistance scaling. The exponent $`\beta _L`$ is related to $`\lambda (2)`$ by:
$$\beta _L=\frac{\mathrm{ln}(1/\lambda (2))}{\mathrm{ln}2}.$$
(15)
## 3 Calculation of the Flicker noise exponent on $`n`$-simplex
It has been shown that the $`4^{th}`$ moment of current distribution is associated with the noise exponent .
Assuming that the cumulant $`S_r^4(I_1,I_2,\mathrm{}..I_n)`$ can be expressed as a homogeneous polynomial of degree $`4`$, the most general polynomial is a linear combination of $`P_1^4`$, $`P_1^2P_2`$, $`P_1P_3`$, $`P_4`$ and $`P_2^2`$. These polynomials are defined as
$`P_1`$ $`=`$ $`I_1+I_2+I_3+\mathrm{}\mathrm{}\mathrm{}\mathrm{}I_n`$
$`P_2`$ $`=`$ $`I_1^2+I_2^2+I_3^2+\mathrm{}\mathrm{}\mathrm{}\mathrm{}I_n^2`$
$`P_3`$ $`=`$ $`I_1^3+I_2^3+I_3^3+\mathrm{}\mathrm{}\mathrm{}\mathrm{}I_n^3\mathrm{and}`$
$`P_4`$ $`=`$ $`I_1^4+I_2^4+I_3^4+\mathrm{}\mathrm{}\mathrm{}\mathrm{}I_n^4`$
However in present case $`P_1=0`$ due to current conservation. Hence $`S_r^4(I_1,I_2,\mathrm{}..I_n)`$ can be written as
$$S_r^4(I_1,I_2,\mathrm{}..I_n)=A_rP_2^2(I_1,I_2,\mathrm{}..I_n)+B_rP_4(I_1,I_2,\mathrm{}..I_n)$$
(16)
The next step is to determine $`S_{r1}^4(I_1,I_2,\mathrm{}\mathrm{}\mathrm{}I_n)`$. To establish a recursion relation between
$`S_r^4(I_1,I_2,\mathrm{}\mathrm{}.I_n)`$ and $`S_{r1}^4(I_1,I_2,\mathrm{}\mathrm{}I_n)`$ we obtained current distribution in an $`n`$-simplex lattice at each node. It is easy to see that the current distribution at each node is
$$\underset{k=1}{\overset{n}{}}\left[I_k+\frac{1}{n1}\underset{\genfrac{}{}{0pt}{}{j=1}{kj}}{\overset{n}{}}I_kI_j\right]$$
by current conservation. In figure 1 we have shown the current along each bond. Therefore, we can write
$$S_r^4(I_1,I_2,\mathrm{}..I_n)=S_{r1}^4(I)+S_{r1}^4(II)+\mathrm{}\mathrm{}\mathrm{}\mathrm{}..+S_{r1}^4(n)$$
where $`S_{r1}(I)`$, $`S_{r1}(II),\mathrm{}\mathrm{}..`$ are the current cumulants of $`(r1)`$th order of $`n`$-simplex lattice. $`I`$, $`II,\mathrm{}\mathrm{}`$ represents shaded region in figure 1. Above equation can be expressed as
$`S_r^4(I_1,I_2,\mathrm{}..I_n)`$ $`=`$ $`A_{r1}P_2^2(I_1,I_2,\mathrm{}..I_n)+B_{r1}P_4(I_1,I_2,\mathrm{}..I_n)`$ (17)
$`+`$ $`A_{r1}P_4(I_1,I_2,\mathrm{}..I_n)+B_{r1}P_2^2(I_1,I_2,\mathrm{}..I_n)`$
which establish the transformation relation between $`r`$ and $`(r1)`$th order. Comparing equations (16) and (17) we obtain the recursion relation between the $`n`$-simplex lattices of the $`r`$ and $`(r1)`$ th order.
$$\left(\begin{array}{c}A_r\\ B_r\end{array}\right)=\frac{1}{n^4}\left(\begin{array}{cc}(n^3+2)n& n^2(n+1)^2\\ 6& (2n+3)n\end{array}\right)\left(\begin{array}{c}A_{r1}\\ B_{r1}\end{array}\right)$$
(18)
The eigenvalues corresponding to the transformation matrix for the $`n`$-simplex lattice are given by
$$\lambda _n^\pm (a=4)=\frac{(n^3+2n+5)n\pm n(n+1)\sqrt{n^42n^3n^2+2n+25}}{2n^4}$$
and the fractal scaling exponent corresponding to largest eigenvalue is
$$D(a=4)=\frac{\mathrm{ln}\lambda _n^+(a=4)}{\mathrm{ln}2}.$$
## 4 Discussion
We have seen that various moments of branch current give rise to different exponents, namely exponent $`b`$ associated with the noise amplitude, $`\beta _L`$ associated with resistance scaling and $`\overline{d}`$ associated with the mass of the fractal. The relation between $`D(4)`$ obtained in Section 3 and the exponent $`b`$ for normalised noise is as follows:
$$\frac{S_R}{R_r^2}\rho _RL_r^b.$$
(19)
Now, $`S_{R_r}_p|I_p|^4`$
$$\frac{_p(I_p)_r^4}{R_r^2}\rho _{R_r}L_r^b.$$
(20)
This gives:
$$\frac{\rho _{R_{r+1}}}{\rho _{R_r}}=\lambda (\rho _{R_r})2^b$$
(21)
or,
$$b=\frac{\mathrm{ln}(1/\lambda (\rho _R))}{\mathrm{ln}2}.$$
(22)
Now,
$$\lambda (\rho _R)=\frac{_p(I_p)_{r+1}^4}{_p(I_p)_r^4}\frac{R_r^2}{R_{r+1}^2}$$
(23)
or,
$$\lambda (\rho _R)=\frac{A_{r+1}}{A_r}\left(\frac{R_r}{R_{r+1}}\right)^2$$
(24)
$$\lambda (\rho _R)=\frac{A_{r+1}}{A_r}\left(\frac{1}{\lambda (R)}\right)^2.$$
(25)
Substituting this in equation(8) gives:
$$b=\frac{\mathrm{ln}\left(\lambda ^2(R)\times \frac{A_r}{A_{r+1}}\right)}{\mathrm{ln}2}.$$
(26)
This expression gives the respective values of the exponent $`b`$ as $`1.1844`$, $`1.0629`$ and $`0.9269`$ for the 3, 4 and 5-simplex respectively. It is clear that the inequality $`\overline{d}b\beta _L`$ is satisfied in each of the three cases.
In the limit $`n`$ goes to infinity $`\lambda (2)=(n+2)/n`$ goes to 1. This is due to the fact that a large number of parallel equi-resistance paths are available for current flow. Such a large number of paths are available that in going from one order to the next we are in effect not altering the equivalent resistance. The exponent $`\beta _L`$ decreases in magnitude as we go to higher dimension, implying that resistance becomes less dependent on the length of the fractal.
With regard to flicker noise, we have seen that the scaling relation becomes increasingly complex as the order of simplex is increased. The noise versus resistance exponent $`Q`$ is defined by the following:
$$Q=2+\frac{t}{k}$$
(27)
where $`t`$ and $`k`$ are given by:
$$R(\mathrm{\Delta }p)^t$$
(28)
and
$$\frac{S_R}{R^2}(\mathrm{\Delta }p)^k.$$
(29)
The experimental measurements \[19-20\] of $`t`$ and $`k`$ were made on 2d-carbon-vax mixtures and found to be $`2.3\pm 0.4`$ and $`5\pm 1`$ respectively. The direct plot of $`S_R`$ versus $`R`$ leads to $`S_rR^Q`$ where $`Q=3.7\pm 0.2`$. The value we obtain for $`Q`$ is in agreement with this as is clear from Table 1.
However, similar measurements on two dimensional films and metallic films have given values of $`Q`$ differing from what we predict. Perhaps instead of taking the $`n`$-simplex lattice, if one considered a 2-d Sierpinski gasket \[13-14\] better results could be expected. For all $`n>a`$, there is a finite dimension matrix whose largest eigenvalue will give the characteristic exponents. The matrix elements are function of $`n`$ and hence eigenvalues will be the well defined function of $`n`$. But for $`n<a`$ the result will be obtained by smaller matrix. The generalization to higher value of $`a`$ and rescaling factor $`b>2`$ is under progress.
## Acknowledgements
We would like to thank Yashwant Singh and Deepak Dhar for many helpful discussions. Financial assistance from the Department of Science and Technology India is acknowledged. One of us(SK) would like to thank INSA-DFG for financial support.
## References
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9. B. Docut and R. Rammal, Phy. Rev. Lett., $`\mathrm{𝟓𝟓}`$, 1148, (1985)
10. I. Zivic, S. Milosevic, and H. E. Stanley, Phy. Rev. E, $`\mathrm{𝟒𝟕}`$, 2340 (1990)
11. D. C. Hong, and H. E. Stanley, J. Phys. A, $`\mathrm{𝟏𝟔}`$, L525 (1983)
12. H. J. Herrmann, D. C. Honmg and H.E. Stanley, J. Phys. A, $`\mathrm{𝟏𝟕}`$, L261 (1984)
13. B. B. Mandelbrot, The Fractal Geometry of Nature, (Freeman, NY 1982)
14. L. Pietronero and E. Tosatti (eds) Fractals in Physics, (North Holland: Amsterdam 1986)
15. D. Kessler, J. Koplick and H. Levine, Adv. in Phys., $`\mathrm{𝟑𝟕}`$, 255 (1988)
16. D. R. Nelson and M. E. Fisher, Ann. Phys., $`\mathrm{𝟗𝟏}`$, 266 (1975)
17. D. Dhar, J. Math. Phys., $`\mathrm{𝟏𝟖}`$, 577 (1977)
18. S. Roux and C. D. Mitescu, Phy. Rev. B , $`\mathrm{𝟑𝟓}`$, 898 (1987)
19. R. Rammal, C. Tannous and A.M.S. Tremblay, Phy. Rev. Lett., $`\mathrm{𝟓𝟒}`$, 1718 (1985)
20. R. Rammal C. Tannous and A. M. S. Tremblay, Phy. Rev. A, $`\mathrm{𝟑𝟏}`$, 2662 (1985)
21. P. Alstrom, D. Stassinopoulos and H. E. Stanley, Physica A, $`\mathrm{𝟏𝟓𝟑}`$, 20 (1988)
22. P. Y. Tong and K. W. Yu, Phys. Lett A, $`\mathrm{𝟏𝟔𝟎}`$, 293 (1991)
Figure 1 : Schematic representation of $`n`$-simplex lattice ($`n=6`$). The current along each bond has been shown.
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# Decoherence and the Final Pointer Basis.
## I Introduction.
We will demonstrate that for a wide set of quantum systems the quantum regime can be consider as the transient phase while the final classical equilibrium regime is the permanent state. We will find a basis where exact matrix decoherence appears for these final states. Therefore we will find a set of final intrinsically consistent histories.
## II Decoherence.
### A Decoherence in the energy.
Let us consider a closed and isolated quantum system with $`N+1`$ dynamical variables and a Hamiltonian endowed with a continuous spectrum and just one bounded ground state. So the discrete part of the spectrum of $`H`$ has only one value $`\omega _0`$ and the continuous spectrum is let say $`0\omega <\mathrm{}`$. Eventually we will give the collective name $`x`$ to both $`\omega _0`$ and $`\omega .`$ Let us assume that it is possible to diagonalize the Hamiltonian $`H`$, together with $`N`$ observables $`O_i`$ ($`i=1,\mathrm{},N)`$. The operators ($`H`$, $`O_1`$,…,$`O_N`$) form a complete set of commuting observables (CSCO). For simplicity we also assume a discrete spectrum for the $`N`$ observables $`O_i`$. Therefore we write
$$H=\omega _0\underset{m}{}|\omega _0,m\omega _0,m|+_0^{\mathrm{}}\omega \underset{m}{}|\omega ,m\omega ,m|d\omega $$
(1)
where $`\omega _0<0`$ is the energy of the ground state, and $`m\{m_1,\mathrm{},m_N\}`$ labels a set of discrete indexes which are the eigenvalues of the observables $`O_1`$,…,$`O_N`$. $`\{|\omega _0,m,|\omega ,m\}`$ is a basis of simultaneous generalized eigenvectors of the CSCO:
$`H|\omega _0,m`$ $`=`$ $`\omega _0|\omega _0,m,H|\omega ,m=\omega |\omega ,m,`$
$`O_i|\omega _0,m`$ $`=`$ $`m_i|\omega _0,m,O_i|\omega ,m=m_i|\omega ,m.`$
The most general observable that we are going to consider in our model reads:
$`O`$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}O(\omega _0)_{mm^{}}|\omega _0,m\omega _0,m^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega O(\omega )_{mm^{}}|\omega ,m\omega ,m^{}|+`$ (5)
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega O(\omega ,\omega _0)_{mm^{}}|\omega ,m\omega _0,m^{}|+`$
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}O(\omega _0,\omega ^{})_{mm^{}}|\omega _0,m\omega ^{},m^{}|+`$
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑\omega 𝑑\omega ^{}O(\omega ,\omega ^{})_{mm^{}}|\omega ,m\omega ^{},m^{}|,`$
where $`O(\omega )_{mm^{}}`$, $`O(\omega ,\omega _0)_{mm^{}}`$, $`O(\omega _0,\omega )_{mm^{}}`$ and $`O(\omega ,\omega ^{})_{mm^{}}`$ are ordinary functions of the real variables $`\omega `$ and $`\omega ^{}`$(these functions must have some mathematical properties in order to develop the theory; these properties are listed in paper ). We will say that these observables belong to a space $`𝒪`$. This space has the basis $`\{|\omega _0,mm^{})`$, $`|\omega ,mm^{})`$, $`|\omega \omega _0,mm^{})`$, $`|\omega _0\omega ^{},mm^{})`$, $`|\omega \omega ^{},mm^{})\}`$:
$`|\omega _0,mm^{})|\omega _0,m\omega _0,m^{}|,|\omega ,mm^{})|\omega ,m\omega ,m^{}|,`$
$$|\omega \omega _0,mm^{})|\omega ,m\omega _0,m^{}|,|\omega _0\omega ^{},mm^{})|\omega _0,m\omega ^{},m^{}|,$$
(6)
$`|\omega \omega ^{},mm^{})|\omega ,m\omega ^{},m^{}|`$
The quantum states $`\rho `$ are measured by the observables just defined, computing the mean values of these observable in the quantum states, i. e. in the usual notation: $`O_\rho =Tr(\rho ^{}O)`$ . These mean values, generalized as in paper , can be considered as linear functionals $`\rho `$ (mapping the vectors $`O`$ on the real numbers), that we can call $`(\rho |O)`$ . In fact, this is a generalization of the usual mean value definition. Then $`\rho 𝒮𝒪^{^{}},`$ where $`𝒮`$ is a convenient convex set contained in $`𝒪^{^{}}`$, the space of linear functionals over $`𝒪`$ , . The basis of $`𝒪^{}`$ (that can also be considered as the co-basis of $`𝒪)`$ is $`\{(\omega _0,mm^{}|`$, $`(\omega ,mm^{}|`$, $`(\omega \omega _0,mm^{}|`$, $`(\omega _0\omega ^{},mm^{}|`$, $`(\omega \omega ^{},mm^{}|\}`$ defined as functionals by the equations:
$`(\omega _0,mm^{}|\omega _0,nn^{})=\delta _{mn}\delta _{m^{}n^{}},(\omega ,mm^{}|\eta ,nn^{})=\delta (\omega \eta )\delta _{mn}\delta _{m^{}n^{}},`$
$`(\omega \omega _0,mm^{}|\eta \omega _0,nn^{})=\delta (\omega \eta )\delta _{mn}\delta _{m^{}n^{}},`$
$`(\omega _0\omega ^{},mm^{}|\omega _0\eta ^{},nn^{})=\delta (\omega ^{}\eta ^{})\delta _{mn}\delta _{m^{}n^{}},`$
$$(\omega \omega ^{},mm^{}|\eta \eta ^{},nn^{})=\delta (\omega \eta )\delta (\omega ^{}\eta ^{})\delta _{mn}\delta _{m^{}n^{}}.$$
(7)
and all other $`(.|.)`$ are zero. Then, a generic quantum state reads:
$`\rho `$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}\overline{\rho (\omega _0)}_{mm^{}}(\omega _0,mm^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega \overline{\rho (\omega )}_{mm^{}}(\omega ,mm^{}|+`$ (11)
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega \overline{\rho (\omega ,\omega _0)}_{mm^{}}(\omega \omega _0,mm^{}|+`$
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega ^{}\overline{\rho (\omega _0,\omega ^{})}_{mm^{}}(\omega _0\omega ^{},mm^{}|+`$
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega {\displaystyle _0^{\mathrm{}}}d\omega ^{}\overline{\rho (\omega ,\omega ^{})}_{mm^{}}(\omega \omega ^{},mm^{}|`$
where
$`\overline{\rho (\omega ,\omega _0)}_{mm^{}}=\rho (\omega _0,\omega )_{m^{}m},\overline{\rho (\omega ,\omega ^{})}_{mm^{}}=\rho (\omega ^{},\omega )_{m^{}m},`$
and $`\overline{\rho (\omega _0)}_{mm}`$ and $`\overline{\rho (\omega )}_{mm}`$ are real and non negative satisfying the total probability condition
$$(\rho |I)=\underset{m}{}\rho (\omega _0)_{mm}+\underset{m}{}_0^{\mathrm{}}𝑑\omega \rho (\omega )_{mm}=1,$$
(12)
where $`I=_m|\omega _0,m\omega _0,m|+_0^{\mathrm{}}𝑑\omega _m|\omega ,m\omega ,m|`$ is the identity operator in $`𝒪`$. Eq. (12) is the extension to state functionals of the usual condition $`Tr\rho ^{}=1`$, used when $`\rho `$ is a density operator.
The time evolution of the quantum state $`\rho `$ reads:
$`\rho (t)`$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}\overline{\rho (\omega _0)}_{mm^{}}(\omega _0,mm^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega \overline{\rho (\omega )}_{mm^{}}(\omega ,mm^{}|+`$ (16)
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega \overline{\rho (\omega ,\omega _0)}_{mm^{}}e^{i(\omega \omega _0)t}(\omega \omega _0,mm^{}|+`$
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega ^{}\overline{\rho (\omega _0,\omega ^{})}_{mm^{}}e^{i(\omega _0\omega ^{})t}(\omega _0\omega ^{},mm^{}|+`$
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega {\displaystyle _0^{\mathrm{}}}d\omega ^{}\overline{\rho (\omega ,\omega ^{})}_{mm^{}}e^{i(\omega \omega ^{})t}(\omega \omega ^{},mm^{}|`$
As we only measure mean values of observables in quantum states, i. e.:
$`O_{\rho (t)}`$ $`=`$ $`(\rho (t)|O)=`$ (17)
$`=`$ $`{\displaystyle \underset{mm^{}}{}}\overline{\rho (\omega _0)}_{mm^{}}O(\omega _0)_{mm^{}}+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega \overline{\rho (\omega )}_{mm^{}}O(\omega )_{mm^{}}+`$ (21)
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega \overline{\rho (\omega ,\omega _0)}_{mm^{}}e^{i(\omega \omega _0)t}O(\omega ,\omega _0)_{mm^{}}+`$
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}\overline{\rho (\omega _0,\omega ^{})}_{mm^{}}e^{i(\omega _0\omega ^{})t}O(\omega _0,\omega ^{})_{mm^{}}+`$
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega {\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}\overline{\rho (\omega ,\omega ^{})}_{mm^{}}e^{i(\omega \omega ^{})t}O(\omega ,\omega ^{})_{mm^{}},`$
using the Riemann-Lebesgue theorem we obtain the limit, for all $`O𝒪`$
$$\underset{t\mathrm{}}{lim}O_{\rho (t)}=O_\rho _{}$$
(22)
where we have introduced the diagonal asymptotic or equilibrium state functional
$$\rho _{}=\underset{mm^{}}{}\overline{\rho (\omega _0)}_{mm^{}}(\omega _0,mm^{}|+\underset{mm^{}}{}_0^{\mathrm{}}d\omega \overline{\rho (\omega )}_{mm^{}}(\omega ,mm^{}|$$
(23)
Therefore, in a weak sense we have:
$$W\underset{t\mathrm{}}{lim}\rho (t)=\rho _{}$$
(24)
Thus, any quantum state goes weakly to a linear combination of the energy diagonal states $`(\omega _0,mm^{}|`$ and $`(\omega ,mm^{}|`$ (the energy off-diagonal states $`(\omega \omega _0,mm^{}|`$, $`(\omega _0\omega ^{},mm^{}|`$ and $`(\omega \omega ^{},mm^{}|`$ are not present in $`\rho _{}`$). This is the case if we observe and measure the system evolution with any possible observable of space $`𝒪`$. Then, from the observational point of view, we have decoherence of the energy levels, even that, from the strong limit point of view the off-diagonal terms never vanish, they just oscillate, since we cannot directly use the Riemann-Lebesgue theorem in the operator equation (16).
### B Decoherence in the other ”momentum” dynamical variables.
Having established the decoherence in the energy levels we must consider the decoherence in the other dynamical variables $`O_i`$, of the CSCO where we are working. We will call these variables ”momentum variables”. For the sake of simplicity we will consider, as in the previous section, that the spectra of these dynamical variables are discrete. As the expression of $`\rho _{}`$ given in eq. (23) involve only the time independent components of $`\rho (t)`$, it is impossible that a different decoherence process take place to eliminate the off-diagonal terms in the remaining $`N`$ dynamical variables. Therefore, the only thing to do is to find if there is a basis where the off-diagonal components of $`\rho (\omega _0)_{mm^{}}`$ and $`\rho (\omega )_{mm^{}}`$ vanish at any time before the equilibrium is reached.
Let us consider the following change of basis
$`|\omega _0,r={\displaystyle \underset{m}{}}U(\omega _0)_{mr}|\omega _0,m,|\omega ,r={\displaystyle \underset{m}{}}U(\omega )_{mr}|\omega ,m,`$
where $`r`$ and $`m`$ are short notations for $`r\{r_1,\mathrm{},r_N\}`$ and $`m\{m_1,\mathrm{},m_N\}`$, and $`\left[U(x)^1\right]_{mr}=\overline{U(x)}_{rm}`$ ($`x`$ denotes either $`\omega _0<0`$ or $`\omega ^+`$).
The new basis $`\{|\omega _0,r,|\omega ,r\}`$ verify the generalized orthogonality conditions
$`\omega _0,r|\omega _0,r^{}`$ $`=`$ $`\delta _{rr^{}},\omega ,r|\omega ^{},r^{}=\delta (\omega \omega ^{})\delta _{rr^{}},`$
$`\omega _0,r|\omega ,r^{}`$ $`=`$ $`\omega ,r|\omega _0,r^{}=0.`$
As $`\overline{\rho (\omega _0)}_{mm^{}}=\rho (\omega _0)_{m^{}m}`$ and $`\overline{\rho (\omega )}_{mm^{}}=\rho (\omega )_{m^{}m}`$, it is possible to choose $`U(\omega _0)`$ and $`U(\omega )`$ in such a way that the off-diagonal parts of $`\rho (\omega _0)_{rr^{}}`$ and $`\rho (\omega )_{rr^{}}`$ vanish, i.e.
$`\rho (\omega _0)_{rr^{}}=\rho _r(\omega _0)\delta _{rr^{}},\rho (\omega )_{rr^{}}=\rho _r(\omega )\delta _{rr^{}}.`$
Therefore, there is a final pointer basis for the observables given by $`\{|\omega _0,rr^{})`$, $`|\omega ,rr^{})`$, $`|\omega \omega _0,rr^{})`$, $`|\omega _0\omega ^{},rr^{})`$, $`|\omega \omega ^{},rr^{})\}`$ and defined as in eq. (6). The corresponding final pointer basis for the states $`\{(\omega _0,rr^{}|`$, $`(\omega ,rr^{}|`$, $`(\omega \omega _0,rr^{}|`$, $`(\omega _0\omega ^{},rr^{}|`$, $`(\omega \omega ^{},rr^{}|\}`$ diagonalizes the time independent part of $`\rho (t)`$ and therefore it diagonalizes the final state $`\rho _{}`$
$$\rho _{}=W\underset{t\mathrm{}}{lim}\rho (t)=\underset{r}{}\rho _r(\omega _0)(\omega _0,rr|+\underset{r}{}_0^{\mathrm{}}d\omega \rho _r(\omega )(\omega ,rr|.$$
(25)
Now we can define the final exact pointer observables
$$P_i=\underset{r}{}P_r^i(\omega _0)|\omega _0,r\omega _0,r|+_0^{\mathrm{}}𝑑\omega \underset{r}{}P_r^i(\omega )|\omega ,r\omega ,r|.$$
(26)
As $`H`$ and $`P_i`$ are diagonal in the basis $`\{|\omega _0,r`$, $`|\omega ,r\}`$, the set $`\{H,P_i,\mathrm{}P_N\}`$ is precisely the complete set of commuting observables (CSCO) related to this basis, where $`\rho _{}`$ is diagonal in the corresponding co-basis for states. For simplicity we define the operators $`P_i`$ such that $`P_r^i(\omega _0)=P_r^i(\omega )=r_i`$, thus
$$P_i|\omega _0,r=r_i|\omega _0,r,P_i|\omega ,r=r_i|\omega ,r.$$
(27)
Therefore $`\{|\omega _0,r`$, $`|\omega ,r\}`$ is the observers pointer basis where there is a perfect decoherence in the corresponding state co-basis. Moreover the generalized states $`(\omega _0,rr|`$ and $`(\omega ,rr|`$ are constants of the motion, and therefore these exact pointer observables have a constant statistical entropy and will be ”at the top of the list” of Zurek’s ”predictability sieve” .
Therefore:
i.- Decoherence in the energy is produced by the time evolution.
ii.- Decoherence in the other dynamical variables can be seen if we choose an adequate basis, namely the final pointer basis.
Our main result is eq. (25): When $`t\mathrm{}`$ then $`\rho (t)\rho _{}`$ and in this state the dynamical variables $`H,P_1,\mathrm{},P_N`$ are well defined. Therefore the eventual conjugated variables to these momentum variables (namely: configuration variables, if they exists) are completely undefined.
In fact, calling by $`𝕃_i`$ the generator of the displacements along the eventual configuration variable conjugated to $`P_i`$, we have $`(𝕃_i\rho _{}|O)=(\rho _{}|𝕃_i^{}O)=(\rho _{}|[P_i,O])=0`$ for all $`O𝒪`$. Then $`\rho _{}`$ is homogeneous in these configuration variables.
From the preceding section we may have the feeling that the process of decoherence must be found in all the physical systems, and therefore, all of them eventually would become classical when $`\mathrm{}0`$. It is not so as explained in .
## III The classical equilibrium limit.
### A Expansion in sets of classical motions.
In this section we will use the Wigner integrals that introduce an isomorphism between quantum observables $`O`$ and states $`\rho `$ and their classical analogues $`O^W(q,p)`$ and $`\rho ^W(q,p)`$ :
$`O^W(q,p)`$ $`=`$ $`{\displaystyle 𝑑\lambda q\frac{\lambda }{2}|O|q+\frac{\lambda }{2}\mathrm{exp}(\frac{i\lambda p}{\mathrm{}})}`$ (28)
$`\rho ^W(q,p)`$ $`=`$ $`{\displaystyle \frac{1}{\pi \mathrm{}}}{\displaystyle }d\lambda (\rho ||q+\lambda q\lambda |)\mathrm{exp}({\displaystyle \frac{2i\lambda p}{\mathrm{}}}).`$ (29)
It is possible to prove that $`𝑑q𝑑p\rho ^W(q,p)=(\rho |I)=1`$, but $`\rho ^W`$ is not in general non negative. It is also possible to deduce that
$$(\rho ^W|O^W)=𝑑q𝑑p\rho ^W(q,p)O^W(q,p)=(\rho |O),$$
(30)
and therefore to the mean value in the classical Liouville space it corresponds the mean value in the quantum Liouville space.
Moreover, calling by $`L`$ the classical Liouville operator, and by $`𝕃`$ the quantum Liouville-Von Neumann operator, we have
$$L\left[\rho ^W(q,p)\right]=\left[𝕃\rho \right]^W(q,p)+O(\mathrm{}),$$
(31)
where $`L\rho ^W(q,p)=i\{H^W(q,p),\rho ^W(q,p)\}_{PB}`$ and
$$(𝕃\rho |O)=(\rho |[H,O]).$$
(32)
Finally, if $`O=O_1O_2`$, where $`O_1`$ and $`O_2`$ are two quantum observables, we have
$$O^W(q,p)=O_1^W(q,p)O_2^W(q,p)+O(\mathrm{}).$$
(33)
We will prove that the distribution function $`\rho _{}^W(q,p)`$, that corresponds to the state functional $`\rho _{}`$ via the Wigner integral is a non negative function of the classical constants of the motion, in our case $`H^W(q,p)`$, $`P_1^W(q,p)`$,…, $`P_N^W(q,p),`$ obtained from the corresponding quantum operators $`H`$, $`P_1`$,…, $`P_N`$.
From eq. (25) we have:
$$\rho _{}=W\underset{t\mathrm{}}{lim}\rho (t)=\underset{r}{}\rho _r(\omega _0)(\omega _0,rr|+\underset{r}{}_0^{\mathrm{}}d\omega \rho _r(\omega )(\omega ,rr|,$$
(34)
so we must compute:
$$\rho _{\omega r}^W(q,p)\pi ^1(\omega ,rr||q+\lambda q\lambda |)e^{2ip\lambda }d\lambda $$
(35)
We know from section II. C, (or we can prove directly from eqs.(25-27)) that
$`(\omega _0,rr|H^n)`$ $`=`$ $`\omega _0^n,(\omega ,rr|H^n)=\omega ^n,`$ (36)
$`(\omega _0,rr|P_i^n)`$ $`=`$ $`r_i^n,(\omega ,rr|P_i^n)=r_i^n,i=1,\mathrm{},N`$ (37)
for $`n=0,1,2,\mathrm{}`$ Using the relation (33) between quantum and classical products of observables and relation (30) between quantum and classical mean values, in the limit $`\mathrm{}0`$ (we will consider that we always take this limit when we refer to classical equations below) we deduce that the characteristic property of the distribution $`\rho _{\omega r}^W(q,p)`$, that corresponds to the state functional $`(\omega ,rr|`$, is:
$$\rho _{\omega r}^W(q,p)[H^W(q,p)]^n𝑑q𝑑p=\omega ^n,\rho _{\omega r}^W(q,p)[P_i^W(q,p)]^n𝑑q𝑑p=r_i^n,$$
(38)
for any natural number $`n.`$ Thus $`\rho _{\omega r}^W(q,p)`$ must be the functional
$$\rho _{\omega r}^W(q,p)=\delta (H^W(q,p)\omega )\delta (P_1^W(q,p)r_1)\mathrm{}\delta (P_N^W(q,p)r_N).$$
(39)
For the distribution $`\rho _{\omega _0r}^W(q,p)`$ corresponding to the state functional $`(\omega _0,rr|`$, we obtain
$$\rho _{\omega _0r}^W(q,p)=\delta (H^W(q,p)\omega _0)\delta (P_1^W(q,p)r_1)\mathrm{}\delta (P_N^W(q,p)r_N).$$
(40)
Therefore, going back to eq. (34) and since the Wigner relation is linear, we have:
$$\rho _{}^W(q,p)=\underset{r}{}\rho _r(\omega _0)\rho _{\omega _0r}^W(q,p)+\underset{r}{}_0^{\mathrm{}}𝑑\omega \rho _r(\omega )\rho _{\omega r}^W(q,p).$$
(41)
Also we obtain $`\rho _{}^W(q,p)0`$, because $`\rho _r(\omega _0)`$ and $`\rho _r(\omega )`$ are non negative.
Therefore, the classical state $`\rho _{}^W(q,p)`$ is a linear combination of the generalized classical states $`\rho _{xr}^W(q,p)`$ (where $`x`$ is either $`\omega _0`$ or $`\omega `$), having well defined values $`x`$, $`r_1`$,…, $`r_N`$ of the classical observables $`H^W(q,p)`$, $`P_1^W(q,p)`$,…, $`P_N^W(q,p)`$ and the corresponding classical canonically conjugated variables completely undefined since $`\rho _{xr}^W(q,p)`$ is not a function of these variables. So we reach, in the classical case, to the same conclusion than in the quantum case (see end of subsection 2. 2). But now all the classical canonically conjugated variables $`a_0,a_1,\mathrm{},a_N`$ do exist since they can be found solving the corresponding Poisson brackets differential equations. We can also expand the densities given in eqs. (39-41) in terms of classical motions as shown in .
## IV Conclusion.
i.- We have shown that the quantum state functional $`\rho (t)`$ evolves to a diagonal state $`\rho _{}`$.
ii.- This quantum state $`\rho _{}`$ has $`\rho _{}^W(q,p)`$ as its corresponding classical density.
iii.- This classical density can be decomposed in sets of classical motions where $`H^W`$, $`P_1^W`$,…, $`P_N^W`$ remain constant. These motions have origins $`a_0(0),a_1(0),\mathrm{},a_N(0)`$ distributed in an homogeneous way.
iv.- From eqs. (39-41) we obtained that
$`\rho _{}^W(q,p)=f(H^W(q,p),P_1^W(q,p),\mathrm{},P_N^W(q,p))0.`$
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# Causal Boundary Entropy From Horizon Conformal Field Theory
## Abstract
The quantum theory of near horizon regions of spacetimes with classical spatially flat, homogeneous and isotropic Friedman-Robertson-Walker geometry can be approximately described by a two dimensional conformal field theory. The central charge of this theory and expectation value of its Hamiltonian are both proportional to the horizon area in units of Newton’s constant. The statistical entropy of horizon states, which can be calculated using two dimensional state counting methods, is proportional to the horizon area and depends on a numerical constant of order unity which is determined by Planck scale physics. This constant can be fixed such that the entropy is equal to a quarter of the horizon area in units of Newton’s constant, in agreement with thermodynamic considerations.
Black holes possess geometric entropy equal to a quarter of the area of their horizon in units of Newton’s constant, known as the Bekenstein-Hawking entropy . Since the discovery of black hole entropy, many attempts were made to identify its microscopic, statistical mechanics origin. Strominger , and more recently Carlip have argued that the statistical origin of black hole entropy is the ensemble of states of a conformal field theory (CFT) describing fluctuations of two dimensional (2D) horizon surfaces. They have used 2D methods to evaluate the density of states of this theory, showing that their entropy is indeed a quarter of the horizon area. Some attempts were made to identify the horizon CFT , and to extend the results to cosmological de Sitter space .
It is widely accepted that geometric entropy must also be attributed to the horizon of de Sitter space . The argument is that de Sitter horizons are event horizons (as are black hole horizons), and therefore a thermodynamic system crossing them is forever removed from an observer’s ken. Therefore the loss of the system’s entropy must be compensated by an increase of geometric entropy in order for the second law to remain valid. In general, a cosmological horizon is not an event horizon, and a system crossing it is not necessarily forever out of view and so, it may be argued, there is no compelling reason to associate an entropy with a cosmological horizon.
In , I have proposed that geometric entropy has to be attributed to cosmological horizons (or, in general, to causal boundaries), whether or not they are event horizons. I have argued that entropy of quantum fluctuations can be lost if the scale of causal connection becomes smaller than their wavelength, for example, in an inflating universe. This is in violation of the second law. The role of proposed geometric entropy is precisely to restore validity of the second law in such situations.
Here I show that the effective quantum theory of near horizon (NH) regions of spacetimes with classical spatially flat Friedman-Robertson-Walker (FRW) geometry is a 2D CFT, appearing due to huge redshifts suffered by horizon fluctuations which allow only massless fluctuations to survive. The central charge of this CFT $`c=\alpha \frac{A^H}{4\pi G_N}`$, is proportional to the horizon area $`A^H`$ in units of Newton’s constant $`G_N`$, and depends on a numerical constant of order unity $`\alpha `$, which is determined by Planck scale physics. The expectation value of the Hamiltonian of the theory $`L_0=\frac{A^H}{8\pi G_N}`$, is also proportional to the horizon’s area in units of Newton’s constant. The asymptotic density of horizon states, and therefore the horizon entropy $`S^H`$, can be obtained using Cardy’s formula $`S^H=2\pi \sqrt{\frac{c}{6}\left(L_0\frac{c}{24}\right)}`$ (see for detailed considerations about application of Cardy’s formula in this context). The cutoff dependent numerical coefficient $`\alpha `$ can be set to $`\alpha =6`$, such that $`S^H=\frac{A^H}{4G_N}`$, in agreement with thermodynamic considerations.
Horizons of FRW spaces are not necessarily event horizons, so our results indicate that it is the existence of causal boundary which is the source of geometric entropy, and that for a causal boundary to have geometric entropy it is not required to hold information forever. Our results strongly support the conjecture that causal boundaries and not only event horizons have geometric entropies proportional to their area, and therefore strengthen considerably the conclusion based on this conjecture that a certain class of singularities are thermodynamically forbidden (see also ).
Our starting point is the 4D Einstein-Hilbert action
$$S^{(4)}=\frac{1}{16\pi G_N}d^4x\sqrt{g}R^{(4)}+S_m,$$
(1)
$`G_N`$ being Newton’s constant, $`g`$ the determinant of the 4D metric $`g_{\mu \nu }`$, $`R^{(4)}`$ is the 4D Ricci scalar, and $`S_m`$ is the matter action. We consider spatially flat FRW solutions $`ds_4^2=dt^2+a^2(t)dr^2+a^2(t)r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right)`$, with expanding scale factors $`a(t)=a(t_0)\left(\frac{t}{t_0}\right)^\beta `$. The matter has an ideal fluid type energy momentum tensor derived from $`S_m`$, given by $`T_\nu ^\mu =diag(\rho ,p,p,p)`$. The energy density $`\rho `$, and pressure $`p`$ are related by a simple equation of state $`p=w\rho `$, which determines the scale factor expansion rate $`\beta =\frac{2}{3(1+w)}`$. de Sitter space, for which the scale factor expands exponentially, should be considered as the limiting case $`w1`$, $`\beta \mathrm{}`$.
FRW spaces have finite causal connection scale $`R_{CC}`$ , generically called “horizon”. Beyond the horizon local interactions are not effective. This causal connection scale is determined by the Hubble parameter $`H=\frac{\dot{a}}{a}`$, and its derivative: $`R_{CC}^1=\sqrt{\mathrm{Max}[\dot{H}+2H^2,\dot{H}]}`$. In , a covariant expression for $`R_{CC}`$ is given, but we will approximate it here simply by $`R_{CC}=H^1`$, a form applicable in almost all situations. For $`w<\frac{1}{3}`$ the expansion is inflationary, so fixed comoving points “inside the horizon”, will in time “exit the horizon”, while for $`w>\frac{1}{3}`$ the expansion is decelerated, so fixed comoving points “outside the horizon” , will in time “enter the horizon”.
Horizons are 2D surfaces that are classically, in the homogeneous and isotropic cases that we are interested in here, well defined spherical shells, which may evolve in time. Quantum mechanically they can fluctuate, so it is reasonable to expect that they can be described effectively by 2D field theories.
We focus on the NH geometry by choosing a 2D metric $`ds_2^2=dt^2+a^2(t)dr^2`$, fixing the position of the horizon $`H^1(t)`$, at some specified time $`t^{}`$, $`H^1(t^{})=d`$, and changing coordinates to Schwartzschild-like coordinates $`R=a(t)r`$, and $`T`$ defined by $`dT=dt+\frac{R/d}{1R/d}dR`$. In the new coordinates
$`ds_2^2=(1H^2R^2)dT^22\left[HR{\displaystyle \frac{R}{d}}{\displaystyle \frac{1H^2R^2}{1\left(\frac{R}{d}\right)^2}}\right]dRdT`$ (2)
$`+\left[1\left({\displaystyle \frac{R}{d}}\right)^2{\displaystyle \frac{1H^2R^2}{\left(1\left(\frac{R}{d}\right)^2\right)^2}}+2HR{\displaystyle \frac{R/d}{1\left(\frac{R}{d}\right)^2}}\right]dR^2.`$ (3)
In the NH “shell” defined by
$$\frac{H^1(t)d}{d}\frac{Rd}{d}1,$$
(4)
(3) reduces to
$$ds_2^22(1R/d)dT^2+\frac{1}{2}\frac{1}{1R/d}dR^2.$$
(5)
One more change of coordinates $`d\rho =\frac{1}{2}\frac{1}{1R/d}dR`$, brings (5) into $`ds_2^2=e^{2\rho /d}\left(dT^2+d\rho ^2\right)`$ in the NH region $`\rho /d1`$, so that $`\gamma _{ab}=e^{2\rho /d}\eta _{ab}`$, $`\sqrt{\gamma }=e^{2\rho /d}`$, and $`\gamma ^{ab}=e^{+2\rho /d}\eta _{ab}`$. For later reference, we note that $`t`$ is a NH lightcone coordinate,
$$t=T\rho .$$
(6)
We would like to obtain an effective 2D field theory of NH geometries, so we parametrize the angular part of the metric with a field $`\mathrm{\Phi }`$, which will eventually determine the horizon surface $`\mathrm{\Phi }=H^1(t)`$, $`ds_4^2=\gamma _{ab}(t,r)dx^adx^b+\mathrm{\Phi }^2(t,r,\theta ,\varphi )\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right)`$. $`\mathrm{\Phi }`$ is allowed to have general dependence on the coordinates. For this metric (1) reduces to
$$S^{(4)}=\frac{1}{8\pi G_N}𝑑t𝑑r𝑑\mathrm{\Omega }_2\sqrt{\gamma }\left\{\gamma ^{ab}_a\mathrm{\Phi }_b\mathrm{\Phi }+\frac{1}{2}\mathrm{\Phi }^2R^{(2)}+1\right\},$$
(7)
where $`d\mathrm{\Omega }_2=\mathrm{sin}\theta d\theta d\varphi `$. In (7) we have dropped the matter action $`S_m`$, since its only function is to determine the classical solution. Note that terms containing angular derivatives are absent in (7) .
Mass terms in action (7) are suppressed in the NH region (4) . Obviously, the last term in (7) is exponentially suppressed by a factor $`e^{2\rho /d}`$ with respect to the first two terms<sup>*</sup><sup>*</sup>*Strictly speaking, certain variations coming from this term are suppressed. The NH effective action is therefore the following,
$$S_{\mathrm{NH}}^{(4)}=\frac{1}{8\pi G_N}𝑑T𝑑\rho 𝑑\mathrm{\Omega }_2\sqrt{\gamma }\left\{\gamma ^{ab}_a\mathrm{\Phi }_b\mathrm{\Phi }+\frac{1}{2}R^{(2)}\mathrm{\Phi }^2\right\}.$$
(8)
We proceed to reduce action (8) to 2D. First, we expand $`\mathrm{\Phi }`$ in spherical harmonics,
$$\mathrm{\Phi }(t,r,\theta ,\varphi )=\underset{l,|m|l_{\mathrm{max}}}{\overset{l_{\mathrm{max}}}{}}\mathrm{\Phi }_{l,m}(t,r)Y_l^m(\theta ,\phi ).$$
(9)
The maximal angular momentum $`l_{\mathrm{max}}`$ in expansion (9) is determined by the short distance cutoff of the theory, as we discuss later. We then substitute (9) into (8), perform the angular integration using the orthogonality property of $`Y_l^m`$’s, and obtain a dimensionally reduced NH 2D effective action,
$`S_{\mathrm{NH}}^{(2)}={\displaystyle \frac{1}{8\pi G_N}}{\displaystyle }dTd\rho \sqrt{\gamma }\times `$ (10)
$`\left\{{\displaystyle \underset{l,|m|l_{\mathrm{max}}}{\overset{l_{\mathrm{max}}}{}}}\gamma ^{ab}_a\mathrm{\Phi }_{l,m}_b\mathrm{\Phi }_{l,m}+{\displaystyle \frac{1}{2}}R^{(2)}{\displaystyle \underset{l,|m|l_{\mathrm{max}}}{\overset{l_{\mathrm{max}}}{}}}\mathrm{\Phi }_{l,m}^2\right\}.`$ (11)
Only a single field $`\mathrm{\Phi }_{l,m}^2`$, couples to the 2D curvature term. This field is simply the area of the horizon shell $`A^H`$, $`A^H=𝑑\mathrm{\Omega }_2\mathrm{\Phi }^2(t,r,\theta ,\varphi )=\underset{l,|m|l_{\mathrm{max}}}{\overset{l_{\mathrm{max}}}{}}\mathrm{\Phi }_{l,m}(t,r)^2.`$ Here we have performed the angular integration using the orthogonality property of $`Y_l^m`$’s. We will be interested in fluctuations of the horizon which keep the area fixed at its (time-dependent) classical value, and therefore we will freeze quantum fluctuations of this mode. Since we are interested in counting states for the case of large $`l_{\mathrm{max}}`$, projecting out a single mode will not compromise the generality of our results. Freezing the quantum fluctuations of the area has the benefit of simplifying quantization of the NH theory enormously. Furthermore, we do not take into account fluctuations of matter sources, assuming that their only effect is to determine the time dependent expectation value of the horizon.
So, when all is said and done, the remaining NH action is simply a sum of actions of independentWe ignore the single overall constraint on their sum, which effectively removes a single field. free scalar fields, minimally coupled to 2D gravity,
$$S_{\mathrm{NH}}^{(2)}=\frac{1}{8\pi G_N}d^2x\sqrt{\gamma }\underset{l,|m|l_{\mathrm{max}}}{\overset{l_{\mathrm{max}}}{}}\gamma ^{ab}_a\mathrm{\Phi }_{l,m}_b\mathrm{\Phi }_{l,m}.$$
(12)
Theory (12) can be quantized in the 2D conformal gauge, using standard DDK arguments. The conformal anomaly of $`\mathrm{\Phi }_{l,m}`$’s induces a kinetic term for the 2D conformal mode, and renormalizes the 2D action such that the full theory is a CFT, whose total central charge vanishes. But the central charge of the Liouville mode cannot be used to calculate the density of states . For state counting purposes it counts as a single field.
To determine the horizon’s entropy we need to compute the effective central charge $`c`$, of the NH CFT, and the expectation value of its Hamiltonian $`L_0`$. We first calculate the total effective central charge $`c`$, which is approximately equal to the sum of individual matter central charges,
$$c\underset{l,|m|l_{\mathrm{max}}}{\overset{l_{\mathrm{max}}}{}}1l_{\mathrm{max}}^2.$$
(13)
In (13) we have neglected contributions from ghosts, from the Liouville mode, ignored the area constraint, and included redundant contributions from a small number of gauge modes, but since we are interested in the case of large $`l_{\mathrm{max}}`$, we are justified in doing so.
The maximal angular momentum $`l_{\mathrm{max}}`$, is determined by Planck scale physics. The smallest angular variations $`\mathrm{\Delta }\phi _{\mathrm{min}}`$ and $`\mathrm{\Delta }\theta _{\mathrm{min}}`$ allowed as fluctuations of a sphere of radius $`d`$, are determined by the short distance cutoff of the theory $`\mathrm{}_{UV}`$,
$$\mathrm{\Delta }\phi _{\mathrm{min}}=\mathrm{}_{UV}/d,\mathrm{\Delta }\theta _{\mathrm{min}}=\mathrm{}_{UV}/d.$$
(14)
Since $`Y_l^me^{im\phi }e^{il\theta }`$, the smallest angular variations $`\mathrm{\Delta }\phi _{\mathrm{min}}`$ and $`\mathrm{\Delta }\theta _{\mathrm{min}}`$ determine the maximal angular momentum,
$$m_{\mathrm{max}}=\frac{C_m}{\mathrm{\Delta }\phi _{\mathrm{min}}},l_{\mathrm{max}}=\frac{C_l}{\mathrm{\Delta }\theta _{\mathrm{min}}},$$
(15)
where $`C_l`$, $`C_m`$ are numerical coefficients of order unity. We may use eqs. (14,15) to estimate the maximal allowed angular momentum $`l_{\mathrm{max}}^2=C_lC_md^2/\mathrm{}_{UV}^2`$. Assuming that the short distance cutoff is some numerical factor of order unity $`k`$, times the Planck length $`\mathrm{}_{UV}=k\sqrt{G_N}`$, and denoting $`\alpha =\frac{C_lC_m}{k^2}`$ we obtain our final expression for the total central charge of the NH theory,
$$c=\alpha \frac{A^H}{4\pi G_N}.$$
(16)
The expectation value of $`L_0`$ is determined by the classical background. Recall that the classical solution is a function of time only $`\mathrm{\Phi }(r,t,\theta ,\phi )=H^1(t)`$, and that $`H^1(t)=d\left(\frac{t}{t^{}}\right)`$. But this means that only the $`l=0`$, $`m=0`$ mode has non-trivial expectation value
$$\mathrm{\Phi }_{0,0}=\sqrt{4\pi }d\left(\frac{t}{t^{}}\right).$$
(17)
To simplify evaluation of $`L_0`$, we go to 2D lightcone coordinates in two steps, first setting $`\tau =\frac{T}{t^{}}`$, $`\sigma =\frac{\rho }{t^{}}`$, and then setting $`x^\pm =\tau \pm \sigma `$. Note that according to (6), $`x^{}=\frac{t}{t^{}}`$. We expand
$$\mathrm{\Phi }_{0,0}=q+px^{}+\underset{n0}{}\frac{1}{n}\alpha _ne^{2in\pi {\scriptscriptstyle \frac{\sigma }{\sigma _{\mathrm{max}}}}x^{}},$$
(18)
where $`\sigma _{\mathrm{max}}`$ determines the range of the 2D coordinate $`\sigma `$. Since $`x^{}=\frac{t}{t^{}}`$, we can compare (18) and (17), and observe that for the classical background only $`p`$ is non-vanishing $`p=\sqrt{4\pi }d`$, while all the $`\alpha _n`$’s and $`q`$ vanish. Since $`L_0=\frac{1}{8\pi G_N}\left[p^2+\underset{n0}{}\alpha _n\alpha _n\right]`$, it follows that $`L_0=p^2/8\pi G_N`$, but $`p^2=A^H`$, so
$$L_0=\frac{A^H}{8\pi G_N}.$$
(19)
In general, CFT’s have two sets of independent modes which are either functions of $`x^+`$, or $`x^{}`$, but the NH theory has only one set of modes . As we show $`\mathrm{\Phi }_{l,m}=\mathrm{\Phi }_{l,m}(x^{})`$, leaving only one Virasoro algebra as symmetry of the NH CFT. Near the horizon, as we have already seen, propagating modes are massless, due to redshift effects. But the same redshift effects allow them to propagate only along outgoing light-like trajectories in the $`x^{}`$ direction. Near black hole horizons, similar redshift effects allow only ingoing modes to propagate. To see this in more detail, we look at the $`x^+`$ derivative
$$_+=\frac{1}{2}_{T}^{}{}_{|t}{}^{}+\frac{1}{2}_{\rho }^{}{}_{|t}{}^{}=_{T}^{}{}_{|t}{}^{}=_{\rho }^{}{}_{|t}{}^{},$$
(20)
where the last equation is obtained using eq.(6). For simplicity, we now set $`t^{}=1`$, so $`x^{}=t`$. Expressing $`T`$ and $`\rho `$ derivatives in terms of $`R`$ derivatives, we find
$$_{\rho }^{}{}_{|t}{}^{}_{T}^{}{}_{|t}{}^{}2\left(1\frac{R}{d}\right)_R,$$
(21)
so smooth functions in the original variables $`(R,t)`$ have vanishing $`x^+`$ derivative in the NH region (4). Further examination shows that all $`x^+`$ derivatives $`_+^n`$, vanish for such functions, so smooth functions in the original variables $`(R,t)`$ are not functions of $`x^+`$ but only functions of $`x^{}`$, as claimedWe will not discuss possible diffeomorphisms anomalies due to the chiral nature of NH theory..
We are now ready to use Cardy’s formula to evaluate the entropy of the horizon,
$`S_{\mathrm{horizon}}`$ $`=`$ $`2\pi \sqrt{{\displaystyle \frac{c}{6}}\left(L_0{\displaystyle \frac{c}{24}}\right)}`$ (22)
$`=`$ $`{\displaystyle \frac{A^H}{4G_N}}\sqrt{{\displaystyle \frac{\alpha }{3}}\left(1{\displaystyle \frac{\alpha }{12}}\right)},`$ (23)
where $`\alpha `$ is defined above (16).
Since coefficient $`\alpha `$ is determined by short distance (Planck scale) physics, we expect it to be universal. It should not be sensitive to the macroscopic, large scale physics which determines the exact nature of classical solutions. If so, we may use the limiting case of de Sitter space to “calibrate” it. In de Sitter space the horizon entropy is known to be $`\frac{A^H}{4G_N}`$ from other considerations. This procedure sets the value of $`\alpha `$ at $`\alpha =6`$, leading to the conclusion that, in general,
$$S_{\mathrm{horizon}}=\frac{A^H}{4G_N}.$$
(24)
The horizon shell, whose entropy we have just calculated, is a thin shell, since its thickness is much smaller than its radius, but its thickness is macroscopic, much larger than the Planck length, so we might have expected that its entropy turns out to be proportional to its volume, in units of Planck volume. The entropy of an arbitrary shell does scale as its volume, but not the NH shell. So what is so special about the horizon? Huge redshifts “squash” NH shells, and make their entropy proportional to their area and not to their volume.
Our results apply also to collapsing pressureless matter, since the such systems can be described by FRW metric. In this case, the horizon whose entropy we have computed is an apparent horizon. The apparent horizon reaches the event horizon when all matter has collapsed to a point.
Our explicit calculations were carried out for 4-dimensional, homogeneous and isotropic, spatially flat geometries. But the essential ingredient in our calculation was the huge redshift near the horizon. This left only massless outgoing 2D modes, which are naturally described by a CFT. Since this ingredient seems to be present whenever causal boundaries form, I believe that similar methods can be applied to higher dimensional spaces, along the lines of , to more general spatially curved spaces, using the covariant definition of causal connection scale , and to string theory along the lines of .
###### Acknowledgements.
I wish to thank S. de Alwis, S. Elitzur, S. Foffa and S. Solodukhin for discussions, and CERN-TH division, where part of this work was carried out, for hospitality.
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# 1 Introduction
## 1 Introduction
The Non-Perturbative (NP) renormalization technique proposed in has been successfully applied to compute renormalization constants of two-fermion and four-fermion operators with Wilson fermions. The standard implementation consists in imposing the RI/MOM renormalization conditions on conveniently defined amputated projected Green functions computed between off-shell quark states of operators of interest, evaluated numerically at fixed momentum and lattice spacing (fixed lattice coupling $`\beta `$). In order to control discretization effects, one should work at renormalization scales $`\mu `$ well below the UV cutoff (i.e. $`\mu 𝒪(a^1)`$). Moreover, any contribution in the renormalization constants due to spontaneous chiral symmetry breaking terms in the Green’s functions has to be avoided to ensure that the non-perturbatively renormalized operators satisfy the standard QCD axial Ward Identities (abbreviated as WIs), i.e. $`\mathrm{\Lambda }_{QCD}\mu `$. Thus, the existence of a “renormalization window” $`\mathrm{\Lambda }_{QCD}\mu 𝒪(a^1)`$ is essential to the reliability of the RI/MOM scheme in practical simulations, as was pointed out in ref. . We note in passing that, as discussed in refs. , the lower bound can be relaxed by matching renormalization conditions between a coarse and a fine lattice in the spirit of ref. ).
This has not always been the case at current values of the lattice coupling. In off-shell correlation functions, at high enough momenta, the dominant contributions are perturbative. For example, the leading non perturbative contribution of the pseudoscalar correlation function can be easily obtained from the Operator Product Expansion (OPE) of the quark propagator, which has been worked out in refs. . Using either the lattice or the Sum Rules determination of the chiral condensate and quark mass values typical of present-day simulations, we estimate that the non perturbative contribution to the pseudoscalar correlation function at $`\mu 2`$ GeV is at most 2%. Therefore, the existence of a “renormalization window’ has been checked in by comparing the $`p^2`$ dependence of the Green’s functions to the logarithmic one predicted by continuum perturbation theory <sup>1</sup><sup>1</sup>1Note that continuum perturbation theory converges much better than the lattice one. Moreover, higher orders have been calculated in the former case, whereas only 1-loop results are available in lattice perturbation theory.. In this sense, at $`\beta =6.0,6.2`$ and $`6.4`$, the scalar density renormalization constant $`Z_S`$ is well behaved in a large $`\mu `$-range, whereas there is no clear evidence of such a satisfactory renormalization window for the RI/MOM renormalization constant of the pseudoscalar density $`Z_P`$ (see ref. for details).
Recently an alternative interpretation of the observed discrepancy between the RI/MOM and perturbative $`Z_P`$ results has been proposed in ref. . It is suggested that the $`p^2`$ dependence of the pseudoscalar correlation function is due to large contributions of the spontaneous chiral symmetry breaking terms, which persist at scales used in present day simulations (i.e. $`\mu a^1`$). In other words, the lower bound of the renormalization window is not adequately satisfied. In order to test this proposal, the authors of ref. analyze the pseudoscalar correlation function, taking into account the leading OPE non-perturbative contribution and ignoring discretization effects. They find that “a stricking and unexpected feature of the lattice data is the very large size of the Goldstone boson contribution to the pseudoscalar vertex”.
In the present work we critically examine the claim of ref. by comparing the RI/MOM determination of the ratio $`Z_P/Z_S`$ to those obtained from two WIs at $`\beta =6.2,6.4`$ (the quenched approximation and Wilson fermions are implied throughout). WI results are unaffected by Goldstone pole contamination but are subject to discretization effects. Thus any discrepancy between two WI results gives an estimate of $`𝒪(a)`$ effects, whereas an eventually big discrepancy between the WI results and the RI/MOM ones would signal a Goldstone pole contamination in $`Z_P`$. We find that individual WI results for $`Z_P/Z_S`$ are characterized by errors of about 5% but they are discrepant to each other at the level of 10-15% due to $`𝒪(a)`$ effects. The RI/MOM result obtained at $`\mu 2`$ GeV is compatible with the WI ones but has a larger error of about 10-15%. Thus, contrary to the claim of ref. , we conclude that the RI/MOM determination of $`Z_P/Z_S`$ (and subsequently of $`Z_P`$) is affected by large discretization errors while any contamination from Goldstone pole contributions cannot be discerned. Analogous conclusions, based on an analysis of the form factors of the improved quark propagator, have been drawn in ref. .
Finally we wish to stress that the WI proposed in this paper (based on off-shell correlation functions and non-degenerate quark masses) has given a very stable estimate of $`Z_P/Z_S`$. In view of this, we prefer to evaluate $`Z_P`$ from the WI result of $`Z_P/Z_S`$ and the RI/MOM value of $`Z_S`$.
## 2 The RI/MOM renormalization scheme and WIs
In this section we review the basic characteristics of the RI/MOM renormalization scheme, implemented with the lattice regularization . After giving the essential definitions, we discuss the compatibility of the RI/MOM scheme with (lattice) WIs at large renormalization scales $`\mu `$, where the presence of the Goldstone pole becomes negligible.
We start with basic definitions. Given the quark propagator $`𝒮(x_1x_2;m)=\psi (x_1)\overline{\psi }(x_2)`$ and its Fourier transform $`𝒮(p;m)=𝑑x\mathrm{exp}(ipx)𝒮(x,m)`$ (here $`m`$ is the quark mass), we define two “projections” (i.e. traces in spin and color space) as
$`\mathrm{\Gamma }_\mathrm{\Sigma }(p;m)={\displaystyle \frac{i}{48}}\text{Tr}\left[\gamma _\mu {\displaystyle \frac{𝒮^1(p;m)}{p_\mu }}\right]`$
$`\mathrm{\Gamma }_m(p;m)={\displaystyle \frac{1}{12}}\text{Tr}\left[{\displaystyle \frac{𝒮^1(p;m)}{m}}\right]`$ (1)
We also consider bilinear quark operators of the form $`O_\mathrm{\Gamma }(x)=\overline{\psi }_1(x)\mathrm{\Gamma }\psi _2(x)`$ where $`\psi _f(x)`$ is the quark field and $`\mathrm{\Gamma }`$ a generic Dirac matrix. For definitiveness we work with two different flavors $`f=1,2`$ with corresponding masses $`m_1`$ and $`m_2`$. Specific non-singlet bilinear operators will be denoted as $`S(x)`$, $`P(x)`$ (scalar and pseudoscalar densities) and $`V_\mu (x)`$, $`A_\mu (x)`$ (vector and axial currents). Given the insertion of the operator $`O_\mathrm{\Gamma }(x=0)`$ in the $`2`$-point fermionic Green’s function
$$G_\mathrm{\Gamma }(p)=𝑑x_1𝑑x_2\mathrm{exp}[ip(x_1x_2)]\psi (x_1)O_\mathrm{\Gamma }(0)\overline{\psi }(x_2)$$
(2)
and the corresponding amputated correlation function $`\mathrm{\Lambda }_\mathrm{\Gamma }(p)=𝒮^1(p)G_\mathrm{\Gamma }(p)𝒮^1(p)`$, the projected amputated Green’s function $`\mathrm{\Gamma }_\mathrm{\Gamma }(p)`$ is defined as
$$\mathrm{\Gamma }_\mathrm{\Gamma }(p)=\frac{1}{12}\text{Tr}\left[\text{I}\text{P}_\mathrm{\Gamma }\mathrm{\Lambda }_\mathrm{\Gamma }(p)\right]$$
(3)
where the trace is over spin and color indices and $`\text{I}\text{P}_\mathrm{\Gamma }`$ is the Dirac matrix which renders the tree-level value of $`\mathrm{\Gamma }_\mathrm{\Gamma }(p)`$ equal to unity (i.e. it projects out the nominal Dirac structure of the Green function $`\mathrm{\Lambda }_\mathrm{\Gamma }(p)`$):
$`\text{I}\text{P}_S=I;\text{I}\text{P}_P=\gamma _5`$
$`\text{I}\text{P}_V={\displaystyle \frac{1}{4}}\gamma _\mu ;\text{I}\text{P}_A={\displaystyle \frac{1}{4}}\gamma _5\gamma _\mu `$ (4)
Everything defined so far is a bare quantity, assumed to be regularized on the lattice <sup>2</sup><sup>2</sup>2This means that the integrals of eq. (2) are really sums ($`a^8_{x_1,x_2}`$) which run over all lattice sites, labelled by $`x_1`$, $`x_2`$, etc. Also note that the $`\mu `$-dependence of renormalized correlation functions is sometimes suppressed; i.e. we use $`\widehat{\mathrm{\Gamma }}(p;\widehat{m})`$ instead of $`\widehat{\mathrm{\Gamma }}(p^2/\mu ^2;\widehat{m}^2/\mu ^2)`$. We opt for the lattice regularization scheme with Wilson fermions. Then the corresponding renormalized quantities in a given mass independent renormalization scheme are formally given by
$`\widehat{m}_f(\mu )`$ $`=`$ $`\underset{a0}{lim}[Z_m(a\mu )m_f(a)]=\underset{a0}{lim}\left[Z_m\left(m_{0f}m_C\right)\right]`$
$`\widehat{𝒮}(p;\widehat{m}_f,\mu )`$ $`=`$ $`\underset{a0}{lim}[Z_\psi (a\mu )𝒮(p;m_f,a)]`$
$`\widehat{\mathrm{\Gamma }}_\mathrm{\Sigma }(p^2/\mu ^2;m_f^2/\mu ^2)`$ $`=`$ $`\underset{a0}{lim}\left[Z_\psi ^1(a\mu )\mathrm{\Gamma }_\mathrm{\Sigma }(ap,am_f)\right]`$
$`\widehat{\mathrm{\Gamma }}_m(p^2/\mu ^2;\widehat{m}_f^2/\mu ^2)`$ $`=`$ $`\underset{a0}{lim}\left[Z_\psi ^1(a\mu )Z_m^1(a\mu )\mathrm{\Gamma }_m(ap;am_f)\right]`$
$`\widehat{\mathrm{\Gamma }}_\mathrm{\Gamma }(p^2/\mu ^2;\widehat{m}_f^2/\mu ^2)`$ $`=`$ $`\underset{a0}{lim}\left[Z_\psi ^1(a\mu )Z_\mathrm{\Gamma }(a\mu )\mathrm{\Gamma }_\mathrm{\Gamma }(ap;am_f)\right]`$ (5)
Note that $`m_{0f}`$ denotes the bare quark mass of a given flavor, $`m_f=m_{0f}m_C`$ is power subtracted and logarithmically divergent and $`\widehat{m}_f`$ is renormalized. Moreover, $`Z_\psi ^{1/2}`$ is the quark field renormalization and $`Z_\mathrm{\Gamma }`$ the operator renormalization. The functional dependence of the above expressions is determined by dimensional arguments and Lorenz invariance.
The renormalization constants in the RI/MOM scheme are defined by imposing the following off-shell conditions in the deep Euclidean region: the wave function renormalization $`Z_\psi `$ is obtained from
$$\widehat{\mathrm{\Gamma }}_\mathrm{\Sigma }(p^2/\mu ^2;\widehat{m}^2/p^2)|_{\begin{array}{c}p^2=\mu ^2\\ \widehat{m}^2=0\end{array}}=\underset{a0}{lim}\underset{m0}{lim}Z_\psi ^1(a\mu )\mathrm{\Gamma }_\mathrm{\Sigma }(a\mu ;a^2m^2)=1$$
(6)
the quark mass renormalization $`Z_m`$ is achieved through
$$\widehat{\mathrm{\Gamma }}_m(p^2/\mu ^2;\widehat{m}^2/p^2)|_{\begin{array}{c}p^2=\mu ^2\\ \widehat{m}^2=0\end{array}}=\underset{a0}{lim}\underset{m0}{lim}Z_\psi ^1(a\mu )Z_m^1(a\mu )\mathrm{\Gamma }_m(a\mu ;a^2m^2)=1$$
(7)
and for the bilinear operator $`Z_\mathrm{\Gamma }`$ is obtained by imposing the renormalization condition
$$\widehat{\mathrm{\Gamma }}_\mathrm{\Gamma }(p^2/\mu ^2;\widehat{m}^2/p^2)|_{\begin{array}{c}p^2=\mu ^2\\ \widehat{m}^2=0\end{array}}=\underset{a0}{lim}\underset{m0}{lim}Z_\psi ^1(a\mu )Z_\mathrm{\Gamma }(a\mu )\mathrm{\Gamma }_\mathrm{\Gamma }(a\mu ;a^2m^2)=1$$
(8)
For any fixed scale $`\mu `$, this procedure removes the UV divergences from all Green functions and thus renormalization is achieved. However, renormalization conditions must also be chosen so that the resulting renormalized operators transform “correctly” under the chiral group; i.e. they should belong to the same chiral representation as the nominal bare operators. This additional requirement, which ensures that chiral symmetry survives renormalization, is only true at large enough scales $`\mu `$. At small scales symmetry violating effects due to spontaneous chiral symmetry breaking in QCD should appear as $`\mathrm{\Lambda }_{QCD}`$-dependent form factors in the renormalized Green functions $`\widehat{\mathrm{\Gamma }}_\mathrm{\Gamma }(p^2/\mu ^2;\widehat{m}^2/p^2;\mathrm{\Lambda }_{QCD}/p^2)`$. Moreover, in practical simulations, usually performed at (light) non-zero quark mass, explicit chiral symmetry breaking form factors, proportional to $`\widehat{m}(\mu )`$ will also be present. Both types of form factors become negligible if the scale $`\mu `$ is adequately large. Therefore, the requirement $`\mathrm{\Lambda }_{QCD}\mu `$ has to be satisfied in principle, so that quantities renormalized in the RI/MOM scheme respect chiral symmetry. Moreover, since one is usually working at fixed UV cutoff in simulations, the requirement $`\mu 𝒪(a^1)`$ must be satisfied in order to control discretization errors. Thus, the reliability of the RI/MOM scheme on the lattice is ensured provided one is working in a renormalization window $`\mathrm{\Lambda }_{QCD}\mu 𝒪(a^1)`$.
The above general statements have been shown to be true in ref. (and generalized for some cases of additive renormalization of dimension-six four-fermion operators in refs. ). In particular it has been shown that the RI/MOM scheme is always compatible with vector WIs, whereas it is only compatible with the axial WIs at large scales $`\mu `$. The key observation is that the vector WIs
$`\widehat{\mathrm{\Gamma }}_S(p;\widehat{m})`$ $`=`$ $`\widehat{\mathrm{\Gamma }}_m(p;\widehat{m})`$
$`\widehat{\mathrm{\Gamma }}_V(p;\widehat{m})`$ $`=`$ $`\widehat{\mathrm{\Gamma }}_\mathrm{\Sigma }(p;\widehat{m})`$ (9)
at momentum $`p^2=\mu ^2`$, are automatically satisfied by the renormalized quantities $`\widehat{\mathrm{\Gamma }}_\mathrm{\Sigma }`$, $`\widehat{\mathrm{\Gamma }}_m`$, $`\widehat{\mathrm{\Gamma }}_V`$ and $`\widehat{\mathrm{\Gamma }}_S`$ determined in the RI/MOM scheme by eqs. (6), (7) and (8). Conversely, if we use the RI/MOM scheme to fix, say, $`\widehat{\mathrm{\Gamma }}_\mathrm{\Sigma }`$ and $`\widehat{\mathrm{\Gamma }}_m`$ (i.e. $`Z_\psi `$ and $`Z_m`$), then $`Z_V`$ and $`Z_S`$ (through the identity $`Z_S=Z_m^1`$), as fixed by the WIs (9), are compatible to those obtained from the RI/MOM condition (8).
The compatibility of axial WIs to RI/MOM is more intricate; the proof would proceed in exactly the same way as in the vector case, if it were not for the extra terms on the l.h.s. of WIs
$`\widehat{\mathrm{\Gamma }}_P(p;\widehat{m})+\widehat{m}{\displaystyle \frac{\widehat{\mathrm{\Gamma }}_P(p;\widehat{m})}{\widehat{m}}}`$ $`=`$ $`\widehat{\mathrm{\Gamma }}_m(p;\widehat{m})`$
$`\widehat{\mathrm{\Gamma }}_A\left(p;\widehat{m}=0\right)+\underset{q_\rho 0}{lim}{\displaystyle \frac{q_\mu }{48}}\text{Tr}\left[\gamma _5\gamma _\rho {\displaystyle \frac{\widehat{\mathrm{\Lambda }}_A^\mu (p+q/2,pq/2;\widehat{m}=0)}{q_\rho }}\right]`$ $`=`$ $`\widehat{\mathrm{\Gamma }}_\mathrm{\Sigma }\left(p;\widehat{m}=0\right)`$ (10)
which is the axial analogue to the vector WIs (9). Note that the second term of the l.h.s. of WIs (10) does not vanish in the chiral limit, due to the presence of a Goldstone pole. As shown in ref. these extra terms are negligible in the deep Euclidean region $`\mathrm{\Lambda }_{QCD}p`$. In this limit, the WIs and the RI/MOM conditions are compatible, just like the vector case. This limit corresponds to the lower bound of the renormalization window requirement, inherent in the RI/MOM scheme. Clearly, all statements made thus far are true up to discretization effects, which are present both in the RI/MOM and the WI implementations on the lattice. We will address this problem at a later stage of this work. What we need to consider at present, is that RI/MOM results are reliable in the renormalization window $`\mathrm{\Lambda }_{QCD}p𝒪(a^1)`$, whereas WI results are reliable in the window $`p𝒪(a^1)`$.
The main aim of the present work is to establish whether there is significant Goldstone pole contamination of the lattice RI/MOM results obtained thus far. In particular, the RI/MOM determination of $`Z_P/Z_S`$ at current $`\beta `$ values and momenta does not display a clear plateau; see refs. -. The problem may be either due to the Goldstone pole problem or, as claimed in ref. , due to discretization effects<sup>3</sup><sup>3</sup>3Moreover, the 1-loop Perturbation Theory PT (and Boosted PT -BPT) result for $`Z_P`$ is in stark disagreement with the RI/MOM estimate. This comparison, however, does not reveal much, as the BPT result also suffers from $`𝒪(g_0^4)`$ errors.. Since the WI determination of the ratio $`Z_P/Z_S`$ does not suffer from Goldstone pole contamination, comparison of the WI and the RI/MOM results for this quantity should offer a way of estimating this effect. For example, the WI result from eqs. (9) and (10) is
$$\frac{Z_P}{Z_S}=\frac{{\displaystyle \frac{\mathrm{\Gamma }_S}{\mathrm{\Gamma }_P}}}{1+{\displaystyle \frac{m}{\mathrm{\Gamma }_P}}{\displaystyle \frac{\mathrm{\Gamma }_P}{m}}}$$
(11)
whereas the RI/MOM result from eq. (8) is
$$\frac{Z_P}{Z_S}=\frac{\mathrm{\Gamma }_S}{\mathrm{\Gamma }_P}$$
(12)
The denominator on the r.h.s. of the WI (11) is missed by the RI/MOM eq. (12) and this results to a Goldstone pole contamination of the latter. This particular WI determination, however, is difficult to implement in practice (e.g. realization of derivatives w.r.t. the quark mass in lattice numerical simulations etc.). In the next section we discuss several WIs, which are equivalent (up to discretization effects) but more useful in practice.
## 3 The ratio $`Z_P/Z_S`$ from WIs
We will now review three methods, based on WIs, for the determination of the scale independent ratio $`Z_P/Z_S`$. A first method consists in computing $`Z_P/Z_S`$ as the ratio of the PCAC current quark mass to the bare (subtracted) quark mass . In standard notation (see refs. , ) the current quark mass for Wilson fermions is obtained from the standard PCAC relation (an axial WI, valid $`x0`$):
$$2\left[m_0\overline{m}(m_0)\right]=Z_A\frac{_0d^3xA_0(x)P(0)}{d^3xP(x)P(0)}$$
(13)
whereas the mass power subtraction $`m_C`$ can be obtained by linearly extrapolating either the square of the pion mass $`m_\pi ^2`$ or $`\left[m_0\overline{m}(m_0)\right]`$ to their vanishing value. Then vector and axial WIs (see refs. ) determine the quark mass renormalization
$$\widehat{m}=Z_S^1\left[m_0m_C\right]=Z_P^1\left[m_0\overline{m}\right]$$
(14)
and thus by computing $`\left[m_0\overline{m}\right]`$ at several $`m_0`$’s the ratio $`Z_P/Z_S`$ is obtained as the slope of
$$m_0\overline{m}=\frac{Z_P}{Z_S}\left[m_0m_C\right].$$
(15)
This determination does not depend on the critical quark mass $`m_C`$. It does, however, depend on the determination of $`Z_A`$. Since the methods relies on the WI (13) on hadronic states, we will label it WIh.
A second method consists in combining eq. (15) with the axial WI
$`(m_{01}+m_{02}2\overline{m})\mathrm{\Gamma }_P(ap;am_1,am_2)`$ $`=`$ $`{\displaystyle \frac{1}{12}}\text{Tr}𝒮^1(ap;am_1)+{\displaystyle \frac{1}{12}}\text{Tr}𝒮^1(ap;am_2)`$ (16)
(considered in the mass-degenerate case), in order to obtain
$$\frac{Z_P}{Z_S}=\frac{\text{Tr}𝒮^1(ap;am)}{12\left[m_0m_C\right]\mathrm{\Gamma }_P(ap;am,am)}$$
(17)
This determination requires the projected amputated correlation functions and the quark propagators in momentum space. We anticipate at this stage, that this method fails in practice. This is due to the large discretization errors of the quark propagator (a detailed investigation of these effects, which remain large even when Clover improvement is implemented beyond tree-level, has been performed in ref. ).
A third method consists in combining eq. (16) for the degenerate mass case with the vector WI
$`(m_{02}m_{01})\mathrm{\Gamma }_S(ap;am_1,am_2)`$ $`=`$ $`{\displaystyle \frac{1}{12}}\text{Tr}𝒮^1(ap;am_1)+{\displaystyle \frac{1}{12}}\text{Tr}𝒮^1(ap;am_2)`$ (18)
so as to eliminate the quark propagators. Using also eq. (15) we obtain
$$\frac{Z_P}{Z_S}=\frac{\left(m_1m_2\right)\mathrm{\Gamma }_S(ap;am_1,am_2)}{m_1\mathrm{\Gamma }_P(ap;am_1,am_1)m_2\mathrm{\Gamma }_P(ap;am_2,am_2)}$$
(19)
The advantage of this determination is that the large $`𝒪(ap)`$ discretization errors, present in the quark propagator (cf. ref. ), are eliminated. Since this determination depends on correlation functions of external quark states, we will label it WIq. Note that this method is characterized by non-degenerate quark masses.
Finally, we mention that a WI with hadronic states has also been used for the computation of the ratio $`Z_P/Z_S`$; see refs. .
## 4 Results
Our results are based on earlier simulations; all technical details can be found in refs. . Here we just mention that the dataset we analyze has been obtained with the Wilson (unimproved) action, at couplings $`\beta =6.2`$ and $`6.4`$ and at several hopping parameters in the range of the strange quark mass. All results shown here are extrapolated to the chiral limit. The RI/MOM results are those of ref. (which agree with the ones of ref. ) whereas the WIh and WIq results are new. We mostly show the Wilson $`\beta =6.2`$ case <sup>4</sup><sup>4</sup>4Note that our $`\beta =6.0`$ dataset has also been analyzed, with qualitatively similar results..
In figure 1 we show the data and linear fit characteristic of the method WIh. The slope of the fitting line is $`Z_P/[Z_SZ_A]`$. In order to extract $`Z_P/Z_S`$ from the slope, we use the RI/MOM estimate for $`Z_A`$, as given in ref. . This is in accordance (with larger errors) to the RI/MOM result of ref. .
In Fig. 2(a) we show our $`\beta =6.2`$ results obtained from WI (17), as a function of the (lattice) scale
$$a^2\overline{\mu }^2=\underset{\nu }{}\mathrm{sin}^2(a\mu _\nu )$$
(20)
Due to enormous systematic errors, our results range over a large interval and no plateau is seen. The tree-level value of eq. (17) is given by
$$\frac{Z_P}{Z_S}=1+\frac{2\underset{\nu }{}\mathrm{sin}^2\left(ap_\nu /2\right)}{am}$$
(21)
where the second term on the r.h.s. is the tree-level $`𝒪(a)`$ effect arising from the propagator sum in the numerator of eq. (17). In Fig. 2(b) we have corrected our results by this factor, in the spirit of KLM tree-level improvement (cf. ref. ). Although this factor accounts for a big part of the discretization errors, the remaining effects are still large and render the method unreliable.
In Fig. 3 we compare the $`\beta =6.2`$ results for the RI/MOM and WIq determinations, as a function of the (lattice) scale $`a^2\overline{\mu }^2`$. It is clear that at small scales there is a large discrepancy due to the sizeable Goldstone pole contamination of the RI/MOM method, which decreases with increasing scale. At large scales it is impossible to discern whether the discrepancy is due to a remnant Goldstone pole contamination or the discretization effects. Moreover, even at larger scales, the renormalization window of the RI/MOM method, compared to the plateau displayed by the WIq results, is rather poor. In ref. , the central values of the RI/MOM results were taken at $`(a\overline{\mu })^20.8`$ and the systematic errors estimated as the spread of the values in the region $`0.8(a\overline{\mu })^21.5`$. This choice takes into account the instability of the renormalization window, yielding a RI/MOM result with large errors. On the contrary, the WI data displays a good plateau already at rather small scales. Our choice $`(a\overline{\mu })^2[0.2,0.8]`$ for the WIq plateau gives very accurate results.
In order to estimate the discretization errors, we collect in Fig. 4, the various determinations of $`Z_P/Z_S`$ at $`\beta =6.2`$ and $`6.4`$. The most accurate estimates are those based on the two WIs (WIh and WIq). They are not, however fully compatible, due to discretization effects (recall that WI results do not suffer from the Goldstone pole problem). This indicates that these effects introduce a systematic error of $`15\%`$ at $`\beta =6.2`$, which decreases to $`10\%`$ at $`\beta =6.4`$. The RI/MOM determination, which is far less accurate due to the instability of the renormalization window, is nevertheless in the same range of values. Thus, it is clear that any remnant Goldstone pole contamination in the RI/MOM result is not superior in magnitude to the uncertainties arising from discretization effects. Also, the fact that the spread of the various determinations of $`Z_P/Z_S`$ decrease with increasing $`\beta `$, indicates that the discrepancies are due to (decreasing) finite cutoff effects.
Our new WI results demonstrate that the discrepancy of $`1015\%`$ for $`Z_P/Z_S`$ carries over to $`Z_P`$, computed (non-perturbatively) from
$$Z_P=\left[\frac{Z_P}{Z_S}\right]^{\mathrm{WI}}Z_S^{\mathrm{RI}/\mathrm{MOM}}.$$
(22)
In practice we use the WIq determination for the ratio $`Z_P/Z_S`$ in the above equation; i.e. combining eq. (19) with the RI/MOM renormalization condition (8) for the scalar density, we have
$$Z_P=Z_\psi \frac{m_1m_2}{m_1\mathrm{\Gamma }_P(a\mu ;am_1,am_1)m_2\mathrm{\Gamma }_P(a\mu ;am_2,am_2)}.$$
(23)
The wave-function renormalization constant $`Z_\psi `$ has been computed in ref. . In Table 1 we collect our results for $`Z_P`$, $`Z_S`$ and their ratio. Consequently, at these $`\beta `$ values, determinations of the quark mass (and chiral condensate) which differ in the choice of non-perturbative renormalization (i.e. WIq, WIh, RI/MOM), would also display the same variations <sup>5</sup><sup>5</sup>5In refs. , a good agreement was observed between the quark mass (and chiral condensate), renormalized with $`Z_S`$ and that renormalized with $`Z_P`$ in the RI/MOM scheme. This is equivalent to the good agreement between the central values of the WIh and RI/MOM estimates of $`Z_P/Z_S`$. It is clear from fig. 4 that this is accidental.. These differences should vanish upon extrapolation to the continuum limit.
It is interesting to note that, upon expressing the first of WIs (18) in terms of renormalized quantities and using the RI/MOM renormalization condition for $`\mathrm{\Gamma }_S`$, the relation
$$\left[\widehat{m}_1(\mu )\widehat{m}_2(\mu )\right]^{\text{RI/MOM}}=\left[\frac{1}{12}\text{Tr}\widehat{𝒮}^1(\mu ;\widehat{m}_1)\frac{1}{12}\text{Tr}\widehat{𝒮}^1(\mu ;\widehat{m}_2)\right]^{\text{RI/MOM}}$$
(24)
can readily be obtained. This gives a non-perturbative quark-mass determination which does not require knowledge of $`Z_S`$ (or $`Z_P`$); cf. ref. . Moreover, it does not suffer from large $`𝒪(ap)`$ discretization effects (they cancel in the difference of quark propagators at any quark mass values), which would otherwise have to be isolated by chiral extrapolation. This determiantion does however require the computation of the wave function renormalization $`Z_\psi `$ in the RI/MOM scheme. Since this is a mass-independent renormalization scheme, the above can only be implemented in the limit $`m^2/\mu ^21`$.
## Acknowledgments
We wish to thank V. Lubicz, G. Martinelli and S. Sint for numerous illuminating discussions. A.V. thanks the DESY Theory Group for its hospitality during the early stages of this work. L. G. has been supported in part under DOE grant DE-FG02-91ER40676.
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# X–ray Emission of Mkn 421: New Clues From Its Spectral Evolution. I. Temporal Analysis
## 1. Introduction
Blazars are radio–loud AGNs characterized by strong variability, large and variable polarization, and high luminosity. Radio spectra smoothly join the infrared-optical-UV ones. These properties are successfully interpreted in terms of synchrotron radiation produced in relativistic jets and beamed into our direction due to plasma moving relativistically close to the line of sight (e.g. Urry & Padovani up95, 1995). Many blazars are also strong and variable sources of GeV $`\gamma `$–rays, and in a few objects the spectrum extends up to TeV energies. The hard X– to $`\gamma `$–ray radiation forms a separate spectral component, with the luminosity peak located in the MeV–TeV range.
The emission up to X–rays is thought to be due to synchrotron radiation from high energy electrons in the jet, while it is likely that $`\gamma `$-rays derive from the same electrons via inverse Compton (IC) scattering of soft (IR–UV) photons –synchrotron or ambient soft photons (e.g. Sikora, Begelman & Rees sbr94, 1994, Ghisellini et al. gg\_sed98, 1998).
The contributions of these two mechanisms characterize the average blazar spectral energy distribution (SED), which typically shows two broad peaks in a $`\nu F_\nu `$ representation (e.g. von Montigny et al. vmon95, 1995; Sambruna, Maraschi & Urry smu96, 1996; Fossati et al. 1998a, ): the energies at which the peaks occur and their relative intensity provide a powerful diagnostic tool to investigate the properties of the emitting plasma, such as electron energies and magnetic field (e.g. Ghisellini et al. gg\_sed98, 1998). Moreover variability studies, both of single band and of simultaneous multifrequencies data, constitute the most effective means to constrain the emission mechanisms at work in these sources as well as the geometry and modality of the energy dissipation. The quality and amount of X–ray data on the brightest sources start to allow us to perform a thorough temporal analysis as function of energy and determine the spectral evolution with good temporal resolution.
In X–ray bright BL Lacs (HBL, from High-energy-peak-BL Lacs, Padovani & Giommi pg95, 1995) the synchrotron maximum (usually) occurs in the soft-X–ray band, and the inverse Compton emission extends in some cases to the TeV band where – thanks to ground based Cherenkov telescopes – four sources have been detected up to now: Mkn 421 (Punch et al. punch92, 1992), Mkn 501 (Quinn et al. quinn96, 1996), 1ES 2344+514 (Catanese et al. catanese\_2344\_98, 1998), and PKS 2155–304 (Chadwick et al. chadwick98, 1998). If the interpretation of the SED properties in terms of synchrotron and IC radiation is correct, a correlation between the X–ray and TeV emission is expected.
Mkn 421 ($`z`$ = 0.031) is the brightest BL Lac object at X–ray and UV wavelengths and the first extragalactic source discovered at TeV energies, where dramatic variability has been observed with doubling times as short as 15 minutes (Gaidos et al. gaidos96, 1996). As such it was repeatedly observed with X–ray satellites, including BeppoSAX. Remarkable X–ray variability correlated with strong activity at TeV energies has been found on different occasions (Macomb et al. macomb95, 1995, 1996, Takahashi et al. takahashi96, 1996, Fossati et al. 1998b, , Maraschi et al. maraschi\_letter, 1999). In particular, the 1998 BeppoSAX data presented here were simultaneous with a large TeV flare detected by the Whipple Observatory (Maraschi et al. maraschi\_letter, 1999).
This paper is the first of two, which present the results of a uniform, detailed spectral and temporal analysis of BeppoSAX observations of Mkn 421 performed during 1997 and 1998. Here we focus on the data reduction and the timing analysis, and also discuss the results on the spectral variability derived from the different properties of the flux variations in different energy bands.
The paper is organized as follows. We briefly summarize the characteristics of BeppoSAX (§2), and introduce the observations studied (§3). We then address the temporal analysis of the variability, considering several energy bands and comparing the light curve features by means of a few simple estimators for the 1997 and 1998 observations (§4). The remarkable flare observed in 1998 is the object of a further deeper analysis, reported in Section §5, focused on timescales and time lags. Section 6 contains a summary of the results of the temporal analysis, preparing the ground for the comprehensive discussion presented in Paper II (Fossati et al. fossati\_II, 2000). There they are considered together with the results of the spectral analysis and thus used to constrain a scenario able to interpret the complex spectral and temporal findings.
## 2. BeppoSAX overview
For an exhaustive description of the Italian/Dutch BeppoSAX mission we refer to Boella et al. (boella97, 1997) and references therein. The narrow field coaligned instrumentation (NFI) on BeppoSAX consists of a Low Energy Concentrator Spectrometer (LECS), three Medium Energy Concentrator Spectrometers (MECS), a High Pressure Gas Scintillation Proportional Counter (HPGSPC), and a Phoswich Detector System (PDS). The LECS and MECS have imaging capabilities in the 0.1–10 keV and 1.3–10 keV energy band, respectively, with energy resolution of 8% at 6 keV. At the same energy, the angular resolution is about 1.2 arcmin (Half Power Radius). In the overlapping energy range the MECS effective area (150 cm<sup>2</sup>) is $``$ 3 times that of the LECS. Furthermore the exposure time for the LECS is limited by stronger operational constraints to avoid UV light contamination through the entrance window (LECS instrument is operated during Earth dark time only). The HPGSPC covers the range 4–120 keV, and the PDS the range 13–300 keV. In the overlapping interval the HPGSPC has a better energy resolution than the PDS, but it is less sensitive of both PDS and MECS. Therefore, HPGSPC data will not be discussed in this paper.
The present analysis is based on the SAXDAS linearized event files for the LECS and the three MECS experiments, together with appropriate background event files, as produced at the BeppoSAX Science Data Center (rev 0.2, 1.1 and 2.0). The PDS data reduction was performed using the XAS software (Chiappetti & Dal Fiume chiappetti\_dalfiume, 1997) according to the procedure described in Chiappetti et al. (chiappetti\_2155, 1999).
## 3. Observations
Mkn 421 has been observed by BeppoSAX in the springs of 1997 and 1998. The journal of observations is given in Table 1.
### 3.1. 1997
The 1997 observation comprised several pointings spanning the interval between April 29<sup>th</sup> and May 7<sup>th</sup>. MECS data are not available for May 7<sup>th</sup> because of the failure of the detector unit 1 on May 6<sup>th</sup>. In this paper we will not consider the LECS data of this last day of the 1997 campaign, because unfortunately LECS data alone do not provide useful spectral and variability information. The net exposure time (excluding the May 7<sup>th</sup> data) was $``$ 52 and $``$ 117 ks for LECS and MECS respectively, while the on–source time coverage computed as the sum of the observations between each T<sub>start</sub> and T<sub>stop</sub> is of about 58 hr.
Results on the first half of the 1997 campaign (April 29<sup>th</sup> to May 1<sup>st</sup>) have been presented by Guainazzi et al. (guainazzi\_mkn421, 1999), where the details and motivation of the observations are given. We re–analyzed those data along with the new ones, applying the same techniques, in order to obtain a homogeneous set of results, necessary for a direct comparison.
### 3.2. 1998
In 1998 BeppoSAX observed Mkn 421 as part of a long monitoring campaign involving BeppoSAX, ASCA (Takahashi, Madejski, & Kubo takahashi99\_veritas, 1999), RossiXTE (Madejski et al., in preparation), and coverage from ground based TeV observatories (Maraschi et al. maraschi\_letter, 1999). The BeppoSAX observation comprised two distinct long pointings, started respectively on April 21<sup>st</sup> and 23<sup>rd</sup>. The total net exposure time for MECS telescopes has been of 64 ks and the LECS one adds up to 50 ks. Hereinafter we will refer to the two 1998 pointings simply as April 21<sup>st</sup> and April 23<sup>rd</sup>, leaving out the year.
The actual on–source time was 55.5 hr. Unfortunately celestial and satellite constraints in late April 1998 were such that in each orbit there are two intervals during which the reconstructed attitude is undefined. One occurs when the source is occulted by the dark Earth, but the second longer one occurs during source visibility periods. Thus there are about 19 min per orbit when the NFIs are pointing at the source, but there is a gap in attitude reconstruction despite the high Earth elevation angle. During these intervals the source position is drifting along the $`\pm `$X detector axis, from a position in the center to about 12 arcmin off, passing the strongback window support. These intervals are excluded by the LECS/MECS event file generation. Using the XAS software it is possible to accumulate MECS images also during such intervals, showing that the source is actually drifting. The drift causes a strong energy dependent modulation, clearly visible in light curves accumulated with XAS, and which would be very difficult to recover, as requires to build a response matrix integrated from different position dependent matrices weighted according to the time spent in each position. We then did not try to recover these data for the analysis.
In the case of PDS the source remains inside the collimator flat–top during the drift, and therefore we used for the accumulation all events taken when the Earth elevation angle was greater than 3 degrees (inclusive of attitude gaps).
## 4. Temporal Analysis
Light curves have been accumulated for different energy bands. The input photon lists were extracted from the full dataset selecting events in a circular region centered on the position of the point source. The extraction radii are 8 and 6 arcmin for LECS and MECS respectively, chosen to be large enough to collect most of the photons in the whole energy range<sup>1</sup><sup>1</sup>1For bright and soft sources – this matters especially because the LECS point spread function (PSF) gets rapidly broad below the Carbon edge, i.e. E $``$ 0.3 keV – it is recommended to select quite large values. : the used values ensure that $`<`$ 5 % of the photons are missed, at all energies.
For MECS data we considered the merged photon list of all the available MECS units, which were 3 for the 1997 observations and only 2 (namely MECS units 2 and 3) in 1998.
The expected background contribution is less than a fraction of a percent (typically of the order of 0.2–0.5 %) and therefore the light curves have not been background subtracted.
In order to maximize the information on the spectral/temporal behavior, the choice of energy bands shall be such that they are as independent as possible. This should take into account the instrumental efficiencies and, at a certain level, the details of the spectral shape. The main trade–off is the number of photons in each band, that should be large enough to obtain statistically meaningful results. Taking into account the main features of the LECS and MECS effective areas (e.g. the Carbon edge at 0.29 keV which provides with independent energy bands below and above this energy) and the steepness of the Mkn 421 spectrum (above a few keV the available photons quickly become insufficient), we determined the following bands: 0.1–0.5 keV and 0.5–2 keV for LECS data, 2–3 and 4–6 keV for MECS data. Their respective “barycentric” energies<sup>2</sup><sup>2</sup>2These have been computed directly from the count spectra and thus they are not model dependent. Moreover, despite of the observed spectral variability, their values do not change more than a few eV. are 0.26, 1.16, 2.35, 4.76 keV, respectively, thus providing a factor $`20`$ leverage, useful to test the energy dependence of the variability characteristics.
### 4.1. 1997
The four resulting light curves for 1997 are presented in Figure 1. The source shows a high degree of flux variability, with possibly a major flare between the third and fourth pointings. The vertical scale is logarithmic allowing a direct comparison of the amplitude of variations at different energies. As anticipated we are not going to consider further the LECS data of May 7<sup>th</sup>, however for sake of completeness we report that the count rate was of 0.37$`\pm `$0.01 and 1.06$`\pm `$0.04 cts/s for the 0.1–0.5 and 0.5–2.0 keV energy bands respectively, comparable to the level measured on May 5<sup>th</sup>.
The comparison of data with the overlaid reference grid shows that the amplitude of variability is larger at higher energies, as commonly observed in blazars at energies above the SED peak (e.g. Ulrich, Maraschi & Urry umu97, 1997). We will discuss this issue more quantitatively in section §4.3.
### 4.2. 1998
The 1998 light curves are shown in Figure 2, (see the caption for details about the axis scales). The overall appearance is somewhat different, being dominated by a single isolated flare at the beginning of the campaign. The variability amplitude is similar to the 1997 one in each corresponding energy band (see also §4.3).
The most striking and important result of the campaign is that in correspondence with the X–ray flare of April 21<sup>st</sup> a sharp TeV flare was detected by the Whipple Cherenkov Telescope, with amplitude of a factor 4 and a halving time of about 2 hr (Maraschi et al. maraschi\_letter, 1999). In Figure 3 the Whipple TeV (E$``$2 TeV) light curve is shown together with the BeppoSAX ones: LECS (0.1–0.5 keV), MECS (4–6 keV) and PDS (12–26 keV) instruments. The peaks in the 0.1–0.5 keV, 4–6 keV and 2 TeV light curves are simultaneous within one hour. Note however that the TeV variation appears to be both larger and faster than the X–ray one. The 12–26 keV light curve shows a broader and slightly delayed peak with respect to lower energy X–rays. However the very limited statistics does not allow us to better quantify this event. A more detailed account on this particular result is presented and discussed by Maraschi et al. (maraschi\_letter, 1999).
It is also worth noticing that the positive (flux) offset of the beginning of April 23<sup>rd</sup> with respect to the end of the April 21<sup>st</sup> suggests the presence of a second flare occurring between the two pointings. This seems to be also confirmed by the RossiXTE All Sky Monitor (ASM) light curve, shown in Fig. 4 together with the MECS one in the (2–10 keV) nominal working energy range of the ASM. Single ASM $``$90 s dwells have been rebinned in 14400 s (i.e. 4 hr) intervals, weighting the contribution of each single dwell on its effective exposure time and on the quoted error. The MECS light curve bin is 1000 s. Indeed the ASM light curve appears to unveil the presence of a second flare occurring in between the two BeppoSAX pointings, with a brightness level similar to that of the detected one.
On the other hand the few Whipple data points at the time of this putative second X–ray flare (at T = (100–120) $`\times 10^3`$ s, see Figure 3) do not indicate any (major) TeV activity: the count rate measured by the Whipple telescope was in fact significantly lower than that measured simultaneously with the first X–ray flare.
### 4.3. Comparison of 1997 and 1998 variability characteristics
We can characterize the variability in different X–ray bands by two commonly used estimators: the fractional root mean square ($`r.m.s.`$) variability parameter F<sub>var</sub>, and the minimum “doubling/halving time” T<sub>short</sub> (for definitions and details see Appendix A, and Zhang et al. zhang\_2155, 1999).
#### 4.3.1 Fractional r.m.s. variability
For each (0.1–0.5, 0.5–2, 2–3 and 4–6 keV) band we computed F<sub>var</sub> for the whole 1997 and 1998 datasets which have the same on–source coverage (i.e. $``$ 55 hr).
Moreover, for the 1997 we considered three (partially overlapping) subsets of data each spanning a time interval comparable to the length of the whole 1998 observation, i.e. $``$ 300 ks. The intervals cover the first 300 ks, a middle section, and the last 300 ks, as marked on Fig. 1.
For each dataset F<sub>var</sub> is computed for 6 different light curve binning intervals, namely 200, 500, 1000, 1500, 2000 and 2500 s, and then the weighted average<sup>3</sup><sup>3</sup>3To the F<sub>var</sub> values obtained for each different time binning, we assigned as weight $`1/\sigma ^2`$, where $`\sigma `$ is their respective one sigma uncertainty. of these values is computed. The results are listed in Table 2 and shown in Figure 5.
If we treat all the datasets as independent and momentarily do not consider the subset–1 of the 1997 light curve (empty triangles in Fig. 5) on which we will comment later, it appears that:
1. In terms of root mean square variability the lower energy flux is less variable than the higher energy one, as already pointed out in section §4.1. This holds independent of both the length of the time interval and the state of the source.
2. There is no relation between the brightness level and the variability amplitude, since although in 1998 the source was at least a factor 2 brighter than in 1997 the corresponding F<sub>var</sub> are not larger.
What is different for the subset–1 of the 1997 light curve ?
First two caveats on F<sub>var</sub>. This quantity is basically only sensitive to the average excursion around the mean flux, and does not carry any information about either duty cycle or the actual flux excursion. Therefore in order for F<sub>var</sub> to meaningfully represent the variability, the “extrema” of the source variability have to be sampled. This is clearly not the case for the first 1997 sub–light curve (see Figure 1). Furthermore, F<sub>var</sub> is estimated from the comparison of the variances of the light curve and the measurements, but while the latter ones probably obey a Gaussian distribution and are much smaller than the former ones, the “probability” for the source to show a significant deviation from the average (i.e. the count rate histogram distribution) is in most cases highly non–Gaussian. For a typical “flaring” light curve this has a “core” at low count rates, resulting from the long(er) time spent by the source in between flares, and a very extended tail to high(er) rates. The implicit assumption of Gaussian variability can thus affect the estimate of F<sub>var</sub> (as for large amplitude and low duty cycle variability).
F<sub>var</sub> is thus a particularly poor indicator especially for observations with a window function like the one of 1997, for which the source temporal covering fraction is of the same order of the single outburst, making likely to miss the flare peak, and the total span of the observation does not significantly exceed the frequency of peaks. We therefore believe this affects the results relative to the first subsample of 1997.
In order to obtain a variability estimator more sensitive to the dynamic range towards higher count rates (flaring) we also considered a modified definition of it, F<sub>var,med</sub>: to study only the width (amplitude) of the brighter tail only data above the median count rate are considered to compute variance and expected variance. Note that the modified definition is still affected by the problem of “representativity” of all the source states. Since the median is always smaller than the average (for all energies, and for all trial time binnings), the values of F<sub>var,med</sub> are larger than F<sub>var</sub>, but in every other respect the results are qualitatively unchanged.
#### 4.3.2 Doubling/halving timescale
Similarly, the “minimum halving/doubling” timescale has been computed for each energy band taking the average of values for the different input light curve binning (but discarding the 200 s one because too strongly subject to spurious results). As in Zhang et al. (zhang\_2155, 1999) we rejected a T<sub>short,ij</sub> if the fractional uncertainty was larger than 20 %.
In order to minimize the contamination by isolated data points we filtered the light curves excluding data lying at more than 3-$`\sigma `$ from the average computed over the 6 nearest neighbors (3 on each side), and for the 500 s binning light curves we did not consider pairs between data closer than 3 positions along the time series. Finally instead of taking the single absolute minimum value of T<sub>short,ij</sub> (over all possible $`i,j`$ pairs) T<sub>short</sub> is determined as the average over the 5 lowest values.
The resulting T<sub>short</sub> are listed at the bottom of Table 2 and shown in Figure 6a. The main findings are:
1. there is not significant difference between 1997 and 1998 observations, at any energy.
2. There is weak evidence that the softest X–ray band exhibits slower variability, while
3. the timescales for the three higher energy bands are indistinguishable, and there is no sign of a trend with energy.
We also distinguish doubling and halving timescales. Unfortunately it has not been possible to obtain an estimate of the doubling time for the lowest energy band for 1998 data, because there are only a few, uncertain, data during the rise of the flare. We find that (see Table 2, and Figure 6b,c):
1. again the softest X–ray band shows longer (halving) timescales;
2. there is marginal evidence that the doubling timescale is shorter than the halving one (i.e. the rise is faster than the fall).
This latter result seems to hold both for 1997 and 1998 light curves and for each single trial binning time, and for each energy (except in 5 cases out of 35)<sup>4</sup><sup>4</sup>4We stress here that although the individual timescales are statistically consistent with each other, the behavior appear to be systematic at all energies.. However it is strongly weakened by the averaging over the several time binnings because T<sub>short</sub> slightly increases with the duration of the time bins, thus yielding a larger uncertainty on the average.
## 5. 1998 detailed temporal analysis
We performed a more detailed analysis of the characteristics of the April 21<sup>st</sup> flare, in particular to determine any significant different behavior among different energy bands. Three are the main goals: i) to measure the flare exponential (or power–law) decay timescales (or power–law index); ii) to evidence possible differences in the flare rise and decay timescales; iii) to look for time lags among the different energies.
We considered the energy ranges already adopted in Section 4.
### 5.1. Decay Timescales
A first estimate of the decay timescales has been presented by Maraschi et al. (maraschi\_letter, 1999), where a simple exponential law has been fit to the decaying phase of the light curves, and a dependence of the e–folding timescale on the energy band has been found. Here we analyze the issue of determining the timescale in more detail and in particular we consider the requirement for the presence of an underlying steady emission.
In particular, the decay timescales for different energy bands have been estimated by fitting the post–flare light curves with an exponential decay superimposed to a constant (quasi–steady) flux eventually constrained by the fit (we immediately stress that an additional non–zero steady emission is necessary in order to obtain a meaningful fit with an exponential decay). In one set of fits the underlying contribution is set to zero. More precisely, we performed the fits with three possible conditions for the level of the steady contribution: un–constrained, constrained and zero.
The analytical expression of the function used in the fits is:
$`\mathrm{F}(t)`$ $`=`$ $`\mathrm{F}_{\mathrm{steady}}+\mathrm{F}_{\mathrm{flaring}}e^{\frac{\mathrm{t}\mathrm{T}_{\mathrm{ref}}}{\tau }}`$
$`=`$ $`\mathrm{F}_{\mathrm{steady}}\left(1+e^{\frac{\mathrm{t}\mathrm{T}_{\mathrm{ref}}}{\tau }}\right)`$
The model parameters are then the decay timescale $`\tau `$, the absolute value of F<sub>steady</sub>, and the ratio $``$ between the flaring and the steady component taken at a reference time T<sub>ref</sub>, which ideally should correspond to the peak of the flare.
The “reference time” T<sub>ref</sub> for April 21<sup>st</sup> flare was T$`{}_{\mathrm{ref}}{}^{}=25\times 10^3`$ s (T$`=0=10924.0`$ TJD), and T$`{}_{\mathrm{ref}}{}^{}=170\times 10^3`$ s for April 23<sup>rd</sup>. Note that the value of $``$ obtained from the fit to the second dataset is only a lower limit on the actual flare amplitude since we do not have an estimate of the time at which the peak of this putative second flare occurred.
In any case, as our focus is actually on the properties of the first, well defined, outburst, we used the parameters obtained from the April 23<sup>rd</sup> data only to constrain the contribution from the steady component. Furthermore, it turns out that the determination of the timescale $`\tau `$ is not affected by the uncertainty on its contribution (see Table 3 and Fig. 8).
We adopted a 1000 s binning for the 0.5–2 keV and 2–3 keV data, which have better statistics, while we the binning time for the 0.1–0.5 keV and 4–6 keV light curves is of 2000 and 1500 s, respectively. The analysis has been performed in two stages, for each energy band:
* Fit to the April 21<sup>st</sup> and 23<sup>rd</sup> datasets independently, yielding values of $`\tau `$, F<sub>steady</sub> and $``$. As an example, in Figure 7 the 2–3 keV light curve is shown together with its best fit models for April 21<sup>st</sup> and 23<sup>rd</sup>.
* Fit of the April 21<sup>st</sup> post–flare phase constraining the level of the underlying steady component to be below the level of April 23<sup>rd</sup> (the actual constraint is the upper bracket for a two parameters 90% confidence interval). We assume that the underlying emission varies on a timescale longer than that spanned by the observations.
As anticipated we also modeled the decay of the April 21<sup>st</sup> flare setting to zero the offset.
The results are summarized in Tables 3 and 4. The best fit $`\tau `$ for all the three cases are plotted in Fig. 8 as a function of energy.
The 0.1–0.5 keV timescale is not affected by the constraint on the baseline flux and the 4–6 keV confidence interval only suffers a minor cut (smaller errorbar), while both the 0.5–2 and 2–3 keV parameters are significantly changed because the best fit un–constrained level of the steady contribution is significantly higher than allowed by the April 23<sup>rd</sup> data (see for instance Fig. 7).
The main findings are:
* the decay timescales depend critically on the presence/ absence of a contribution by a non variable component. This is true not only for the value themselves (for the case without baseline are between a factor 2 and 4 longer) but also for the relationship between the timescales and energy. In fact:
+ In the cases with baseline, the timescales range between 30 and 45$`\times 10^3`$ s, and do not show a clear (if any) relationship with the energy, rather suggesting an achromatic post–flare evolution. In fact, according to a $`\chi ^2`$ test, the values of $`\tau `$ for the 4 energies are consistent with coming from the same distribution (the so called “null hypothesis”).
+ On the contrary, in the case of pure exponential decay the timescales follow a weak inverse relationship with the energy, $`\tau E^{0.18\pm 0.02}`$ (both when considering $`\tau `$ or $`E`$ as the dependent variable). The $`\chi ^2`$ probability in favor of the “null hypothesis” is in this case negligible.
* the amplitude of the variability (as measured by $``$) positively correlates with the energy, as reported in many cases for HBLs (e.g. Sambruna et al. sambruna\_exosat, 1994, Ulrich et al. umu97, 1997).
#### 5.1.1 Is there a steady component ?
Variability timescales and their energy dependence carry precious information on the physics of the source: the puzzling dependence of the results on the hypothesis of the presence of a steady component requires further analysis.
Here we only discuss a simple test of the robustness of the results of previous section, that however provides useful hints. This consists in determining $`\tau `$ for different choices of the end time (T<sub>end</sub>) of the interval over which the fit is performed. The underlying idea is that if the shape of the decay is close to the one we are assuming (exponential with/without an offset) one can expect that the model parameters do not change much as a function of T<sub>end</sub>. On the contrary, if the chosen model provides only a poor description of the actual decay characteristics, the parameters will likely take different values depending on T<sub>end</sub>, unless of course we postulate a process whose timescale changes with time.
For each of the three cases (constrained, unconstrained, zero–baseline) we tried 5 different choices of the end time, namely T$`{}_{\mathrm{end}}{}^{}=60`$, 70, 80, 90 and 100$`\times 10^3`$ seconds (from T<sub>0</sub>), the last one corresponding to the entire duration of the data coverage on April 21<sup>st</sup>.
To evaluate the likelihood that the 5 values of $`\tau `$ obtained for these choices of T<sub>end</sub> are drawn from the same (normal) distribution, we considered the probability $`𝒫`$ of the $`\chi ^2`$ computed with respect to their weighted<sup>5</sup><sup>5</sup>5To each measure we assigned as $`\sigma `$ (and in turn as weight $`1/\sigma ^2`$), the mean between its $`+`$ and $``$ uncertainties. average<sup>6</sup><sup>6</sup>6It should be noted that the weighted average has the property of minimizing the $`\chi ^2`$ for a dataset, and thus the resulting values of probability that the data do belong to a unique distribution are upper limits..
In Table 4 we summarize the results obtained under each of the different hypotheses. There are a few points to note:
1. in the zero–baseline hypothesis the timescales obtained changing T<sub>end</sub> are not consistent with being several different realizations from a unique distribution. Furthermore the clear trend of (decreasing) $`\tau `$ with (decreasing) T<sub>end</sub> supports the possibility that the variation of $`\tau `$ is significant.
2. On the contrary in both cases with baseline the values of $`\tau `$’s for the 0.1–0.5, 0.5–2 and 2–3 keV bands are very close for all values of T<sub>end</sub>. For the 4–6 keV band the scatter is somewhat larger, likely due to the fact that at this energy there is the largest amplitude of variability, i.e. the smallest contribution of the putative baseline. The same kind of effect of lowered sensitivity to the details of a putative baseline, is responsible for the good result of the 4–6 keV band in the case without baseline.
3. Even in the case without baseline the values of $`\tau `$ for the various energy bands become very similar for T$`{}_{\mathrm{end}}{}^{}=60\times 10^3`$ s. Therefore even if the failure of the test by the pure exponential decay model could be interpreted as the signature of a mechanism whose characteristic timescale changes with time (unlike a true exponential decay), there is anyway an indication that in the earlier stages of the decay the evolution is not energy dependent.
On the basis of this simple consistency check we can therefore infer that if the decay were exponential there is evidence for the presence of an underlying non–variable component. Thus the post–flare phase does not show any dependence on energy, and the timescale is of the order of 30–40 ks.
### 5.2. Power Law Decay
As an alternative description of the decay of April 21<sup>st</sup> flare we consider a power law, where F$`{}_{\mathrm{flaring}}{}^{}(\mathrm{T}/\mathrm{T}_{\mathrm{ref}})^\eta `$, again leaving the possibility of an offset from a steady component. The results are summarized in Table 3 and Figure 9.
A fundamental difference with respect to the exponential decay case is that an additional steady component is never required: the power law fits yield only (not very stringent) upper limits on F<sub>steady</sub> and only for the 0.1–0.5 keV band the formal best fit value for F<sub>steady</sub> is different from zero (F$`{}_{\mathrm{steady}}{}^{}=0.26`$ cts/s).
Again there is only a weak relationship of the decay properties with energy, with a marginal indication that the flux decrease is faster at higher energies. The values for the 4 energy bands are consistent with a common (weighted average) value of $`\eta =0.583`$, while for the subset comprising only the three higher energy datasets the probability for the “null hypothesis” is reduced to about 0.07.
### 5.3. Rise vs. Fall
A further interesting question is whether the rise and the decay phases of the flare are characterized by the same timescale (i.e. if the flare is symmetric or asymmetric) and in particular how these properties might correlate with energy.
The observational coverage of the rise is not good enough to apply the same direct technique used for the decay (§5.1), to constrain the timescale of an exponential rise, with all the parameters left free. Thus the flux level of the steady component is fixed at the best fit value obtained for the decay, or set to 0.
We are mainly interested in the comparison of the properties of the outburst just before and after the top of the flare (see discussion in Paper 2), and the interval of about 17–18 ks preceding the flare is long enough for this purpose.
We thus focused on the “tip” of the flare, considering data on the decay side only up to the time when the flux reaches approximately the level at which the observation starts. The duration of the post–peak intervals needed to decrease the flux to the pre–flare level are of about 15–20, 25, 30–35 and 40 ks for the four energy bands. Since the uncertainties on the decay time of the light curves are not larger than 5 ks, the differences in duration and the trend (longer time at higher energy) are real.
The differences between rise and fall timescales, estimated with fits to the light curves, are shown in Table 5, and plotted in Figure 10. The flare is symmetric for the softest energy band, while it is definitely asymmetric in the other three bands, with an asymmetry which appears to increase systematically with energy: at higher energies the rise phase is increasingly faster than the decay (the significance is reported in Table 5).
It is worth stressing that these numbers indicate only relative changes in the timescales before and after the peak of the flare. On the other hand, since we found that the decay timescales are very similar once the baseline flux is taken into account (§5.1), this asymmetry probably means that the rise times get shorter with increasing energy.
### 5.4. Time Lag
We performed a detailed cross correlation analysis using two different techniques suited to unevenly sampled time series: the Discrete Correlation Function (DCF, Edelson & Krolik edelson\_krolik88, 1988) and the Modified Mean Deviation (MMD, Hufnagel & Bregman hufnagel\_bregman92, 1992). Moreover, Monte Carlo simulations, taking into account ”flux randomization” (FR) and ”random subset selection” (RSS) of the data series (see Peterson et al. peterson\_etal98, 1998), were used to statistically determine the significance of the time lags between different X–ray bands obtained with the DCF and MMD. We refer for the relevant details of such analysis to Zhang et al. (zhang\_2155, 1999).
We binned the light curves over 300 s in the 0.1–1.5 and 3.5–10 keV bands (whose effective barycentric energies are E $`0.8`$ and E $`5`$ keV), and a trial time step of 720 s is adopted for both the DCF and MMD. DCF amplitude versus time lag is plotted in Figure 11. A negative value means that variations in the 3.5–10 keV band light curve lag those occurring in the 0.1–1.5 keV one (i.e. hard lag). The best Gaussian fits for both DCF and MMD result in negative time lags of $`2.8\pm 0.2`$ (DCF) and $`2.1\pm 0.3`$ (MMD) ks, indicating that the medium energy X–ray photons lag the soft ones.
The Cross Correlation Peak Distribution (CCPD) obtained from the FR/RSS Monte Carlo simulations is shown in Figure 12 both for the DCF and the MMD methods. The average lags resulting from CCPD are $`2.7_{1.2}^{+1.9}`$ ks for DCF, and $`2.3_{0.7}^{+1.2}`$ ks for MMD (1 $`\sigma `$ confidence intervals, two sided with respect to the average), confirming the significance of the above results with high ($`>`$ 90 %) confidence.
The total integral probabilities for a negative lag (that would be the actual measure of the confidence of the “discovery” of the hard lag) are $`95.0`$% (DCF) and $`98.7`$% (MMD).
## 6. Summary of Variability Properties
We presented a comprehensive temporal analysis of the flux variability characteristics in several energy bands, of BeppoSAX observations of Mkn 421. The primary results of this study are the following:
1. The fractional r.m.s. variability is higher at higher energies.
2. The fractional r.m.s. variability does not change with the brightness state of the source.
3. The minimum halving/doubling time is longer for the softest energy band.
4. The minimum halving/doubling time at a given energy does not change with the brightness state of the source.
5. There is a hint that the doubling timescale is shorter than the halving timescale.
The findings 1 and 2, on F<sub>var</sub>, confirm the well known (although not fully explained) phenomenology of HBLs (see Ulrich et al. umu97, 1997 and references therein).
The results on the variability timescales (findings 3, 4, 5) show good agreement with those found with the more thorough analysis of the 1998 flare, that can be summarized as:
1. The flare decay is consistent with being achromatic, both if modeled as an exponential decay with an additional contribution from a steady component and in the case of a power law decay.
2. The results on timescales are at odds with the simplest possibility of interpreting the decay phase, i.e. that it is driven by the radiative cooling of the emitting electrons. This would give rise, in the simplest case, to a dependence of the timescale on energy, $`\tau \mathrm{E}^{1/2}`$. The tracks corresponding to this relationship are overlaid to the data in Fig. 8, and it is clear that they can not be reconciled.
3. The harder X–ray photons lag the soft ones, with a delay of the order of few ks. This finding is opposite to what is normally found in the best monitored HBL where the soft X–rays lag the hard ones (e.g. Takahashi et al. takahashi96, 1996 for Mkn 421; Urry et al. urry\_2155, 1993, Zhang et al. zhang\_2155, 1999, and Kataoka et al. kataoka\_2155, 2000 for PKS 2155–304; Kohmura et al. kohmura, 1994 for H0323$`+`$022). The latter behavior is usually interpreted in terms of cooling of the synchrotron emitting particles.
4. Possible asymmetry in the rise/decay phases: the flare seems to be symmetric at the energies corresponding (roughly) to the peak of the synchrotron component (see the spectral analysis in Paper II), while it might have a faster rise at higher energies. This could be connected to the observed hard–lag.
5. Finally, the characteristics of the April 21<sup>st</sup> flare suggest the presence of a quasi–stationary emission contribution, which seems to be dominated by a highly variable peaked spectrum.
## 7. Conclusions
BeppoSAX has observed Mkn 421 in 1997 and 1998. We analyzed and interpreted the combined spectral and temporal evolution in the X–ray range. During these observations the source has shown a large variety of behaviors, providing us a great wealth of information, but at the same time revealing a richer than expected phenomenology.
In this paper we have presented the first part of the analysis, focused on the study of the variability properties.
The fact that the higher energy band lags the softer one (with a delay of the order of 2–3 ks) and the energy dependence of the shape of the light curve during the flare (with faster flare rise time at higher energies) provide strong constraints on any possible time dependent particle acceleration mechanism. In particular, if we are indeed observing the first direct signature of the ongoing particle acceleration, progressively “pumping” electrons from lower to higher energies, the measure of the delay between the peaks of the light curves at the corresponding emitted frequencies would provide a tight constraint on the timescale of the acceleration process.
The decomposition of the observed spectrum into two components (a quasi stationary one and a peaked, highly variable one) might allow us to determine the nature and modality of the energy dissipation in relativistic jets.
In Paper II we complement these findings with those of the time resolved spectral analysis and develop a scenario to interpret the complex spectral and temporal phenomenology.
We are grateful to the BeppoSAX Science Data Center (SDC) for their invaluable work and for providing standardized product data archive, and to the RossiXTE ASM Team. We thank Gianpiero Tagliaferri and Paola Grandi for their contribution to our successful BeppoSAX program, and for useful comments and the anonymous referee for useful suggestions which have improved the clarity of the paper. AC, MC and YHZ acknowledge the Italian MURST for financial support. This research was supported in part by the National Science Foundation under Grant No. PHY94–07194 (AC). Finally, GF thanks Cecilia Clementi for providing tireless stimulus.
## Appendix A A. Definitions of F<sub>var</sub> and T<sub>short</sub>
The fractional $`r.m.s.`$ variability amplitude is a useful parameter to characterize the variability in unevenly sampled light curves. It is defined as the square root of the so called excess variance (e.g., Nandra et al. nandra\_97, 1997). This parameter, also known as the true variance (Done et al. done\_92, 1992), is computed by taking the difference between the variance of the overall light curve and the variance due to measurement error, normalized by dividing by the average squared flux (count rate).
We consider a dataset $`F_i(t_i)`$ ($`i=1,N`$), with an uncertainty $`\sigma _i`$ assigned to each point.
The fractional r.m.s. variability parameter is then defined as:
$$\mathrm{F}_{\mathrm{var}}\frac{\left(\sigma _F^2\mathrm{\Delta }_F^2\right)}{F}^{1/2},$$
(A1)
where
$`\sigma _F^2`$ $``$ $`{\displaystyle \frac{1}{N1}}{\displaystyle \underset{i=1}{\overset{N}{}}}(F_iF)^2`$ (A2)
$`\mathrm{\Delta }_F^2`$ $``$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\sigma _i^2.`$ (A3)
The definition of the “shortest variability timescale” is the following:
$`\mathrm{T}_{\mathrm{short}}`$ $``$ $`\underset{i,j}{\mathrm{min}}(\mathrm{T}_{\mathrm{short},ij})`$ (A4)
$`\mathrm{T}_{\mathrm{short},ij}`$ $``$ $`\left|{\displaystyle \frac{F_{ij}\mathrm{\Delta }T_{ij}}{\mathrm{\Delta }F_{ij}}}\right|,`$ (A5)
where $`F_{ij}(F_i+F_j)/2`$, $`\mathrm{\Delta }F_{ij}F_iF_j`$, and $`\mathrm{\Delta }T_{ij}T_iT_j`$.
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# 1 Cosmological constant
## 1 Cosmological constant
As is well known, the cosmological constant appears as a constant in the Einstein equations :
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=8\pi G_NT_{\mu \nu }+\lambda g_{\mu \nu },$$
(1)
where $`G_N`$ is Newton’s constant and $`T_{\mu \nu }`$ is the energy-momentum tensor. The cosmological constant $`\lambda `$ is thus of the dimension of an inverse length squared. It was introduced by Einstein in order to build a static universe model, its repulsive effect compensating the gravitational attraction, but, as we now see, constraints on the expansion of the Universe impose for it a very small upper value.
It is more convenient to work in the specific context of a Friedmann universe, with a Robertson metric :
$$ds^2=dt^2a^2(t)\left[\frac{dr^2}{1kr^2}+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)\right],$$
(2)
where $`a(t)`$ is the cosmic scale factor. Implementing energy conservation into the Einstein equations then leads to the Friedmann equation which gives an expression for the Hubble parameter $`H`$ :
$$H^2\frac{\dot{a}^2(t)}{a^2(t)}=\frac{1}{3}\left(\lambda +8\pi G_N\rho \right)\frac{k}{a^2},$$
(3)
where, using standard notations, $`\dot{a}`$ is the time derivative of the cosmic scale factor, $`\rho =T_{}^{0}{}_{0}{}^{}`$ is the energy density and the term proportional to $`k`$ is a spatial curvature term (see (2)). Note that the cosmological constant appears as a constant contribution to the Hubble parameter.
Evaluating each term of the Friedmann equation at present time allows for a rapid estimation of an upper limit on $`\lambda `$. Indeed, we have $`H_0=h_0\times 100\mathrm{km}.\mathrm{s}^1\mathrm{Mpc}^1`$ with $`h_0`$ of order one, whereas the present energy density $`\rho _0`$ is certainly within one order of magnitude of the critical energy density $`\rho _c=3H_0^2/(8\pi G_N)=h_0^22.10^{26}`$ kg.m<sup>-3</sup>; moreover the spatial curvature term certainly does not represent presently a dominant contribution to the expansion of the Universe. Thus, (3) implies the following constraint on $`\lambda `$ :
$$|\lambda |H_0^2.$$
(4)
In other words, the length scale $`\mathrm{}_\mathrm{\Lambda }|\lambda |^{1/2}`$ associated with the cosmological constant must be larger than $`H_0^1=h_0^1.10^{26}`$ m, and thus a macroscopic distance.
This is not a problem as long as one remains classical. Indeed, $`H_0^1`$ provides a natural macroscopic scale for our present Universe. The problem arises when one tries to combine gravity with the quantum theory. Indeed, from the Newton’s constant and the Planck constant $`\mathrm{}`$ one can construct a mass scale or a length scale
$`m_P`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{}c}{8\pi G_N}}}=2.4\times 10^{18}\mathrm{GeV}/\mathrm{c}^2,`$
$`\mathrm{}_P`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{m_Pc}}=8.1\times 10^{35}\mathrm{m}`$
The above constraint now reads :
$$\mathrm{}_\mathrm{\Lambda }|\lambda |^{1/2}\frac{1}{H_0}10^{60}\mathrm{}_P.$$
(5)
In other words, there are more than sixty orders of magnitude between the scale associated with the cosmological constant and the scale of quantum gravity.
A rather obvious solution is to take $`\lambda =0`$. This is as valid a choice as any other in a pure gravity theory. Unfortunately, it is an unnatural one when one introduces any kind of matter. Indeed, set $`\lambda `$ to zero but assume that there is a non-vanishing vacuum (i.e. groundstate) energy: $`<T_{\mu \nu }>=<\rho >g_{\mu \nu }`$, then the Einstein equations (1) read
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=8\pi G_NT_{\mu \nu }+8\pi G_N<\rho >g_{\mu \nu },$$
(6)
The last term is interpreted as an effective cosmological constant :
$$\lambda _{eff}=8\pi G_N<\rho >\frac{\mathrm{\Lambda }^4}{m_P^2}.$$
(7)
Generically, $`<\rho >`$ receives a non-zero contribution from symmetry breaking: for instance, the scale $`\mathrm{\Lambda }`$ would be typically of the order of $`100`$ GeV in the case of the gauge symmetry breaking of the Standard Model or $`1`$ TeV in the case of supersymmetry breaking. But the constraint (5) now reads:
$$\mathrm{\Lambda }10^{30}m_P10^3\mathrm{eV}.$$
(8)
It is this very unnatural fine-tuning of parameters (in explicit cases $`<\rho >`$ and thus $`\mathrm{\Lambda }`$ are functions of the parameters of the theory) that is referred to as the cosmological constant problem, or more accurately the vacuum energy problem.
## 2 The role of supersymmetry
If the vacuum energy is to be small, it may be easier to have it vanishing through some symmetry argument. Global supersymmetry is the obvious candidate. Indeed, the supersymmetry algebra
$$\{Q_r,\overline{Q}_s\}=2\gamma _{rs}^\mu P_\mu $$
(9)
yields the following relation between the Hamiltonian $`H=P_0`$ and the supersymmetry generators $`Q_r`$:
$$H=\frac{1}{4}\underset{r}{}Q_r^2,$$
(10)
and thus the vacuum energy $`<0|H|0>`$ is vanishing if the vacuum is supersymmetric ($`Q_r|0>=0`$).
Unfortunately, supersymmetry has to be broken at some scale since its prediction of equal mass for bosons and fermions is not observed in the physical spectrum. Then $`\mathrm{\Lambda }`$ is of the order of the supersymmetry breaking scale, that is a few hundred GeV to a TeV.
However, the right framework to discuss these issues is supergravity i.e. local supersymmetry since locality implies here, through the algebra (9), invariance under “local” translations that are the diffeomorphisms of general relativity. In this theory, the graviton, described by the linear perturbations of the metric tensor $`g_{\mu \nu }(x)`$, is associated through supersymmetry with a spin 3/2 field, the gravitino $`\psi _\mu `$. One may write a supersymmetric invariant combination of terms in the action:
$$𝒮=d^4x\sqrt{g}\left[3m_P^2m_{3/2}^2m_{3/2}\overline{\psi }_\mu \sigma ^{\mu \nu }\psi _\nu \right],$$
(11)
where $`\sigma _{\mu \nu }=[\gamma _\mu ,\gamma _\nu ]/4`$. If the first term is made to cancel the vacuum energy, then the second term is interpreted as a mass term for the gravitino. We thus see that the criterion for spontaneous symmetry breaking changes from global supersymmetry (non-vanishing vacuum energy) to local supersymmetry or supergravity (non-vanishing gravitino mass). It is somewhat a welcome news that a vanishing vacuum energy is not tied to a supersymmetric spectrum. On the other hand, we have lost the only rationale that we had to explain a zero cosmological constant.
Let us recall for future use that, in supergravity, the potential for a set of scalar fields $`\varphi ^i`$ is written in terms of the Kähler potential $`K(\varphi ^i,\overline{\varphi }^{\overline{i}})`$ (the normalisation of the scalar field kinetic terms is simply given by the Kähler metric $`K_{i\overline{j}}=^2K/\varphi ^i\overline{\varphi }^{\overline{j}}`$) and of the superpotential $`W(\varphi ^i)`$, an holomorphic function of the fields:
$$V=e^{K/m_P^2}\left[\left(W_i+\frac{K_i}{m_P^2}W\right)K^{i\overline{j}}\left(\overline{W}_{\overline{j}}+\frac{\overline{K}_{\overline{j}}}{m_P^2}\overline{W}\right)3\frac{|W|^2}{m_P^2}\right]+\mathrm{D}\mathrm{terms}$$
(12)
where $`K_i=K/\varphi ^i`$, etc. and $`K^{i\overline{j}}`$ is the inverse metric of $`K_{i\overline{j}}`$. Obviously, the positive definiteness of the global supersymmetry scalar potential is lost in supergravity.
## 3 Observational results
Over the last years, there has been an increasing number of indications that the Universe is presently undergoing accelerated expansion. This appears to be a strong departure from the standard picture of a matter-dominated Universe. Indeed, the standard equation for the conservation of energy,
$$\dot{\rho }=3(p+\rho )H,$$
(13)
allows to derive from the Friedmann equation (3), written in the case of a universe dominated by a component with energy density $`\rho `$ and pressure $`p`$ :
$$\frac{\ddot{a}}{a}=\frac{4\pi G_N}{3}(\rho +3p).$$
(14)
Obviously, a matter-dominated ($`\rho 0`$) universe is decelerating. One needs instead a component with a negative pressure.
A cosmological constant is associated with a contribution to the energy-momentum tensor as in (6)(7):
$$T_\nu ^\mu =\mathrm{\Lambda }^4\delta _\nu ^\mu =(\rho ,p,p,p)$$
(15)
The associated equation of motion is therefore
$$p=\rho .$$
(16)
It follows from (14) that a cosmological constant tends to accelerate expansion.
The discussion of data is thus often expressed in terms of the energy density $`\mathrm{\Lambda }^4`$ stored in the vacuum versus the energy density $`\rho _M`$ in matter fields (baryons, neutrinos, hidden matter,…). It is customary to normalize with the critical density (corresponding to a flat Universe) :
$$\mathrm{\Omega }_\mathrm{\Lambda }=\frac{\mathrm{\Lambda }^4}{\rho _c},\mathrm{\Omega }_M=\frac{\rho _M}{\rho _c},\rho _c=\frac{3H^2}{8\pi G_N}.$$
(17)
The relation
$$\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }=1,$$
(18)
a prediction of many inflation scenarios, is found to be compatible with recent Cosmic Microwave Background measurements <sup>2</sup><sup>2</sup>2This follows from the fact that the first acoustic peak is expected at an “angular” scale $`\mathrm{}200/\sqrt{\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }}`$ .. It is striking that independent methods based on the measurement of different observables on rich clusters of galaxies all point towards a low value of $`\mathrm{\Omega }_M1/3`$ : mass-to-light ratio, baryon mass to total cluster mass ratio (the total baryon density in the Universe being fixed by primordial nucleosynthesis), cluster abundance. This necessarily implies a non-vanishing $`\mathrm{\Omega }_\mathrm{\Lambda }`$ (non-vanishing cosmological constant or a similar dynamical component).
There are indeed some indications going in this direction from several types of observational data. One which has been much discussed lately uses supernovae of type Ia as standard candles<sup>3</sup><sup>3</sup>3 by calibrating them according to the timescale of their brightening and fading.. Two groups, the Supernova Cosmology Project and the High-$`z`$ Supernova Search have found that distant supernovae appear to be fainter than expected in a flat matter-dominated Universe. If this is to have a cosmological origin, this means that, at fixed redshift, they are at larger distances than expected in such a context and thus that the Universe is accelerating.
More precisely, the relation between the flux $`f`$ received on earth and the absolute luminosity $``$ of the supernova depends on its redshift $`z`$, but also on the geometry of spacetime. Traditionally, flux and absolute luminosity are expressed on a log scale as apparent magnitude $`m_B`$ and absolute magnitude $`M`$ (magnitude is $`2.5\mathrm{log}_{10}`$ luminosity + constant). The relation then reads
$$m_B=5\mathrm{log}(H_0d_L)+M5\mathrm{log}H_0+25.$$
(19)
The last terms are $`z`$-independent, if one assumes that supernovae of type Ia are standard candles; they are then measured by using low $`z`$ supernovae. The first term, which involves the luminosity distance $`d_L`$, varies logarithmically with $`z`$ up to corrections which depend on the geometry. Expanding in $`z`$ <sup>4</sup><sup>4</sup>4 Of course, since supernovae of redshift $`z1`$ are now being observed, an exact expression must be used to analyze data. The more transparent form of (20) gives the general trend., one obtains :
$$H_0d_L=cz\left[1+\frac{1q_0}{2}z+\mathrm{}\right],$$
(20)
where $`q_0a\ddot{a}/\dot{a}^2`$ is the deceleration parameter. This parameter is easily obtained from (14): in a spatially flat Universe with only matter and a cosmological constant (cf. (16)), $`\rho =\rho _M+\mathrm{\Lambda }^4`$ and $`p=\mathrm{\Lambda }^4`$ which gives
$$q_0=\mathrm{\Omega }_M/2\mathrm{\Omega }_\mathrm{\Lambda }.$$
(21)
This allows to put some limit on $`\mathrm{\Omega }_\mathrm{\Lambda }`$ on the model considered here (see Fig. 1).
Let us note that the combination (21) is ‘orthogonal’ to the combination $`\mathrm{\Omega }_T\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }`$ measured in CMB experiments (see footnote preceding page). The two measurements are therefore complementary: this is sometimes referred to as ‘cosmic complementarity’.
Of course, such type of measurement is sensitive to many possible systematic effects (evolution besides the light-curve timescale correction, etc.), and this has fueled a healthy debate on the significance of present data. This obviously requires more statistics and improved quality of spectral measurements. A particular tricky systematic effect is the possible presence of dust that would dimmer supernovae at large redshift.
Other results come from gravitational lensing. The deviation of light rays by an accumulation of matter along the line of sight depends on the distance to the source
$$r=_t^{t_0}\frac{dt}{a(t)}=\frac{1}{a(t_0)H_0}\left(z\frac{1}{2}(1+q_0)z^2+\mathrm{}\right)$$
(22)
and thus on the cosmological parameters $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$. As $`q_0`$ decreases (i.e. as the Universe accelerates), there is more volume and more lenses between the observer and the object at redshift $`z`$. Several methods are used: abundance of multiply-imaged quasar sources , strong lensing by massive clusters of galaxies (providing multiple images or arcs) , weak lensing .
## 4 Quintessence
From the point of view of high energy physics, it is however difficult to imagine a rationale for a pure cosmological constant, especially if it is nonzero but small compared to the typical fundamental scales (electroweak, strong, grand unified or Planck scale). There should be physics associated with this form of energy and therefore dynamics. For example, in the context of string models, any dimensionful parameter is expressed in terms of the fundamental string scale $`M_s`$ and vacuum expectation values of scalar fields. The physics of the cosmological constant is then the physics of the corresponding scalar fields.
Introducing dynamics generally modifies the equation of state (16) to the more general form with negative pressure :
$$p=w\rho ,w<0.$$
(23)
Let us recall that $`w=0`$ corresponds to non-relativistic matter (dust) whereas $`w=1/3`$ corresponds to radiation. A network of light, nonintercommuting topological defects on the other hand gives $`w=n/3`$ where $`n`$ is the dimension of the defect i.e. $`1`$ for a string and $`2`$ for a domain wall. Finally, the equation of state for a minimally coupled scalar field necessarily satisfies the condition $`w1`$.
Experimental data may constrain such a dynamical component just as it did with the cosmological constant. For example, in a spatially flat Universe with only matter and an unknown component $`X`$ with equation of state $`p_X=w_X\rho _X`$, one obtains from (14) with $`\rho =\rho _M+\rho _X`$, $`p=w_X\rho _X`$ the following form for the deceleration parameter
$$q_0=\frac{\mathrm{\Omega }_M}{2}+(1+3w_X)\frac{\mathrm{\Omega }_X}{2},$$
(24)
where $`\mathrm{\Omega }_X=\rho _X/\rho _c`$. Supernovae results give a constraint on the parameter $`w_X`$ as shown in Fig. 2.
Similarly, gravitational lensing effects are sensitive to this new component through (22).
A particularly interesting candidate in the context of fundamental theories is the case of a scalar<sup>5</sup><sup>5</sup>5 A vector field or any field which is not a Lorentz scalar must have settled down to a vanshing value. Otherwise, Lorentz invariance would be spontaneously broken. field $`\varphi `$ slowly evolving in a runaway potential which decreases monotonically to zero as $`\varphi `$ goes to infinity . This is often referred to as quintessence. This can be extended to the case of a very light field (pseudo-Goldstone boson) which is presently relaxing to its vacuum state . We will discuss the two situations in turn.
### 4.1 Runaway quintessence
A runaway potential is frequently present in models where supersymmetry is dynamically broken. Indeed, supersymmetric theories are characterized by a scalar potential with many flat directions, i.e. directions $`\varphi `$ in field space for which the potential vanishes. The corresponding degeneracy is lifted through dynamical supersymmetry breaking, that is supersymmetry breaking through strong interaction effects. In some instances (dilaton or compactification radius), the field expectation value $`<\varphi >`$ actually provides the value of the strong interaction coupling. Then at infinite $`\varphi `$ value, the coupling effectively goes to zero together with the supersymmetry breaking effects and the flat direction is restored: the potential decreases monotonically to zero as $`\varphi `$ goes to infinity.
Dynamical supersymmetry breaking scenarios are often advocated because they easily yield the large scale hierarchies necessary in grand unified or superstring theories in order to explain the smallness of the electroweak scale with respect to the fundamental scale. Let us take the example of supersymmetry breaking by gaugino condensation in effective superstring theories. The value $`g_0`$ of the gauge coupling at the string scale $`M_s`$ is provided by the vacuum expectation value of the dilaton field $`s`$ (taken to be dimensionless by dividing by $`m_P`$) present among the massless string modes: $`g_0^2=<s>^1`$. If the gauge group has a one-loop beta function coefficient $`b`$, then the running gauge coupling becomes strong at the scale
$$\mathrm{\Lambda }M_se^{1/2bg_0^2}=M_se^{s/2b}.$$
(25)
At this scale, the gaugino fields are expected to condense. Through dimensional analysis, the gaugino condensate $`<\overline{\lambda }\lambda >`$ is expected to be of order $`\mathrm{\Lambda }^3`$. Terms quadratic in the gaugino fields thus yield in the effective theory below condensation scale a potential for the dilaton :
$$V\left|<\overline{\lambda }\lambda >\right|^2e^{3s/b}.$$
(26)
The $`s`$-dependence of the potential is of course more complicated and one usually looks for stable minima with vanishing cosmological constant. But the behavior (26) is characteristic of the large $`s`$ region and provides a potential slopping down to zero at infinity as required in the quintessence solution. A similar behavior is observed for moduli fields whose vev describes the radius of the compact manifolds which appear from the compactification from 10 or 11 dimensions to 4 in superstring theories .
Let us take therefore the example of an exponentially decreasing potential. More explicitly, we consider the following action
$$𝒮=d^4x\sqrt{g}\left[\frac{m_P^2}{2}R\frac{1}{2}^\mu \varphi _\mu \varphi V(\varphi )\right],$$
(27)
which describes a real scalar field $`\varphi `$ minimally coupled with gravity and the self-interactions of which are described by the potential:
$$V(\varphi )=V_0e^{\lambda \varphi /m_P},$$
(28)
where $`V_0`$ is a positive constant.
The energy density and pressure stored in the scalar field are respectively :
$$\rho _\varphi =\frac{1}{2}\dot{\varphi }^2+V(\varphi ),p_\varphi =\frac{1}{2}\dot{\varphi }^2V(\varphi ).$$
(29)
We will assume that the background (matter and radiation) energy density $`\rho _B`$ and pressure $`p_B`$ obey a standard equation of state
$$p_B=w_B\rho _B.$$
(30)
If one neglects the spatial curvature ($`k0`$), the equation of motion for $`\varphi `$ simply reads
$$\ddot{\varphi }+3H\dot{\varphi }=\frac{dV}{d\varphi },$$
(31)
with
$$H^2=\frac{1}{3m_P^2}(\rho _B+\rho _\varphi ).$$
(32)
This can be rewritten as
$$\dot{\rho }_\varphi =3H\dot{\varphi }^2.$$
(33)
We are looking for scaling solutions i.e. solutions where the $`\varphi `$ energy density scales as a power of the cosmic scale factor: $`\rho _\varphi a^{n_\varphi }`$ or $`\dot{\rho }_\varphi /\rho _\varphi =n_\varphi H`$. In this case, one easily obtains from (29) and (33) that the $`\varphi `$ field obeys a standard equation of state
$$p_\varphi =w_\varphi \rho _\varphi ,$$
(34)
with
$$w_\varphi =\frac{n_\varphi }{3}1.$$
(35)
Hence
$$\rho _\varphi a^{3(1+w_\varphi )}.$$
(36)
If one can neglect the background energy $`\rho _B`$, then (32) yields a simple differential equation for $`a(t)`$ which is solved as :
$$at^{2/[3(1+w_\varphi )]}.$$
(37)
Since $`\dot{\varphi }^2=(1+w_\varphi )\rho _\varphi `$, one deduces that $`\varphi `$ varies logarithmically with time. One then easily obtains from (31,32) that
$$\varphi =\varphi _0+\frac{2}{\lambda }\mathrm{ln}(t/t_0).$$
(38)
and<sup>6</sup><sup>6</sup>6 under the condition $`\lambda ^26`$ ($`w_\varphi 1`$ since $`V(\varphi )0`$).
$$w_\varphi =\frac{\lambda ^2}{3}1,$$
(39)
It is clear from (39) that, for $`\lambda `$ sufficiently small, the field $`\varphi `$ can play the role of quintessence.
But the successes of the standard big-bang scenario indicate that clearly $`\rho _\varphi `$ cannot have always dominated: it must have emerged from the background energy density $`\rho _B`$. Let us thus now consider the case where $`\rho _B`$ dominates. It turns out that the solution just discussed with $`\rho _\varphi \rho _B`$ and (39) is a late time attractor only if $`\lambda ^2<3(1+w_B)`$. If $`\lambda ^2>3(1+w_B)`$, the global attractor turns out to be a scaling solution with the following properties:<sup>7</sup><sup>7</sup>7See ref. for the case where the scalar field is non-minimally coupled to gravity.
$`\mathrm{\Omega }_\varphi {\displaystyle \frac{\rho _\varphi }{\rho _\varphi +\rho _B}}`$ $`=`$ $`{\displaystyle \frac{3}{\lambda ^2}}(1+w_B)`$ (40)
$`w_\varphi `$ $`=`$ $`w_B`$ (41)
The second equation (41) clearly indicates that this does not correspond to a quintessence solution (23).
The semi-realistic models discussed earlier tend to give large values of $`\lambda `$ and thus the latter scaling solution as an attractor. For example, in the case (26) where the scalar field is the dilaton, $`\lambda =3/b`$ with $`b=C/(16\pi ^2)`$ and $`C=90`$ for a $`E_8`$ gauge symmetry down to $`C=9`$ for $`SU(3)`$. Moreover , on the observational side, the condition that $`\rho _\varphi `$ should be subdominant during nucleosynthesis (in the radiation-dominated era) imposes to take rather large values of $`\lambda `$. Typically requiring $`\rho _\varphi /(\rho _\varphi +\rho _B)`$ to be then smaller than $`0.2`$ imposes $`\lambda ^2>20`$.
Although not quintessence, such attractor models with a fixed fraction $`\mathrm{\Omega }_\varphi `$ as in (40) have interest of their own , in particular for structure formation if $`\lambda [5,6]`$. It has been proposed recently to make the prefactor $`V_0`$ in (28) a trigonometric function in $`\varphi `$. This allows for some modulation around the previous attractor in an approximately oscillatory way: $`\mathrm{\Omega }_\varphi `$ could then have been small at the time of nucleosynthesis and be much larger at present times. Finally, very recently , such models have been coupled to a system of a Brans-Dicke field and a dynamical field characterizing the cosmological constant, with a diverging kinetic term, to provide a relaxation mechanism for the cosmological constant .
Ways to obtain a quintessence component have been proposed however. Let me sketch some of them in turn.
One is the notion of tracker field<sup>8</sup><sup>8</sup>8 Somewhat of a misnomer since in this solution, as we see below, the field $`\varphi `$ energy density tracks the radiation-matter energy density before overcoming it (in contradistinction with (40)). One should rather describe it as a transient tracker field. . This idea also rests on the existence of scaling solutions of the equations of motion which play the role of late time attractors, as illustrated above. An example is provided by a scalar field described by the action (27) with a potential
$$V(\varphi )=\lambda \frac{\mathrm{\Lambda }^{4+\alpha }}{\varphi ^\alpha }$$
(42)
with $`\alpha >0`$. In the case where the background density dominates, one finds an attractor scaling solution $`\varphi a^{3(1+w_B)/(2+\alpha )}`$, $`\rho _\varphi a^{3\alpha (1+w_B)/(2+\alpha )}`$. Thus $`\rho _\varphi `$ decreases at a slower rate than the background density ($`\rho _Ba^{3(1+w_B)}`$) and tracks it until it becomes of the same order at a given value $`a_Q`$. More precisely :
$`\varphi `$ $`=`$ $`m_P\sqrt{{\displaystyle \frac{\alpha (2+\alpha )}{3(1+w_B)}}}\left({\displaystyle \frac{a}{a_Q}}\right)^{3(1+w_B)/(2+\alpha )},`$ (43)
$`\rho _\varphi `$ $``$ $`\lambda {\displaystyle \frac{\mathrm{\Lambda }^{4+\alpha }}{m_P^\alpha }}\left({\displaystyle \frac{a}{a_Q}}\right)^{3\alpha (1+w_B)/(2+\alpha )}.`$ (44)
One finds
$$w_\varphi =1+\frac{\alpha (1+w_B)}{2+\alpha }.$$
(45)
Shortly after $`\varphi `$ has reached for $`a=a_Q`$ a value of order $`m_P`$, it satisfies the standard slowroll conditions:
$`m_P|V^{}/V|`$ $``$ $`1,`$ (46)
$`m_P^2|V^{\prime \prime }/V|`$ $``$ $`1,`$ (47)
and therefore (45) provides a good approximation to the present value of $`w_\varphi `$. Thus, at the end of the matter-dominated era, this field may provide the quintessence component that we are looking for.
Two features are particularly interesting in this respect. One is that this scaling solution is reached for rather general initial conditions, i.e. whether $`\rho _\varphi `$ starts of the same order or much smaller than the background energy density . The second deals with the central question in this game: why is the $`\varphi `$ energy density (or in the case of a cosmological constant, the vacuum energy density) emerging now? Since $`\varphi `$ is of order $`m_P`$ in this scenario, it can be rephrased here into the following: why is $`V(m_P)`$ of the order of the critical energy density $`\rho _c`$? Using (44), this amounts to a constraint on the parameters of the theory:
$$\mathrm{\Lambda }\left(H_0^2m_P^{2+\alpha }\right)^{1/(4+\alpha )}.$$
(48)
For example, this gives for $`\alpha =2`$, $`\mathrm{\Lambda }10`$ MeV, not such an unnatural value.
Let us note here the key difference between this tracking scenario and the preceding one<sup>9</sup><sup>9</sup>9 I wish to thank M. Joyce for discussions on this point.. Whereas the exponential potential model accounts for a fixed fraction $`\mathrm{\Omega }_\varphi `$ in the attractor solution (and thus $`\varphi `$ is a tracker in the strict sense), the final attractor in the tracker field solution corresponds to scalar field dominance ($`\mathrm{\Omega }_\varphi 1`$). It is the scale $`\mathrm{\Lambda }`$ which allows to regulate the time at which the scalar field starts to emerge and makes it coincide with present time. The welcome property is that the required value for $`\mathrm{\Lambda }`$ falls in a reasonable range from a high energy physics point of view. On the other hand, we will see below that the fact that the present value for $`\varphi `$ is of order $`m_P`$ is a source of problems.
Models of dynamical supersymmetry breaking easily provide a model of the type just discussed . Let us consider supersymmetric QCD with gauge group $`SU(N_c)`$ and $`N_f<N_c`$ flavors, i.e. $`N_f`$ quarks $`Q^i`$ (resp. antiquarks $`\overline{Q}_i`$), $`i=1\mathrm{}N_f`$, in the fundamental $`𝐍_𝐜`$ (resp. anti-fundamental $`\overline{𝐍}_𝐜`$ of $`SU(N_c)`$. At the scale of dynamical symmetry breaking $`\mathrm{\Lambda }`$ where the gauge coupling becomes strong<sup>10</sup><sup>10</sup>10It is given by an expression such as (25) where $`g_0`$ is the value of the gauge coupling at the large scale $`M_s`$ and $`b`$ the one-loop beta function coefficient for gauge group $`SU(N_c)`$., boundstates of the meson type form : $`\mathrm{\Pi }_{}^{i}{}_{j}{}^{}=Q^i\overline{Q}_j`$. The dynamics is described by a superpotential which can be computed non-perturbatively using standard methods :
$$W=(N_cN_f)\frac{\mathrm{\Lambda }^{(3N_cN_f)/(N_cN_f)}}{\left(\mathrm{det}\mathrm{\Pi }\right)^{1/(N_cN_f)}}.$$
(49)
Such a superpotential has been used in the past but with the addition of a mass or interaction term (i.e. a positive power of $`\mathrm{\Pi }`$) in order to stabilize the condensate. One does not wish to do that here if $`\mathrm{\Pi }`$ is to be interpreted as a runaway quintessence component. For illustration purpose, let us consider a condensate diagonal in flavor space: $`\mathrm{\Pi }_{}^{i}{}_{j}{}^{}\varphi ^2\delta _j^i`$ (see for a more complete analysis). Then the potential for $`\varphi `$ has the form (42), with $`\alpha =2(N_c+N_f)/(N_cN_f)`$. Thus,
$$w_\varphi =1+\frac{N_c+N_f}{2N_c}(1+w_B),$$
(50)
which clearly indicates that the meson condensate is a potential candidate for a quintessence component.
Another possibility for the emergence of the quintessence component out of the background energy density might be attributed to the presence of a local minimum (a “bump”) in the potential $`V(\varphi )`$: when the field $`\varphi `$ approaches it, it slows down and $`\rho _\varphi `$ decreases more slowly ($`n_\varphi `$ being much smaller as $`w_\varphi `$ temporarily becomes closer to $`1`$, cf.(35)). If the parameters of the potential are chosen carefully enough, this allows the background energy density, which scales as $`a^{3(1+w_B)}`$ to become subdominant. The value of $`\varphi `$ at the local minimum provides the scale which allows to regulate the time at which this happens. This approach can be traced back to the earlier work of Wtterich and has recently been advocated by Albrecht and Skordis in the context of an exponential potential. They argue quite sensibly that, in a “realistic” string model, $`V_0`$ in (28) is $`\varphi `$-dependent: $`V_0(\varphi )`$. This new field dependence might be such as to generate a bump in the scalar potential and thus a local minimum. Since
$$\frac{1}{V}\frac{dV}{d\varphi }=\frac{V_0^{}(\varphi )}{V_0(\varphi )}\frac{\lambda }{m_P},$$
(51)
it suffices that $`m_PV_0^{}(\varphi )/V_0(\varphi )`$ becomes temporarily larger than $`\lambda `$ in order to slowdown the redshift of $`\rho _\varphi `$: once $`\rho _\varphi `$ dominates, an attractor scaling solution of the type (38,39) is within reach, if $`\lambda `$ is not too large. As pointed out by Albrecht and Skordis, the success of this scheme does not require very small couplings.
One may note that, in the preceding model, one could arrange the local minimum in such a way as to completely stop the scalar field, allowing for a period of true inflation . The last possibility that I will discuss goes in this direction one step further. It is known under several names: deflation , kination , quintessential inflation . It is based on the remark that, if a field $`\varphi `$ is to provide a dynamical cosmological constant under the form of quintessence, it is a good candidate to account for an inflationary era where the evolution is dominated by the vacuum energy. In other words, are the quintessence component and the inflaton the same unique field?
In this kind of scenario, inflation (where the energy density of the Universe is dominated by the $`\varphi `$ field potential energy) is followed by reheating where matter-radiation is created by gravitational coupling during an era where the evolution is driven by the $`\varphi `$ field kinetic energy (which decreases as $`a^6`$). Since matter-radiation energy density is decreasing more slowly, this turns into a radiation-dominated era until the $`\varphi `$ energy density eventually emerges as in the quintessence scenarios described above.
Finally, it is worth mentioning that, even though the models discussed above all have $`w_\varphi 1`$, models with lower values of $`w_\varphi `$ may easily be constructed. One may cite models with non-normalized scalar field kinetic terms , or simply models with non-minimally coupled scalar fields . Indeed, it has been argued by Caldwell that such a “phantom” energy density component fits well the present observational data.
### 4.2 Pseudo-Goldstone boson
There exists a class of models very close in spirit to the case of runaway quintessence: they correspond to a situation where a scalar field has not yet reached its stable groundstate and is still evolving in its potential.
More specifically, let us consider a potential of the form:
$$V(\varphi )=M^4v\left(\frac{\varphi }{f}\right),$$
(52)
where $`M`$ is the overall scale, $`f`$ is the vacuum expectation value $`<\varphi >`$ and the function $`v`$ is expected to have coefficients of order one. If we want the potential energy of the field (assumed to be close to its vev $`f`$) to give a substantial fraction of the energy density at present time, we must set
$$M^4\rho _cH_0^2m_P^2.$$
(53)
However, requiring that the evolution of the field $`\varphi `$ around its minimum has been overdamped by the expansion of the Universe until recently imposes
$$m_\varphi ^2=\frac{1}{2}V^{\prime \prime }(f)\frac{M^4}{f^2}H_0^2.$$
(54)
Let us note that this is one of the slowroll conditions familiar to the inflation scenarios.
From (53) and (54), we conclude that $`f`$ is of order $`m_P`$ (as the value of the field $`\varphi `$ in runaway quintessence) and that $`M10^3`$ eV (not surprisingly, this is the scale $`\mathrm{\Lambda }`$ typical of the cosmological constant, see (8)). The field $`\varphi `$ must be very light: $`m_\varphi h_0\times 10^{60}m_Ph_0\times 10^{33}\mathrm{eV}`$. Such a small value is only natural in the context of an approximate symmetry: the field $`\varphi `$ is then a pseudo-Goldstone boson.
A typical example of such a field is provided by the axion field (QCD axion or string axion ). In this case, the potential simply reads :
$$V(\varphi )=M^4\left[1+\mathrm{cos}(\varphi /f)\right].$$
(55)
## 5 Quintessential problems
However appealing, the quintessence idea is difficult to implement in the context of realistic models . The main problem lies in the fact that the quintessence field must be extremely weakly coupled to ordinary matter. This problem can take several forms :
$``$ we have assumed until now that the quintessence potential monotonically decreases to zero at infinity. In realistic cases, this is difficult to achieve because the couplings of the field to ordinary matter generate higher order corrections that are increasing with larger field values, unless forbidden by a symmetry argument. For example, in the case of the potential (42), the generation of a correction term $`\lambda _dm_P^{4d}\varphi ^d`$ puts in jeopardy the slowroll constraints on the quintessence field, unless very stringent constraints are imposed on the coupling $`\lambda _d`$. But one typically expects from supersymmetry breaking $`\lambda _dM_S^4/m_P^4`$ where $`M_S`$ is the supersymmetry breaking scale .
Similarly, because the $`vev`$ of $`\varphi `$ is of order $`m_P`$, one must take into account the full supergravity corrections. One may then argue that this could put in jeopardy the positive definiteness of the scalar potential, a key property of the quintessence potential. This may point towards models where $`<W>=0`$ (but not its derivatives, see (12)) or to no-scale type models: in the latter case, the presence of 3 moduli fields $`T^i`$ with Kähler potential $`K=_i\mathrm{ln}(T^i+\overline{T}^i)`$ cancels the negative contribution $`3|W|^2`$ in (12).<sup>11</sup><sup>11</sup>11 Moreover, supergravity corrections may modify some of the results. For example, the presence of a (flat) Kähler potential in (12) induces exponential field-dependent factors. A more adequate form for the inverse power law potential (42) is thus $`V(\varphi )=\lambda e^{\varphi ^2/2M_P^2}\mathrm{\Lambda }^{4+\alpha }/\varphi ^\alpha `$. The exponential factor is not expected to change much the late time evolution of the quintessence energy density. Brax and Martin argue that it changes the equation of state through the value of $`w_\varphi `$.
$``$ the quintessence field must be very light . If we return to our example of supersymmetric QCD in (42), $`V^{\prime \prime }(m_P)`$ provides an order of magnitude for the mass-squared of the quintessence component:
$$m_\varphi \mathrm{\Lambda }\left(\frac{\mathrm{\Lambda }}{m_P}\right)^{1+\alpha /2}H_010^{33}\mathrm{eV}.$$
(56)
using (48). Similarly, we have seen that the mass of a pseudo-Goldstone boson that could play the rôle of quintessence is typically of the same order. This field must therefore be very weakly coupled to matter; otherwise its exchange would generate observable long range forces. Eötvös-type experiments put very severe constraints on such couplings.
Again, for the case of supersymmetric QCD, higher order corrections to the Kähler potential of the type
$$\kappa (\varphi _i,\varphi _j^{})\left[\beta _{ij}\left(\frac{Q^{}Q}{m_P^2}\right)+\overline{\beta }_{ij}\left(\frac{\overline{Q}\overline{Q}^{}}{m_P^2}\right)\right]$$
(57)
will generate couplings of order 1 to the standard matter fields $`\varphi _i`$, $`\varphi _j^{}`$ since $`<Q>`$ is of order $`m_P`$. In order to alleviate this problem, Masiero, Pietroni and Rosati have proposed a solution much in the spirit of the least coupling principle of Damour and Polyakov : the different functions $`\beta _{ij}`$ have a common minimum close to the value $`<Q>`$, which is most easily obtained by assuming “flavor” independence of the functions $`\beta _{ij}`$. This obviously eases the Eötvös experiment constraints. In the early stages of the evolution of the Universe, when $`Qm_P`$, couplings of the type (57) generate a contribution to the mass of the $`Q`$ field which, being proportional to $`H`$, does not spoil the tracker solution.
$``$ it is difficult to find a symmetry that would prevent any coupling of the form $`\beta (\varphi /m_P)^nF^{\mu \nu }F_{\mu \nu }`$ to the gauge field kinetic term. Since the quintessence behavior is associated with time-dependent values of the field of order $`m_P`$, this would generate, in the absence of fine tuning, corrections of order one to the gauge coupling. But the time dependence of the fine structure constant for example is very strongly constrained : $`|\dot{\alpha }/\alpha |<5\times 10^{17}\mathrm{yr}^1`$. This yields a limit :
$$|\beta |10^6\frac{m_PH_0}{<\dot{\varphi }>},$$
(58)
where $`<\dot{\varphi }>`$ is the average over the last $`2\times 10^9`$ years.
A possible solution is to implement an approximate continuous symmetry of the type: $`\varphi \varphi +`$ constant . This symmetry must be approximate since it must allow for a potential $`V(\varphi )`$. Such a symmetry would only allow derivative couplings, an example of which is an axion-type coupling $`\stackrel{~}{\beta }(\varphi /m_P)F^{\mu \nu }\stackrel{~}{F}_{\mu \nu }`$. If $`F_{\mu \nu }`$ is the color $`SU(3)`$ field strength, QCD instantons yield a mass of order $`\stackrel{~}{\beta }\mathrm{\Lambda }_{_{QCD}}^2/m_P`$, much too large to satisfy the preceding constraint. In any case, supersymmetry would relate such a coupling to the coupling $`\beta (\varphi /m_P)F^{\mu \nu }F_{\mu \nu }`$ that we started out to forbid.
The very light mass of the quintessence component points towards scalar-tensor theories of gravity, where such a dilaton-type (Brans-Dicke) scalar field is found. This has triggered some recent interest for this type of theories. Attractor scaling solutions have been found for non-minimally coupled fields . However, as discussed above, one problem is that scalar-tensor theories lead to time-varying constants of nature. One may either put some limit on the couplings of the scalar field or use the attractor mechanism towards General Relativity that was found by Damour and Nordtvedt . This mechanism exploits the stabilsation of the dilaton-type scalar through its conformal coupling to matter. Indeed, assuming that this scalar field $`\varphi `$ couples to matter through an action term $`𝒮_m(\psi _m,f(\varphi )g_{\mu \nu }))`$, then its equation of motion takes the form:
$$\frac{2}{3\varphi _{}^{}{}_{}{}^{2}}\varphi ^{\prime \prime }+(1w_B)\varphi ^{}=(13w_B)\frac{d\mathrm{ln}f(\varphi )}{d\varphi },$$
(59)
where $`\varphi ^{}=d\varphi /d\mathrm{ln}a`$. This equation can be interpreted as the motion of a particle of velocity-dependent mass $`2/(3\varphi _{}^{}{}_{}{}^{2})`$ subject to a damping force $`(1w_B)\varphi ^{}`$ in an external force deriving from a potential $`v_{\mathrm{eff}}(\varphi )=(13w_B)\mathrm{ln}f(\varphi )`$. If this effective potential has a minimum, the field quickly settles there. Bartolo and Pietroni have recently proposed a model of quintessence (they add a potential $`V(\varphi )`$) using this mechanism: the quintessence component is first attracted to General Relativity and then to a standard tracker solution.
Scalar-tensor theories of gravity naturally arise in the context of higher-dimensional theories and we will return to such scenarios in the next section where we discuss these theories.
All the preceding shows that there is extreme fine tuning in the couplings of the quintessence field to matter, unless they are forbidden by some symmetry. This is somewhat reminiscent of the fine tuning associated with the cosmological constant. In fact, the quintessence solution does not claim to solve the cosmological constant (vacuum energy) problem described above. If we take the example of a supersymmetric theory, the dynamical cosmological constant provided by the quintessence component clearly does not provide enough amount of supersymmetry breaking to account for the mass difference between scalars (sfermions) and fermions (quarks and leptons): at least $`100`$ GeV. There must be other sources of supersymmetry breaking and one must fine tune the parameters of the theory in order not to generate a vacuum energy that would completely drown $`\rho _\varphi `$.
In any case, the quintessence solution shows that, once this fundamental problem is solved, one can find explicit fundamental models that effectively provide the small amount of cosmological constant that seems required by experimental data.
## 6 Extra spacetime dimensions
The old idea of Kaluza and Klein about compact extra dimensions has received a new twist with the realisation, motivated by string theory , that such extra dimensions may only be felt by gravitational interactions . In other words, our 4-dimensional world of quarks, leptons and gauge interactions may constitute a hypersurface in a higher-dimensional Universe. Such a hypersurface is called a brane in modern jargon: certain types of branes (Dirichlet branes) appear as solitons in open string theories . In what follows, we will mainly consider 4-dimensional branes to which are confined observable matter as well as standard non-gravitational gauge interactions. The part of the Universe which is not confined to the brane is called the bulk (which for simplicity we will take to be 5-dimensional).
In this framework, the very notion of a cosmological constant takes a new meaning and there has been recently a lot of activity to try to unravel it. The hope is that the cosmological constant problem itself may receive a different formulation, easier to deal with.
If we think of the cosmological constant as some vacuum energy, one has the choice to add it to the brane or to the bulk. The consequences are quite different:
$``$ If we introduce a vacuum energy $`\lambda _b>0`$ on the brane, it creates a repulsive gravitational force outside (i.e. in the bulk). Indeed, a result originally obtained by Ipser and Sikivie in the case of a domain wall may be adapted here as follows: let $`p`$ and $`\rho `$ be the pressure and energy density on the brane, then if $`\rho +3p`$ is positive (resp. negative), a test body may remain in the bulk stationary to the brane if it accelerates away from (resp. towards) the brane.<sup>12</sup><sup>12</sup>12 One may note that, if the expansion in the brane is standard, then, according to (14), the expansion in the brane is decelerating (resp. accelerating). In the case of a positive cosmological constant, $`\rho =p=\lambda _b`$ and $`\rho +3p=2\lambda _b<0`$.
Projected back to our 4-dimensional brane-world, this yields a different behaviour from the one seen in a standard 4-dimensional world. For example, the vacuum energy contributes to the Hubble parameter describing the expansion of the brane world in a (non-standard) quadratic way : $`H^2=\lambda _b^2/(36M^6)+\mathrm{}`$, where $`M`$ is the fundamental 5-dimensional scale.
$``$ If we introduce a vacuum energy $`\lambda _B`$ ) in the 5-dimensional bulk (this $`\lambda _B`$ is then of mass dimension 5), this will induce a potential for the modulus field whose vev measures the radius of the compact dimension. Let us call for simplicity $`R`$ this modulus, which is often referred to as the radion. Then in the case of a single compact dimension, $`V(R)=\lambda _BR`$ .
The contribution of this bulk vacuum energy to the square of the Hubble parameter on the brane is standard (linear) : $`H^2=\lambda _B/(6M^3)+\mathrm{}`$
Allowing both types of vacuum energies allows to construct static solutions with a cancelling effect in the bulk. Indeed, if one imposes the condition :
$$\lambda _B=\frac{\lambda _b^2}{6M^3},$$
(60)
the effective 4-dimensional cosmological constant, i.e. the constant term in the Hubble parameter $`H`$, vanishes.
A striking property of this type of configuration is that it allows to localize gravity on the brane. This is the so-called Randall-Sundrum scenario (see also for earlier works). The 5-dimensional Einstein equations are found to allow for a 4-dimensionally flat solution with a warp factor (i.e. an overall fifth dimension-dependent factor in front of the four-dimensional metric) :
$$ds^2=e^{|y\lambda _b|/(3M^3)}\eta _{\mu \nu }dx^\mu dx^\nu +dy^2$$
(61)
if the condition (60) is satisfied. Let us note that this condition ensures that the bulk is anti-de Sitter since $`\lambda _B<0`$. If $`\lambda _b>0`$, one finds a single normalisable massless mode of the metric which is interpreted as the massless 4-dimensional graviton. The wave function of this mode turns out to be localized close to the brane, which gives an explicit realisation of 4-dimensional gravity trapping. There is also a continuum of non-normalisable massive modes (starting from zero mass) which are interpreted as the Kaluza-Klein graviton modes.
Of course, the Randall-Sundrum condition (60) is another version of the standard fine tuning associated with the cosmological constant. One would like to find a dynamical justification to it.
Some progress has recently been made in this direction . The presence of a scalar field in the bulk, conformally coupled to the matter on the brane allows for some relaxation mechanism that screens the 4-dimensional cosmological constant from corrections to the brane vacuum energy. Let us indeed consider such a scalar field, of the type discussed above in the context of scalar-tensor theories. The action is of the following form :
$`𝒮`$ $`=`$ $`{\displaystyle d^5x\sqrt{g^{(5)}}\left[\frac{M^3}{2}R^{(5)}\frac{1}{2}^N\varphi _N\varphi V(\varphi )\right]}`$ (62)
$`+𝒮_m(\psi _m,g_{\mu \nu }f(\varphi ))`$
where the fields $`\psi _m`$ are matter fields localized on the brane, located at $`y=0`$, and $`g_{\mu \nu }`$ is the 4-dimensional metric ($`N`$ are 5-dimensional indices whereas $`\mu `$, $`\nu `$ are 4-dimensional indices). We will be mostly interested in the 4-dimensional vaccuum energy so that we can write the 4-dimensional matter action as :
$$𝒮_m=d^4x\sqrt{g^{(4)}}\lambda _bf^2(\varphi ).$$
(63)
Five-dimensional Einstein equations projected on the brane, provide the following Friedmann equation:
$$H^2=\frac{1}{18M^6}\lambda _b^2f^2(\varphi )\left[f^2(\varphi )3M^2f_{}^{}{}_{}{}^{2}(\varphi )\right]+\frac{1}{3}V(\varphi ).$$
(64)
The other equations, including the $`\varphi `$ equation of motion, ensure that this vanishes, irrespective of the precise value of $`\lambda _b`$, for the following metric :
$$ds^2=e^{\alpha (y)}dx^\mu dx^\nu +dy^2,$$
(65)
where the derivative of the function $`\alpha (y)`$ with respect to $`y`$ is fixed on the brane by junction conditions (assuming a symmetry $`yy`$)
$$\alpha ^{}(0)=\frac{\lambda }{3M^5}f^2(\varphi )|_{\mathrm{y}=0}.$$
(66)
In other words, the cosmological constant is, to a first order, not sensitive to the corrections to the vacuum energy coming from the Standard Model interactions.
For specific values of the potential, such a dynamics localizes the gravity around the brane. For example, with vanishing potential, the solution of the equations is obtained for
$$f(\varphi )=e^{\varphi /(M\sqrt{3})}.$$
(67)
One obtains a flat 4-dimensional spacetime (indeed, in this case, this is the unique solution ) although the vacuum energy may receive non-vanishing corrections.
The price to pay is the presence of a singularity close to the brane. It remains to be seen what is the interpretation of this singularity, how it should be treated and whether this reintroduces fine tuning . Also a full cosmological treatment, i.e. including time dependence, is needed.
Presumably, supersymmetry plays an important role in this game if one wants to deal with stable solutions. Supersymmetry indeed may prove to be in the end the rationale for the vanishing of the cosmological constant. The picture that would emerge would be one of a supersymmetric bulk with vanishing cosmological constant and with supersymmetry broken on the brane (remember that supersymmetry is related to translational invariance) . Models along these lines have been discussed recently by Gregory, Rubakov and Sibiryakov : the four-dimensional gravity is localized on the brane due to the existence of an unstable graviton boundstate. Presumably in such models one does not recover the standard theory of gravity.
## 7 Conclusion
The models discussed above are many. This is not a surprise since the cosmological constant problem, although it has attracted theorists for decades, has not received yet a convincing treatment. What is new is that one expects in a not too distant future a large and diversified amount of observational data that should allow to discreminate among these models. One may mention the MAP and PLANCK satellites on the side of CMB measurements. The SNAP mission should provide, on the other hand, large numbers of type Ia supernovae which should allow a better handle on this type of measurements and a significant increase in precision. But other methods will also give complementary information: lensing, galaxy counts , gravitational wave detection , …
Acknowledgments:
I wish to thank Christophe Grojean, Michael Joyce, Reynald Pain, James Rich and Jean-Philippe Uzan for discussions and valuable comments on the manuscript. I thank the Theory Group of Lawrence Berkeley National Lab where I found ideal conditions to finish writing these lecture notes.
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# TRI-PP-00-23 MKPH-T-00-03 FZJ-IKP(TH)-2000-03 Radiative pion capture by a nucleon
## I Introduction
Radiative pion capture by a nucleon is one of the obvious reactions to use as a testbed for heavy baryon chiral perturbation theory (HBChPT). For charged pions, the reaction begins at $`O(p)`$, which is leading order in HBChPT, and it is known that the $`O(p^3)`$ result for s-wave multipole is in reasonable agreement with most measurements . The p-wave multipoles however seem never to have been calculated. This is in contrast to the neutral pion case where both s- and p-wave multipoles have been extensively discussed . A calculation beyond the s-wave provides insight into the convergence of the chiral expansion and also serves to determine some of the HBChPT parameters that are required for other reactions, such as radiative muon capture by a nucleon, where the existing experimental data are in surprising disagreement with theoretical expectations . Thus an investigation of the p-wave multipoles in the charged case is a useful thing to do and is the primary aim of this work.
In the present work, the only explicit fields in the chiral Lagrangian are the pions and nucleons. Other physical particles will enter the calculation through their implicit contributions to the Lagrangian’s parameters (LEC’s). For some reactions it is advantageous to include the $`\mathrm{\Delta }(1232)`$ explicitly, as done for example in Ref. , and it is possible that this could be a useful approach for radiative pion capture as well, once one goes away from threshold. However, it is consistent to absorb such resonances in the LEC’s and we shall see that for the present reaction a reasonable fit to the data can be obtained when the $`\mathrm{\Delta }(1232)`$ is left implicit in the HBChPT parameters.
Experimental data for the $`\pi ^{}p\gamma n`$ differential cross section was reported fifteen years ago from a TRIUMF experiment at beam energies of $`T_\pi =27.4`$ and 39.3 MeV . A recent TRIUMF experiment has taken data at $`T_\pi =9.88`$, 14.62 and 19.85 MeV . There is also very recent data for the inverse reaction $`\gamma pn\pi ^+`$ taken very near threshold at $`T_\gamma `$ 153 MeV corresponding to $`T_\pi 3`$ MeV. In this study, we will not attempt to apply HBChPT to energies above 40 MeV.
There are at least two modern theoretical discussions of radiative charged-pion capture (both discussions actually address the inverse reaction: charged-pion photoproduction). One is an HBChPT study of the s-wave at threshold by Bernard, Kaiser and Meißner, and another is a dispersion theoretical analysis of s- and p-waves by Hanstein, Drechsel and Tiator. The present work goes beyond threshold and also explicitly computes the p-wave multipoles. The comparison of our work to the threshold results of Ref. is found to be quite interesting and to provide a useful constraint on our results.
In section II, we establish the general expressions for kinematics, multipoles and the differential cross section. Section III discusses the HBChPT calculation and section IV presents and discusses our results, both at threshold and in general. Section V contains a summary of what has been learned from this effort, and what the next steps could be.
## II Kinematics and Multipoles
In radiative charged-pion capture by a nucleon, a low energy $`\pi ^\pm `$ with four-momentum $`q^\mu =(E_\pi ,\stackrel{}{q})`$ in the center-of-mass system gets absorbed by a slowly moving nucleon of mass $`m_N`$. In the final state, one observes a recoiling nucleon and a low energy photon with polarization four-vector $`ϵ^\mu =(ϵ_0,\stackrel{}{ϵ})`$ and four-momentum $`k^\mu =(\omega ,\stackrel{}{k})`$. The pion’s center-of-mass energy is related to $`s`$, the square of the total energy in the center of mass, and to $`T_\pi `$, the kinetic energy in the lab frame by
$$E_\pi =\frac{s+m_\pi ^2m_N^2}{2\sqrt{s}}=\frac{m_\pi ^2+m_N(m_\pi +T_\pi )}{\sqrt{(m_N+m_\pi )^2+2m_NT_\pi }},$$
(1)
where $`m_\pi `$ and $`m_N`$ are respectively the pion and nucleon masses. The analogous formulas for the photon energy in the center of mass are
$$\omega =\frac{sm_N^2}{2\sqrt{s}}=\frac{m_NT_\gamma }{\sqrt{m_N^2+2m_NT_\gamma }},$$
(2)
where $`T_\gamma `$ is the corresponding laboratory gamma energy for the inverse process. All energy dependence will be expressed via the pion energy in the center-of-mass system. For the energy of the final state photon we therefore employ
$$\omega =E_\pi \frac{m_\pi ^2}{2m_N}+\frac{E_\pi m_\pi ^2}{2m_N^2}+𝒪(1/m_N^3).$$
(3)
The differential cross section for the pion capture process in the center-of-mass frame is,
$$\frac{\mathrm{d}\sigma ^{\pi N\gamma N}}{\mathrm{d}\mathrm{\Omega }_\gamma }=\frac{\omega }{|\stackrel{}{q}|}\frac{1}{2}\underset{\mathrm{pol}^{}\mathrm{s}}{}\left|\right|^2,$$
(4)
and that for the inverse (photoproduction) reaction is
$$\frac{\mathrm{d}\sigma ^{\gamma N\pi N}}{\mathrm{d}\mathrm{\Omega }_\pi }=\frac{|\stackrel{}{q}|}{\omega }\frac{1}{4}\underset{\mathrm{pol}^{}\mathrm{s}}{}\left|\right|^2,$$
(5)
where $``$ is the amplitude defined below. Notice that Eqs. (4) and (5) explicitly contain the average over initial and sum over final spins and polarizations and that the two cross sections are related by the usual detailed balance relation.
Essentially all previous work has dealt with the inverse, photoproduction, process, $`\gamma N\pi N`$ and the conventions for that process are by now well established. Thus in the Coulomb gauge with $`ϵ_0=0`$ and the transversality condition $`\stackrel{}{ϵ}\stackrel{}{k}=0`$ the amplitude for that process can be written in terms of the T-matrix as
$`^{\gamma N\pi N}={\displaystyle \frac{m_N}{4\pi \sqrt{s}}}Tϵ`$ $`=`$ $`F_1(E_\pi ,x)i\chi ^{}\stackrel{}{\sigma }\stackrel{}{ϵ}\chi +F_2(E_\pi ,x)\chi ^{}\stackrel{}{\sigma }\widehat{q}\stackrel{}{\sigma }\left(\widehat{k}\times \stackrel{}{ϵ}\right)\chi +`$ (7)
$`F_3(E_\pi ,x)i\chi ^{}\stackrel{}{\sigma }\widehat{k}\stackrel{}{ϵ}\widehat{q}\chi +F_4(E_\pi ,x)i\chi ^{}\stackrel{}{\sigma }\widehat{q}\stackrel{}{ϵ}\widehat{q}\chi ,`$
where $`\sigma ^i`$ is a Pauli matrix in spin space between the two-component spinors of the incoming/outgoing nucleon ($`\chi /\chi ^{}`$), $`ϵ`$ is the photon polarization vector and $`x=\mathrm{cos}\theta `$ corresponds to the cosine of the angle between the photon and the pion momenta.
Furthermore, each structure amplitude $`F_i(E_\pi ,x)`$ ($`i`$=1,2,3,4) can be decomposed into three isospin channels ($`a`$=1,2,3)
$$F_i^a(E_\pi ,x)=F_i^{()}(E_\pi ,x)iϵ^{a3b}\tau ^b+F_i^{(\mathit{0})}(E_\pi ,x)\tau ^a+F_i^{(+)}(E_\pi ,x)\delta ^{a3},$$
(8)
where $`\tau ^a`$ denotes a Pauli matrix in isospin space. The physical structure amplitudes are then obtained from the linear combinations
$`F_i^{\gamma n\pi ^{}p}`$ $`=`$ $`\sqrt{2}\left[F_i^{(\mathit{0})}F_i^{()}\right],`$ (9)
$`F_i^{\gamma p\pi ^+n}`$ $`=`$ $`\sqrt{2}\left[F_i^{(\mathit{0})}+F_i^{()}\right].`$ (10)
The full physics content of this process is encoded in the four structure amplitudes $`F_i`$, which are complicated functions of $`E_\pi `$ and $`\theta `$, and in the amplitude of Eq. (7), the square of which is used to get the cross section. However it may be more intuitive to discuss the underlying physics in terms of a multipole decomposition. The HBChPT formalism which we are employing in the following sections involves an expansion in terms of the pion energy divided by a scale of approximately 1 GeV, i.e. it is only reliable in a kinematic region of low energy pions. With this in mind we restrict the multipoles we consider to s- and p-waves only. They can be found from the $`F`$-amplitudes via
$`E_{0+}\left(E_\pi \right)`$ $`=`$ $`{\displaystyle _1^1}𝑑x\left\{{\displaystyle \frac{1}{2}}F_1(E_\pi ,x){\displaystyle \frac{1}{2}}xF_2(E_\pi ,x)+{\displaystyle \frac{1}{6}}\left[1P_2(x)\right]F_4(E_\pi ,x)\right\},`$ (11)
$`M_{1+}\left(E_\pi \right)`$ $`=`$ $`{\displaystyle _1^1}𝑑x\left\{{\displaystyle \frac{1}{4}}xF_1(E_\pi ,x){\displaystyle \frac{1}{4}}P_2(x)F_2(E_\pi ,x)+{\displaystyle \frac{1}{12}}\left[P_2(x)1\right]F_3(E_\pi ,x)\right\},`$ (12)
$`M_1\left(E_\pi \right)`$ $`=`$ $`{\displaystyle _1^1}𝑑x\left\{{\displaystyle \frac{1}{2}}xF_1(E_\pi ,x)+{\displaystyle \frac{1}{2}}F_2(E_\pi ,x)+{\displaystyle \frac{1}{6}}\left[1P_2(x)\right]F_3(E_\pi ,x)\right\},`$ (13)
$`E_{1+}\left(E_\pi \right)`$ $`=`$ $`{\displaystyle _1^1}dx\{{\displaystyle \frac{1}{4}}xF_1(E_\pi ,x){\displaystyle \frac{1}{4}}P_2(x)F_2(E_\pi ,x)+{\displaystyle \frac{1}{12}}[1P_2(x)]F_3(E_\pi ,x)`$ (15)
$`+{\displaystyle \frac{1}{10}}[xP_3(x)]F_4(E_\pi ,x)\},`$
with the $`P_i(x),i2`$ being Legendre polynomials.
The formulas above are those conventionally defined for the photoproduction reaction $`\gamma N\pi N`$, whereas we are interested particularly in the capture process $`\pi N\gamma N`$. The cross sections for these two processes are related trivially by the detailed balance equation arising from Eqs.(4) and (5). The relation between the amplitudes is however more complicated, arising from time reversal and depending explicitly on the phases of the parts of the amplitude. In our conventions we find (up to a possible overall, and thus irrelevant phase)
$$^{\pi N\gamma N}=[^{\gamma N\pi N}]^{}.$$
(16)
If we apply Eq. (16) to Eq. (7) to get the amplitude for pion capture the structure functions $`F_i`$ attract various phases and a complex conjugate and the order of the structures corresponding to $`F_2`$ is reversed. Putting the $`F_2`$ structures back in the original order generates extra terms and makes some of the coefficients of the four independent structures linear combinations of the $`F_i`$. Thus if we were to define the amplitude for the pion capture reaction to be of the original general form of Eq. (7) then the $`F_i`$ for pion capture will be linear combinations, complex conjugated, with various phase changes, of the $`F_i`$ for photoproduction. An alternative, and probably more sensible choice, is to define the amplitude for the capture reaction via the action of Eq. (16) on the definition used for the photoproduction direction. This eliminates the problem of linear combinations, but still leaves the two sets of $`F_i`$ related by a complex conjugate and various phase changes.
The third alternative, which is the one we adopt, is to just do the calculation for the photoproduction direction in the first place, and then make the connection to the pion capture direction at the level of the cross section. This has the advantage of keeping a close connection with the conventions and the large body of previous work dealing with photoproduction. Thus the formulas for the $`F_i`$ which we quote, and more importantly those for the multipoles, are actually for the $`\gamma N\pi N`$ direction. This means for example that our numerical results for the multipoles can be compared directly and without ambiguity with the dispersion relation calculation for photoproduction of Ref. , even though the parameters are being fixed primarily by the pion capture data.
## III The HBChPT Calculation
The HBChPT Lagrangian is ordered in powers of momenta and pion masses, which are small compared to both the chiral scale, $`4\pi F`$, and the nucleon mass, $`m_N`$,
$$_{\pi N}=_{\pi N}^{(1)}+_{\pi N}^{(2)}+_{\pi N}^{(3)}+\mathrm{}.$$
(17)
The lowest-order Lagrangian is
$$_{\pi N}^{(1)}=\overline{N}_v(iv+g_ASu)N_v$$
(18)
where
$`N_v(x)`$ $`=`$ $`\mathrm{exp}[im_{0N}vx]{\displaystyle \frac{1}{2}}(1+v/)\psi (x),`$ (19)
$`S_\mu `$ $`=`$ $`{\displaystyle \frac{i}{2}}\gamma _5\sigma _{\mu \nu }v^\nu ,`$ (20)
$`u_\mu `$ $`=`$ $`iu^{}(_\mu ir_\mu )uiu(_\mu i\mathrm{}_\mu )u^{},`$ (21)
$`_\mu `$ $`=`$ $`_\mu +\mathrm{\Gamma }_\mu iv_\mu ^{(s)},`$ (22)
$`\mathrm{\Gamma }_\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[u^{}(_\mu ir_\mu )u+u(_\mu i\mathrm{}_\mu )u^{}\right],`$ (23)
with $`m_{0N}`$ and $`g_A`$ being the lowest-order nucleon mass and axial coupling respectively. The external photon field is included via $`r_\mu =\mathrm{}_\mu =(e/2)\tau ^3A_\mu `$, and $`u`$ is a nonlinear representation of the pion fields, for example
$$u=\mathrm{exp}\left[\frac{i}{2F_0}\left(\begin{array}{cc}\pi ^0& \sqrt{2}\pi ^+\\ \sqrt{2}\pi ^{}& \pi ^0\end{array}\right)\right].$$
(24)
The parameter $`F_0`$ corresponds to the pion decay constant in the chiral limit (normalized so that the physical value $`F=92.4`$ MeV).
The higher-order Lagrangians $`_{\pi N}^{(n)}`$ will be written in the notation of Ecker and Mojžiš and are exactly the same as those used in Ref. . Results for the multipoles in the present work depend on four combinations of parameters from $`_{\pi N}^{(3)}`$, namely $`b_{10}`$, $`b_{19}`$, $`b_{21}^r(\mu )`$ and $`2b_{22}^r(\mu )+b_{23}`$, where $`\mu `$ is the renormalization scale. The numerical values of $`b_{19}`$ and $`b_{23}`$ were determined in Ref. . The three remaining parameters, $`b_{10}`$, $`b_{21}^r(\mu )`$ and $`b_{22}^r(\mu )`$, will be determined in the present work.
The calculation requires an evaluation of tree-level and one-pion-loop diagrams, which can be organized into four classes depending on whether the radiated photon is emitted from the initial nucleon, the final nucleon, the pion, or from the $`\pi NN`$ vertex. The calculation was performed in a general gauge (and is fully gauge invariant). While this meant more work, the ability to check gauge invariance provided a very important tool for eliminating errors in what was an algebraically complex calculation. The result was then reduced to the special case of $`vϵ=0`$. In this gauge, only one of the four classes of diagrams has any dependence on the unknown HBChPT parameters, $`b_{10}`$, $`b_{21}^r(\mu )`$ and $`b_{22}^r(\mu )`$, namely photon emission from the $`\pi NN`$ vertex.
Adding all contributions together gives the amplitude of Eq. (7) with the structure amplitudes, $`F_i(E_\pi ,x)`$, given explicitly in Appendix A. Although only charged-pion processes are discussed in this work, the calculation was actually performed for general isospin. We have verified that the $`\pi ^0`$ amplitudes agree with Ref. .
## IV Results
### A The differential cross section
Using our calculation from the previous section with Eq. (4) or Eq. (5) and the $``$ of Eq. (7) and the F’s of the Appendix, we can immediately compute the differential cross section. At $`O(p)`$ and $`O(p^2)`$ the result is completely determined, whereas at $`O(p^3)`$ it depends on three unknown parameters, which will now be determined via a least-squares fit to the experimental data.
Ref. provides 11 measurements of the differential cross section for $`\pi ^{}p\gamma n`$ at $`T_\pi =9.88`$, 14.62 and 19.85 MeV and Ref. provides an additional 16 measurements at $`T_\pi =27.4`$ and 39.3 MeV. A further 8 measurements, these for the inverse reaction $`\gamma pn\pi ^+`$ very near threshold ($`T_\pi 3`$ MeV), come from Ref. . We have performed fits to several subsets of this set of data, as well as to the complete set. A comparison of these fits allows us to check for consistency among the data sets and also for a possible breakdown of the HBChPT form as $`T_\pi `$ increases.
The values of the three fitted parameters are given in Table I. It is reassuring to see that within the uncertainties all of the various data sets lead to the same numerical values for these parameters, though the fit becomes more stable and the uncertainties smaller as we increase the number of data points included in the fit. It should also be noted that each of our least-squares fits actually finds two sets of parameters, characterized by nearly identical values of $`b_{21}^r`$ and $`b_{22}^r`$ but quite different values of $`b_{10}`$, depending on where the least-squares routine begins in parameter space. This presumedly reflects the fact that the cross section is quadratic in the $`b_i`$’s and that the data is not sufficiently good to distinguish the two solutions. We refer to these two minima in parameter-space as “A” and “B”, and then label our solution sets as A($`n`$) and B($`n`$), where $`n`$ is the number of experimental measurements used in the fit. For the various subsets of pion capture data, A($`n`$) and B($`n`$) give essentially indistinguishable $`\chi ^2`$ values and differential cross sections. Addition of the very low energy photoproduction data of Ref. produces a small improvement in the $`\chi ^2`$ of the A(35) solution relative to the B(35) one. The two solutions can be distinguished, however, by their quite different values of $`b_{10}`$ and also by the different individual p-wave multipoles, as will be discussed below.
The results of our best fits to the cross section data are shown in Fig. 1, along with the parameter-free $`O(p)`$ and $`O(p^2)`$ calculations and the experimental data. As these plots indicate, the $`O(p)`$ calculation disagrees with the data. $`O(p^2)`$ contributions reduce the discrepancy, but do not eliminate it. The $`O(p^3)`$ terms are necessary for a good fit to the data. The $`O(p)`$ terms clearly dominate (note the suppressed zero in the plots), but the contributions of $`O(p^2)`$ and $`O(p^3)`$ are comparable at most angles. For $`\gamma n\pi ^{}p`$ the two contributions seem to add, whereas for $`\gamma p\pi ^+n`$ they have opposite signs and tend to cancel. The fact that the $`O(p^2)`$ and $`O(p^3)`$ terms are more or less equal may raise some concern that the HBChPT expansion has not yet fully converged at $`O(p^3)`$. This point can also be made from Table I, which gives the values of the three parameters that were determined in the fits. For a nicely converging chiral expansion that just contains pions and nucleons as effective degrees of freedom, one probably would have expected each of the $`b_i`$ to acquire values near unity. The fact that we find values somewhat larger than this perhaps can be seen as an indication of the role of explicit matter fields like the $`\mathrm{\Delta }`$ isobar and vector mesons. The discussion of such issues, however, has to be delayed to a future communication. Here we only provide the first step and fix the contact terms numerically at the scale $`\mu =m_N`$. Note however that the value of $`b_{10}`$ obtained from the A (but not the B) solution is quite consistent in magnitude with the value of the parameter $`b_P`$, which is a linear combination of $`b_{10}`$ and $`b_9`$, obtained in Ref. by fitting $`\pi ^0`$ data.
### B Threshold Results
Expressions are simplified somewhat at threshold, that is in the limit in which the pion kinetic energy $`T_\pi `$ goes to zero. Using $`M_{1+}=\omega |\stackrel{}{q}|m_{1+}`$, $`M_1=\omega |\stackrel{}{q}|m_1`$, and $`E_{1+}=\omega |\stackrel{}{q}|e_{1+}`$, the multipoles, given for the photoproduction process $`\gamma +N\pi +N`$, follow directly from Eqs. (11-15) and the expressions for the $`F`$’s given in the Appendix. The resulting expressions are given purely in terms of physical quantities. For the $`(\mathit{0})`$ isospin channel we obtain
$`E_{0+}^{(\mathit{0})}(m_\pi )`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi (m_N+m_\pi )}}{\displaystyle \frac{eG_A}{2F}}\left[{\displaystyle \frac{m_\pi }{2m_N}}+{\displaystyle \frac{m_\pi ^2}{4m_N^2}}(\mu _p+\mu _n)\right],`$ (25)
$`m_{1+}^{(\mathit{0})}(m_\pi )`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi (m_N+m_\pi )}}{\displaystyle \frac{eG_A}{2F}}\left[{\displaystyle \frac{(\mu _p+\mu _n)}{6m_\pi m_N}}{\displaystyle \frac{1}{12m_N^2}}+{\displaystyle \frac{(\mu _p+\mu _n)}{6m_N^2}}+{\displaystyle \frac{2b_{10}}{3G_A(4\pi F)^2}}\right],`$ (26)
$`m_1^{(\mathit{0})}(m_\pi )`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi (m_N+m_\pi )}}{\displaystyle \frac{eG_A}{2F}}\left[{\displaystyle \frac{(\mu _p+\mu _n)}{3m_\pi m_N}}+{\displaystyle \frac{7}{24m_N^2}}{\displaystyle \frac{(\mu _p+\mu _n)}{3m_N^2}}+{\displaystyle \frac{2b_{10}}{3G_A(4\pi F)^2}}\right],`$ (27)
$`e_{1+}^{(\mathit{0})}(m_\pi )`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi (m_N+m_\pi )}}{\displaystyle \frac{eG_A}{2F}}{\displaystyle \frac{1}{24m_N^2}},`$ (28)
and for the $`()`$ isospin channel
$`E_{0+}^{()}(m_\pi )`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi (m_N+m_\pi )}}{\displaystyle \frac{eG_A}{2F}}[1+{\displaystyle \frac{m_\pi ^2}{8m_N^2}}{\displaystyle \frac{m_\pi ^2}{4m_N^2}}(\mu _p\mu _n)+{\displaystyle \frac{\pi ^2m_\pi ^2}{4(4\pi F)^2}}`$ (30)
$`{\displaystyle \frac{m_\pi ^2}{G_A(4\pi F)^2}}(2b_{19}2b_{21}^r(\mu )2b_{22}^r(\mu )b_{23}+G_A\mathrm{ln}{\displaystyle \frac{m_\pi ^2}{\mu ^2}})],`$
$`m_{1+}^{()}(m_\pi )`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi (m_N+m_\pi )}}{\displaystyle \frac{eG_A}{2F}}[{\displaystyle \frac{1}{6m_\pi ^2}}{\displaystyle \frac{1}{12m_\pi m_N}}{\displaystyle \frac{(\mu _p\mu _n)}{6m_\pi m_N}}+{\displaystyle \frac{5}{48m_N^2}}`$ (33)
$`{\displaystyle \frac{(\mu _p\mu _n)}{6m_N^2}}+{\displaystyle \frac{2G_A^2}{3(4\pi F)^2}}{\displaystyle \frac{2G_A^2\pi }{3(4\pi F)^2}}+{\displaystyle \frac{G_A^2\pi ^2}{12(4\pi F)^2}}`$
$`+{\displaystyle \frac{1}{6G_A(4\pi F)^2}}(2b_{19}4b_{22}^r(\mu )2b_{23}2G_A^3\mathrm{ln}{\displaystyle \frac{m_\pi ^2}{\mu ^2}})],`$
$`m_1^{()}(m_\pi )`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi (m_N+m_\pi )}}{\displaystyle \frac{eG_A}{2F}}[{\displaystyle \frac{1}{3m_\pi ^2}}+{\displaystyle \frac{1}{6m_\pi m_N}}{\displaystyle \frac{(\mu _p\mu _n)}{6m_\pi m_N}}+{\displaystyle \frac{1}{24m_N^2}}`$ (36)
$`{\displaystyle \frac{(\mu _p\mu _n)}{6m_N^2}}{\displaystyle \frac{4G_A^2}{3(4\pi F)^2}}{\displaystyle \frac{2G_A^2\pi }{3(4\pi F)^2}}+{\displaystyle \frac{G_A^2\pi ^2}{3(4\pi F)^2}}`$
$`{\displaystyle \frac{1}{3G_A(4\pi F)^2}}(2b_{19}4b_{22}^r(\mu )2b_{23}2G_A^3\mathrm{ln}{\displaystyle \frac{m_\pi ^2}{\mu ^2}})],`$
$`e_{1+}^{()}(m_\pi )`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi (m_N+m_\pi )}}{\displaystyle \frac{eG_A}{2F}}\left[{\displaystyle \frac{1}{6m_\pi ^2}}+{\displaystyle \frac{1}{12m_\pi m_N}}{\displaystyle \frac{5}{48m_N^2}}{\displaystyle \frac{b_{19}}{3G_A(4\pi F)^2}}\right].`$ (37)
To make contact with previous work, observe that the $`O(p)`$ and $`O(p^2)`$ parts of these expressions are just what one would obtain from an expansion of the usual Born graphs using pseudovector coupling. The $`O(p^3)`$ parts contain higher order pieces of the expansion of the Born graphs, loop contributions, and contributions from the part of the Lagrangian involving the LEC’s.
The numerical values of the threshold multipoles at each order in HBChPT are displayed in Table II. The $`O(p^3)`$ results are given for both solutions, A($`n`$) and B($`n`$). Again the results are essentially the same within errors for any of the subsets of data used, though the fit is most accurate when the full 35 points are included. The $`m_{1+}`$ and $`m_1`$ multipoles differ dramatically between A($`n`$) and B($`n`$), as they have an important dependence on $`b_{10}`$ which is quite different for the two solutions. $`e_{1+}`$ is constant, as it depends only on the parameter $`b_{19}`$ which was fixed from muon capture and $`E_{0+}`$ is nearly constant as it depends only on the parameters $`b_{22}^r,b_{21}^r,b_{19}`$ and $`b_{23}`$ which are all essentially the same for the two fits.
Also shown in Table II are the results of a dispersion theory calculation by Hanstein, Drechsel and Tiator. For the electric multipoles $`E_{0+}`$ and $`e_{1+}`$ the agreement with the HBChPT results is quite good for both the $`\pi ^+`$ and $`\pi ^{}`$ cases. For the magnetic multipoles $`m_{1+}`$ and $`m_1`$ the agreement with A(35) is good, albeit not spectacular. One must recognize however that there are uncertainties in the dispersion relations results also, which were quoted only for the $`E_{0+}`$ multipole. For the B(35) fit however the HBChPT and dispersion results for these multipoles are quite different. Thus comparison with the dispersion relation results strongly favors the A(35) solution over the B(35) one.
One can gain some further insight via a more detailed comparison with the dispersion relation results. Observe first that Eqs. (25-37) give the eight observable multipole amplitudes in terms of four parameters: $`b_{10}`$, $`b_{19}`$, $`b_{21}^r(\mu )`$ and $`2b_{22}^r(\mu )+b_{23}`$. This means that four parameter-free relations exist among the multipoles in the $`O(p^3)`$ HBChPT calculation. For example, Table III gives a set of four quantities which are independent of these four parameters, along with their values as obtained from HBChPT and dispersion theory. For these four quantities the convergence of the HBChPT expansion is good and the results agree quite well with the dispersion relation predictions of Ref. .
This idea can be carried a step further by looking at combinations of the multipoles which depend on only one or only a few of the $`b_i`$’s. Such results are tabulated in Table IV. The multipole $`e_{1+}^{()}`$ depends, in fact only weakly, on $`b_{19}`$ and one can see from the table that the HBChPT results converge well and agree with the dispersion theory result. The next two entries, $`m_{1+}^{()}`$ and $`m_1^{()}`$ depend in addition on the combination $`2b_{22}^r+b_{23}`$ and also agree with the dispersion relation results. The next entry $`E_{0+}^{()}`$ depends in addition on $`b_{21}^r`$ and the following one, $`E_{0+}^{()}+3m_\pi ^2(m_{1+}^{()}e_{1+}^{()})`$ depends only on $`b_{21}^r`$. Both show good convergence and good agreement with the dispersion theory. Finally the last two entries $`m_{1+}^{(\mathit{0})}`$ and $`m_1^{(\mathit{0})}`$ depend only on $`b_{10}`$. Here the B(35) solution is clearly ruled out by comparison with the dispersion relation results. The A(35) solution agrees moderately well, especially since the dispersion results come from taking the difference of two large numbers, and so probably have significant uncertainties. As found before however, the convergence of the magnetic multipoles is not as good as for the electric multipoles.
One can summarize the results of this evaluation of the threshold multipoles and comparison with the dispersion relation calculation of Ref. as follows. Generally the HBChPT calculation produces results for the multipoles for the physical processes that converge and that agree with the dispersion relation calculation. Likewise the various LEC’s seem to be well determined. The second solution, B(35), which could not be distinguished from the other one on the basic of $`\chi ^2`$ alone, seems to be ruled out by comparison with dispersion relation results. The weakest link appears to be in the convergence of the HBChPT expansion for the magnetic multipoles, which is not as good as that for the electric multipoles, and in the detailed combinations of multipoles depending on $`b_{10}`$ alone.
To improve the calculation it might be interesting to extend it to one higher order, which can be done still within the context of a one-loop calculation. Thus one could see if the $`O(p^4)`$ terms indicated real convergence. One might also think about including the $`\mathrm{\Delta }(1232)`$ as an explicit degree of freedom. In the present calculation $`\mathrm{\Delta }`$ effects are included implicitly in the LEC’s, which is a perfectly consistent approach. One alternatively could extract them explicitly along the lines of Ref. . Very preliminary estimates seem to indicate that such effects are relatively small in the very near threshold region we are considering, but it might be worth doing a full calculation.
Finally, as somewhat of a side issue, we note that an alternative representation of the near-threshold differential cross section which is often used is
$`{\displaystyle \frac{\omega }{|\stackrel{}{q}|}}{\displaystyle \frac{\mathrm{d}\sigma ^{\gamma N\pi N}}{\mathrm{d}\mathrm{\Omega }_\pi }}`$ $`=`$ $`A+Bx+Cx^2`$ (38)
$`A`$ $`=`$ $`|E_{0+}|^2+{\displaystyle \frac{1}{2}}|P_2|^2+{\displaystyle \frac{1}{2}}|P_3|^2`$ (39)
$`B`$ $`=`$ $`2\mathrm{R}\mathrm{e}(E_{0+}P_1^{})`$ (40)
$`C`$ $`=`$ $`|P_1^2|{\displaystyle \frac{1}{2}}|P_2|^2{\displaystyle \frac{1}{2}}|P_3|^2`$ (41)
$`P_1`$ $`=`$ $`3E_{1+}+M_{1+}M_1`$ (42)
$`P_2`$ $`=`$ $`3E_{1+}M_{1+}+M_1`$ (43)
$`P_3`$ $`=`$ $`2M_{1+}+M_1`$ (44)
However, this near-threshold result differs somewhat from the general result we have used. It is obtained by expanding the original amplitude, e.g. the pion pole contributions, and keeping terms only through $`x^2`$, which is sufficient to give the cross section in terms of s- and p-wave multipoles. In contrast we used the square of the full HBChPT amplitude to get the cross section, and only later after fitting the data extracted the s- and p-wave multipoles.
## V Summary and Outlook
We have investigated the radiative capture of a charged pion by a nucleon using heavy baryon chiral perturbation theory and have obtained explicit expressions for the amplitude and for the s- and p-wave multipoles, expressed, as is more conventional, as amplitudes for the inverse photoproduction process. Up to $`O(p^3)`$, these expressions depend upon three parameters that were determined by fitting to data for $`\pi ^{}`$ capture by a proton and for very near threshold photoproduction. Two satisfactory fits were obtained, which were indistinguishable, based only on comparison with the data.
Using the LEC’s obtained from these fits, the eight s- and p-wave multipoles (four for the $`\pi ^+`$ case and four for the $`\pi ^{}`$ case) were calculated and compared with results previously obtained from dispersion theory . In general the agreement was good for one of the fits, A(35), whereas there were significant differences when the other fit was used. This same result held for combinations of the multipoles depending on just a few of the parameters. We thus conclude that the A(35) fit gives an acceptable result, and thus that the three parameters determined in that fit, $`b_{10},b_{21}^r,b_{22}^r`$ and given in Table I are available for future studies of other reactions.
In general the convergence of the HBChPT expansion was very good for the electric multipoles, but somewhat less good for the magnetic ones. This suggests that it might be valuable to consider extending the present work to $`O(p^4)`$ or to include explicit $`\mathrm{\Delta }(1232)`$ fields in the chiral Lagrangian.
## Acknowledgments
We are grateful to Dave Hutcheon for providing us with unpublished TRIUMF data, to Elie Korkmaz for the data from the SAL experiment, and to Dave Hutcheon, Lothar Tiator, and Olaf Hanstein for helpful conversations. This work was supported in part by the Natural Sciences and Engineering Research Council of Canada.
## A Structure amplitudes
Up to $`O(p^3)`$ in HBChPT, the structure amplitudes of Eq. (7), corresponding to the photoproduction process $`\gamma +N\pi +N`$, are found to be
$`F_1^{(\mathit{0})}(E_\pi ,x)`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi \sqrt{s}}}{\displaystyle \frac{eG_A}{2F}}\{{\displaystyle \frac{1}{2m_N}}[E_\pi +x|\stackrel{}{q}|(\mu _p+\mu _n)]+{\displaystyle \frac{2xE_\pi |\stackrel{}{q}|b_{10}}{G_A(4\pi F)^2}}`$ (A2)
$`+{\displaystyle \frac{1}{4m_N^2}}[|\stackrel{}{q}|^2{\displaystyle \frac{1}{2}}xE_\pi |\stackrel{}{q}|+(2E_\pi ^2m_\pi ^2+xE_\pi |\stackrel{}{q}|2x^2|\stackrel{}{q}|^2)(\mu _p+\mu _n)]\},`$
$`F_1^{()}(E_\pi ,x)`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi \sqrt{s}}}{\displaystyle \frac{eG_A}{2F}}\{1{\displaystyle \frac{x|\stackrel{}{q}|}{2m_N}}(\mu _p\mu _n)`$ (A9)
$`+{\displaystyle \frac{1}{4m_N^2}}\left[+E_\pi ^2{\displaystyle \frac{m_\pi ^2}{2}}(2E_\pi ^2m_\pi ^2+xE_\pi |\stackrel{}{q}|2x^2|\stackrel{}{q}|^2)(\mu _p\mu _n)\right]`$
$`{\displaystyle \frac{2m_\pi ^2b_{19}}{G_A(4\pi F)^2}}+{\displaystyle \frac{2E_\pi ^2}{G_A(4\pi F)^2}}\left(b_{21}^r(\mu ){\displaystyle \frac{G_A}{2}}(1+G_A^2)\mathrm{ln}{\displaystyle \frac{m_\pi ^2}{\mu ^2}}\right)`$
$`+{\displaystyle \frac{E_\pi (E_\pi x|\stackrel{}{q}|)}{G_A(4\pi F)^2}}(2b_{22}^r(\mu )+b_{23}+G_A^3\mathrm{ln}{\displaystyle \frac{m_\pi ^2}{\mu ^2}})+{\displaystyle \frac{1}{4(4\pi F)^2}}[\pi ^2m_\pi ^2\text{}`$
$`8E_\pi |\stackrel{}{q}|\mathrm{ln}\left({\displaystyle \frac{E_\pi +|\stackrel{}{q}|}{m_\pi }}\right)+4i\pi m_\pi ^2\mathrm{ln}\left({\displaystyle \frac{E_\pi +|\stackrel{}{q}|}{m_\pi }}\right)4m_\pi ^2\left(\mathrm{ln}\left({\displaystyle \frac{E_\pi +|\stackrel{}{q}|}{m_\pi }}\right)\right)^2`$
$`+4i\pi E_\pi |\stackrel{}{q}|\text{}]+{\displaystyle \frac{xG_A^2E_\pi |\stackrel{}{q}|}{(4\pi F)^2}}[2{\displaystyle \frac{2|\stackrel{}{q}|}{E_\pi }}\mathrm{ln}\left({\displaystyle \frac{E_\pi +|\stackrel{}{q}|}{m_\pi }}\right)+{\displaystyle \frac{\pi ^2m_\pi ^2}{4E_\pi ^2}}{\displaystyle \frac{2\pi m_\pi }{E_\pi }}`$
$`+{\displaystyle \frac{m_\pi ^2}{E_\pi ^2}}\left(\mathrm{ln}\left({\displaystyle \frac{E_\pi +|\stackrel{}{q}|}{m_\pi }}\right)\right)^2]\},`$
$`F_2^{(\mathit{0})}(E_\pi ,x)`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi \sqrt{s}}}{\displaystyle \frac{eG_A}{2F}}\left\{{\displaystyle \frac{E_\pi |\stackrel{}{q}|}{8m_N^2}}{\displaystyle \frac{x|\stackrel{}{q}|^2}{4m_N^2}}(\mu _p+\mu _n)+{\displaystyle \frac{2E_\pi |\stackrel{}{q}|b_{10}}{G_A(4\pi F)^2}}\right\},`$ (A10)
$`F_2^{()}(E_\pi ,x)`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi \sqrt{s}}}{\displaystyle \frac{eG_A}{2F}}\{{\displaystyle \frac{|\stackrel{}{q}|}{2m_N}}(\mu _p\mu _n)+{\displaystyle \frac{|\stackrel{}{q}|}{4m_N^2}}[E_\pi (E_\pi x|\stackrel{}{q}|)(\mu _p\mu _n)]`$ (A12)
$`+{\displaystyle \frac{G_A^2E_\pi |\stackrel{}{q}|}{2(4\pi F)^2}}[{\displaystyle \frac{\pi ^2m_\pi ^2}{E_\pi ^2}}{\displaystyle \frac{4\pi m_\pi }{E_\pi }}2\pi i{\displaystyle \frac{|\stackrel{}{q}|}{E_\pi }}+2\pi i{\displaystyle \frac{m_\pi ^2}{E_\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{E_\pi +|\stackrel{}{q}|}{m_\pi }}\right)]\},`$
$`F_3^{(\mathit{0})}(E_\pi ,x)`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi \sqrt{s}}}{\displaystyle \frac{eG_A}{2F}}\{{\displaystyle \frac{|\stackrel{}{q}|}{2m_N}}(\mu _p+\mu _n)+{\displaystyle \frac{E_\pi |\stackrel{}{q}|}{8m_N^2}}[32(\mu _p+\mu _n)]`$ (A14)
$`+{\displaystyle \frac{x|\stackrel{}{q}|^2}{2m_N^2}}(\mu _p+\mu _n){\displaystyle \frac{2E_\pi |\stackrel{}{q}|b_{10}}{G_A(4\pi F)^2}}\},`$
$`F_3^{()}(E_\pi ,x)`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi \sqrt{s}}}{\displaystyle \frac{eG_A}{2F}}\{+{\displaystyle \frac{|\stackrel{}{q}|}{(E_\pi x|\stackrel{}{q}|)}}+{\displaystyle \frac{|\stackrel{}{q}|}{2m_N}}(\mu _p\mu _n){\displaystyle \frac{m_\pi ^2|\stackrel{}{q}|}{4m_N^2(E_\pi x|\stackrel{}{q}|)}}`$ (A18)
$`+{\displaystyle \frac{|\stackrel{}{q}|}{m_N^2}}\left[{\displaystyle \frac{E_\pi }{4}}+{\displaystyle \frac{m_\pi ^2}{8(E_\pi x|\stackrel{}{q}|)}}+{\displaystyle \frac{1}{4}}(E_\pi 2x|\stackrel{}{q}|)(\mu _p\mu _n)\right]{\displaystyle \frac{2G_A^2E_\pi |\stackrel{}{q}|}{(4\pi F)^2}}`$
$`{\displaystyle \frac{2m_\pi ^2|\stackrel{}{q}|b_{19}}{G_A(4\pi F)^2(E_\pi x|\stackrel{}{q}|)}}+{\displaystyle \frac{E_\pi |\stackrel{}{q}|}{G_A(4\pi F)^2}}\left(2b_{22}^r(\mu )+b_{23}+G_A^3\mathrm{ln}{\displaystyle \frac{m_\pi ^2}{\mu ^2}}\right)`$
$`+{\displaystyle \frac{G_A^2E_\pi |\stackrel{}{q}|}{(4\pi F)^2}}[{\displaystyle \frac{2|\stackrel{}{q}|}{E_\pi }}\mathrm{ln}\left({\displaystyle \frac{E_\pi +|\stackrel{}{q}|}{m_\pi }}\right){\displaystyle \frac{\pi ^2m_\pi ^2}{4E_\pi ^2}}{\displaystyle \frac{m_\pi ^2}{E_\pi ^2}}\left(\mathrm{ln}\left({\displaystyle \frac{E_\pi +|\stackrel{}{q}|}{m_\pi }}\right)\right)^2+{\displaystyle \frac{2\pi m_\pi }{E_\pi }}]\},`$
$`F_4^{(\mathit{0})}(E_\pi ,x)`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi \sqrt{s}}}{\displaystyle \frac{eG_A}{2F}}\left\{{\displaystyle \frac{|\stackrel{}{q}|^2}{2m_NE_\pi }}+{\displaystyle \frac{1}{m_N^2}}\left[{\displaystyle \frac{E_\pi ^2}{4}}{\displaystyle \frac{m_\pi ^4}{4E_\pi ^2}}{\displaystyle \frac{x|\stackrel{}{q}|^3}{2E_\pi }}{\displaystyle \frac{|\stackrel{}{q}|^2}{4}}(\mu _p+\mu _n)\right]\right\},`$ (A19)
$`F_4^{()}(E_\pi ,x)`$ $`=`$ $`{\displaystyle \frac{m_N}{4\pi \sqrt{s}}}{\displaystyle \frac{eG_A}{2F}}\{{\displaystyle \frac{|\stackrel{}{q}|^2}{E_\pi (E_\pi x|\stackrel{}{q}|)}}{\displaystyle \frac{|\stackrel{}{q}|^2}{2m_NE_\pi }}[1+{\displaystyle \frac{m_\pi ^2}{E_\pi (E_\pi x|\stackrel{}{q}|)}}]`$ (A22)
$`{\displaystyle \frac{|\stackrel{}{q}|^2}{4m_N^2}}\left[1+{\displaystyle \frac{m_\pi ^2}{E_\pi ^2}}{\displaystyle \frac{2x|\stackrel{}{q}|}{E_\pi }}(\mu _p\mu _n){\displaystyle \frac{3m_\pi ^2}{2E_\pi (E_\pi x|\stackrel{}{q}|)}}+{\displaystyle \frac{m_\pi ^4}{E_\pi ^3(E_\pi x|\stackrel{}{q}|)}}\right]`$
$`+{\displaystyle \frac{2m_\pi ^2|\stackrel{}{q}|^2b_{19}}{G_A(4\pi F)^2E_\pi (E_\pi x|\stackrel{}{q}|)}}\},`$
where $`|\stackrel{}{q}|=\sqrt{E_\pi ^2m_\pi ^2}`$, $`m_N`$ is the renormalized nucleon mass, and $`m_\pi `$ is the renormalized pion mass.
Note that all of the parameters in these expressions have been renormalized. The calculation was performed using the bare Lagrangian parameters, which were then converted to renormalized parameters as follows:
$`2a_7`$ $`=`$ $`\mu _p+\mu _n,`$ (A23)
$`4a_6`$ $`=`$ $`\mu _p\mu _n+{\displaystyle \frac{4\pi G_A^2m_\pi m_N}{(4\pi F)^2}},`$ (A24)
$`F_0`$ $`=`$ $`F\left\{1{\displaystyle \frac{m_\pi ^2}{F^2}}\left[l_4^r(\mu ){\displaystyle \frac{1}{(4\pi )^2}}\mathrm{ln}\left({\displaystyle \frac{m_\pi ^2}{\mu ^2}}\right)\right]\right\},`$ (A25)
$`g_A`$ $`=`$ $`G_A{\displaystyle \frac{4a_3G_Am_\pi ^2}{m_N^2}}+{\displaystyle \frac{G_A^3m_\pi ^2}{(4\pi F)^2}}{\displaystyle \frac{4m_\pi ^2}{(4\pi F)^2}}\left[b_{17}^r(\mu ){\displaystyle \frac{G_A}{4}}(1+2G_A^2)\mathrm{ln}\left({\displaystyle \frac{m_\pi ^2}{\mu ^2}}\right)\right].`$ (A26)
$`\mu _p2.79`$ and $`\mu _n1.91`$ are the magnetic moments of the proton and neutron, respectively. The expression for the bare pion decay constant $`F_0`$ in terms of the renormalized $`F`$ and for the bare $`g_A`$ in terms of the physical $`G_A1.26`$ depend somewhat on the explicit form of the Lagrangian used, and are derived, for example, in Ref. .
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# Soliton dynamics in 1D quantum antiferromagnets
## I Introduction
In the past few decades it has become well-established that the physical properties of some magnetic materials, TMMC, CsNiF<sub>3</sub> (caesium nickel fluride) and CuCl<sub>2</sub> 2NC<sub>5</sub>H<sub>5</sub> (dicloro-bis-piridine copper II), for instance, have essentially one dimensional character above their transition temperature. In those kind of materials the distance between magnetic ions along a given direction (magnetic chain direction) is shorter than in the other directions. In such an arrangement the intrachain coupling constant is typically more than two orders of magnitude stronger than the interchain coupling constant. Therefore, the system can be considered as a set of weakly interacting magnetic chains . Due to the relative simplicity of obtaining solitonic or solitary-wave solutions in 1D systems, these quasi-one-dimensional magnets turn out to be the paradigm for the study of the influence of the non-liner modes (solitons) in the dynamical properties of such systems at finite temperatures. Although all real magnetic materials investigated are not perfectly one dimensional, the assumption of the 1D behavior is shown to be in good agreement with the experimental results (see Ref.1 and the references therein).
In magnetic materials solitons or solitary-waves can be regarded as ‘kinks’ or ‘twists’ in the spin space moving with constant speed and carrying a constant topological charge defined by the values of the spin variables at infinity. For low enough temperatures, when the linear modes (spin-waves) are not excited, the magnetic system can be represented in first approximation by a gas of non-interacting solitons. Using this idea, Mikeska calculated the soliton contribution to the dynamical structure factor of the classical one-dimensional magnets. From both works we learn that the assumption of ballistic motion for solitons is the origin of the ‘central peak’ behavior observed in neutron scattering experiments. A different situation could be found from the quantum field theory point of view when the temperature is raised. In this case the spin-wave (SW) modes are excited, therefore not all of the degrees of freedom of the system contribute to the soliton formation and a residual interaction (which couples the center of mass of the soliton to the spin-wave modes) shows up. In practice, the specific form of this kind of interaction is obtained via the collective coordinate method in the quantization process of the classical hamiltonian.
The soliton-SW coupling may result in a dissipative regime to the soliton motion depending on the form of the potential generated by the presence of the non-linear excitation. As it is known, the equation of motion for the spin variable in the TMMC below and above $`T_N`$ are a 2SG and a SG equation respectively. This fact makes the TMMC a suitable probe to investigate the appearance of a dissipative regime in the soliton motion provided that, below and above a certain Néel transition temperature $`T_N`$, the classical equation of motion for the spin variables are substantially different. The main purpose of this work will be therefore the analysis of the magnetic soliton motion above and below $`T_N`$ and the possible influence of the dissipative regime, found for $`T<T_N`$, on the solitonic contribution to the dynamical structure factor. In doing that, we will use the method developed in Ref. 9 for the analysis of the dissipative dynamics of solitons. As pointed out before, this formalism is based on the collective coordinate method and allows us to transform the original hamiltonian of the spin degree of freedom into one of a particle (the soliton) coupled to an infinite set of linear-modes (SW).
Starting from the interacting soliton-SW hamiltonian it is possible to obtain a Brownian-like equation of motion for the soliton center of mass via the Feynman-Vernon formalism. This effective equation of motion is written in terms of a damping constant that depends on the phase shifts of the scattering problem that emerges by the presence of the non-linear excitation coupled to the SW. Therefore, the analysis of the scattering properties of the 2SG potential (for $`T<T_N`$) and the SG potential (for $`T>T_N`$) allows us to calculate the mobility as a function of the temperature and the external magnetic field. As it will be shown, above $`T_N`$ the SG solitons have infinite mobility in agreement with the ballistic motion used to understand the neutron scattering experiment for $`H/T10kOe/K`$. On the other hand, for low temperatures or low magnetic field, when the spin equation of motion for TMMC have a 2SG form, the soliton mobility is finite, changing the form of the dynamical structure factor considerably.
To begin with, in Sec. II we review the models currently applied to the spin dynamics of the TMMC compound, above and below its transition temperature, and also the corresponding classical equations of motion. In Sec. III we summarize the obtainment of the quantized soliton-SW hamiltonian and the damping parameter of the soliton Brownian motion is calculated as a function of the temperature and the external magnetic field. Sec. IV is devoted to the study of the influence of the soliton damped motion on the dynamical structure factor and, finally, our conclusions are presented in Sec. V.
## II The model for TMMC
The antiferromagnet TMMC has extensively been studied from the theoretical and experimental points of view. The hamiltonian describing the interacting 3D array of classical spins in this material can be written as
$$=\underset{j}{}H_j\frac{1}{2}J_{}\underset{ii^{}}{}\underset{j}{}S_{i,j}S_{i^{},j},$$
(1)
where
$$H_j=\underset{k}{}\left\{J_{||}S_{j,k}S_{j,k+1}+A(S_{j,k}^z)^2g\mu _BBS_{j,k}^x\right\}.$$
(2)
The hamiltonian $`H_j`$ describes the nearest neighbour intrachain interaction between spins with an easy plane anisotropy ($`A>0`$) placed in an external magnetic field ($`B`$) in the x direction. The spins will be treated as classical vectors of lengh $`S`$ and the constants $`J_{}`$ and $`J_{}`$, both positive, correspond to the antiferromagnetic and ferromagnetic exchange coupling constants, respectively. The second term in the r.h.s. of (1) represents an interchain interaction between the spins, completing the description of the 3D spin arrangement. Finally, the following values of material parameters will be used: $`J_{}=13.4K`$, $`S=5/2`$, $`A/J_{}=0.010.02`$, $`J_{}/J_{}=1.510^5`$ and $`g=2.01`$.
In order to start the classical description of the spin dynamics it is convenient to look at two main different situations, namely, temperatures below and above the transition temperature. For temperatures below $`T_N`$ the system described by (1) displays a long range magnetic order, therefore the staggered spontaneous magnetization is not zero and the system can be described in the mean field approximation as a set of non-interacting antiferromagnetic chains with an additional spontaneous magnetization in the y direction. Explicitly,
$`H`$ $`=`$ $`{\displaystyle \underset{i}{}}\{J_{||}S_iS_{i+1}+A(S_i^z)^2g\mu _BBS_i^x`$ (4)
$`g\mu _BB_{}^{MF}(1)^iS_i^y\},`$
where
$$B_{}^{MF}=\eta J_{}(1)^iS_i^y/g\mu _B,$$
(5)
and $`\eta `$ accounts for the presence of neighbouring chains in the model. In the specific case of TMMC, $`\eta =6`$. The intrachain mean field $`B_{}^{MF}`$ is usually replaced by its saturation value $`B_{}^S22.3`$Oe which results from the substitution of $`S_i^{(y)}`$ in (5) by its maximum value.
At this point, we can carry on the classical description of the spin dynamics. In order to do that it is convenient to change the spin variables to the following form
$`S_{e,o}`$ $`=`$ $`\pm S[\mathrm{sin}(\mathrm{\Theta }\pm \theta )\mathrm{cos}(\mathrm{\Phi }\pm \phi ),`$ (7)
$`\mathrm{sin}(\mathrm{\Theta }\pm \theta )\mathrm{sin}(\mathrm{\Phi }\pm \phi ),\mathrm{cos}(\mathrm{\Theta }\pm \theta )],`$
where $`e`$ and $`o`$ stands for even and odd sites within a chain.
Using the representation (7) a $`\mathrm{\Phi }`$-dependent part of the hamiltonian (4) can be obtained (see Ref. 3 for details). Explicitly,
$`H^\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{1}{2}}J_{}S^2{\displaystyle }dz[{\displaystyle \frac{1}{c^2}}(_t\mathrm{\Phi })^2+(_z\mathrm{\Phi })^2`$ (9)
$`{\displaystyle \frac{1}{4}}b^2\mathrm{sin}^2\mathrm{\Phi }2b_{}\mathrm{sin}\mathrm{\Phi }],`$
where
$$c^2=4+\frac{2A}{J_{}},b=\frac{g\mu _BB}{J_{}S},b_{}=\frac{g\mu _BB_{}^{MF}}{J_{}S}.$$
(10)
It should be stressed that the hamiltonian (9) is an approximated description of the real TMMC system. To reproduce the experimental results, magnon-mass and solitonic energy, for instance, quantum effects and the out of plane component of the magnetization must be taken into account. To go on with the classical description of the $`\mathrm{\Phi }`$-dependent part of the original hamiltonian (1) the equation of motion associated to (9),
$$\frac{1}{c^2}_{tt}\mathrm{\Phi }=_{zz}\mathrm{\Phi }\frac{b^2}{8}\mathrm{sin}2\mathrm{\Phi }b_{}\mathrm{sin}\mathrm{\Phi }$$
(11)
has to be solved. Equation (11) is not completely integrable, however, it has solitonic solutions in the form of 2$`\pi `$-kinks(antikinks) moving with velocity $`u`$. Explicitly,
$`\mathrm{cos}\mathrm{\Phi }`$ $`=`$ $`\pm 2{\displaystyle \frac{\sqrt{\alpha }}{1+\alpha \mathrm{sinh}^2y}}\mathrm{sinh}y,`$ (12)
$`\mathrm{sin}\mathrm{\Phi }`$ $`=`$ $`1{\displaystyle \frac{2}{1+\alpha \mathrm{sinh}^2y}},`$ (13)
where
$$\alpha =\frac{b_{}}{b_{}+b^2/4},y=(zz(t))\sqrt{\frac{b_{}+b^2/4}{1u^2/c^2}},$$
(14)
and the position of the soliton center of mass $`z(t)`$ is given by
$$z(t)=z_0+ut.$$
(15)
On the other hand, for temperatures above $`T_N`$ the value of the $`b_{}`$ is very small. In fact, in this situation $`b_{}`$ can be set equal to zero and, a well-known solitonic solution for equation (11) can be found: the $`\pi `$-kink(antikink) solution for the SG equation
$$\mathrm{sin}\mathrm{\Phi }_s(z,t)=\pm \mathrm{tanh}\left[(1u^2/c^2)^{1/2}(zz(t))b/2\right].$$
(16)
As it can be seen, the model for TMMC in the continuum approximation leads us to different kinds of solitonic solutions depending on the temperature. A 2SG soliton solution given by (12) and (13) for $`T<T_N`$ and, a SG solution (16) for temperatures above $`T_N`$. As it was already mentioned, from the classical point of view, these soliton solutions will move with constant velocity througout the sample. However, looking at the soliton dynamics from the quantum field theory perspective, the interaction with the spin waves can transform this ballistic regime into a dissipative one. In the next section we shall be aiming at the investigation of the mobility of the two types of solitons, below and above $`T_N`$.
## III Soliton Mobility
The quantum dynamics of our spin system (9) can be analyzed by studying the quantum mechanics of the field theory described by the action
$$S[\mathrm{\Phi }]=J_{}S^2\left\{\frac{1}{2c^2}(_t\mathrm{\Phi })^2\frac{1}{2}(_z\mathrm{\Phi })^2+U(\mathrm{\Phi })\right\}𝑑t𝑑z,$$
(17)
where
$$U(\mathrm{\Phi })=\frac{b^2}{8}\mathrm{sin}2\mathrm{\Phi }+b_{}\mathrm{sin}\mathrm{\Phi }.$$
(18)
To quantize the system described by (17) we need to evaluate
$$G(t)=\text{tr}𝒟\mathrm{\Phi }\mathrm{exp}\frac{i}{\mathrm{}}S[\mathrm{\Phi }]$$
(19)
where the functional integral has the same initial and final configurations and tr means to evaluate it over all such configurations. As the functional integral in (19) is impossible to be evaluated for a potential energy density as in (18) we must choose an approximation to do it. Since the magnetic moments at the manganese sites in the TMMC are large (5/2), the semi-classical limit will be chosen as the appropriate one in our case. Within the functional integral formalism of quantum mechanics, the semi-classical limit is simply the stationary phase method applied to (19) around the solitonic solutions (12), (13) or (16) in which we are interested. When this is done we are left with an eigenvalue problem that reads
$$\left\{\frac{d^2}{dz^2}+U^{\prime \prime }(\mathrm{\Phi }_s)\right\}\psi _n(zz_0)=k_n^2\psi _n(zz_0),$$
(20)
where $`\mathrm{\Phi }_s`$ is denoting the soliton-like solution around which we are expanding $`\mathrm{\Phi }(z,t)`$ and $`\psi _n(zz_0)`$ are the spin wave modes in the presence of the soliton.
Now one can easily show that $`d\mathrm{\Phi }_s/dz`$ is a solution of (20) with $`k_n=0`$. The existence of this mode is related to the translation invariance of the system and causes the divergence of the functional integral in (19) in the semi-classical limit (Gaussian approximation). The way out of this problem is the so-called collective coordinate method . This method consists basically in expanding the field configurations about $`\mathrm{\Phi }_s(z)`$ as
$$\mathrm{\Phi }(z,t)=\mathrm{\Phi }_s(zz_0(t))+\underset{n=1}{\overset{\mathrm{}}{}}c_n\psi _n\left(zz_0(t)\right),$$
(21)
but regarding the $`c`$-number $`z_0`$ as a position operator. Using expansion (21), the second quantized version of (9) can be written as
$`H={\displaystyle \frac{1}{2M_s}}(P{\displaystyle \underset{mn}{}}\mathrm{}g_{mn}b{}_{n}{}^{+}b_{m}^{})^2+{\displaystyle \mathrm{}\mathrm{\Omega }_nb{}_{n}{}^{+}b_{n}^{}}.`$ (22)
where $`\mathrm{\Omega }_nck_n`$.
In the hamiltonian (22), $`P`$ stands for the momentum canonically conjugated to $`z_0`$,
$$M_s=\frac{2J_{}S^2a}{c^2}_{\mathrm{}}^+\mathrm{}𝑑zU(\mathrm{\Phi }_s(z))$$
(23)
is the soliton mass and the coupling constant $`g_{mn}`$ is given by
$$g_{mn}=\frac{1}{2i}\left[\sqrt{\frac{\mathrm{\Omega }_m}{\mathrm{\Omega }_n}}+\sqrt{\frac{\mathrm{\Omega }_n}{\mathrm{\Omega }_m}}\right]𝑑z\psi _m(z)\frac{d\psi _n(z)}{dz}.$$
(24)
The operators $`b^+`$ and $`b`$ are respectively the creation and annihilation operators of the excitations of the magnetic system (magnons) in the presence of the soliton. In fact, the term
$$\underset{mn}{}\mathrm{}g_{mn}b{}_{n}{}^{+}b_{m}^{},$$
(25)
can be interpreted as the total linear momentum of the magnons of the system and therefore, we are left with a problem in which the momentum associated to the soliton is now coupled to the magnons’ momenta. This effective model suggests that, as the population of magnons is a temperature-dependent quantity, the mobility of the soliton will be strongly related to the temperature of the system and its dynamics (determined by (22)) will be non trivial.
At this point we are ready to study the mobility of the wall because we have been able to map that problem into the hamiltonian (22), which on its turn has been recently used to study the mobility of polarons, heavy particles in 1D environments and skyrmions in 2D electronic systems. The main result obtained in those calculations can be summarized as follows. The damping function $`\gamma (t)`$ (basically the inverse of the mobility) is given by
$`\gamma (t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{2M}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}d\omega d\omega ^{}\{S(\omega ,\omega ^{})(\omega \omega ^{})\times `$ (27)
$`[n(\omega )n(\omega ^{})]\mathrm{cos}(\omega \omega ^{})t\},`$
where
$$n(\omega )=\frac{1}{e^{\beta \mathrm{}\omega }1}$$
(28)
is the Bose function and,
$$S(\omega ,\omega ^{})=\underset{mn}{}|g_{mn}|^2\delta (\omega \mathrm{\Omega }_n)\delta (\omega ^{}\mathrm{\Omega }_m)$$
(29)
is the so-called scattering function.
In the long time limit $`\gamma (t)`$ can, to a good approximation, be written as
$$\gamma (t)\overline{\gamma }(T)\delta (t),$$
(30)
where $`\delta (t)`$ is the Dirac delta function and $`\overline{\gamma }(T)`$ is given by
$$\overline{\gamma }(T)=\frac{1}{2\pi M_s}_0^{\mathrm{}}𝑑E(E)\frac{\beta Ee^{\beta E}}{(e^{\beta E}1)^2}.$$
(31)
In (31), $`(E)`$ is the reflection coefficient of the “potential” $`U^{\prime \prime }(\mathrm{\Phi }_s)`$ in the Schrödinger-like equation (20). For simplicity we will express the reflection coefficient $`(E)`$ in terms of the even and odd scattering phase shifts as
$$(k)=\mathrm{sin}^2\left(\delta ^e(k)\delta ^o(k)\right).$$
(32)
At this point we can perform the calculation of the soliton mobility in TMMC for temperatures above and below $`T_N`$.
### A Soliton mobility for $`T<T_N`$
To calculate the soliton mobility below $`T_N`$ we need the explicit form of the potential $`U^{\prime \prime }(\mathrm{\Phi }_s)`$ involved in (31). For the case of the 2SG soliton this potential can be written as
$`U^{\prime \prime }(\mathrm{\Phi }_s^{2SG})`$ $`=`$ $`{\displaystyle \frac{1}{\lambda ^2}}[12\text{sech}^2({\displaystyle \frac{z}{\lambda }}+\rho )2\text{sech}^2({\displaystyle \frac{z}{\lambda }}\rho )`$ (34)
$`+2\text{sech}({\displaystyle \frac{z}{\lambda }}+\rho )\text{sech}({\displaystyle \frac{z}{\lambda }}\rho )],`$
where
$$\lambda =\frac{1}{b_{}+b^2/4},\mathrm{cosh}\rho =\frac{1}{\sqrt{\alpha }}.$$
(35)
The second and third terms in the r.h.s. of (34) are the potentials of the noninteracting $`\pi `$-solitons located at $`z/\lambda =\pm \rho `$ whereas the last term describes the interaction of the two $`\pi `$-solitons at $`z/\lambda =\pm \rho `$ respectively.
For all finite values of $`\lambda `$ and $`\rho `$, the system is translationally invariant and, consequently, the potential (34) has a zero energy state that is given by
$$\psi _0\text{sech}(\frac{z}{\lambda }+\rho )+\text{sech}(\frac{z}{\lambda }\rho ),$$
(36)
which is nothing but the Goldstone mode of the $`2\pi `$-soliton for finite transverse magnetization and finite external field.
In order to evaluate the expression (31) for the damping constant we need the even and odd phase shifts associated to the potential (34). Unfortunately, their analytical evaluation is very complicated for all finite values of $`\lambda `$ and $`\rho `$, and in what follows we will only study the situation of weak external fields ($`b_{}b^2/2`$). In this case ($`\rho 1`$) the Shrödinger-like equation (20) can be written as
$$\left\{\frac{d^2}{dz^2}+V(z)\right\}\psi _n(z)=\kappa _n^2\psi _n(z),$$
(37)
where
$$\kappa _n^2=k_n^2\frac{1}{\lambda ^2}\frac{\rho ^2}{\lambda ^2},$$
(38)
and the potential (34) is now reduced to the sum of two contributions, one coming from the spontaneous staggered magnetization and, the other from the presence of the weak external field. Explicitly,
$$V(z)=V_0(z)+(\frac{\rho }{\lambda })^2V_1(z),$$
(39)
with
$$V_0(z)=2\text{sech}^2\left(\frac{z}{\lambda }\right)$$
(40)
and
$$V_1(z)=8\mathrm{tanh}^2\left(\frac{z}{\lambda }\right)\text{sech}^2\left(\frac{z}{\lambda }\right).$$
(41)
The calculation of the even and odd phase shifts for a potential of the form (39)-(41) is reported in Ref. 10 and here we will only show the fundamental results of the numerical solution of the Schrödinger-like equation (37). Fig. 1 and Fig. 2 show the even and odd parity phase shifts for different values of the external field.
The values of $`\delta _e`$ and $`\delta _o`$ for $`k=0`$ are in agreement with the 1D version of the Levinson’s theorem which establishes that
$`\delta ^e(k=0)=\pi (n^e{\displaystyle \frac{1}{2}}),`$ (42)
$`\delta ^o(k=0)=\pi n^o,`$ (43)
where $`n^e`$ and $`n^o`$ are the number of even and odd parity bound states.
As it can be seen in Fig. 1 the even phase shift is $`\pi /2`$ at the origin. This behavior is in complete agreement with the existence of an even bound state corresponding to the Goldstone mode. On the other hand, the odd phase shift $`\delta _o(0)=\pi `$, indicates the presence of an odd bound state. This result was previously obtained by Kivshar et al. in the study of the small-amplitude modes around the localized solution of the 2SG equation and shows that there is always an odd bound state in this kind of system. Therefore, the spectrum of (39) is composed by: i) the $`\psi _0`$ solution (36) corresponding to the translation mode of the soliton (Goldstone mode) ii) an internal mode which appears when the system is perturbed by the external magnetic field and iii) the $`\psi _k`$ solutions which constitute the continuum modes and correspond to magnons.
In order to find the damping coefficient we must compute the reflection coefficient $`(k)`$. This can be done by inserting the numerical results of the even and odd phase shifts into the general expression (32). In Fig. 3, we have plotted $`(k)`$ for different values of the perturbation parameter $`\rho `$ for the whole range of $`k`$. As it can be seen the major contribution for the reflection coefficient comes from the low energy states, in agreement with the well behaved potentials (40) and (41).
Having done that, one can immediately integrate the function $`(k)`$ in expression (31) which finally allows us to describe the damping coefficient as a function of the temperature (see Fig.4). It is important to notice that we have not considered the odd bound state of the potential (39) in computing the damping coefficient because in evaluating the scattering matrix (29), only elastic terms are taken into account (see for instance Ref. 9, Ref. 12 or Ref. 13).
As it can be seen, the damping coefficient is linear for high temperatures. This result can be obtained directly from (31). In fact, for $`T`$ high enough the damping constant can be approximated by
$$\overline{\gamma }(T)\frac{1}{2\pi M_s\beta }_0^{\mathrm{}}𝑑E\frac{(E)}{E}T$$
(44)
which is linear in $`T`$, independently of the explicit form of $`(E)`$. In the low temperature regime we can write
$$\overline{\gamma }(T)\frac{1}{2\pi M_s}_0^{\mathrm{}}𝑑E(E)\beta Ee^{\beta E},$$
(45)
where $`E`$ always presents a gap determined by the presence of the magnetic field and/or the spontaneous staggered magnetization. Here we shall not attempt to write an approximate expression for (45) because the correct behavior of the reflection coefficient was only numerically determined. As it is shown in Fig. 4, for low enough temperatures, the damping coefficient drops exponentially to zero due to the existence of the gap. As the temperature increases the damping coefficient rises following a power law behavior until it becomes linear for high enough temperatures. This strong temperature dependence of the damping parameter, for $`T`$ below the transition temperature, will influence directly the correlation function between the magnetic solitons.
### B Soliton mobility for $`T>T_N`$
To perform the calculation of the $`\pi `$-soliton mobility for $`T>T_N`$ we simply set to zero the $`b_{}`$ in the hamiltonian (9) and therefore, the equation of motion for the $`\mathrm{\Phi }`$-dependent part of the spin degree of freedom (11) becomes a SG equation with solitonic solution (16). In this case the potential involved in the Schrödinger-like equation (20) which determines the fluctuations around the soliton solution have the form
$$U^{\prime \prime }(z)=\xi ^2(12\text{sech}^2\xi z),$$
(46)
where $`\xi =b/2`$. The spectrum of (46) contains a bound state with zero energy
$$\psi _0=\sqrt{\frac{\eta }{2}}\text{sech}(\xi z),k_0^2=0,$$
(47)
which constitutes the translation mode of the soliton (Goldstone mode), and a continuum of quasiparticles modes (magnons) given by
$$\psi _n(x)=\frac{1}{\sqrt{L}}\left[\frac{k_n+i\xi \mathrm{tanh}(\xi z)}{k_n+i\xi }\right]e^{ik_nz},$$
(48)
where
$$k_n=\frac{2n\pi }{L}\frac{\delta (k_n)}{L},\delta (k)=\mathrm{arctan}\left[\frac{2\xi k}{k^2\xi ^2}\right].$$
(49)
As it was already mentioned, the reflection coefficient $``$ for a general symmetric potential can be expressed in terms of the corresponding even and odd phase shifts by the relation (32). Re-expressing (48) in terms of parity eigenstates it is easy to prove that the potential (46) belongs to the class of reflectionless potentials because its phase shifts are given by
$$\delta ^{e,o}(k)=\mathrm{arctan}(\xi /k),$$
(50)
that do not distinguish between odd and even parities. Therefore no matter how high the temperature rises above $`T_N`$ the damping coefficient is always zero and as a direct consequence the ballistic regime for the soliton results.
As we have seen, the solitonic solutions in TMMC have different regimes for $`T`$ below and above $`T_N`$. Below the transition temperature the $`2\pi `$-solitons behave like a Brownian particle with a finite damping parameter. On the other hand, for $`T`$ above $`T_N`$ the $`\pi `$-solitons have infinite mobility corresponding to the ballistic regime. The next section is devoted to studying the influence of the changes of the solitonic solutions mobility in the dynamical properties of TMMC.
## IV Dynamical Structure Factor
In this section we will investigate the dependence of the dynamical properties of TMMC with respect to the temperature and the magnetic field. For $`T`$ below the transition temperature, this will be done through the computation of the dynamical structure factor of a dilute gas of 2$`\pi `$-solitons in a dissipative regime. With this result, we can analyze the main differences with the assumption of ballistic regime using by Holyst in the same situation.
In a general form, the longitudinal and transverse dynamical structure factors with respect to the external field $`B`$ can be defined as
$$𝒮^{||()}=\frac{1}{(2\pi )^2}𝑑t𝑑ze^{i(qz\omega t)}S^{x(y)}(0,0)S^{x(y)}(z,t),$$
(51)
where $`S^{x(y)}`$ corresponds to the spin component in the x(y) direction.
To begin with, let us recall the main results for the longitudinal dynamical structure factor reported in Ref. 3. Using the model of non-interacting 2$`\pi `$-soliton gas in the ballistic regime $`𝒮^{||}(q,\omega )`$ can be written approximately as
$$𝒮^{||}(q,\omega )=n_{2\pi }S^2|F_{2\pi }^x(q)|^2\frac{p(\omega /q)}{2\pi q},$$
(52)
where
$$p(\omega /q)=\sqrt{\frac{\beta E_{2\pi }}{2\pi c^2}}\mathrm{exp}\frac{\beta E_{2\pi }\omega ^2}{2c^2q^2},$$
(53)
$`E_{2\pi }`$ $`=`$ $`2Bg\mu _BS[\sqrt{1+4b_{}/b^2}`$ (55)
$`+4b_{}b^2\mathrm{sinh}^1(bb_{}^{1/2}/2)]`$
and
$$F_{2\pi }^x(q)=\frac{i\pi d_\pi }{2}\frac{\mathrm{sin}(qd_\pi \sigma )}{\mathrm{cosh}(q\pi d_\pi \sqrt{1\alpha }/8)},$$
(56)
$$\sigma =\frac{\sqrt{1\alpha }}{8}\mathrm{ln}\left(\frac{2}{\alpha }1+\frac{2}{\alpha }\sqrt{1\alpha }\right),d_\pi =\frac{8}{b}.$$
(57)
The correlations described by (52) are induced by single kinks moving from the origin to the position $`z`$ in a time interval 0 to $`t`$. As expected, the Maxellian velocity distribution used to describe the $`2\pi `$-kink gas is directly reflected in the Gaussian dependence of the longitudinal structure factor with the frequency.
On the other hand, to get a better idea of the changes in the dynamical properties when we cross the transition temperature, it is convenient to calculate the dynamical structure factor for $`T`$ above the transition temperature. As it was demostrated before, above $`T_N`$, the $`\pi `$-solitons moves without dissipation and, therefore, the ballistic regime is valid. Using again the model of a dilute gas of solitons, the dynamical structure factor can be written as
$`𝒮^{||}(q,\omega )`$ $`=`$ $`{\displaystyle \frac{S^2}{(2\pi )^{3/2}}}\sqrt{{\displaystyle \frac{E_\pi \beta }{c^2q^2}}}|F^{}(q)|^2\mathrm{exp}({\displaystyle \frac{E_\pi \beta \omega ^2}{2c^2q^2}}),`$ (58)
where
$$F^{}(q)=\frac{2\pi }{b}\text{sech}(q\pi /b)\text{and}E_\pi =Bg\mu _BS.$$
(59)
As it can be seen, the dependence with the frequency remains almost unchanged no matter what the temperature is. At the same time, it should be noticed that once the intensity of the central peak and the density of kinks in (53) and (58) are proportional, the intensity of the central peak for $`T<T_N`$ will be lower. This is a consequence of the smaller number of solitons for temperatures below the transition temperature.
With the previous results for $`T`$ above and below $`T_N`$ in mind, we can go further on and study the influence of the dissipative regime in the dynamical properties of the $`2\pi `$-kink gas. As it was shown before, below the transition temperature the $`2\pi `$-solitons move in a dissipative regime. Therefore, the position of the center of mass for those kind of excitations as a function of time can be written as
$$z(t)=z_0+\frac{v_0}{\gamma (T)}(1\mathrm{exp}\gamma (T)t),$$
(60)
where $`z_0`$ is the initial position, $`v_0`$ is the initial velocity and $`\gamma (T)`$ is the temperature-dependent damping parameter. Now, to calculate the dynamical structure factor we will use the $`2\pi `$-soliton solutions (12)-(14) with $`z(t)`$ given by (60). Following the same procedure that led to equation (52) and, after some calculations, $`𝒮^{||}(q,\omega )`$ can be written as
$$𝒮^{||}(q,\omega )=\frac{2n_{2\pi }S^2}{\pi }|F_{2\pi }^x(q)|^2\mathrm{\Gamma }(q,\omega ),$$
(61)
where
$$\mathrm{\Gamma }=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n}{n!}\left(\frac{q^2}{2\beta E_{2\pi }\gamma ^2}\right)^n\underset{m=0}{\overset{2n}{}}(1)^mC_m^{2n}\frac{2m\gamma }{m^2\gamma ^2+\omega ^2}$$
(62)
and
$$C_m^{2n}=\frac{(2n)!}{m!(2nm)!}.$$
(63)
As it can be seen, the dissipative regime for the magnetic solitons in the case in which $`T<T_N`$, changes considerably the behavior of the longitudinal dynamical structure factor. Although the expression (61) is valid for all finite values of $`q`$, we couldn’t perform the entire sum to get a closed expression. Therefore, it is helpful to study the behavior of (61)-(63) for small momentum in order to compare it with the ballistic behavior result (52). Assuming that
$$q\frac{\gamma }{c}\sqrt{2\beta E_{2\pi }},$$
(64)
the dynamical structure factor $`𝒮^{||}(q,\omega )`$ can be written as
$$𝒮^{||}(q,\omega )=\frac{2\alpha }{\pi }|F_{2\pi }^x(q)|^2\mathrm{\Lambda }(q,\omega )$$
(65)
where
$`\mathrm{\Lambda }(q,\omega )`$ $`=`$ $`2\pi \delta (\omega )\mathrm{exp}{\displaystyle \frac{q^2c^2}{\beta E_{2\pi }\gamma ^2}}+`$ (67)
$`{\displaystyle \frac{q^2c^2}{\beta E_{2\pi }\gamma ^2}}\left[{\displaystyle \frac{\gamma }{\gamma ^2+\omega ^2}}{\displaystyle \frac{\gamma }{4\gamma ^2+\omega ^2}}\right].`$
Within the approximation of small momentum, the behavior of $`𝒮^{||}(q,\omega )`$ with the frequency, changes from the ‘Gaussian’ central peak to a ‘Lorentzian’ dependence. Therefore, as the temperature is lowered below $`T_N`$, the central peak behavior is replaced by a smoother-one in the frequency domain. This result is a direct consequence of the dissipative regime of the $`2\pi `$-soliton and, as the damping constant $`\gamma `$ can be controlled by changing the temperature and the magnetic field, a possible indication of a non-ballistic regime has been found. Another quantity that can be computed in order to get a better idea of the influence of the dissipative motion of magnetic solitons in TMMC is the $`T_1`$ time of NMR. This problem is currently being investigated by one of us.
## V Conclusions
In this paper we have analyzed the possibility of identifying two different regimes of motion for the magnetic solitons in the TMMC antiferromagnet. We were able to show that above the transition temperature $`T_N`$, the $`\pi `$-soliton moves without dissipation, even from the field theoretical point of view. This result is in complete agreement with the ballistic regime adopted to understand the experimental data reported in Ref. 11.
On the other hand, for $`T`$ below the transition temperature the $`2\pi `$-solitons in the system move with finite mobility. Therefore, a dissipative equation of motion has to be used in the description of the soliton’s center of mass motion. This difference in the regime of motion is directly reflected in the longitudinal structure factor and, therefore, can be used as an indication of the finite mobility of the solitonic solutions below the transition temperature. The results presented here could be directly compared to the experimental data one may obtain when testing the TMMC antiferromagnet in this temperature regime.
Although the formulation used to compute the damping parameter is valid for all values of the external magnetic field, we have restricted ourselves to the study of very weak fields. However, our formulation can be used to study situations with arbitrarily stronger magnetic fields and to other magnetic materials that support solitonic solutions without any major qualitative difference. For instance, we could treat systems modelled by the 1-D Dzyaloshinski-Moriya antiferromagnet which can naturally be described by a 2DSG hamiltonian independently of the temperature.
## VI Acknowledgment
AVF wishes to thank Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP) for financial support, whereas AOC kindly acknowledges partial support from Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq).
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# The 𝜋, 𝐾⁺, and 𝐾⁰ electromagnetic form factors
## I Introduction
The light pseudoscalar mesons play an important role in understanding low-energy QCD. They are the lightest observable hadronic bound states of a quark and an anti-quark, and are the Goldstone bosons associated with chiral symmetry breaking. Their static properties such as the mass and decay constants have been studied extensively . Dynamic properties and scattering observables are much less understood theoretically, but therefore not less important to calculate within QCD. In this respect, the elastic electromagnetic form factors of the pion and kaon are very interesting: the probe is well understood, there are accurate data for $`F_\pi `$ at low $`Q^2`$ to confront theoretical calculations with, and the charge radii $`r_\pi ^2`$, $`r_{K^+}^2`$, and $`r_{K^0}^2`$ are experimentally known. Currently, there are several experiments at JLab to determine both the pion and the kaon form factor in the range $`0.5<Q^2<3\mathrm{GeV}^2`$ to better accuracy , which could help to discriminate between different model calculations.
To calculate these form factors, we use an approach based on the Dyson–Schwinger equations \[DSEs\], which form an excellent tool to study nonperturbative aspects of hadron properties in QCD . The approach is consistent with quark and gluon confinement , generates dynamical chiral symmetry breaking , and is Poincaré invariant. It is straightforward to implement the correct one-loop renormalization group behavior of QCD , and obtain agreement with perturbation theory in the perturbative region. Provided that the relevant Ward identities are preserved in the truncation of the DSEs, the corresponding currents are conserved. Axial current conservation induces the Goldstone nature of the pions and kaons ; electromagnetic current conservation produces the correct hadronic charge without fine-tuning.
We obtain the meson Bethe–Salpeter amplitudes \[BSAs\] and the quark-photon vertex as the solutions of respectively the homogeneous and inhomogeneous Bethe–Salpeter equation \[BSE\] in ladder truncation. The required dressed quark propagators are obtained from solutions of the quark DSE in rainbow truncation. Non-analytic effects from vector mesons are automatically taken into account, because these vector $`q\overline{q}`$ bound states appear as poles in the quark-photon vertex solution . We employ a realistic model for the effective quark-antiquark coupling that has been shown to reproduce the pion and kaon masses and decay constants as well as the masses and decay constants for the vector mesons $`\rho `$, $`\varphi `$ and K to within 10% . The model parameters are all fixed in previous work and constrained only by $`m_\pi `$, $`m_K`$, $`f_\pi `$ and $`\overline{q}q`$. The produced pion charge radius is within 2% of the experimental value . Here, we use the same approach, without parameter adjustment, to calculate the neutral and charged kaon form factors and charge radii; we also extend our previous pion form factor calculations to the spacelike $`Q^2`$-domain anticipated for future JLab data.
In Sec. II we review the formulation that underlies a description of the pion and kaon charge form factors within a modeling of QCD through the DSEs. Within the impulse approximation, we outline the manner in which a ladder-rainbow dynamics for the propagators, BSAs and quark-photon vertex preserves the meson electromagnetic current. We further indicate the additional terms needed for current conservation, if one goes beyond rainbow-ladder truncation for the DSEs. In Sec. III we discuss the details of the model and present our numerical results for the form factors. Concluding remarks are given in Sec. IV.
## II Pseudoscalar electromagnetic form factors
The 3-point function describing the coupling of a photon with momentum $`Q`$ to a pseudoscalar meson with initial and final momenta $`P_\pm =P\pm Q/2`$ respectively can be written as the sum of two terms
$`\mathrm{\Lambda }_\nu ^{a\overline{b}}(P,Q)`$ $`=`$ $`\widehat{Q}^a\mathrm{\Lambda }_\nu ^{a\overline{b}a}(P,Q)+\widehat{Q}^{\overline{b}}\mathrm{\Lambda }_\nu ^{a\overline{b}\overline{b}}(P,Q),`$ (1)
with $`\widehat{Q}`$ the electric charge of the (anti-)quark, $`\frac{2}{3}`$ for the $`u`$-quark, and $`\frac{1}{3}`$ for the $`d`$\- and $`s`$-quarks, and with $`\mathrm{\Lambda }^{a\overline{b}a}`$ and $`\mathrm{\Lambda }^{a\overline{b}\overline{b}}`$ describing the coupling of a photon to the quark and anti-quark inside the meson respectively. The meson form factor is defined as
$$\mathrm{\Lambda }_\nu ^{a\overline{b}}(P,Q)=2P_\nu F(Q^2),$$
(2)
and the corresponding charge radius as $`r^2=6F^{}(Q^2)`$ at $`Q^2=0`$. Analogously, we can define a form factor for each of the two terms on the RHS of Eq. (1)
$$\mathrm{\Lambda }_\nu ^{a\overline{b}\overline{b}}(P,Q)=2P_\nu F_{a\overline{b}\overline{b}}(Q^2).$$
(3)
Current conservation dictates that each of the form factors $`F_{a\overline{b}\overline{b}}(Q^2)`$ and $`F_{a\overline{b}a}(Q^2)`$ are 1 at $`Q^2=0`$.
### A Impulse Approximation
Using dressed quark propagators, bound state BSAs, and the dressed $`qq\gamma `$-vertex, form factors can be calculated in impulse approximation. We denote by $`\mathrm{\Gamma }_\mu ^a(q,q^{};Q)`$ the quark-photon vertex describing the coupling of a photon with momentum $`Q`$ to a quark with final and initial momenta $`q`$ and $`q^{}=qQ`$ respectively and flavor $`a`$. With this notation, the vertices in Eq. (1) take the form<sup>*</sup><sup>*</sup>*We use Euclidean metric $`\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu }`$, $`\gamma _\mu ^{}=\gamma _\mu `$ and $`ab=_{i=1}^4a_ib_i`$.
$`\mathrm{\Lambda }_\nu ^{a\overline{b}\overline{b}}(P,Q)=2N_c{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4k}{(2\pi )^4}}\mathrm{Tr}[S^a(q)\mathrm{\Gamma }^{a\overline{b}}(q,q_+;P_{})`$ (4)
$`\times `$ $`S^b(q_+)i\mathrm{\Gamma }_\nu ^b(q_+,q_{};Q)S^b(q_{})\overline{\mathrm{\Gamma }}^{a\overline{b}}(q_{},q;P_+)],`$ (5)
where $`q=k+\frac{1}{2}P`$, $`q_\pm =k\frac{1}{2}P\pm \frac{1}{2}Q`$, $`P_\pm =P\pm \frac{1}{2}Q`$, and analogously for $`\mathrm{\Lambda }_\nu ^{a\overline{b}a}`$. The notation $`^\mathrm{\Lambda }`$ denotes a translationally-invariant regularization of the integral, with $`\mathrm{\Lambda }`$ the regularization mass-scale, which can be removed at the end of all calculations by taking the limit $`\mathrm{\Lambda }\mathrm{}`$. $`S(q)`$ is the dressed quark propagator and $`\mathrm{\Gamma }^{a\overline{b}}(q,q^{};P)`$ is the meson BSA, with $`P^2=m^2`$ the on-shell meson momentum, and $`q`$ and $`q^{}=qP`$ the quark and anti-quark momenta respectively.
Both $`S(q)`$ and $`\mathrm{\Gamma }^{a\overline{b}}(q,q^{};P)`$ are solutions of their respective DSEs
$`S(p)^1=Z_2i/p+Z_4m(\mu )`$ (6)
$`+`$ $`Z_1{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4q}{(2\pi )^4}}g^2D_{\mu \nu }(pq){\displaystyle \frac{\lambda ^i}{2}}\gamma _\mu S(q)\mathrm{\Gamma }_\nu ^i(q,p),`$ (7)
and
$$\mathrm{\Gamma }^{a\overline{b}}(p,p^{};Q)=^\mathrm{\Lambda }\frac{d^4q}{(2\pi )^4}K(p,q;Q)\chi ^{a\overline{b}}(q,q^{};Q),$$
(8)
where $`D_{\mu \nu }(k)`$ is the renormalized dressed-gluon propagator, $`\mathrm{\Gamma }_\nu ^i(q,p)`$ is the renormalized dressed quark-gluon vertex, $`K`$ is the renormalized $`\overline{q}q`$ scattering kernel that is irreducible with respect to a pair of $`\overline{q}q`$ lines, and $`\chi ^{a\overline{b}}(q,q^{};Q)=S^a(q)\mathrm{\Gamma }^{a\overline{b}}(q,q^{};Q)S^b(q^{})`$ is the BS wave function.
The solution of Eq. (6) is renormalized according to $`S(p)^1=i/p+m(\mu )`$ at a sufficiently large spacelike $`\mu ^2`$, with $`m(\mu )`$ the renormalized quark mass at the scale $`\mu `$. In Eq. (6), $`S`$, $`\mathrm{\Gamma }_\mu ^i`$ and $`m(\mu )`$ depend on the quark flavor, although we have not indicated this explicitly. The renormalization constants $`Z_2`$ and $`Z_4`$ depend on the renormalization point and the regularization mass-scale, but not on flavor: in our analysis we employ a flavor-independent renormalization scheme.
The meson BSAs $`\mathrm{\Gamma }^{a\overline{b}}(q,q^{};P)`$ are normalized according to the canonical normalization condition
$`P_\mu =N_c{\displaystyle \frac{}{P_\mu }}{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4q}{(2\pi )^4}}\{\mathrm{Tr}[\overline{\mathrm{\Gamma }}^{a\overline{b}}(\stackrel{~}{q}^{},\stackrel{~}{q};Q)`$ (11)
$`\times S^a(q+\eta P)\mathrm{\Gamma }^{a\overline{b}}(\stackrel{~}{q},\stackrel{~}{q}^{};Q)S^b(q+(\eta 1)P)]+`$
$`{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4k}{(2\pi )^4}}\mathrm{Tr}\left[\overline{\chi }^{a\overline{b}}(\stackrel{~}{k}^{},\stackrel{~}{k};Q)K(\stackrel{~}{k},\stackrel{~}{q};P)\chi ^{a\overline{b}}(\stackrel{~}{q},\stackrel{~}{q}^{};Q)\right]\},`$
at the mass shell $`P^2=Q^2=m^2`$, with $`\stackrel{~}{q}=q+\eta Q`$, $`\stackrel{~}{q}^{}=q+(\eta 1)Q`$, and similarly for $`\stackrel{~}{k}`$ and $`\stackrel{~}{k}^{}`$. We use the conventions where $`f_\pi =92\mathrm{MeV}`$, and $`\eta `$ describes the momentum partitioning between the quark and anti-quark. Note that physical observables should be independent of this parameter.
For pseudoscalar bound states the BSA is commonly decomposed into
$`\mathrm{\Gamma }(k+\eta P,k+(\eta 1)P;P)=`$ (12)
$`=`$ $`\gamma _5[iE(k^2;kP;\eta )+/PF(k^2;kP;\eta )`$ (14)
$`+/kG(k^2;kP;\eta )+\sigma _{\mu \nu }k_\mu P_\nu H(k^2;kP;\eta )],`$
with the invariant amplitudes $`E`$, $`F`$, $`G`$ and $`H`$ being Lorentz scalar functions of $`k^2`$ and $`kP=kP\mathrm{cos}\theta `$. Subsequently, each invariant amplitude can be expanded in $`kP`$ based on Chebyshev polynomials
$$f(k^2,kP;P^2)=\underset{i=0}{\overset{\mathrm{}}{}}U_i(\mathrm{cos}\theta )(kP)^if_i(k^2;P^2).$$
(15)
For charge-parity eigenstates, such as the pion, the invariant amplitudes $`E`$, $`F`$, $`G`$, and $`H`$ have a well-defined charge-parity if one chooses $`\eta =\frac{1}{2}`$. Therefore, these amplitudes are either entirely even ($`E`$, $`F`$, and $`H`$) or odd ($`G`$) in $`kP`$, and one needs only the even (or odd) Chebyshev moments to completely describe these amplitudes, which makes this a convenient decomposition.
### B The quark-photon vertex
The quark-photon vertex is the solution of the renormalized inhomogeneous BSE
$`\mathrm{\Gamma }_\mu ^a(p_+,p_{};Q)`$ $`=`$ $`Z_2\gamma _\mu +{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4q}{(2\pi )^4}}K(p,q;Q)`$ (17)
$`\times S^a(q_+)\mathrm{\Gamma }_\mu ^a(q_+,q_{};Q)S^a(q_{}),`$
with $`p_\pm =p\pm \frac{1}{2}Q`$ and $`q_\pm =q\pm \frac{1}{2}Q`$, and with the same kernel $`K`$ as the homogeneous BSE for meson bound states. Because of gauge invariance, it satisfies the Ward–Takahashi identity \[WTI\]
$$iQ_\mu \mathrm{\Gamma }_\mu ^a(p_+,p_{};Q)=S_a^1(p_+)S_a^1(p_{}).$$
(18)
Solutions of the homogeneous version of Eq. (17) at discrete timelike momenta $`Q^2`$ define vector meson bound states with masses $`m_V^2=Q^2`$. It follows that $`\mathrm{\Gamma }_\mu ^a(p;Q)`$ has poles at those locations, and behaves like
$$\mathrm{\Gamma }_\mu ^a(p_+,p_{};Q)\frac{\mathrm{\Gamma }_\mu ^{a\overline{a}V}(p_+,p_{};Q)f_Vm_V}{Q^2+m_V^2},$$
(19)
in the vicinity of these bound states, where $`\mathrm{\Gamma }_\mu ^{a\overline{a}V}`$ is the $`a\overline{a}`$ vector meson BSA, and $`f_V`$ the electroweak decay constant .
For the photon coupled to $`u`$\- and $`d`$-quarks, this results in a $`\rho `$-meson pole at $`Q^2=0.6\mathrm{GeV}^2`$. For the photon coupled to $`s`$-quarks, the first pole is located around $`Q^2=1.0\mathrm{GeV}^2`$ at the $`\varphi `$-mass. At the level of the ladder approximation, which is commonly used in practical calculations, there is no width generated for the vector meson, and the vertex has real poles. One would have to incorporate the open $`\pi \pi `$ channel in the ladder BSE kernel to produce a vector meson width; for the vertex, this would generate an imaginary part beyond the threshold for pion production, $`Q^2<4m_\pi ^2`$ in the timelike region.
The full vertex $`\mathrm{\Gamma }_\mu ^a`$ can be decomposed into 4 longitudinal components and 8 transverse components. The longitudinal components do not contribute to the form factors. In Ref. it was shown that only 5 of the 8 transverse components are important for the pion form factor, in the momentum range $`0.3<Q^2<1.0\mathrm{GeV}^2`$ the remaining 3 components contribute less than 1%. We expect that this will also be the case for the kaon form factor, and use the Dirac amplitudes $`T_1`$ to $`T_5`$ of Ref. only.
### C Charge conservation
At $`Q=0`$ the quark-photon vertex is completely specified by the differential Ward identity
$$i\mathrm{\Gamma }_\mu ^b(p,p;0)=\frac{}{p_\mu }S_b^1(p).$$
(20)
If this is inserted in Eq. (5), one finds after a change of integration variables $`kk\frac{1}{2}P`$
$`\mathrm{\Lambda }_\nu ^{a\overline{b}\overline{b}}(P,0)=2P_\mu F_{a\overline{b}\overline{b}}(0)=2N_c{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4q}{(2\pi )^4}}`$ (22)
$`\mathrm{Tr}\left[\overline{\mathrm{\Gamma }}^{a\overline{b}}(q^{},q;P)S^a(q)\mathrm{\Gamma }^{a\overline{b}}(q,q^{};P){\displaystyle \frac{S^b(qP)}{P}}\right],`$
with $`q^{}=qP`$. Comparing this expression with Eq. (11) with $`\eta =0`$, we recognize that the physical result $`F(Q^2=0)=1`$ follows directly from the canonical normalization condition for $`\mathrm{\Gamma }^{a\overline{b}}`$ with a BSE kernel $`K`$ independent of the meson momentum $`P`$. For the ladder truncation of the kernel, which we consider in our calculation in the next section, this is the case.
With a general momentum partitioning parameter $`\eta `$, the relation between the normalization condition and electromagnetic current conservation is not so obvious. However, using a different $`\eta `$ in loop diagrams (without external quark lines) is equivalent to a shift in integration variables. For processes that are not anomalous, loop integrals are independent of a shift of integration variables, provided that such a shift is performed consistently, and that all approximations employed respect Poincaré invariance. In performing such a shift, one has to take special care of the BSAs. The vertex function $`\mathrm{\Gamma }(q,q^{};P)`$, as function of the incoming and outgoing quark momenta, does not depend on $`\eta `$; it is only in commonly used decompositions in terms of Lorentz invariant amplitudes such as Eq. (12), where $`\eta `$ becomes relevant. The amplitudes $`E`$, $`F`$, $`G`$, and $`H`$ are scalar functions of $`k^2`$ and $`kP`$, which do depend on the choice for $`\eta `$. Under a change of $`\eta `$, some of the different Dirac structures (e.g. the amplitudes $`F`$ and $`G`$) will mix, as will the Chebyshev moments, $`f_i`$ in Eq. (15). Therefore, the results will be independent of the momentum routing in the loop integrals if and only if all Dirac amplitudes and their dependence on $`kP`$ are properly taken into account. Previously it has been shown that under these conditions the decay constants are indeed independent of $`\eta `$ .
Use of a bare quark-photon vertex, in combination with dressed propagators, in Eq. (5), clearly violates charge conservation and leads to $`F_\pi (0)1`$. With the Ball–Chiu Ansatz , which is commonly used in DSE studies of electromagnetic interactions , the electromagnetic current is explicitly conserved, $`F(Q^2=0)=1`$. However, the behavior of the form factor away from $`Q^2=0`$ is not constrained by current conservation, and in the present model, use of the Ball–Chiu Ansatz leads to a value for $`r_\pi ^2`$ which is about 50% too small . With the quark-photon vertex as the solution of the ladder BSE, together with quark propagators from the rainbow DSE, we satisfy all constraints from current conservation, and the calculated value of $`r_\pi ^2`$ is within 5% of the experimental value .
### D Beyond rainbow-ladder truncation
If one goes beyond the rainbow-ladder truncation for the DSEs for the propagators, BSAs and quark-photon vertex, one has to go beyond impulse approximation for the form factors in order to ensure current conservation. For example, one could include higher-order $`\alpha _s`$ corrections to the rainbow-ladder DSE and BSE kernels, as depicted in Fig. 1. Following the general procedure developed in Ref. , one can show that both the WTI, Eq. (18), and the differential Ward identity, Eq. (20), are preserved in the truncation indicated in Fig. 1, as is the axial-vector WTI, which is important for the Goldstone nature of the pions.
The resulting BSE kernel $`K(q,p;P)`$ now becomes dependent on the meson momentum $`P`$, which means that the second term of the normalization condition, Eq.(11), is nonzero. To be specific, with the choice $`\eta =0`$, this introduces the four extra terms in the normalization condition, diagrammatically depicted in Fig. 2. These four additional diagrams can be generated from the BSE kernel in the bottom part of Fig. 1 by taking the derivative with respect to the meson momentum $`P`$, where $`P`$ flows through one quark propagator only.
Since taking the derivative with respect to $`P`$ is equivalent to the insertion of a zero-momentum photon according to the differential WTI, Eq. (18), it is obvious which diagrams have to be added to the impulse approximation to ensure current conservation, see Fig. (3). In the limit $`Q0`$ these four additional diagrams become identical to the four additional diagrams in Fig. (2), provided that the vertex satisfies the differential WTI. Of course, there are similar contributions to $`\mathrm{\Lambda }^{a\overline{b}a}`$, which can be identified with terms in the normalization condition with $`\eta =1`$.
Also, simple addition of contributions due to pion and kaon loops to Eq. (5), in combination with a ladder-rainbow truncation for the DSEs, will generally violate current conservation. Current conservation requires a consistent treatment of the kernels for both the DSE and BSE equations and the approximation for the photon-hadron coupling. At present it is not clear how to incorporate meson loops self-consistently in such an approach, but we expect corrections coming from such loops to be small in the spacelike region. In Ref. it was demonstrated that the quark core can generate most of the pion charge radius, and that pion loops contribute less then 15% to $`r_\pi ^2`$. For larger values of $`Q^2`$ the effect from meson loops reduces even further, and for $`Q^2>1\mathrm{GeV}^2`$ we expect the contribution of such loops to be negligible.
## III Model calculations
For the BSE we use a ladder truncation
$$K(p,q;P)𝒢(k^2)D_{\mu \nu }^{\mathrm{free}}(k)\frac{\lambda ^a}{2}\gamma _\mu \frac{\lambda ^a}{2}\gamma _\nu ,$$
(23)
where $`D_{\mu \nu }^{\mathrm{free}}(k=pq)`$ is the free gluon propagator in Landau gauge. The resulting BSE is consistent with a rainbow truncation $`\mathrm{\Gamma }_\nu ^a(q,p)\gamma _\nu \lambda ^a/2`$ for the quark DSE, Eq. (6), in the sense that the combination produces vector and axial-vector vertices satisfying the respective WTIs. In the axial case, this ensures that in the chiral limit the ground state pseudoscalar mesons are the massless Goldstone bosons associated with chiral symmetry breaking . In the vector case, this ensures electromagnetic current conservation.
The model is completely specified once a form is chosen for the “effective coupling” $`𝒢(k^2)`$. We employ the Ansatz
$`{\displaystyle \frac{𝒢(k^2)}{k^2}}`$ $`=`$ $`{\displaystyle \frac{4\pi ^2Dk^2}{\omega ^6}}\mathrm{e}^{k^2/\omega ^2}`$ (25)
$`+{\displaystyle \frac{4\pi ^2\gamma _m(k^2)}{\frac{1}{2}\mathrm{ln}\left[\tau +\left(1+k^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2\right)^2\right]}},`$
with $`\gamma _m=12/(332N_f)`$ and $`(s)=(1\mathrm{exp}\frac{s}{4m_t^2})/s`$. This Ansatz preserves the one-loop renormalization group behavior of QCD, and ensures that we reproduce perturbation theory in the perturbative region. The first term of Eq. (25) implements the strong infrared enhancement in the region $`0<k^2<1\mathrm{GeV}^2`$ which is a phenomenological requirement for sufficient dynamical chiral symmetry breaking to produce an acceptable strength for the quark condensate . We use $`m_t=0.5\mathrm{GeV}`$, $`\tau =\mathrm{e}^21`$, $`N_f=4`$, $`\mathrm{\Lambda }_{\mathrm{QCD}}=0.234\mathrm{GeV}`$, and a renormalization point $`\mu =19\mathrm{GeV}`$, well in the perturbative region . The remaining parameters, $`\omega =0.4\mathrm{GeV}`$ and $`D=0.93\mathrm{GeV}^2`$, are fitted to give a good description of the chiral condensate, $`m_{\pi /K}`$ and $`f_\pi `$. The subsequent values for $`f_K`$ and the masses and decay constants of the vector mesons $`\rho ,\varphi ,K^{}`$ are in agreement with the experimental data , see Table I.
### A Results for $`u\overline{u}u`$, $`u\overline{s}u`$, and $`u\overline{s}\overline{s}`$ form factors
The pion and kaon form factors are given by
$`F_\pi (Q^2)`$ $`=`$ $`{\displaystyle \frac{2}{3}}F_{u\overline{d}u}(Q^2)+{\displaystyle \frac{1}{3}}F_{u\overline{d}\overline{d}}(Q^2),`$ (26)
$`F_{K^+}(Q^2)`$ $`=`$ $`{\displaystyle \frac{2}{3}}F_{u\overline{s}u}(Q^2)+{\displaystyle \frac{1}{3}}F_{u\overline{s}\overline{s}}(Q^2),`$ (27)
$`F_{K^0}(Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{3}}F_{d\overline{s}d}(Q^2)+{\displaystyle \frac{1}{3}}F_{d\overline{s}\overline{s}}(Q^2),`$ (28)
where the quark and anti-quark charges are evident. We work in the SU(2) isospin limit, where the strong interaction does not discriminate between $`u`$\- and $`d`$-quarks, so for the pion we simply have $`F_\pi (Q^2)=F_{u\overline{u}u}(Q^2)`$. Thus there are only three independent form factors, $`F_{u\overline{u}u}(Q^2)`$, $`F_{u\overline{s}u}(Q^2)`$, and $`F_{u\overline{s}\overline{s}}(Q^2)`$, which are shown in Fig. 4. Our estimate of the numerical error in these calculations is less than 1% for $`F_{u\overline{u}u}(Q^2)`$, and 2% for the other two form factors.
For the pion we use only the leading terms of the expansion of the BSAs and the quark-photon vertex in $`kP`$, Eq. (15). Higher order terms do not change the results more than 1% in this momentum regime, although they are needed at larger values of $`Q^2`$. For the kaon we have to use more terms in the expansion, even at $`Q^2=0`$, to obtain independence from the parameter $`\eta `$, and to ensure current conservation. With terms up to order $`(kP)^1`$ only, there is a spread in our results of more than 10% at $`Q^2=0`$ (from 0.94 to 1.06) if we change $`\eta `$ between 0 and 1. Including the next two terms reduces that spread to less then 3%, illustrating that the result of a loop integral is independent of this unphysical parameter $`\eta `$, provided that all relevant Dirac structures and the dependence on $`kP`$ are properly taken into account.
The results for $`F_{u\overline{u}u}`$ and $`F_{u\overline{s}u}`$ are remarkably close to each other, indicating that the flavor of the spectator quark matters very little. Within our numerical errors, they are almost indistinguishable on the $`Q^2`$ domain shown. There is a slight difference in the slope of these form factors: $`r_{u\overline{u}u}^2=0.45\mathrm{fm}^2`$ versus $`r_{u\overline{s}u}^2=0.47\mathrm{fm}^2`$. These results are in good agreement with the pion charge radius, $`r_\pi ^2=0.46\mathrm{fm}^2`$, obtained in Ref. using all eight Dirac amplitudes of the quark-photon vertex.
The result for $`F_{u\overline{s}\overline{s}}`$ is quite different in that it has a significantly smaller slope characterized by a radius parameter $`r_{u\overline{s}\overline{s}}^2=0.21\mathrm{fm}^2`$. This is due to the larger mass of the strange quark, and as a consequence the neutral kaon charge radius $`r_{K^0}^2`$ will be negative. A similar effect was observed for the neutron form factor, where the heavier mass of the $`0^+(ud)`$-diquark compared to the $`d`$ quark mass leads to a negative charge radius . Our result is also consistent with the qualitative aspects of the vector meson dominance \[VMD\] picture: the lowest-mass bound state pole in the $`ss\gamma `$-vertex is the $`\varphi `$, at $`Q^2=1.0\mathrm{GeV}^2`$, which is significantly further from the photon point than is the $`\rho `$ pole in the $`uu\gamma `$-vertex at $`Q^2=0.6\mathrm{GeV}^2`$. This observation, as well as the difference between $`r_{u\overline{s}\overline{s}}^2`$ and $`r_{u\overline{s}u}^2`$, is consistent with the larger mass of the strange quark.
### B Results for the meson form factors
The results in this model for the pion form factor at low $`Q^2`$, in particular the pion charge radius, were presented previously . The obtained charge radii for the kaon are presented in Table II, and are in reasonable agreement with the experimental data, without any readjustment of the model. In Fig. 5 we show our result for the charged kaon, which is in good agreement with the available data.
Finally, in Fig. 6 we present $`Q^2F(Q^2)`$ for $`\pi `$ and K<sup>0,±</sup> for a larger $`Q^2`$ range to anticipate data that may be forthcoming from experiments at JLab and possibly other facilities in the future. In this momentum range, even for $`F_\pi (Q^2)`$ the dependence on $`kP`$ becomes important, and terms up to $`(kP)^3`$ in Eq. (15) are required to produce a converged result at $`Q^2=13\mathrm{GeV}^2`$. Higher-order terms do not change the results by more than 1% in this momentum range. Our estimate is that the net numerical accuracy for $`F_{u\overline{u}u}`$, $`F_{u\overline{s}u}`$, and $`F_{u\overline{s}\overline{s}}`$ is about 2-3% at these values of $`Q^2`$. This translates to a similar level of accuracy for $`F_\pi `$ and $`F_{K^+}`$, and to a somewhat larger relative error, about 5%, for $`F_{K^0}`$, which is the difference of $`F_{u\overline{s}\overline{s}}`$ and $`F_{u\overline{s}u}`$. At $`Q^2>3\mathrm{GeV}^2`$, higher-order Chebyshev moments may be necessary, but current numerical methods prevent their accurate determination at large $`Q^2`$.
Over the entire spacelike momentum range considered, $`F_\pi (Q^2)<F_{K^+}(Q^2)`$, and $`Q^2F(Q^2)`$ rises with $`Q^2`$ until $`Q^2=3\mathrm{GeV}^2`$ for all three form factors. In this momentum range our results for both the pion and the $`K^+`$ form factor can be fitted quite well by a simple monopole $`m^2/(Q^2+m^2)`$, with a mass $`m^2=0.53\mathrm{GeV}^2`$ for the pion and $`m^2=0.61\mathrm{GeV}^2`$ for the $`K^+`$. A VMD model, two monopoles for the two form factors $`F_{u\overline{s}u}`$ and $`F_{u\overline{s}\overline{s}}`$ in Eqs. (27) and (28) with the physical $`\rho `$ and $`\varphi `$ masses respectively, does not reproduce our results for the kaon form factors very well. For example, at $`Q^2=1\mathrm{GeV}^2`$, VMD overshoots our $`F_{K^+}`$ calculation by almost 10%, whereas the monopole fit is within 2% of our result. This difference between VMD and our calculations grows with $`Q^2`$.
Above $`Q^23.5\mathrm{GeV}^2`$ the monopole fits begin to deviate significantly from our results and $`Q^2F(Q^2)`$ starts to decrease. In a more realistic model, that takes meson loop corrections into account self-consistently, it is very well conceivable that this turn-over happens at somewhat lower values of $`Q^2`$: meson loops are expected to contribute up to 15% to $`r_\pi ^2`$ , but their contribution to the form factor decreases rapidly with increasing spacelike momenta. In the presence of meson loop corrections the contribution to the form factor from the impulse approximation has to be smaller than in our calculation in order to maintain agreement with the low-$`Q^2`$ data. Therefore it is not unlikely that at intermediate momenta in the present approach we overestimate the form factors, which may explain the difference between the data points at $`Q^2=3.3\mathrm{GeV}^2`$ and our calculated results. More accurate results from JLab, in combination with realistic model calculations that include meson loop corrections self-consistently, may be able to resolve this question.
At asymptotically large $`Q^2`$, factorized pQCD predicts that the form factor behaves like $`Q^2F(Q^2)c`$, with $`c=16\pi f_\pi ^2\alpha _s(Q^2)`$. Since our truncation and the Ansatz, Eq. (25), is constructed so as to preserve asymptotic freedom, we are guaranteed to recover the leading power-law asymptotic behavior. An explicit verification of this behavior, and calculation of the constant $`c`$, is not readily available within our present framework since numerical accuracy at large $`Q^2`$ is problematic. However, it is clear from our results that at $`Q^24\mathrm{GeV}^2`$ the form factor has not yet reached its asymptotic value, and it is unlikely that experiments can access the true asymptotic region in the near future. In simplified models such as that of Ref. however, it is straightforward to demonstrate that the impulse approximation does indeed lead to the power-law behavior predicted by pQCD.
## IV Summary
We calculate the pion and kaon electromagnetic form factors within the DSE approach. The method is completely Poincaré invariant, and the only approximation made is a self-consistent truncation of the set of DSEs, which respects the relevant vector and axial-vector WTIs. The employed quark propagators, the meson BSAs, and the quark-photon vertex are solutions of their DSEs in rainbow-ladder truncation with all parameters fixed previously by fitting the chiral condensate, $`m_{\pi /K}`$ and $`f_\pi `$. We include all relevant Dirac amplitudes for the BSAs and their dependence upon $`kP`$. The electromagnetic current is explicitly conserved in this approach, and there is no fine-tuning needed to obtain $`F_\pi (0)=1=F_{K^+}(0)`$ and $`F_{K^0}(0)=0`$. We also demonstrate explicitly that our results are (within numerical accuracy) independent of the momentum partitioning of the BSAs. The obtained pion and kaon form factors are in good agreement with the available data over the entire $`Q^2`$ range considered, and the calculated charge radii are within the error bars of their experimental values.
These charge radii are somewhat larger than those obtained in a previous study that was framed in terms of semi-phenomenological representations for BSAs and confined quark propagators within the impulse approximation. The main difference with that work is that here we use numerical solutions of truncated DSEs for all the elements needed in Eq. (5), and that all our parameters were fixed previously. In comparison with theoretical calculations based on other methods, it is interesting to note that our results are very similar to those obtained in Ref. , in particular for the neutral kaon.
At intermediate values of $`Q^2`$ our calculations are qualitatively similar to those obtained in both Ref. and Ref. . Up to about $`Q^2=3\mathrm{GeV}^2`$, both $`F_\pi `$ and $`F_{K^+}`$ can be fitted quite well by a monopole form, with monopole masses of $`m^2=0.53\mathrm{GeV}^2`$ and $`m^2=0.61\mathrm{GeV}^2`$ respectively. At large $`Q^2`$ the DSE approach does reproduce the pQCD power-law behavior , but this behavior does not occur until well beyond the $`Q^2`$ range considered in our present calculations and accessible at current accelerators.
## ACKNOWLEDGMENTS
We acknowledge useful conversations and correspondence with C.D. Roberts, D. Jarecke and S.R. Cotanch. This work was funded by the National Science Foundation under grant No. PHY97-22429, and benefited from the resources of the National Energy Research Scientific Computing Center.
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# Isotropization of Bianchi type models and a new FRW solution in Brans-Dicke theory
## 1 INTRODUCTION
The universe is nowadays at big scales homogeneous and isotropic as measured in the CMBR by the COBE satellite , and must also has been having these properties since, at least, the era of nucleosynthesis . In order to explain the isotropy of the universe from theoretical anisotropic models, many authors have considered the Bianchi models that can in principle evolve to a Friedmann-Robertson-Walker (FRW) cosmology. It has been shown that some Bianchi models in General Relativity (GR) tend to their isotropic solutions, up some extent , and even that they can explain the level of anisotropy measured by COBE . Motived by these facts, we have been working in Brans-Dicke (BD) theory to investigate if Bianchi universes are able to isotropize as the universe evolves, and if its evolution can be inflationary. In previous investigations we have shown that anisotropic, Bianchi type I, V, and IX models tend to isotropize as models evolve . However, this may happen for some restrictive values of $`\omega `$ in the cases of Bianchi type I and IX, and only Bianchi type V model can accomplish an isotropization mechanism within BD current constraints on $`\omega `$. It has been also shown that the isotropization mechanism in the Bianchi type V model can be inflationary, without the presence of any cosmological constant, when small values for the coupling constant $`\omega `$ are considered , as in the case of some induced gravity (IG) models . We have recently shown, however, that the isotropization mechanism can be attained with sufficient amount of e-folds of inflation, only for negative values of $`\omega `$. In the present report we review and generalize some of the main results and present a new $`k0`$ FRW solution in BD theory.
This paper is organized as follows. In section 2 the BD Bianchi type I, V and IX equations are presented. In section 3 we review the main results on these models and present a new solution to curve $`k0`$ FRW cosmologies. Finally, conclusions are in section 4.
## 2 ANISOTROPIC EQUATIONS FOR BIANCHI MODELS
In previous investigations we have used scaled variables, in terms of which our solutions have been given, therefore following we use them: the scaled field $`\psi \varphi a^{3(1\nu )}`$, a new cosmic time parameter $`d\eta a^{3\nu }dt`$, $`()^{}\frac{d}{d\eta }`$, the ‘volume’ $`a^3a_1a_2a_3`$, and the Hubble parameters $`H_ia_{i}^{}{}_{}{}^{}/a_i`$ corresponding to the scale factors $`a_i=a_i(\eta )`$ for $`i=1,2,3`$. We assume comoving coordinates and a perfect fluid with barotropic equation of state, $`p=\nu \rho `$, where $`\nu `$ is a constant. Using these definitions one obtains the cosmological equations for Bianchi type I, V and IX models (in units with $`G=c=1`$):
$$(\psi H_i)^{}\psi a^{6\nu }R_{ij}=\frac{8\pi }{3+2\omega }[1+(1\nu )\omega ]\rho a^{3(1+\nu )}\mathrm{for}i=1,2,3.$$
(1)
$`H_1H_2+H_1H_3+H_2H_3+[1+(1\nu )\omega ]\left(H_1+H_2+H_3\right){\displaystyle \frac{\psi ^{}}{\psi }}`$
$`(1\nu )[1+\omega (1\nu )/2](H_1+H_2+H_3)^2{\displaystyle \frac{\omega }{2}}\left({\displaystyle \frac{\psi ^{}}{\psi }}\right)^2{\displaystyle \frac{R_j}{2}}a^{6\nu }`$
$`=\mathrm{\hspace{0.17em}\; 8}\pi {\displaystyle \frac{\rho a^{3(1+\nu )}}{\psi }},`$ (2)
$$\psi ^{\prime \prime }+(\nu 1)a^{6\nu }R_j\psi =\frac{8\pi }{3+2\omega }[2(23\nu )+3(1\nu )^2\omega ]\rho a^{3(1+\nu )},$$
(3)
where a column sum is given by $`R_j\mathrm{\Sigma }_iR_{ij}`$, where $`j=`$I, V or IX and
$$\begin{array}{ccccc}& & ^I& ^V& ^{IX}\\ & & & & \\ & & 0& 2/a_1^2& [a_1^4a_2^4a_3^4+2a_2^2a_3^2]/(2a^6)\\ & & & & \\ & R_{ij}=& 0& 2/a_1^2& [a_2^4a_3^4a_1^4+2a_1^2a_3^2]/(2a^6)\\ & & & & \\ & & 0& 2/a_1^2& [a_3^4a_1^4a_2^4+2a_1^2a_2^2]/(2a^6)\end{array}$$
(4)
For the Bianchi type V model one has the additional constraint: $`H_2+H_3=\mathrm{\hspace{0.17em}\; 2}H_1`$, implying that $`a_2`$ and $`a_3`$ are inverse proportional functions, $`a_2a_3=a_1^2`$; note that the mean Hubble parameter, $`H\frac{1}{3}(H_1+H_2+H_3)`$, is for this Bianchi type model $`H=H_1`$.
Additionally, the continuity equation yields: $`\rho a^{3(1+\nu )}=\mathrm{const}.M_\nu `$, $`M_\nu `$ being a dimensional constant depending on the fluid present. The vacuum case is attained when $`M_\nu =0`$.
The system of ordinary differential equations, Eqs. (1-3), can be once integrated to get<sup>1</sup><sup>1</sup>1A similar equation, that is valid only for the Bianchi type V model, was derived in Ref. . Now we generalize that result in such a way that Eq. (2) is valid for Bianchi models I, V and IX, as well for FRW cosmologies.:
$`\psi \psi ^{\prime \prime }{\displaystyle \frac{2}{3(1\nu )}}\psi _{}^{}{}_{}{}^{2}{\displaystyle \frac{2(13\nu )}{3(1\nu )}}[m_\nu (13\nu )\eta +\delta ]\psi ^{}+[2+(1\nu )(1+3\nu )\omega ]m_\nu \psi +`$
$`{\displaystyle \frac{2}{3(1\nu )}}[23\nu +{\displaystyle \frac{3}{2}}(1\nu )^2\omega ][m_\nu (13\nu )\eta +\delta ]^2+(1\nu )(h_1^2+h_2^2+h_3^2)=0,`$ (5)
where $`\delta `$ is an integration constant, $`m_\nu \frac{8\pi M_\nu }{3+2\omega }`$, and the Hubble rates are written as follows (similar to the Bianchi type I model deduced in Ref. ):
$$H_i=\frac{1}{3}\left(H_1+H_2+H_3\right)+\frac{h_i}{\psi }=\frac{\psi ^{}(13\nu )m_\nu \eta \delta +3(1\nu )h_i}{3(1\nu )\psi },$$
(6)
where the $`h_i`$’s are some unknown functions of $`\eta `$ that determine the anisotropic character of the solutions. If $`h_i=0`$ for $`i=1,2`$, $`3`$ simultaneously, no anisotropy is present, which is the case of FRW cosmologies. Furthermore, Bianchi models obey the condition
$$h_1+h_2+h_3=0$$
(7)
to demand consistency with Eq. (6). For the Bianchi type V model one has additionally that $`h_1=0`$, since $`H_1=H`$ as mentioned above.
Equations (6) and (7) imply that the mean Hubble parameter is determined by $`\psi `$ alone:
$$3H=H_1+H_2+H_3=\frac{1}{1\nu }\left[\frac{\psi ^{}}{\psi }\frac{(13\nu )m_\nu \eta +\delta }{\psi }\right].$$
(8)
In order to analyze the anisotropic character of the solutions, we consider the anisotropic shear, $`\sigma (H_1H_2)^2(H_2H_3)^2(H_3H_1)^2`$. $`\sigma =0`$ is a necessary condition to obtain a FRW cosmology since it implies $`H_1=H_2=H_3`$, cf. Ref. . If the sum of the squared differences of the Hubble expansion rates tends to zero, it would mean that the anisotropic scale factors tend to a single function of time which is, certainly, the scale factor of a FRW solution.
The anisotropic shear becomes, using Eqs. (6) and (7), $`\sigma (\eta )=\frac{3(h_1^2+h_2^2+h_3^2)}{\psi ^2}`$, or the dimensionless shear parameter , using Eq. (8):
$$\frac{\sigma }{H^2}=\frac{27(1\nu )^2(h_1^2+h_2^2+h_3^2)}{[\psi ^{}(13\nu )m_\nu \eta \delta ]^2}.$$
(9)
If the above equations admit solutions such that $`\sigma /H^20`$ as $`\eta \mathrm{}`$ ($`t\mathrm{}`$), then one has time asymptotic isotropization solutions, similar to solutions found for the Bianchi models in GR . In fact, one does not need to impose an asymptotic, infinity condition, but just that $`\eta \eta _{}`$, where $`\eta _{}`$ is yet some arbitrary value to warrant that $`\sigma /H^2`$ can be bounded from above.
## 3 ANISOTROPIC AND ISOTROPIC SOLUTIONS
The problem to find solutions of Bianchi and FRW models in the BD theory has been reduced to solve the coupled system of equations (2), (6) and (1). Let us present in the following subsections the known and new solutions.
### 3.1 Bianchi type I model
For this Bianchi model the known, the general solution is found in which the $`h_i`$’s are constants, then Eq. (2) is decoupled from Eqs. (1) and (6), and the solution is $`\psi =A_I\eta ^2+B_I\eta +C_I`$, where the constants $`A_I,B_I,C_I`$ are reported elsewhere . This model can be solved in a general way since the curvature sum column $`R_I=0`$, then Eq. (3) can be directly integrated. This solution represents a particular solution of Eq. (2). Direct substitution of $`\psi `$ into Eq. (6) gives the Hubble rates, and into Eq. (9) shows that solutions isotropize as time evolves, that is, $`\sigma /H^20`$ as $`\eta \mathrm{}`$, see Ref. . However, the isotropization mechanism is only possible for solutions such that $`\mathrm{\Delta }_IB_I^24A_IC_I`$
$`\mathrm{\Delta }_I={\displaystyle \frac{2(23\nu )+3(1\nu )^2\omega }{3(1\nu )^2(3+2\omega )}}`$ $`[{\displaystyle \frac{(1\nu )^2B_I^2}{2(23\nu )+3(1\nu )^2\omega }}2(13\nu )\delta B`$ (10)
$`+[2(23\nu )+3(1\nu )^2\omega ]\delta ^2`$
$`+3(1\nu )^2(h_1^2+h_2^2+h_3^2)]`$
is negative . Then, some restrictions on $`\omega `$ apply. For instance, in Dehnen’s IG theory $`\omega 1`$, then the isotropization mechanism is not possible in this Bianchi model .
### 3.2 Bianchi type V model
This Bianchi model is more complicated because curvature terms are different than zero. Still, it is possible to find particular solutions of Eq. (2), since for this model the $`h_i`$’s are constants too, and this equation is decoupled as is the case of Bianchi type I model. The general solution of this model should be obtained through the general solution of Eq. (2), yet unknown. We have found a particular solution that is again of the form $`\psi =A_V\eta ^2+B_V\eta +C_V`$, where the constants $`A_V,B_V,C_V`$ are reported in Refs. . This solution is such that
$$\mathrm{\Delta }_V=\frac{8(13\nu )^2}{18\nu +(1+3\nu )^2\omega }h_2^2$$
(11)
is negative for $`\omega >18\nu /(1+3\nu )^2`$, so the solution tends to isotropic solution within BD theory constraints , $`\omega 500`$, that is $`\sigma /H^20`$ as $`\eta \mathrm{}`$. For this Bianchi type model, Dehnen’s IG theory can achieve an isotropization mechanism . An inflationary behavior may be observed in type V models, but to get enough e-foldings of inflation ($`N68`$) one must demand that $`\omega \frac{3}{2}`$ in consistency with previous results .
### 3.3 Bianchi type IX model
Bianchi type IX model is the most complicated to solve, since curvature terms involve quartic polynomials of the scale factors, see Eq. (4). For this Bianchi model it implies, by imposing the condition that $`h_i`$’s are constants, severe algebraic constraints on the scale factors, so it seems more likely that $`h_i`$’s are functions. This explains why no totally anisotropic ($`H_1H_2H_3`$) solution has been found yet. In Ref. we have analyzed the case when the polynomial solution for $`\psi `$ is valid. In this case, unfortunately, we could not found explicitly the values of the constants $`A_{IX},B_{IX},C_{IX}`$. If this solution is valid, however, one has that $`h_1^2+h_2^2+h_3^2=D\eta ^2+F\eta +G`$, where $`D,F,G`$ are constants. Accordingly, Eq. (9) indicates that the solution must tend, as time evolves, to the positive curvature FRW solution, i.e. one has again that $`\sigma /H^20`$ as $`\eta \mathrm{}`$. However, a definitive answer will arrive by obtaining explicitly the values $`A_{IX},B_{IX},C_{IX}`$.
The only possible solutions for Bianchi type I and V models imply that the $`h_i`$’s are constants, whereas for type IX they are unknown functions of $`\eta `$. An explanation of this fact resides in the property that Bianchi type I and V models have curvature terms of FRW type, whereas type IX has a very much complicated form, see Eq. (4). So the things, it seems that the most general solution of Eq. (2) with $`h_i`$’s constants would give general solutions for Bianchi models I and V. The particular quadratic-polynomial solution of Eq. (2) represents in the case of Bianchi type I model its the general solution, whereas for Bianchi model V it is only a particular solution. Then, other particular solutions, possibly of non-polynomial nature, are expected to be found for Eq. (2) that will reveal new aspects of Bianchi type V model. Finally, for Bianchi type IX model the quadratic-polynomial can be a possible solution, not yet confirmed. However, solutions with $`h_i(\eta )`$ valid for the Bianchi type IX model are almost impossible to find because of the complexities involved in the curvature terms.
### 3.4 FRW solutions
It turns out that Eq. (2) is also valid for the FRW models when the anisotropic parameters vanish, i.e. $`h_1=h_2=h_3=0`$. Solutions of this equation solve FRW cosmologies in BD theory . The known solutions for $`\psi `$ are quadratic polynomials in $`\eta `$ as well. The general flat ($`k=0`$) solution is a particular solution of Eq. (2) in which the coefficient of the quadratic polynomial term, $`A`$, is equal to the corresponding coefficient ($`A_I`$) of the Bianchi I case. For curved ($`k0`$) FRW cosmologies the known particular solution is such that the coefficient of $`A`$ of the quadratic polynomial term is equal to the corresponding coefficient ($`A_V`$) in the Bianchi V case. In this way, one can see a correspondence between anisotropic and isotropic solutions.
We have found a new solution of Eq. (2) that is valid for $`k0`$ FRW cosmologies. The new solution is:
$$\psi =m_\nu \left(\frac{13\nu }{1+3\nu }\right)^2\left[\kappa _1+(13\nu )\eta \right]\left[\kappa _1+2\eta +\kappa _2\left[\kappa _1+(13\nu )\eta \right]^{\frac{2}{13\nu }}\right],$$
(12)
where $`\kappa _1,\kappa _2`$ are arbitrary integration constants. This is the general solution of curved FRW cosmologies when the following relationships are valid:
$`\delta `$ $`=`$ $`0`$
$`\omega `$ $`=`$ $`{\displaystyle \frac{18\nu }{(1+3\nu )^2}}.`$ (13)
Though the latter relationship constrains the range of possible values of $`\omega `$ and $`\nu `$, one can find values of physical interest, e.g. $`\omega =1`$ that has some interest in string effective theories, see for instance Ref. . Moreover, when $`\nu 1/3`$ one obtains the GR limit $`\omega \mathrm{}`$. Finally, one gets the dust model ($`\nu _{\mathrm{dust}}=0`$) in the limit when $`\omega 0`$, like in Dehnen’s IG theory .
The Hubble parameter is given by:
$$H=\frac{\kappa _1+\frac{19\nu }{13\nu }\eta +\kappa _2\left[\kappa _1+(13\nu )\eta \right]^{\frac{2}{13\nu }}}{\left[\kappa _1+(13\nu )\eta \right]\left[\kappa _1+2\eta +\kappa _2\left[\kappa _1+(13\nu )\eta \right]^{\frac{2}{13\nu }}\right]},$$
(14)
from which one can find the scale factor:
$$a=\left[\frac{k}{\kappa _2}\right]^{\frac{1}{2(13\nu )}}\frac{\left[\kappa _1+2\eta +\kappa _2\left[\kappa _1+(13\nu )\eta \right]^{\frac{2}{13\nu }}\right]^{\frac{1}{2(13\nu )}}}{\left[\kappa _1+(13\nu )\eta \right]^{\frac{3\nu }{(13\nu )^2}}}.$$
(15)
The BD field ($`\varphi =\psi a^{3(1\nu )}`$) is obtained through Eqs. (12) and (15):
$$\varphi =m_\nu \left(\frac{13\nu }{1+3\nu }\right)^2\left(\frac{\kappa _2}{k}\right)^{\frac{3(1\nu )}{2(13\nu )}}\frac{\left[\kappa _1+(13\nu )\eta \right]^{\frac{1+3\nu }{(13\nu )^2}}}{\left[\kappa _1+2\eta +\kappa _2\left[\kappa _1+(13\nu )\eta \right]^{\frac{2}{13\nu }}\right]^{\frac{1+3\nu }{2(13\nu )}}}.$$
(16)
Eqs. (15) and (16) imply that the sign of $`\kappa _2`$ is equal to the sign of $`k`$ for most values of $`\nu `$. For open ($`k=1`$) models this implies that $`\kappa _2`$ must be positive, and for closed ($`k=+1`$ ) models $`\kappa _2`$ must be negative which allows the solutions to (re)collapse: The value of $`\kappa _2`$ determines the time of maximum expansion, so it is very related to the mass ($`m_\nu `$) of the model. On the other hand, $`\kappa _1`$ represents a $`\eta `$-time shift.
Because of the mathematical form of Eq. (14) it is not possible to have an inflationary era that lasts for a sufficient time period to solve the horizon and flatness problems. To show the behavior of the models we have plotted the above-given formulae for different values of $`\omega `$, the curvature constant $`k`$, and integration constants $`\kappa _1`$ and $`\kappa _2`$. We have chosen in all figures that $`\varphi ^{}_{\eta =0}\varphi _0^{}=0`$ as initial condition. Given a specific value for $`\omega `$ implies two possible values of $`\nu `$ consistent with Eq. (3.4). For $`\omega =1`$, a value that makes BD theory to resemble string effective theories , it corresponds $`\nu _1=(2\sqrt{3})/30.0893`$ (that represents a quasi dust model) and $`\nu _2=(2+\sqrt{3})/31.2440`$. In figures 1 and 2 we plotted the scale factor and BD field for $`\nu _1`$. Figure 1 shows an open model with $`\kappa _1=\kappa _2=k=1`$, whereas figure 2 is a closed model with $`\kappa _1=k=1`$ and $`\kappa _2=0.001`$. In closed models, the smaller $`\kappa _2`$ is, the later in time they will recollapse. However, for models with $`2/(13\nu )<1`$ one can find an upper limit on $`\kappa _2`$ such that models never recollapse (this happens, for instance, when $`\nu =0.3719`$, making $`\omega =500`$, and if $`\kappa _2<1`$); the effect of negative pressure avoids recollapse.
One can compute the asymptotic limit for $`k=1`$ models, when $`\eta \mathrm{}`$. There are two limit cases, when the quantity $`2/(13\nu )`$ is greater or smaller than 1, cf. terms in the equations above. In the former case one has that
$`a_{\mathrm{asimp1}}`$ $`=`$ $`\left(k\right)^{\frac{1}{2(13\nu )}}\left[(13\nu )\eta \right]^{\frac{1}{13\nu }}`$
$`\varphi _{\mathrm{asimp1}}`$ $`=`$ $`m_\nu \kappa _2\left({\displaystyle \frac{13\nu }{1+3\nu }}\right)^2\left(k\right)^{\frac{3(1\nu )}{2(13\nu )}}=\mathrm{const}.,`$ (17)
whereas in the latter case one gets
$`a_{\mathrm{asimp2}}`$ $`=`$ $`\left[{\displaystyle \frac{2k}{\kappa _2}}\right]^{\frac{1}{2(13\nu )}}\left(13\nu \right)^{\frac{3\nu }{(13\nu )^2}}\eta ^{\frac{19\nu }{2(13\nu )^2}}`$
$`\varphi _{\mathrm{asimp2}}`$ $`=`$ $`2m_\nu {\displaystyle \frac{(13\nu )^{3+\frac{9\nu (1\nu )}{(13\nu )^2}}}{(1+3\nu )^2}}\left({\displaystyle \frac{\kappa _2}{2k}}\right)^{\frac{3(1\nu )}{2(13\nu )}}\eta ^{\frac{(1+3\nu )^2}{2(13\nu )^2}}.`$ (18)
The plots of figure 1 tends to Eqs. (3.4); one observes that $`\varphi `$ tends to a constant value that can be fix to be $`G^1`$ through the right choice of the constants $`m_\nu `$ and $`\kappa _2`$.
Figures 3 and 4 show an open and closed model, respectively, with $`\omega =500`$. Again, there are two possible values of $`\nu `$, $`\nu _3=0.2988`$ and $`\nu _4=0.3719`$. In these figures we have chosen the value of $`\nu _3`$. In figure 3 we have chosen $`\kappa _1=\kappa _2=k=1`$, whereas in figure 4 $`\kappa _1=k=1`$, $`\kappa _2=0.8`$. The plots of figure 3 do not follow Eq. (3.4) nor (3.4) since $`2/(13\nu )`$ is numerically very close to one, and the its asymptotic behavior lies somewhere in between.
The relationship between the cosmic time $`t`$ and $`\eta `$, $`d\eta a^{3\nu }dt`$, seems to be too complicated to be integrated in a closed form. However, in the asymptotic limits given by Eqs. (3.4) and (3.4) one obtains that
$`t_{\mathrm{asimp1}}`$ $`=`$ $`\left(k\right)^{\frac{3\nu }{2(13\nu )}}(13\nu )^{\frac{1}{13\nu }}\eta ^{\frac{1}{13\nu }}`$ (19)
$`t_{\mathrm{asimp2}}`$ $`=`$ $`{\displaystyle \frac{2\left[\frac{2k}{\kappa _2}\right]^{\frac{3\nu }{2(13\nu )}}(13\nu )^{2\left(\frac{3\nu }{13\nu }\right)^2}}{29\nu (1+\nu )}}\eta ^{\frac{29\nu (1+\nu )}{2(13\nu )^2}}.`$ (20)
In both cases the functions are monotonic, and for $`\nu _1`$ and $`\nu _3`$ used in our plots, time grows as $`\eta `$ grows. Therefore, our time parametrization is appropriate. In this limit, one can express our solutions in $`\eta `$ in terms of $`t`$ to get that:
$`a_{\mathrm{asimp1}}`$ $`=`$ $`\sqrt{k}t`$
$`\varphi _{\mathrm{asimp1}}`$ $`=`$ $`m_\nu \left({\displaystyle \frac{13\nu }{1+3\nu }}\right)^2\kappa _2\left(k\right)^{\frac{3(1\nu )}{2(13\nu )}}=\mathrm{const}.`$ (21)
and
$`a_{\mathrm{asimp2}}`$ $`=`$ $`\left[{\displaystyle \frac{2k}{\kappa _2}}\right]^{\frac{13\nu }{29\nu (1+\nu )}}\left(13\nu \right)^{\frac{2+3\nu (83\nu )(13\nu )}{2(13\nu )^4}}\left(19\nu (1+\nu )/2\right)^{\frac{19\nu }{29\nu (1+\nu )}}t^{\frac{19\nu }{29\nu (1+\nu )}}`$
$`\varphi _{\mathrm{asimp2}}`$ $`=`$ $`{\displaystyle \frac{2^{\frac{(1+3\nu )(2+3\nu )}{29\nu (1+\nu )}}m_\nu }{(1+3\nu )^2}}\left[{\displaystyle \frac{\kappa _2}{k}}\right]^{\frac{3(1+\nu )(13\nu )}{29\nu (1+\nu )}}\left(13\nu \right)^{\frac{421\nu 27\nu ^2}{29\nu (1+\nu )}}`$ (22)
$`\left(29\nu (1+\nu )\right)^{\frac{(1+3\nu )^2}{29\nu (1+\nu )}}t^{\frac{(1+3\nu )^2}{29\nu (1+\nu )}}.`$
Eq. (3.4) is a particular known flat solution in BD theory or the $`k=1`$ vacuum solution of GR. Eq. (3.4) is the flat space, Nariai solution in which $`\omega `$ is given by Eq. (3.4).
## 4 CONCLUSIONS
We have presented a set of differential equations written in rescaled variables that let us integrate a general equation to get Eq. (2), valid for Bianchi type I, V and IX models, as well as for FRW models. This equation is coupled to Eqs. (1) and (6), but for Bianchi type I and V models the anisotropic parameters ($`h_i`$’s) are constants (since their curvature terms are of FRW type) and Eq. (2) is decoupled. This property allows one to find the general solution of Bianchi type I model and a particular solution of type V; both solutions are quadratic polynomials. This solution let the models isotropize as time evolves, however, this can happens only for some parameter ($`\omega `$, $`\nu `$) range. The polynomial solution may also be valid for Bianchi type IX, but it is not proved yet. If it were, isotropization would be also guaranteed.
We have found a new solution of Eq. (2) valid for curved ($`k0`$) FRW cosmologies, that is, a solution with $`h_i=0`$. This is the general solution subject to the constraint given by Eq. (3.4). Accordingly, we have analyzed two cases of physical interest: the case when $`\omega =1`$, having some interest in string cosmology, that implies an equation of state of a quasi dust model ($`\nu {}_{}{}^{>}\mathrm{\hspace{0.17em}0}`$), and the case when $`\omega =500`$, consistent with current BD local experimental constraints , implying that $`\nu 1/3`$. Although in a different manner, both the scale factor and $`\varphi `$ grow as a function of the parametrized time $`\eta `$ in all figures presented here. One can find specific values of the constants $`\nu ,\kappa _2`$ to have ”closed” models without recollapse. For this to happen $`2/(13\nu )`$ must be less than 1, i.e., $`\nu `$ must be negative. This is a known effect of negative pressures.
The new solution is non-inflationary and for asymptotic times is of power-law type for the scale factor. On the one hand, when the quantity $`2/(13\nu )`$ is greater than 1 the BD field tends to a constant value, then BD’s dynamics evolves similar to GR’s. One the other hand, when $`2/(13\nu )`$ is smaller than 1, both scale factor and BD field behave asymptotically with a power-law, and the solution is equal to Nariai solution for flat space.
Further solutions of Eq. (2) are in order, which can be either within FRW cosmologies or Bianchi type V or IX models. In particular, the general solution with $`h_i`$’s constants should provide the general solution of both curved FRW cosmologies and Bianchi type V model.
Finally, our results could be also of interest for $`\omega (\varphi )`$-theories, where the value of the coupling parameter in some early cosmological era could have been rather different than its value nowadays, $`\omega 500`$.
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# X-ray observations of the starburst galaxy NGC 253:
## 1 Introduction
A hot gaseous component in the interstellar medium with a temperature around 10<sup>6</sup> K is expected to originate from SNRs in the disk of spiral galaxies. The medium might emerge via galactic fountains into the halo of these galaxies (e.g. Spitzer 1956; Cox & Smith 1974; Bregman 1980a,b; Corbelli & Salpeter 1988). Supernovae and winds from massive stars in a central starburst might even drive a large-scale outflow that can shock heat and accelerate ambient interstellar or circumgalactic gas in form of a galactic super-wind (e.g. Heckman et al. 1990). This component (e.g. Cox & Reynolds 1987) was detected by Snowden et al. (1994) in the plane and halo of the Milky Way during the all sky survey and in pointed observations performed by the Röntgen observatory satellite (ROSAT). Hot interstellar gas in the LMC originally detected by observations of the Einstein observatory (e.g. Wang et al. 1991) was confirmed in the ‘first light’ observations with ROSAT (Trümper et al. 1991). Outside the local group, Einstein upper limits to the diffuse emission from hot gas were determined for edge-on galaxies (Bregman & Glassgold 1982) and for the large face-on galaxy M101 (McCammon & Sanders 1984). For the starburst galaxies M82, NGC 253, and NGC 3628 however, Watson et al. (1984), Fabbiano & Trinchieri (1984), Fabbiano (1988) and Fabbiano et al. (1990) resolved extended emission that was attributed to gaseous clouds ejected from the starburst nuclei, with temperatures in the $`10^6`$ K range.
Deep PSPC and HRI observations of the prototypical starburst galaxy NGC 253, due to its low Galactic foreground absorption and big optical extent, are ideal exploiting the ROSAT virtues (Trümper 1983, Aschenbach 1988, Pfeffermann et al. 1987). The low Galactic foreground absorption (N$`{}_{\mathrm{H}}{}^{}=\mathrm{1.3\hspace{0.17em}10}^{20}`$ cm<sup>-2</sup>, Dickey & Lockman 1990) allows soft X-rays from NGC 253 to reach the detector with little attenuation. Pietsch (1992) reported first results based on some early ROSAT PSPC observations. Read et al. (1997) presented X-ray parameters of NGC 253 in a homogeneous analysis of archival ROSAT PSPC data of nearby spiral galaxies, and Dahlem et al. (1998) investigated archival ROSAT and ASCA data in an X-ray mini-survey of nearby edge-on starburst galaxies including NGC 253.
In this paper we present a detailed analysis of the diffuse X-ray emission of NGC 253, characterizing in detail the different emission components of the area close to the nucleus (“nuclear source” and “X-ray plume”), of the disk, and the halo hemispheres. The point source contributions were separated with the help of the point source catalog presented by Vogler & Pietsch (1999, Paper I). Galaxy parameters used throughout this paper are summarized in Table 1. Our analysis is restricted to ROSAT data. Results of ASCA or BeppoSAX are hampered by the limited spatial resolution of these instruments which does not allow the separation of the very different surface brightness components, present in NGC 253. We compare our findings with results from other X-ray investigations (including ASCA and BeppoSAX) and with other wavelengths and discuss the emission components in view of starburst driven super-wind models.
## 2 Observations and point source analysis
NGC 253 was observed with the ROSAT HRI and PSPC for 57.7 ks and 22.8 ks, respectively. For the analysis all available ROSAT observations of the PSPC and HRI have been merged. Details of the observations and methods used to derive the X-ray point source catalog are described in Paper I. The analysis of extended emission components is based on the corrected event files and images created for Paper I’ and on additional procedures using the ESO-MIDAS/EXSAS (ESO-MIDAS 1997, Zimmermann et al. 1997) software package. While the higher resolution of the HRI detector allows us to resolve the bright nuclear and X-ray plume area, the lower instrumental background, the energy resolution and the high collecting area in the 0.1–0.4 keV band makes the PSPC best suited for studying large-scale diffuse emission in the disk and halo of the galaxy.
Separating diffuse emission components from point source contributions is an ambitious task keeping in mind the limited spatial resolution of the ROSAT PSPC detector. To lose as little area as possible, point sources contained in the source catalog of Paper I were cut out with radii optimized to the point spread function (PSF) of the telescope/detector at the relevant energy (see Table 2). Details on rejected sources and corresponding cut radii are mentioned below when describing the methods used to characterize diffuse emission. To keep the uncertainty low, the derived diffuse fluxes and luminosities are not extrapolated for the full area in the standard procedure. To improve the estimates for the nuclear area we use the higher resolution data of the HRI (see Sect. 3.2.1).
The point source catalog of Paper I contains 73 sources in the NGC 253 field, 32 of which are associated with the disk of the galaxy. Though 27 of these sources are detected with the HRI (some being resolvable with the PSPC), the remaining 5 PSPC-only detected sources are most likely due to fluctuations within the diffuse X-ray emission (see Figs. 2 and 4 of Paper I). The area around the nucleus is resolved into a bright, mildly absorbed point source (X33, most likely a black hole X-ray binary) and a more highly absorbed, slightly extended nuclear source embedded in bright diffuse X-ray emission from the X-ray plume (already detected in Einstein observations, see Fabbiano & Trinchieri 1984, Fig. 2). The halo of NGC 253 is filled with diffuse, filamentary X-ray emission (cf. Fig. 1). While four of the sources detected within this region are most likely background active galactic nuclei (AGN, as suggested by time variability arguments and proposed optical identification), two PSPC-only detected sources (X24 and X27 of Paper I) are most likely spurious detections caused by local enhancements in the diffuse emission. X10, another PSPC-only detected extended source in the NW halo, might be a background cluster of galaxies, as it shows a harder spectrum than the surrounding diffuse emission.
## 3 Analysis of diffuse emission and results
Several methods are used to characterize the diffuse X-ray emission components of NGC 253. Iso-intensity contour maps exhibit the point sources and diffuse emission with high resolution for the HRI. The maps are split up in several energy bands, however with lower spatial resolution for the PSPC. The images in different PSPC energy bands can be combined to a “true” colour image of the galaxy. We further investigate the emission components with respect to spatial distribution (box profiles along major and minor axis and azimuthal profiles for different energy bands and the X-ray plume region) and spectral behavior (e.g. integrated spectra of nuclear component, X-ray plume, disk, and halo, temperature profiles of the halo). In the following subsections we describe the procedures used for the analysis and we present a first discussion of results.
### 3.1 Iso-intensity contour and ”true” colour maps
The central area is best imaged with a contour plot from the HRI data. We have used an image with a bin size of 1″, selecting only raw channels 2–8 to reduce the background caused by UV emission and cosmic rays, and smoothed it with a Gaussian filter of 5″ FWHM, corresponding to the on-axis HRI PSF (cf. Fig. 3 of Paper I and Fig. 7). The image shows two bright point-like sources, the central source X34, and X33, embedded in a complex bright diffuse emission structure that is elongated in NW–SE direction with a maximal width of 45″ (560 pc) protruding 60″ (750 pc) from the nuclear source X34 to the SE, and 40″ (500 pc) to the NW. While the source $``$ 20″ south of the nucleus (X33) is a time variable point source (FWHM $``$5″, see Paper I), the source close to the nucleus (X34) is extended (FWHM $``$23″ and $``$19″, i.e. 290 pc and 240 pc in east/west and north/south directions, respectively), with the intensity maximum offset by 8$`\stackrel{}{.}`$3 to the SE of the nucleus, possibly representing the hottest part of the gas out-flowing from the nuclear starburst region (cf. Paper I).
For the PSPC data, contour plots were obtained from images constructed by the superposition of sub-images with 5″ pixel size in the 8 standard bands (R1 to R8, cf. Snowden et al. 1994). All sub-images were corrected for exposure, vignetting, and dead time, and were smoothed with a Gaussian filter of a FWHM corresponding to the on-axis point spread function (PSF) of that particular energy band. The FWHM values used a range from 52″ to 24″. In Fig. 1 contour plots are overlaid on grey scale images of the X-ray emission for the broad (0.1–2.4 keV, upper left), soft (0.1–0.4 keV, upper right), hard1 (0.5–0.9 keV, lower left), and hard2 (0.9–2.0 keV, lower right) band, respectively.
Visual inspection of the sub-images reveals several components of diffuse emission:
* Emission outside of the disk is absent in the hard2 map, and only present in the NW hemisphere in the hard1.
* The soft emission in both halo hemispheres suggests a filamentary structure that cannot be fully resolved due to the limitations of the ROSAT PSPC PSF.
* The diffuse emission in the soft band in the SE shows a flat maximum (projected distance of about 45″ SE of the NGC 253 nucleus) and seems to floats on the disk like a spectacle-glass (diameter of 16′ $``$ 12 kpc), from which two “horns” protrude (separation at the connection with the diffuse disk emission about 5′ $``$ 3.7 kpc) reaching into the SE halo out to a projected distance of more than 7′ (5.2 kpc), being turned inward. The horn emanating from the northern half of the disk is more pronounced.
* The soft emission in the NW (extent 12′ $``$ 9 kpc parallel to the disk and 8$`\stackrel{}{.}`$8 $``$ 6.6 kpc in perpendicular direction) is separated from the emission in the SE by $``$ 1$`\stackrel{}{.}`$2 $``$ 900 pc and also exhibits an inward turned horn-like structure. Again, the horn emanating from the northern half of the disk is more pronounced.
* The emission in the hard bands shows – besides point sources – a strong extended filamentary component from the inner disk that protrudes from a hardly resolved bright nucleus and traces the inner spiral arms. In addition, there is a less structured diffuse component of the disk. From inspection of the sub-bands R6 (0.9–1.3 keV) and R7 (1.3–2.0 keV) we find that the diffuse components vanish in the R7 band.
* Both hard bands reveal emission protruding into the NW halo along the minor axis. While in the hard2 band this component can only be traced for about 1′ $``$ 750 pc (mainly in the area with highly reduced emission in the soft band), it extends along the northern horn as far as the soft emission in the hard1 band.
This complex structure manifests itself in the true colour image (Fig. 2) built up from the PSPC soft (red), hard1 (green) and hard2 (blue) band images displayed in Fig. 1. The output image is made up of “pointers” to a special colour look-up table, representing 216 colour blends spanning the possible combinations of 6 intensity levels for each of the basic colours red, green and blue. With this method soft sources appear in red, with a “red” intensity proportional to the intensity in the soft image, harder sources in green and blue. The point sources in the disk area are hard, the greenish colour of the diffuse disk emission suggests a hotter or more strongly absorbed gaseous component compared to that in the halo. The NW rim of the disk is bluer than the SE part which is most probably due to the fact that in the NW one mainly detects hot gas absorbed in the disk while the emission in the SE is a superposition of the softer halo emission in the near halo hemisphere and the disk component from the SE. The greenish colour connecting to the red NW halo again suggests emission from highly absorbed hot gas. Due to the cut values used, the nuclear area is over-exposed.
### 3.2 Separating diffuse emission components
To quantify the findings of the last subsection, we now derive spatial distributions, spectra and luminosities for the individual diffuse components. Simple power law (POWL), thermal bremsstrahlung (THBR), or thin thermal plasma (THPL) models (Raymond & Smith 1977 and updates) and combinations are fitted to the different components as detailed below. Fits using MEKAL thin thermal plasma models (Mewe et al. 1985 and updates) did not give significantly different results. For the thin thermal plasma model we assume cosmic abundances as the starburst should have enriched the interstellar material with processed material. The statistics of the data and the restricted energy range of the ROSAT PSPC prohibit more detailed modeling.
#### 3.2.1 Diffuse emission from nuclear area and X-ray plume
As demonstrated in Paper I, the nuclear source (X34) peaks at $`\alpha (2000.0)=00^h47^m33\stackrel{s}{.}40,\delta (2000.0)=25\mathrm{°}17^m23\stackrel{}{.}1`$, with a 90% error radius (including systematics) of $`2\stackrel{}{.}5`$, and that it is clearly extended. Outside a radius of $``$ 15″ there is a continuous transition to the diffuse emission along the minor axis (cf. Fig. 3 of Paper I). The angular distribution of the emission components resolved with the HRI with respect to the position of the nucleus can be viewed in the azimuthal plot of the surface brightness in the central area (radius 1′, Fig. 3), integrated over sectors of 5°. In order to avoid the bright emission, the nuclear source was cut out with a radius of 15″. Above a background level of $`25.5\pm 0.8`$ counts per sector (determined in sectors 0°–90° and 210°–300°) $`10^2`$ cts s<sup>-1</sup> arcmin<sup>-2</sup> , three excesses can be identified:
1. A bright peak ($`393\pm 24`$ integrated counts above background) at 195° represents the point source X33. The number of counts extracted for X33 by this method is consistent with the results reported in Sect. 3.2 of Paper I ($`359\pm 21`$ cts) derived with a different method.
2. A broad, slightly asymmetric double-peak structure (FWHM $``$75°, $`546\pm 34`$ excess counts in sector 100°–175°) centered at position angle 145° is due to the extended emission to the SE along the minor axis. The two peaks at position angles $``$120° (FWHM of $``$35°, $`238\pm 21`$ excess counts in sector 100°–135°) and $``$150° (FWHM of $``$45°, $`308\pm 24`$ excess counts in sector 140°–175°) may reflect a cone-like emission structure along the minor axis of NGC 253 (position angle 142°), with an opening angle of $``$30° originating from the galaxy nucleus.
3. A broad peak (FWHM $``$35°, $`128\pm 20`$ excess counts in sector 305°–345°) centered at position angle $``$325° is due to the extended emission to the NW along the minor axis and pointing directly opposite to component two. According to Fig. 3 of Paper I the emission is less extended than that originating from the nuclear area. This may be explained by higher absorption in this direction (see discussion below).
To further characterize the diffuse nuclear emission components, we made use of the spectral capabilities of the PSPC. This investigation, however, is hampered by the larger PSF of the PSPC compared to the HRI. Especially in the soft energy band (FWHM of PSF $``$ 40″) the emission components separated by the HRI are difficult to resolve,and therefore only crude spectral parameters and luminosities can be derived. A better estimate of the luminosities of the nuclear source (X34) and the X-ray plume can be obtained by converting count rates from the ROSAT HRI detector to fluxes using the PSPC-derived spectral parameters (see below). Detailed results for the point source X33 south of the nucleus were already discussed in Paper I.
Here we present source spectra for the extended nuclear source (X34) and the X-ray plume. The spectrum of X34 was extracted with a cut diameter of 40″. To suppress contributions from the nearby source X33, a 70° sector in the direction of the source was excluded from the integration. We have used both, the background from the field outside the galaxy, used for all other spectral investigations, and a locally defined background. The spectrum of the extended emission along the SE minor axis (X-ray plume) was extracted from a sector with opening angle 90° outside the nucleus excluding X33 (with a circle radius of 26″, corresponding to twice the expected FWHM of the PSF of a hard source). Table 3 summarizes the regions for source and background and gives net counts and hardness ratios HR1 and HR2 (see Paper I for definitions). The hardness ratios indicate a highly absorbed and rather steep spectrum for the nuclear source, while the X-ray plume is less absorbed and much softer (cf. Pietsch et al. 1998). We then used simple absorbed spectral models to fit the data in the different regions, as indicated in Table 4. Reduced $`\chi ^2`$, number of degrees of freedom (DOF), absorption column, photon index for power law and temperature for thermal spectra are given. Raw spectra were rebinned to obtain at least the signal to noise level per bin given in Col. 1 of Table 4. We give fluxes and luminosities outside the Galaxy (f$`{}_{}{}^{\mathrm{exgal}}{}_{\mathrm{x}}{}^{}`$ and L$`{}_{}{}^{\mathrm{exgal}}{}_{\mathrm{x}}{}^{}`$, measured N<sub>H</sub> minus Galactic N<sub>H</sub>), as well as intrinsic to the source (f$`{}_{}{}^{\mathrm{intr}}{}_{\mathrm{x}}{}^{}`$ and L$`{}_{}{}^{\mathrm{intr}}{}_{\mathrm{x}}{}^{}`$, i.e. absorption zero). Rough errors for the fluxes and luminosities in Table 4 are indicated by the statistical errors of the net counts (see Table 3). Additional uncertainties arise from the poor knowledge of the spectrum, as can be seen by comparing the fluxes derived for different models.
When the field background is used, simple power law (POWL), thermal bremsstrahlung (THBR), or thin thermal plasma (THPL) spectra with cosmic abundances (Raymond & Smith 1977 and updates) do not allow a proper description of the nuclear spectrum (X34). If, however, the spectrum is corrected for the local background, the first two models fit the nuclear spectrum quite well, and also a THPL approximation cannot be totally rejected. All models indicate high absorption (N<sub>H</sub> of a few times 10<sup>21</sup> to more than 10<sup>22</sup> cm<sup>-2</sup>) and rather soft spectra (photon index 2.9 or temperatures from 0.4 to 1.2 keV). This suggests that with the local background the soft extended emission above the nuclear area is properly subtracted. We have therefore attempted to take it into account by fitting a more complex model to the data derived using the field background. We have assumed a two-component model composed of a THBR spectrum with the temperature fixed to 1.2 keV (the value determined using local background) and a THPL spectrum using cosmic abundance with fixed Galactic absorption. This resulted in an acceptable fit. For the THBR component flux and absorption are consistent with the values for the local background spectrum, and for the THPL component we derive a temperature of 0.44 keV.
For the X-ray plume, POWL and THBR spectra yield $`\chi ^2/\nu 1.7`$ and therefore are hardly acceptable, while a THPL with cosmic abundance clearly has to be rejected ($`\chi ^2/\nu >4`$). Assuming a two-component model of a THBR with the temperature fixed to 1.2 keV - the value fitted for the nucleus using local background - and additionally a THPL spectrum with cosmic abundance, fixed to Galactic absorption, leads to an acceptable fit. The THBR component is much less absorbed than for the nucleus and the temperature of the THPL with $``$ 0.3 keV a bit cooler than for the nuclear spectrum.
As mentioned above, the spatial resolution of the ROSAT HRI allows us to separately estimate the contribution from X34, the X-ray plume, and the point source X33. From the numbers given in Paper I we determine ROSAT HRI count rates of ($`12.2\pm 0.5`$) and ($`27.4\pm 1.1`$) $`\mathrm{\hspace{0.17em}10}^3`$ cts s<sup>-1</sup> for the nuclear source (local background subtracted) and the X-ray plume, respectively. Using the X34\_local THBR model parameters from the PSPC, the count rates for the nuclear source give an intrinsic luminosity of $`\mathrm{1.16\hspace{0.17em}10}^{39}`$ erg s<sup>-1</sup>, slightly higher than the PSPC value. For the X-ray plume, the soft and hard component of the PSPC-based two-component model lead to contributions to the HRI count rates that are a factor of 1.3 higher for the THBR than for the THPL component. Also the soft component of the X34 spectrum - that is even hotter than the X-ray plume component - has to be added to the HRI flux. Therefore we estimate that the THBR and THPL component contribute with a ratio of 3 to 4 and we assume a THPL temperature of 0.4 keV (average of X34 and X-ray plume values). With these assumptions we calculate intrinsic luminosities of 4.6 and 4.0 $`\mathrm{\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup> for the THBR and THPL components of the X-ray plume spectrum, respectively. These values are about a factor of four above those derived with the PSPC. The difference is caused by the PSPC extraction strategy: We used only part of the PSF area of X34 for the spectrum to reject contributions for X33. Also the rejection radii for X33 and X34 strongly reduce the area for the X-ray plume spectrum.
However, we have to place a caveat here. It is clear from the structured HRI and energy-resolved PSPC images that more than two components contribute to the spectra. They originate from the nuclear area, X-ray plume, disk and halo and certainly are seen through different amounts of absorbing matter. Using a local background for the nuclear spectrum allows us to reduce the influence of the surrounding components. Nevertheless, the assumed spectral models will be a simplification. The fact that a thermal bremsstrahlung or power law spectrum fit best for the absorbed nuclear component and also for the absorbed component of the X-ray plume, should not be over-interpreted but may only reflect that several thin thermal components of differing temperature add up to the spectrum. This may be an artefact of the limited energy coverage and resolution of the ROSAT PSPC, and can hopefully be resolved with the next generation of X-ray instruments.
#### 3.2.2 Diffuse emission from disk and halo
To investigate the diffuse emission from the disk and halo we make use of the nearly homogeneous sensitivity of the PSPC for low surface brightness features in the area within the ring of the window support structure (21′ radius). We examine box profiles along the galaxy axes, azimuthal profiles, and spectra integrated over different areas. To suppress contributions from point sources we follow a different strategy in creating profiles for the individual energy bands and in extracting spectral files.
When creating profiles in the soft band, we cut out only sources outside the disk and X15 using a cut radius of $`1\times `$ FWHM of the PSF at 0.30 keV, since no source from the point source catalog coincides with a point source in the soft band in the disk area (within corrected D<sub>25</sub>). Furthermore, X24, X27, and X30 (see Fig. 2) in the NW halo are not cut out as they are most likely “spurious” sources (local enhancements in the diffuse emission, X24 and X27), or were not visible during the PSPC observations (X30).
In the hard band, we cut out sources both in and outside of the disk. Outside we followed the same procedure as for the soft band, inside we cut out all sources from the catalog, with the exception of the extended nuclear source (X33) and the PSPC only detected sources X50, X51, X52, X57, and X62 (see Fig. 2), most likely representing local enhancements in the diffuse emission in the disk. For the cut radius in the hard band we used $`1\times `$ FWHM of the PSF at 1.0 keV.
The spatial variation of the surface brightness distribution along the major axis of NGC 253 (Fig. 4) was derived by integrating counts in the soft and hard band in boxes of 30″ width along this axis. The boxes cover $`\pm 130\mathrm{}`$ perpendicular to the axis which corresponds to the maximum extent along the minor axis of the inclination corrected D<sub>25</sub> ellipse of NGC 253 (see Fig. 2). It is clear from Figs. 1 and 2 that by using this size of box perpendicular to the major axis of the edge-on galaxy, one only traces emission components from the disk and/or from the halo immediately above or below the disk. The neglected emission from the outer halo hemispheres is covered in the profile along the minor axis (boxes of 30″ width along the minor axis and the corrected D<sub>25</sub> of 18$`\stackrel{}{.}`$8 in the direction of the major axis, Fig. 4), and in profiles along axes parallel to the major axis but offset to the SE and NW such that the covered strips are adjacent to the major axis profile (boxes of $`30\mathrm{}\times 480\mathrm{}`$, Fig. 5). The background surface brightness is determined from the minor axis count rates at distances $`>12\mathrm{}`$. In this way, we utilize the high count statistics due to the biggest box size used for the different profiles and avoid a possible reduction of the extragalactic background due to absorption by an H i disk extending further along the major axis than the galaxy’s X-ray emission, as observed in some other galaxies (e.g. M101, Snowden & Pietsch 1995; NGC 55, Barber et al. 1996; NGC 4559, Vogler et al. 1997).
Along the major axis, the soft- and hard-band profiles show a bright nuclear component on top of a distribution that declines more or less symmetrically with distance from the nucleus (Fig. 4). In the hard band the distribution is composed of two distinct components with extents of 4$`\stackrel{}{.}`$5 and 10′ (3.4 kpc and 7.5 kpc) from the nucleus. By way of contrast the soft-band distribution is made up of just one component with an extent of $``$6′ (4.5 kpc), which can be accounted for by absorption due to interstellar material within the disk. This is why we detect more emission from a plane above the disk that is facing the line-of-sight (SE disk corona). The hard band profile is less effected by absorption and traces emission from within, above and below the disk.
Along the minor axis, the soft and hard band profiles differ drastically (Fig. 4). The hard-band profile is dominated by a bright core, positioned symmetrically about the nucleus that can be traced out to 4′ (3 kpc) to the SE, while to the NW is masked by a second component, with an exponential decline out to $``$12′ (9 kpc). In the soft band the maximum is shifted to the SE by $``$45″ relative to the hard maximum, and it is not symmetric. To the SE, the bright core component drops within 3′ (2.2 kpc) and then decreases more or less exponentially out to $``$12′ (9 kpc) from the nucleus. In the NW, the profile first drops within 1′ (750 pc) from the nucleus to just above the background level. Then, starting at a distance of 2′ (1.5 kpc), the emission gradually recovers to a maximum at $``$4′ (3 kpc), before it finally fades exponentially to background level at a similar distance from the nucleus as in the SE. Apart from the trough, the intensity in the NW is about 1.5 times that in the SE at similar offsets from the nucleus.
Parallel to the major axis to the NW and SE (Fig. 5), the soft profiles show emission above the background out to projected distances from the nucleus of $``$6′ (4.5 kpc). The profiles are double-peaked (broad maxima at offsets from the nucleus of $``$3$`\stackrel{}{.}`$5 $``$ 2.6 kpc), reflecting the horn-like structure of the soft halo emission seen in Figs. 1 and 2. In the hard band, excess emission is clearly seen in the NW out to similar distances as in the soft band. In the SE hard profile there is only a slight excess at a projected distance of $``$2$`\stackrel{}{.}`$5 (1.9 kpc) from the nucleus to the NE, at the base of the corresponding horn-like structure in the images.
To further characterize the diffuse emission components from the disk and halo, we analyzed PSPC spectra, integrated over areas summarized in Table 3. In preparing the spectral files, we removed photons from source areas that were rejected from the soft and hard profiles (see above) using cut radii of $`1\times `$ FWHM of PSF at 0.3 or 1.0 keV (depending on wether the sources were cut out in both bands or just in the hard one). Results of fitting simple models are also given in Table 4.
For the disk spectrum, POWL and THBR models yield $`\chi ^2/\nu 1.4`$ and therefore are hardly acceptable, while a single-temperature THPL with cosmic abundance clearly has to be rejected ($`\chi ^2/\nu >5`$). However, a two-temperature cosmic abundance THPL model with the absorbing column of the softer component fixed to the Galactic foreground results in an acceptable fit. In this model, one would assign the unabsorbed 0.20 keV ($`\mathrm{2.3\hspace{0.17em}10}^6`$ K) component to the hot gas above the disk and the absorbed $``$0.7 keV ($`\mathrm{8\hspace{0.17em}10}^6`$ K) component to the hot interstellar medium within the disk. A even better approximation to reality would certainly require models with more temperatures and absorbing columns. However, due to the limited statistics and spatial and spectral resolution of the ROSAT PSPC data, more detailed modeling proves impossible.
For the NW halo spectrum, POWL and THBR model approximations result in barely acceptable fits not just from $`\chi ^2/\nu `$ but also from systematics in the residuals between models and data. A cosmic abundance THPL model approximation is not acceptable ($`\chi ^2/\nu =3.8`$). If restricted to the energy range 0.1 – 0.5 keV, a cosmic abundance THPL model with absorption fixed to the Galactic foreground and a temperature of 0.15 keV results in an acceptable fit (see Fig. 6). Acceptable fits in the full 0.1 – 2.4 keV range can be achieved by adding additional emission components such as a second (hotter) thin thermal gas component or Gaussian lines at around 0.8 and 1.1 keV to such a spectrum.
For the SE halo spectrum we obtain similar results. However, due to the poor photon statistics, the findings are less significant. Opposite to the NW halo, also a two temperature THPL model doesn’t give an acceptable fit. Without acceptable fits we cannot provide confidence regions for the fitted temperatures to compare them to the NW halo hemisphere. This problem can be overcome in using the X-ray hardness ratios as a crude information on the spectral shape. For the SE halo, HR1 is smaller than for the NW halo ($`\mathrm{\Delta }`$HR1 = -0.42$`\pm `$0.10) while the HR2 values coincide within the errors (see Table 3). This clearly indicates a significantly softer spectrum in the SE halo (see Fig.1 in Pietsch et al. 1998).
To make sure that the problems with fitting simple models to the halo spectra were not caused by the special selection of the halo regions we subdivided both halo areas into five boxes of 2′ width along the minor axis. Within the limits of the reduced statistics we obtained similar results. Specifically the residuals at 0.7 keV and above were present in all regions in the NW halo. This is consistent with the detection of hard-band emission from the entire NW halo region that was already reported following the brightness profile analysis above. No significant temperature change could be established within the individual halo hemispheres.
## 4 Discussion
In the following we will compare the diffuse emission components of NGC 253 presented in Sect. 3 to previously reported X-ray results for this galaxy. The individual diffuse emission components are compared to observations of NGC 253 at other wavelength regimes, as well as to results from other spiral galaxies. We also derive parameters for the dense interstellar material in the disk of NGC 253 from its apparent shadowing of X-ray emission in the NW halo. The findings are discussed in view of starburst and super-wind models.
### 4.1 X-ray luminosity of NGC 253 emission components
As described in Paper I and in the previous section, the complex X-ray emission of NGC 253 can be separated into contributions from point sources and diffuse emission. The diffuse emission originates from the nuclear area, the X-ray plume, the disk and from both halo hemispheres. In Table 5 we summarize the contributions of these components to the ROSAT PSPC count rate and give luminosities in the ROSAT band, derived from the best fitting spectra (see Table 4). For the halo emission we use parameters derived for the two temperature thin thermal plasma models. L$`{}_{}{}^{\mathrm{abs}}{}_{\mathrm{x}}{}^{}`$ is the absorbed luminosity as measured at the detector surface, L$`{}_{}{}^{\mathrm{exgal}}{}_{\mathrm{x}}{}^{}`$ the luminosity corrected for Galactic foreground absorption of 1.3 10<sup>20</sup> cm<sup>-2</sup>, and L$`{}_{}{}^{\mathrm{intr}}{}_{\mathrm{x}}{}^{}`$ the intrinsic luminosity assuming no absorption at all. While the results for the first two are rather robust against model uncertainties due to the low Galactic foreground absorption, the values for the intrinsic luminosity – especially for the highly absorbed components – are very sensitive to changes in the model parameters and therefore have to be taken with care.
It is evident from the overall count rate budget in Table 5, that to first order, all components contribute similar amounts. Due to the spatial resolution of the PSPC, it was possible to separately characterize each component’s spectrum. As discussed below, most of the components were already suggested in Einstein observations. However, they could not be investigated in detail due to lack of statistics, and spatial and spectral resolution of the Einstein IPC. In previous publications, ROSAT data for NGC 253 were analyzed in investigations of samples of spiral and nearby edge-on starburst galaxies. The special merits of ASCA and BeppoSAX observations of the galaxy are the improved spectral resolution and range to higher energies covered. However, due to the comparatively low spatial resolution, the missing response below 0.5 keV and the limited statistics, not all components identified above can be spectrally resolved. Within these limitations, the results of the other investigations support ours.
A detailed analysis of the Einstein HRI data of NGC 253 by Fabbiano & Trinchieri (1984) already revealed emission from several point sources in the galaxy as well as diffuse emission from the nucleus, along the minor axis to the SE and from the inner disk. Assuming a 5 keV thermal bremsstrahlung spectrum with Galactic foreground absorption for the nucleus and disk and a 0.5 keV thermal plasma spectrum for the X-ray plume, they derived luminosities in the 0.2–4 keV band (corrected to the NGC 253 distance of 2.58 Mpc used here) of (8, 3.3 and 10)$`\times 10^{38}`$ erg s<sup>-1</sup>, respectively. Taking into account the difference in the model spectra and energy band, these values are consistent with our results. The analysis of the Einstein IPC observations of the galaxy by Fabbiano (1988) demonstrated that the emission profile from the inner disk of NGC 253 along the major axis closely follows the radio continuum emission. In addition extended emission was found from the northern side of the galaxy and attributed to gaseous clouds ejected from the starburst nucleus (luminosity $`\mathrm{1.1\hspace{0.17em}10}^{39}`$erg s<sup>-1</sup>). From the southern halo no emission was detected with the Einstein instruments. This, however, is not surprising since the collecting area of the Einstein IPC was very low in the ROSAT PSPC soft band, where all emission for this component is detected. The limited spatial and spectral resolution as well as the lack of statistics hindered detailed spectral investigations with the Einstein IPC.
Earlier analysis of ROSAT data did not discern the different emission components of NGC 253 in greater detail and therefore the results were still preliminary. The galaxy was e.g. analyzed as part of a sample of nearby spiral galaxies observed with the ROSAT PSPC (Read et al. 1997). For the sample galaxies, one- ore two-component spectra were fitted to the emission as a whole and to the point sources and diffuse emission, individually. For NGC 253, the authors found an integral and diffuse luminosity of (8.1, 5.9)$`\times 10^{39}`$ erg s<sup>-1</sup> (0.1-2 keV, corrected for a distance of 2.58 Mpc), respectively. The fraction of diffuse emission of 74% for the luminosity escaping the galaxy compares well to $`80`$% quoted in Table 5. A one-component thermal plasma only poorly fits the diffuse emission, leading to an absorption compatible with the Galactic foreground, a temperature of 0.47 keV and heavy-element abundances of 0.02 solar. A two-component model, comprising a thin thermal plasma and an absorbed, hard (10 keV) unresolved source component, improved the fit to the diffuse component and indicated that most ($``$90%) of the ’diffuse’ emission is truly cool (0.39 keV), low-metallicity (0.08 solar), diffuse gas, while the rest could be attributed to highly absorbed, hard sources. However, the fit is still not good ($`\chi ^2/\nu 3.7`$), and the authors argued that “a much more complex model, beyond the scope of this work, may be necessary to explain the halo emission from NGC 253 (and, indeed, other starbursts), as a large temperature gradient is believed to exist within the halo of NGC 253”. In our detailed analysis, we could not establish a temperature gradient within the halo as postulated by Read et al.. However we found differing temperatures for the diffuse emission components from the disk, the region immediately above the disk, and the individual halo hemispheres.
Ptak et al. (1997) report on the complex X-ray spectrum of NGC 253 as measured with the ASCA instruments in the energy range 0.5–10 keV, which shows strong emission lines from O, Ne, Fe, Mg, S, and Si above the continuum. Unfortunately, with ASCA it is not possible to spatially resolve point sources from diffuse emission or the different diffuse emission components. The integral spectrum can be fitted by two components, a “soft” component described by a temperature of 0.8 keV and an absorbed “hard” component with a photon index of 2.0 or a temperature of 7 keV. They find that different models (with different continua) yield absolute abundances that differ by more than an order of magnitude, while relative abundances are more robust and suggest an under-abundance of Fe (inferred from the Fe-L complex) relative to $`\alpha `$-burning elements. The authors also try to derive element abundances from the individual line intensities and argue that for the hard component, they have to be significantly sub-solar (if thermal), or that there is a significant non-thermal or non-equilibrium contribution. The ASCA spectral fit is confirmed by BeppoSAX observations that, similar to ASCA, do not spatially resolve the components, though were, for the first time, able to detect the Fe K line at 6.7 keV (Persic et al. 1998).
ROSAT and ASCA observations of NGC 253 are also included in an X-ray mini-survey of nearby edge-on starburst galaxies (Dahlem et al. 1998). A ROSAT HRI image of the central area of NGC 253 is superimposed onto an H$`\alpha `$ image (see Pietsch (1994) for an overlay onto an optical one), PSPC images in three energy bands, with and without point sources, indicate diffuse emission up to the highest ROSAT band. Spectral results include a joint ROSAT and ASCA spectral study, and an investigation of individual areas with the ROSAT PSPC. Detailed discussions are deferred to a specific paper on NGC 253. In modeling the joint ROSAT ASCA spectrum, integrated over the galaxy as a whole, Dahlem et al. find that the difficulty of measuring multiple absorbing columns causes the largest uncertainty. As an example, they point out that it is impossible to measure the absolute or relative abundances or Fe with the integral spectrum, because N<sub>H</sub> (which is at least a few times 10<sup>21</sup> cm<sup>-2</sup> in the direction of the core) trades off directly against the Fe abundance – stronger absorption for Fe L energies than for Si and S lines at higher energies. They therefore conclude that for NGC 253 the absolute and relative abundances derived from the integrated ASCA and PSPC X-ray spectrum are not reliable indicators of the physical properties of the gas. They also briefly discuss spectral modeling of compact sources, core, disk and halo emission. However, the exact procedure is difficult to reproduce. For the core, they do not separate the components of X33, X34, and the X-ray plume, and therefore, the results cannot be directly compared. For the spectrum of the diffuse emission of the disk after point source subtraction, they reduce a local background, and they achieve an acceptable fit for a model that consists of an absorbed power law and a thin thermal (0.25 keV) plasma with solar abundance (which corresponds to our two-component fit). The halo spectrum and three concentric sub-spectra are accumulated, averaging over both halo hemispheres. Their halo spectrum is well fitted with a single temperature thin thermal plasma model with very low abundance. However, they prefer a two-temperature plasma, with temperatures of 0.14 and 0.65 keV and solar abundance with a flux ratio of 4:1 and find no temperature dependence over the halo. The hard spectral component is present in all three halo regions and not just close to the disk. These results are confirmed by our findings for the NW halo. Due to the higher count rate, this hemisphere dominates the averaged spectra and probably masks in their analysis the significantly differing intensity and shape, we found in our analysis of the halo hemispheres.
### 4.2 Emission from the area of the starburst nucleus and X-ray plume of NGC 253
Following to Sect. 3.2.1 the X-ray emission from the nuclear region of NGC 253 can be separated into at least three components:
* the bright point source X33 about 25″ south of the nucleus, most likely a moderately absorbed black hole X-ray binary within NGC 253 discussed in detail in Paper I
* highly absorbed slightly extended emission SE of the nucleus (X34)
* extended cone-like emission along the minor axis mainly to the SE of the nucleus – the X-ray plume
To compare the nuclear emission components to the structures that are thought to trace the nuclear outflow of NGC 253, we superimposed in Fig. 7 the HRI contours of the central area (contours as in Fig. 3 of Paper I) onto a continuum-subtracted H<sub>α</sub> \+ \[N ii $`\lambda \lambda `$6548, 6583 emission line image kindly provided by H. Schulz (cf. Schulz & Wegner 1992). X34 coincides with the central H$`\alpha `$ peak, and the cone-like diffuse emission covers the ”H$`\alpha `$ fan”. Based on Einstein HRI data with lower statistical significance this fact was already noted by McCarthy et al. (1987). They also pointed out that the H$`\alpha `$ bright emission regions to the SE, near the end of the X-ray emission, are most likely ordinary H ii regions in the disk of NGC 253. In the following we discuss the diffuse central X-ray components individually.
#### 4.2.1 Source X34: extended emission from the starburst nucleus of NGC 253?
While Fabbiano & Trinchieri (1984), using Einstein HRI observations, identified the extended source with the nucleus, the ROSAT HRI observations (cf. Paper I) demonstrate that it is offset by 5$`\stackrel{}{.}`$4 to the east and by 6$`\stackrel{}{.}`$4 to the south from the position of the nucleus as defined by radio observations (Ulvestad & Antonucci 1997, $`\alpha (2000.0)=00^h47^m33\stackrel{s}{.}10,\delta (2000.0)=25\mathrm{°}17^m17\stackrel{}{.}4`$). This total offset of 8$`\stackrel{}{.}`$3 to the SE clearly exceeds the uncertainty in the X-ray position of 2$`\stackrel{}{.}`$5, including statistical and systematic errors. The maximum of the emission of X34 is also clearly separated from the radio-bright SNRs and H ii regions (e.g. Ulvestad & Antonucci 1997), which are centered on the nucleus, and from the bright near-infrared emission originating from dense dust clouds and molecular material in the same region (Sams et al. 1994). We therefore conclude that X34 is not emission originating from the nucleus of NGC 253 but represents the position, where emission from gas, ejected along the minor axis, can penetrate into our line of sight through the dense absorbing interstellar medium surrounding the nucleus.
The spectral results for X34 support this view, too. The N<sub>H</sub> values of $`(34)\times 10^{21}`$ cm<sup>-2</sup> derived for the X34 spectrum (subtracting a local background) are more than a factor of 10 below column densities of $`\mathrm{6.3\hspace{0.17em}10}^{22}`$ cm<sup>-2</sup> expected for emission from the NGC 253 nucleus. Such a high absorption is put forward by the visual extinction of A$`{}_{\mathrm{v}}{}^{}>35`$ mag estimated for the nucleus by Prada et al. (1999) based on Br$`\gamma `$ velocity curves along the SW side of the major axis, and the conversion to N<sub>H</sub> of $`\mathrm{1.79\hspace{0.17em}10}^{21}\times `$A<sub>v</sub> cm<sup>-2</sup> (Predehl & Schmitt 1995). A source with a power-law spectrum with photon index 2.9 as measured for X34 would be suppressed in the ROSAT 0.1–2.4 keV band by a factor of $`>5000`$ compared to zero absorption, and still by a factor of $`70`$ compared to the N<sub>H</sub> derived from the X34 fit. If the spectrum of X34 was intrinsically flatter (e.g. power-law of photon index 1.9), the attenuation in the ROSAT band due to absorption might be reduced by just a factor of ten, but would still be substantial. Therefore, if hot gas is ejected from the nuclear area along the minor axis of the galaxy as a super-wind into the halo hemisphere facing us, soft band X-ray emission from this component will be heavily absorbed close to the nucleus and should become less absorbed along its path through the interstellar medium.
The X-ray absorption is usually characterized by the hydrogen column density N<sub>H</sub> made up mostly of H i and H<sub>2</sub>. To calculate the absorption as function of X-ray energy one then uses the effective absorption cross section per hydrogen atom as e.g. tabulated by Morrison & McCammon (1983) based on atomic cross sections and cosmic abundance. These abundances are derived in the solar neighborhood, and even within the Galaxy there is evidence for an radial abundance gradient for elements heavier than helium (Shaver et al. 1983). The X-ray absorption above 0.5 keV is primarily produced by heavier elements, mostly by oxygen and iron and therefore dependent on the metal abundance of the interstellar medium which may significantly differ from solar in various regions of a starburst galaxy like NGC 253. Direct radio measurements of the H i column in the direction of the NGC 253 nucleus were hampered twofold, by the limiting beam resolution of $`>`$30″, and by H i absorption. Therefore, the H i column of $`\mathrm{1.5\hspace{0.17em}10}^{21}`$ cm<sup>-2</sup> for the nucleus of NGC 253 obtained by the interpolation of the H i maps of Puche et al. (1991) and Koribalski et al. (1995), can only be interpreted as a lower limit. Mauersberger et al. (1996) investigate the distribution of molecular gas in the direction of the NGC 253 nucleus and derive based on a conservatively low $`N(\mathrm{H}_2)/I_{\mathrm{CO}}`$ value, a maximum H<sub>2</sub> column density of $`\mathrm{4\hspace{0.17em}10}^{22}`$ cm<sup>-2</sup> which (if half in the front and half in the back of the nucleus) contributes this amount to the N<sub>H</sub> towards the nucleus. The N<sub>H</sub> determined in this way, nicely compares to the value based on the Br$`\gamma `$ velocity curves above. Fits to broad band nuclear X-ray spectra with good statistics that will be derived with XMM-Newton or Chandra, should offer another method to determine the absorption column and possibly even abundances of the interstellar medium towards the nucleus of NGC 253.
With decreasing X-ray energy, the maximum of the emission should be further separated from the position of the nucleus. This effect is indicated in the PSPC images of the different energy bands in Fig. 1. However, one needs observations with high spatial resolution (as good as or better than those provided by the ROSAT HRI) for several energy bands to trace the increasing absorption down to the nucleus and investigate its nature. As can be seen from the numbers given above, an active nucleus with an intrinsic X-ray luminosity of more than a few times $`10^{40}`$ erg s<sup>-1</sup>, would still have been undetectable with ROSAT.
The spatial displacement of X34 from the nucleus can be used to determine the height above the plane of the galaxy from which the emission originates. We assume a positional displacement along the minor axis and take into account the inclination of NGC 253 (78.5°, Pence 1980). An offset of $`8\stackrel{}{.}3\pm 2\stackrel{}{.}5`$ then transforms to a height of ($`100\pm 30`$) pc above the nucleus. The extended nuclear source ($`20`$″ FWHM) can be interpreted as hot gas being part of a super-wind from the nuclear region. The THBR fit indicates a temperature of $`\mathrm{1.3\hspace{0.17em}10}^7`$ K. Its extent of $``$ 250 pc is similar in size to the bright nuclear radio emission detected with the VLA at wavelengths ranging from 1.3 to 20 cm, where a large number of compact radio sources were revealed, embedded within the diffuse radio structure (16″ $``$ 200 pc along major axis, Ulvestad & Antonucci 1997, and references therein), and more extended by a factor of two than the bright nuclear near infrared emission (Sams et al. 1994).
ASCA and BeppoSAX observations (see Sect. 4.1) do not provide the resolution to spatially resolve the nucleus. However, if the hard component fitted to ASCA spectra is identified with emission from the nucleus, its intrinsic luminosity should not exceed $`\mathrm{5\hspace{0.17em}10}^{39}`$ erg s<sup>-1</sup> (Ptak et al. 1997). As discussed before, this assumption oversimplifies the situation since several sources with different absorptions contribute to the ASCA hard spectrum, the highly absorbed nuclear component representing just the one suffering the highest absorption. Its intrinsic luminosity may therefore be well in excess of the number given above and does not really restrict the limits derived from the ROSAT observations. The same arguments hold for the BeppoSAX observations, that show extended emission best modeled by a hot thermal plasma ($`kT6`$ keV) with the inclusion of an Fe K line and have been tentatively identified with a starburst-driven galactic super-wind (Persic et al. 1998). There are, however, also some similarities to the X-ray emission from the Galactic Ridge (Cappi et al. 1999).
The situation should change with the detectors on board the next generation of X-ray observatories, XMM-Newton and Chandra, that will allow observations at higher energies with good spatial resolution, strongly reducing the effects of absorption, and allowing the nucleus to become directly visible. With the help of X-ray variability studies it should then be possible to decide on the nature of the radio-bright nuclear source which, according to Ulvestad & Antonucci (1997), can still be either an AGN or a very compact supernova remnant. Recently, time variability was detected in the highly absorbed hard X-ray component of the nearby starburst galaxy M82 using ASCA data (Ptak & Griffiths 1999, Matsumoto & Tsuru 1999) which leads to the conclusion that M82 hosts a low luminosity AGN. M82 shows X-ray emission from a strong galactic wind and the halo (e.g. Strickland et al. 1997). If indeed such a low-luminosity AGN were also confirmed for NGC 253, the other nearby prototypical starburst galaxy with galactic wind and X-ray emission from the halo, this might strengthen the view of an intrinsic connection of these phenomena as proposed by Pietsch et al. (1998), based on observations of the starburst galaxy NGC 3079, a galaxy which also hosts an active nucleus and a pronounced X-ray halo.
Forbes et al. (2000) present new HST data of the central region of NGC 253 and compare it to other wavelengths. They find that the majority of optical/IR/mm sources are young star clusters which trace a $`50`$ pc ring that defines the inner edge of a cold gas torus. In X-rays they identify the extended nuclear ROSAT HRI emission with the nucleus and argue that not all of the X-ray emission can be associated with the AGN, suggested from the radio, nor with an ultra-luminous supernova. Instead, they associate the emission with the out-flowing super-wind and find the size scale consistent with the idea of collimation by the gas torus. This analysis in principle does conform with our above interpretation. The main difference is that by careful positioning (see Paper I) we demonstrate that the emission in the ROSAT band is not originating from the nucleus but from a height above the plane of NGC 253 where the extinction is already significantly reduced and transparent to radiation in the ROSAT band. At this height, the X-ray beam which may be collimated in the inner disk by the cold gas torus as suggested above, may already have widened by a factor $`4`$ as indicated from the extent reported in Sect. 3.1.
#### 4.2.2 X-ray plume along SE minor axis
Demoulin & Burbidge (1970) were the first to suggest from spectrographic observations of NGC 253 that gas might be flowing out of the nucleus and outside the equatorial plane. In 1978, Ulrich confirmed the outflow along the minor axis with the help of velocity curves derived from measurements of optical lines using slit spectra across the center of the galaxy with different orientations (to distances of 25″ from the nucleus). She explained the origin of the outflow in terms of a nuclear starburst. McCarthy et al. (1987) and Schulz & Wegner (1992) used long-slit spectra and narrow-band images and interpreted them as emission from the surface of a kiloparsec-sized outflow cone driven by the starburst wind. The cone opening angle of 65 ° and outflow speed along the cone of 339 km s<sup>-1</sup> given by Heckman et al. (1990), compares well with values determined by Schulz & Wegner (50° and 390 km s<sup>-1</sup>, respectively).
The overlay of the ROSAT HRI contours onto the H$`\alpha `$ image (Fig. 7) shows a rather close correspondence of the diffuse X-ray emission and the H$`\alpha `$ structures connected with the outflow. The X-ray emission can be traced to a similar distance ($``$ 1′) from the nucleus to the SE as the weak H$`\alpha `$ and shows a clear cone-like structure with an opening angle of $``$30°, and brightened limbs in the ROSAT band (see Sect. 3.2.1) originating from the nucleus, which is also observed in the H$`\alpha `$ images but less pronounced. The opening angle defined by the maxima in Fig. 3 is smaller than that based on optical radial velocities. This can be explained by the near edge on viewing geometry, and as expected, the width of $``$75° of the emission structure as determined from Fig. 3 in Sect. 3.2.1 compares more favorably. The cone-like structure can be understood in terms of models of galactic super-winds driven by a nuclear starburst (e.g. Tomisaka & Ikeuchi 1988, Suchkov et al. 1994). According to these models the cone surface may represent the interaction zone of the emerging wind with the interstellar medium surrounding the nucleus. In H$`\alpha `$ and X-rays we do not see the wind itself which is expected to be hotter than this medium and radiating predominantly at energies outside the ROSAT band. In this picture the relatively sharp cut-off of the X-ray emission to the SE can be understood as representing the scale height of the dense interstellar medium in the galactic disk.
According to the models, a nuclear starburst should produce a bipolar flow along the minor axis. On the other hand, circumnuclear and disk material – as discussed above – will suppress soft X-rays from the starburst wind emerging into the far, NW hemisphere. The HRI images still reveal some emission NW of the nucleus. The morphology, however, is not that of a hollow-cone as in the SE. Therefore, we cannot decide whether this emission leaks through from the back of NGC 253 through areas of reduced absorption (inter-arm region) in the galactic disk, or whether it originates in front of the disk. In principle one could try to settle this question using the spectral capabilities of the PSPC, as we did for the nuclear source and X-ray plume. Unfortunately, the emission is not strong and extended enough and is too close to the bright central sources to allow this kind of analysis.
However, the existence of a flow of gas into the NW halo hemisphere is demonstrated by the observation of an OH-plume pointing in this general direction, with a strong component to the north (Turner 1985). No OH emission (originating from molecular gas, radio 18 cm) is detected in the SE, adding to the arguments that the outflow into the NW halo must even be stronger than that to the SE. This might be explained by a less strong blocking to the NW caused by inhomogeneities in the interstellar medium or an asymmetric position of the starburst with respect to the nucleus. The intensities and spectra of the X-ray emission from the halo hemispheres (see below) further support this picture.
The extended source above the nucleus (X34) and the X-ray plume can be directly compared to the X-ray emission from the nuclear super-bubble in the edge-on starburst LINER/Seyfert 2 galaxy NGC 3079. There the central component extends over 1.7 kpc with a luminosity of $`\mathrm{7\hspace{0.17em}10}^{39}`$ erg s<sup>-1</sup> and is not only detected in H$`\alpha `$ images but also in the 20 cm radio continuum (see Pietsch et al. 1998). Similar to the NGC 253 X-ray plume, the central bright X-ray and H$`\alpha `$ emission is only seen in the east due to absorption in interstellar medium of this galaxy. Another example of X-ray sources that may reflect emission from hot gas in super-bubbles emerging from galactic nuclei are reported by Read et al. (1995) for the merging galaxy pair NGC 4038/9 – the Antennae. There, the sources are offset by $``$1 kpc from the nuclei, are best fitted by absorbed, low temperature plasma models indicating intrinsic luminosities of a $`(13)\times 10^{40}`$ erg s<sup>-1</sup>.
### 4.3 Diffuse emission from the NGC 253 disk and the flat SE halo component immediately above (coronal emission)
In Sect. 3.2.2, we identified at least three components to the diffuse emission of the NGC 253 disk and the flat SE halo using the surface brightness distributions along the major and minor axes (Fig. 4) that are seen to differ in spatial extent:
1. a bright component that originates from the nuclear region was discussed in the previous section and is connected to the extended source X34 and the X-ray plume
2. a soft component most likely originating from hot gas immediately above the layer of the dense interstellar medium of the disk as can be inferred from the shift of its maximum to the SE by $``$45″. It rises to a maximum height of 2.2 kpc above the galaxy plane with an extent of $`\pm `$4.5 kpc along the major axis (coronal emission above the disk)
3. in the hard band, a bright central and a fainter extended component (extents of $`\pm `$3.4 kpc and $`\pm `$7.5 kpc, respectively) becomes visible centered on the plane of the galaxy and most likely originating from within the disk
Components 2 and 3 can clearly be identified in the overlays of the soft and hard PSPC X-ray contours on an optical image of NGC 253 (Fig. 9). Component 2 may be further subdivided into emission connecting component 1 with the horn-like structure in the outer halo hemisphere (as indicated by the curvature of the soft band contours close to the nuclear area), reflecting the strong galactic wind emanating from the starburst nucleus, and a component floating on the disk like a spectacle-glass. The latter most likely originates from hot gas fueled from galactic fountains (see below). The proposed location and viewing geometry of the different halo components is sketched in Fig. 10.
The presence of these bright components hinders the detection of X-ray emission – if present – from the bar of NGC 253 (seen e.g. in infrared data with position angle 70°, extent $``$150″; Forbes & DePoy 1992). On the other hand, there is a clear correspondence in the brightness distribution of the diffuse hard band emission with the spiral structure in the inner and the outer disk of NGC 253 (Fig. 9). South of the nucleus the inner spiral arm is located to the north and the brighter part of the outer spiral arm to the south within the projected elliptical shape of the disk. North of the nucleus, the spiral arm moves outward from south to north (compare optical and H i analysis by Pence 1980 and Koribalski et al. 1995, respectively).
The X-ray surface brightness distributions (Fig. 4) can be compared to the emission at other wavelengths. The soft band profile originates from above the disk, therefore one expects little similarities with emission at optical continuum and infrared wavelengths. The hard band, however, traces diffuse emission in the disk, and correlations with other wavelengths are expected. Thus, for Einstein IPC data Fabbiano (1988) reported a good correlation with the radio profile along the major axis. In Fig. 8 we compare the background corrected ROSAT PSPC hard-band profiles along the major and minor axes to optical, infrared, and radio profiles.
The optical (blue light) profiles were extracted from Pence (1980, Fig. 8) and normalized to the wings of the major axis X-ray profile (7$`\stackrel{}{.}`$5 offset) and to the maximum of the minor axis profile. Along the major axis the optical brightness shows a bright plateau out to a galacto-centric distance of 10′ and coincides nicely with the faint extended X-ray component. The bright central and the nuclear component seem less pronounced or totally missing, which may be explained by extinction effects. Along the minor axis, X-ray and optical brightness profiles are comparable within the disc. However, the optical data do not indicate halo emission, as can be most clearly followed in the X-ray profiles in the NW.
Infrared profiles were published by Scoville et al. (1985). We have compared 2.2 $`\mu `$m data (their Fig. 1) obtained with a 10″ beam from scans along the major axis, and normalized their profile to match our bright central component ($``$2$`\stackrel{}{.}`$5 offset). The NIR profile follows the X-ray brightness of the bright central component and of the wings quite well. The larger X-ray extent to the SW of the bright central component may be explained by the differences in the infrared beam size and the X-ray slit length. As these systematic effects should be much stronger along the minor axis, we did not include these data in the figure. The excess 2.2 $`\mu `$m radiation, if compared with colours normally seen in disk galaxies, cannot be entirely explained by heavy extinction, but is most likely a contribution from clouds of hot dust associated with star formation regions.
While the radio 1.46 GHz radio emission (Hummel et al. 1984) emphasizes the disk emission, the 330 MHz emission (Carilli et al. 1992, see Fig. 11) discloses the true halo emission. As we are here comparing disk distributions, we compare the 1.46 GHz profile (Hummel et al. 1984) collected with a beam size of 68″$`\times `$36″. It is normalized to match the bright central component at the same offset as the infrared. A similar profile was reported by Klein et al. (1983) based on 10.7 GHz observations. While in 1.46 GHz, the nuclear component along the major axis is more extended than the X-ray or infrared emission, the central component is similar in width to the infrared. The extended wings, however, decline faster than X-rays and infrared. A detailed comparison of radio and X-ray properties and especially of the distribution along the minor axis will be given in a separate paper (Ehle et al., in preparation).
How can we explain the similarities and the differences in the different wavelength regimes? The strong emission from the central source in all wavelengths reflects the presence of the strong nuclear starburst. In addition, the bright inner disk emission in the radio, infrared and X-rays is most likely caused by enhanced star formation in this region. The star-forming regions, however, are located in H ii regions well within the disk and are most likely heavily obscured in the optical. This may explain why the bright central disk is not as obvious in this band compared to other wavelengths. However, in a detailed analysis of the morphology of dark lanes and filaments in the dust-rich disk, Sofue et al. (1994) identified arcs with heights of about 100 to 300 pc, connecting together two or more dark clouds, as well as loops and bubbles expanding into the disk-halo interface with diameters of a few hundred pc to $``$1 kpc, and vertical dust filaments, almost perpendicular to the galactic plane and extending almost coherently 1 to 2 kpc into the halo. Sofue et al. propose a boiling disk model, in which the filamentary structures develop due to star-forming activity in the disk combined with the influence of magnetic fields. There is clearly enhanced activity in the inner disk. This picture of the boiling disk, with indications of outflow of hot gas into the low halo (corona of the disk), connects naturally with the relatively soft X-ray emission floating on top of the disk of the galaxy in the SE like a spectacle-glass. In galaxies with starforming activity distributed over several areas in the disk, galactic fountains may be the mechanism to fuel the halo with hot gas (see e.g. the X-ray halo of NGC 891, Bregman & Pildis 1994). In NGC 253, however, the superwind from the nuclear starbursts clearly dominates.
Analyzing the X-ray spectra, a two-temperature thin thermal plasma model was sufficient to describe the integral disk spectrum. The low-temperature component (temperature 0.2 keV, extra-Galactic luminosity $`\mathrm{7.8\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup>) could be fitted assuming no absorption from within NGC 253, and can therefore be identified with the soft-band emission above the disk as deduced from the surface brightness profiles attributed to the disk corona. The component with the higher temperature of 0.7 keV is heavily absorbed ($`\mathrm{9.5\hspace{0.17em}10}^{21}`$ cm<sup>-2</sup> in addition to Galactic foreground) as expected if it originates from sources within the disk, and has an intrinsic luminosity of $`\mathrm{1.2\hspace{0.17em}10}^{39}`$ erg s<sup>-1</sup>. The spectrum certainly represents a mixture of sources. Unresolved point sources below the point source detection threshold (X-ray binaries, SNRs) will contribute, as well as diffuse emission from H ii regions and from the hot component of the interstellar medium. In principle, these components might be separated utilizing their characteristic spectra. However, due to the location in the disk, their spectra suffer from different levels of absorption so that observations with much better statistics and spatial and spectral resolution are required to resolve them.
Mean physical parameters for the flat SE halo component immediately above the disk can be inferred from the above results if we make some assumptions about the geometry of the emission. Here we choose the simple geometry of a disk corona, an circular cylinder with a thickness of 1 kpc and a diameter of 9 kpc with the central area cut out to a radius of 0.8 kpc. This approximates the extraction area of the disk spectrum and the extent of the soft emission along the major axis. The results are not strongly dependent on these geometrical parameters.
In this coronal model, the hot gas with the fitted temperature of 0.2 keV $`\mathrm{2.3\hspace{0.17em}10}^6`$ K (cooling coefficient of $`\mathrm{2.4\hspace{0.17em}10}^{22}`$ erg cm<sup>3</sup> s<sup>-1</sup> according to Raymond et al. (1976)) and the unabsorbed luminosity of $`\mathrm{7.8\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup> is distributed in a volume of $`\mathrm{1.8\hspace{0.17em}10}^{66}`$ cm<sup>3</sup>, assuming a relative volume filling factor $`\eta _1`$ of the hot gas in clouds. Assuming thermal cooling and ionization equilibrium (Nulsen et al. 1984), the electron density $`n_\mathrm{e}`$, mass $`m_{\mathrm{gas}}`$ and cooling time $`\tau `$ of the coronal gas are given in Table 6.
Neglecting effects of differing filling factors (which, however, would only contribute with the square root) the density of the soft and hard components in the outer halo hemispheres are lower by factors of 2 to 3 and more than 5, respectively, than that of the disk coronal hot gas. The contributing mass in the hard components of the halo is similar to that in the disk corona while it is higher by a factor of 2 to 3 in the soft halo components. The cooling time of the soft halo components is a factor of 2 and that of the hard components more than a factor of 10 longer than in the disc corona. The scale height and properties of the coronal emission can be explained by galactic fountains (e.g. Breitschwerdt & Komossa 1999 and references therein), created due to the heating of the hot intercloud medium by explosions of supernovae while the emission in the outer halo most likely originates from the galactic super-wind (e.g. Heckman et al. 1990, see Sect. 4.5).
### 4.4 Shadowing of halo emission by the disk
An eye-catching structure in the ROSAT PSPC soft band image (Figs. 1 and 9) and the surface brightness profile along the minor axis (Fig. 4) is the depression/gap in the diffuse emission between the galaxy disk and the NW halo hemisphere. The explanation for this novel effect is straight forward. The diffuse emission from the corona and outer halo hemisphere on the far side of the galaxy is shadowed by the dense interstellar medium of the intervening disk. In addition, in the gap area no emission from the other hemisphere contributes, due to the near edge-on view (see Fig. 10). Therefore this shadowing effect not only unequivocally reveals the geometry of the system (i.e. the NW edge is the near side of NGC 253), but also allows us to determine lower limits to the column density of the intervening material.
To estimate the effects of absorption, we could assume for the unabsorbed flux in the NW an identical profile as observed in the SE and calculate the amount of absorption from the reduction factor determined from the profile really measured. However, we already pointed out in Sect. 3.2.2 that the fluxes in the unabsorbed parts in the NW are $``$1.5 times brighter than in the SE at similar nuclear distances. Also, the spectrum in the outer halo is significantly harder. In Table 7 we put together reduction factors with respect to pure Galactic foreground absorption for the 0.1–0.4 keV and 0.5–2.0 keV ROSAT PSPC bands, calculated for thin thermal plasma spectra, with temperatures of 0.2, 0.3 and 0.4 keV and increasing additional shadowing columns of cold gas (N<sub>H</sub> of 2.5 to $`40\times 10^{20}`$ cm<sup>-2</sup>). While, for a given absorbing column, the reduction factors for soft band emission do not vary strongly for the temperature range investigated, for a fixed temperature there is strong variation with increasing absorption (factors of $``$1.2 and $`>1000`$, respectively). The corresponding hard band variations are stronger for changing temperatures, but much weaker with absorption (2 and 10, respectively).
The theoretically expected reduction factors can be compared to factors determined from the minor axis profiles (Fig. 4), by dividing count rates extrapolated from larger off-axis angles to the inner disk by the measured count rates (Table 8). At an off-axis angle of 4′ the factors still indicate absorbing columns of $`\mathrm{2\hspace{0.17em}10}^{20}`$ cm<sup>-2</sup> and at 3′ of $`\mathrm{4\hspace{0.17em}10}^{20}`$ cm<sup>-2</sup> using the soft-band information. The less reliable extrapolation of the hard band would indicate even higher absorption.
These results can be compared with H i measurements. Puche et al. (1991) present data obtained with a beam of 68″ diameter. Their lowest contour of $`\mathrm{2.4\hspace{0.17em}10}^{20}`$ cm<sup>-2</sup> is offset from the nucleus along the minor axis to the NW by 3′. With a beam diameter of 30″ the lowest contour of $`\mathrm{4.8\hspace{0.17em}10}^{20}`$ cm<sup>-2</sup> in the map of Koribalski et al. (1995) shows the same offset. As explained in detail in Sect. 4.2.1, N<sub>H</sub> is determined by contributions from H i and H<sub>2</sub>. The H<sub>2</sub> distribution of NGC 253 can be inferred from CO maps (Houghton et al. 1997). While these maps have less resolution than the H i maps, they more or less coincide in general extent and imply a significant contribution to N<sub>H</sub>. Using our X-ray absorption technique, we can determine absorption columns in the NW outer disk with similar accuracy. Unfortunately, the X-ray shadowing method is not generally applicable to measure the density of the interstellar medium in galaxies, but NGC 253 reflects a case of special luck. Galactic foreground N<sub>H</sub> needs to be rather low ($`\mathrm{1.5\hspace{0.17em}10}^{20}`$ cm<sup>-2</sup>), the galaxy inclination has to be close to edge on, and – the most important requirement – the galaxy halo must emit soft X-rays to allow this kind of study.
Even for our NGC 253 analysis one has to keep in mind that the count rates are averaged along the disk at varying absorption depths and X-ray brightnesses. Data with better statistics as expected from the upcoming XMM-Newton observatory will allow the determination of the absorbing column along the NW outer disk with good spatial resolution. Such data should also clearly uncover the geometry of the outflow to the NW. Its soft emission is hidden by the disk, however it shines through at energies above 0.5 keV (see the hard1 and hard2 images in Fig. 1), indicating that unfortunately the hemisphere with the more spectacular outflow is shadowed by the NGC 253 disk.
### 4.5 Extended soft X-ray halo of NGC 253
The detection of hot gas in the halo of NGC 253 extending to projected distances of 9 kpc from the galactic plain, allows us to investigate a component of the interstellar medium that cannot be accessed in optical observations. The X-ray structure (Sect. 3.2.2) clearly indicates that the emission is not filling the halo homogeneously. Instead, the soft image resembles a horn-like structure in both hemispheres, as would be expected if the emission regions were shaped like a hollow-cone. While the hemispheres look rather similar in the broad band image, they differ significantly in the hard band. This behavior is also reflected in the spectral parameters, the NW halo hemisphere being significantly hotter than the SE one. In addition to this difference between the two halo hemispheres, the northern most “horn” in both hemispheres is more pronounced.
Analogously to the corona above the disk (see Sect. 4.3), we can determine mean physical properties for the emitting medium. For simplicity we assume that the hot gas in both hemispheres is filling a cylinder of 4.5 kpc radius and 7 kpc height. It is evident from the images, that the filling factors $`\eta _2`$ and $`\eta _3`$ for the halo hemispheres differ. As discussed in Sect. 3.2.2, no simple model results in a good approximation for the individual halo spectra. Nevertheless, to derive halo gas parameters (cf. Table 6), we made use of the temperatures (0.19 keV $`\mathrm{2.2\hspace{0.17em}10}^6`$ K and 0.12 keV $`\mathrm{1.4\hspace{0.17em}10}^6`$ K) and corresponding unabsorbed luminosities ($`\mathrm{1.0\hspace{0.17em}10}^{39}`$ erg s<sup>-1</sup> and $`\mathrm{5\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup>), derived from a thin thermal plasma model fit to the NW and SE halo hemispheres, respectively. The significantly higher luminosity and temperature in the NW with respect to the SE halo hemisphere nearly cancel each other and only lead to slightly higher electron densities and masses, while the cooling time stays the same.
The ROSAT halo parameters in the NW are consistent with the parameters derived from Einstein observations (Fabbiano 1988), once corrected to our assumed distance and to the ROSAT measured temperature. The sum of our components differs, however, significantly from the parameters derived by Read et al. (1997) using integral properties of the diffuse emission of NGC 253; they find lower gas masses and shorter cooling times as well as allowing for lower electron densities in the radiating plasma.
If the hot gas in the outer halo represents nuclear and disk material transported into the halo by the super-wind with a velocity of $``$350 km s<sup>-1</sup> as modeled for the optical emission cone close to the nucleus (cf. Sect. 4.2.2), we can calculate from the halo extent of 9 kpc a lower limit for the age of the wind of $`\mathrm{2.5\hspace{0.17em}10}^7`$ y. This time is comparable to the age of the rapid massive star formation in the nucleus of NGC 253, derived from evolutionary synthesis models ($`23\times 10^7`$ y, e.g. Engelbracht et al. 1998), it is shorter, however, by at least an order of magnitude than the cooling time of the hot gas in the halo calculated above. From these considerations it is possible that the hot halo gas was heated in the starburst nucleus and transported into the halo by the wind. However, one could also consider that the wind is heating ambient gas in the halo.
The first scenario has been proposed by Tomisaka & Ikeuchi (1988). Their model of a galactic scale bipolar flow assumes a disk component to represent the interstellar medium of the starburst galaxy, which for some time confines the hot matter within a cool shell, which is further heated by supernova explosions of massive stars formed in the nuclear starburst. Eventually, the shell breaks up and releases the hot gas into the halo. Suchkov et al. (1994) improved the model by including an ambient two-component disk-halo interstellar medium and argued that this two-component representation is crucial for adequate modeling of starbursts. In their models, the bulk of the soft thermal X-ray emission from starbursts arises in the wind-shocked material of the halo and disk gas transported into the halo, rather than in the hot wind material itself. The super-wind itself is too hot and thin, to be visible in the ROSAT band and reaches speeds of $``$2000 km s<sup>-1</sup>. The models allow to produce X-ray geometries and filamentary structures similar to those observed in NGC 253 and propose a temperature distribution of the heated material of $`25\times 10^6`$ K. However, they neglect the role of cosmic rays, magnetic fields, and delayed recombination (Breitschwerdt & Komossa 1999).
The spectra from the halo hemispheres cannot be fitted by simple thin thermal plasma models. However, two temperature models or a one temperature model with additional emission components above 0.7 keV give reasonable fits to the data (Sect. 3.2.2). This confirms on a smaller scale results by e.g. Dahlem et al. 1998 and Strickland & Stevens 1998 who found observational evidence that the spectra of the diffuse soft X-ray emission in galactic winds or superbubbles suggest more than a single temperature. The temperatures and luminosities of the observed NGC 253 halo components are consistent with luminosities and the temperature distribution proposed by the models for the heated material in the halo.
Breitschwerdt & Schmutzler (1999) showed that the assumption of cooling via collisional ionization equilibrium used in the thin thermal plasma models are not correct if the dynamical time scale of the plasma, as for instance in an out-flowing wind, is much shorter than the intrinsic time scales (e.g. recombination, collisional excitation, ionization etc.). Therefore, the dynamical and thermal evolution of the plasma has to be treated self-consistently. First results show that the expected spectra of cooling wind material in the halo should exhibit enhanced emission in lines at energies above $``$0.7 keV (e.g. from O vii, O viii, and highly ionized Fe), compared to equilibrium models, as well as a lower plasma temperature due to adiabatic expansion. Since recombination of highly ionized species is delayed, the resulting X-ray line emission reflects the “initial temperature” of the plasma close to the disk, when it was much hotter. In these non-equilibrium models one cannot determine a single cooling timescale, “equilibrium cooling” due to expansion is likely to be much shorter than the equilibrium value given in Table 6 and the time scale for line recombination longer. Numerical simulations for a galactic outflow including non-equilibrium X-ray emission in a self-consistent fashion have been folded through the ROSAT PSPC instrumental response (Breitschwerdt & Freyberg 2000). Best fitting spectral models are two temperature thin thermal plasma or a one temperature model with an additonal Gaussian component similar to our finding. Breitschwerdt & Freyberg argue that the two fitted temperatures are not physical but mimic non-equilibrium X-ray emission.
Keeping in mind the super-wind origin of the halo emission, one can consider several explanations for the asymmetries in the halo X-ray emission. The asymmetry may either be caused by differing wind parameters in both hemispheres, originating for instance from a slight asymmetry of the driving starburst, or by a different structure of the ambient halo medium MacLow & McCray 1988). These explanations may well explain the asymmetry between the SE and NW halo, they also may account for the stronger “horns” in the north than in the south of both hemispheres. One can, however, think of other possibilities to generate the asymmetry of the “horns”. There could be more ambient material on the N side of the galaxy to shock-exite (perhaps due to an earlier fountain episode?), or even the interaction of the halo gas due to the relative motion with the intergalactic medium of the Sculptor group of galaxies, of which NGC 253 is a member (e.g. Puche & Carignan 1988), could be responsible. Unfortunately, neither the inter-galaxy medium nor the relative motions are known for the Sculptor group galaxies to strengthen or reject the latter hypothesis.
The X-ray halo has no equivalent in the optical continuum. The optical continuum halo mainly follows the elliptical shape of the galactic disk to distances along the major axis of 22.5 kpc (Pence 1980, Beck et al. 1982) and is most likely made up of late type stars. This situation is different for the H$`\alpha `$ emission.
Radio images at 330 MHz (Carilli et al. 1992) show a synchrotron emitting halo which is more extended perpendicular to the disk and in general matches the soft X-ray halo (cf. Fig. 11). At radio frequencies of 1.5 GHz, the emission is box-shaped similar to the disk corona emission (Hummel et al. 1984, Carilli et al. 1992, Beck et al. 1994). Carilli et al. present evidence for outflow from the disk in a ”spur” which – corrected for the NGC 253 distance assumed in this paper – rises 2.5 kpc above the plane at about the position of the NE X-ray horn. Detailed comparison of the morphology of the halo emission is hampered by the poorer resolution of the 0.33 GHz radio (60″ and outermost contour 120″) with respect to the soft X-ray contours (40″). However, there is indication from the overlay that the radio emission does not show the horn-like structure which dominates the X-ray morphology. So in X-rays, the curvature of the NE horn points at a connection to the nucleus, while the NE radio spur seems to be directly connected to the disk and offset to the north. Radio and X-ray data with better resolution and statistics are urgently needed to clarify these differences, that can not be understood in terms of models put forward for the common origin of the X-ray and radio halo emission. For the NGC 253 radio halo, Beck et al. (1994) extracted magnetic fields and rotation measures using multi-frequency observations of the radio continuum emission. X-ray, radio and H$`\alpha `$ measurements will be compared in detail in a separate paper (Ehle et al., in preparation) to determine the importance of magnetic and thermal effects for the interstellar medium in the galaxy disk and halo.
### 4.6 Comparison with diffuse X-ray emission detected in other spiral galaxies
With the help of the ROSAT satellite, diffuse emission from the disk and halo of several late-type spiral galaxies was separated, irrespective of their face-on or edge-on orientation. References for some galaxies analyzed in detail and used for comparison to our NGC 253 results have already been given or are given below. In addition, Read et al. (1997) present a homogeneous, however less deep, analysis of 17 nearby galaxies partly overlapping with our sample, and determine rough parameters for the diffuse emission. Detailed analysis of close to face-on galaxies (e.g. M 101, Snowden & Pietsch 1995; NGC 4258, Pietsch et al. 1994; NGC 4449, Vogler & Pietsch 1997; M 83, Ehle et al. 1998; NGC 7793, Read & Pietsch 1999) clearly show diffuse emission. The separation in disk and halo components proved difficult and was mainly based on the absorbing columns derived from spectral fits. The separation is much easier for (close to) edge-on galaxies due to the fact that in these systems disk and extended halo components spatially separate. Diffuse emission has not been detected from the halo of all “normal” edge-on galaxies. We list some examples with halo detection indicated in brackets: NGC 4565 (+), NGC 4656 (?), NGC 5907 (-) (Vogler et al. 1996); NGC 4631 (+) (Wang et al. 1995, Vogler & Pietsch 1996); NGC 4559 (-) (Vogler et al. 1997). However, the galaxies showing X-ray emission from the halo most convincingly, are really the nearby edge-on starburst galaxies like NGC 253, M82, NGC 3079, and NGC 3628.
While the halo of NGC 253 presents a spectacular image and clearly stands out compared to other normal galaxies, its properties do not at all beat other nearby edge-on starburst galaxies. In Table 9 we compare the z extent of the halo above the disk as well as the temperature T, absorption corrected luminosity (0.1-2.4 keV) L<sub>X</sub> of the halo gas and inferred gas mass $`m_{\mathrm{gas}}`$. From Table 9 one would gather that in view of the halo extent, NGC 253 is not on the low side of the sample. Recently, however, Lehnert et al. (1999) reported emission at a distance of 11.6 kpc for the smaller M82 halo, which they interpret as being due to shock-heating of a massive ionized cloud by the starburst super-wind, which would clearly exceed the NGC 253 halo extent. Apart from NGC 3628, the temperatures of the halo gas of the other galaxies are higher by a factor of 2 to 3 and the luminosities by up to a factor of 10 implying gas masses larger by up to a factor of 20. Therefore, it is not the extreme halo properties of NGC 253, that have pushed it to the front of galaxy halo investigations, but its proximity and low Galactic foreground absorption. It is interesting to note that for the four galaxies in the above sample, besides starburst activity, LINER activity has also been discussed, pointing to the presence of an active nucleus. All four galaxies are members of dense groups, and at least for M82, NGC 3079 and NGC 3628 it is clear that interactions within the group have most likely initiated the starburst (and AGN) activity. For NGC 253, we cannot be sure that a recent encounter with neighboring galaxies of the Sculptor group has initiated the activity. As an alternative explanation for the starburst, bar activity has been put forward. However again, a ‘flyby’ of another Sculptor group galaxy might have caused the bar, which in turn may have caused the starburst.
## 5 Summary and conclusions
Following the detailed analysis of deep ROSAT HRI and PSPC observations for X-ray point sources (Paper I), we here have characterized the diffuse X-ray emission of this edge-on starburst galaxy. After subtraction of the point-source contribution we determined the geometry and spectra of the diffuse components and discuss our results in view of observations at other wavelengths and from other galaxies. In detail, X-ray, radio, and H$`\alpha `$ measurements will be compared in a separate paper (Ehle et al., in preparation). Our main results can be summarized as follows:
* The diffuse X-ray luminosity of NGC 253 is distributed in about equal emission components from nuclear area, disk, and halo and contributes $``$80% to the total X-ray luminosity of NGC 253 (L$`{}_{\mathrm{X}}{}^{}=\mathrm{5\hspace{0.17em}10}^{39}`$ erg s<sup>-1</sup>, corrected for foreground absorption). The starburst nucleus itself is highly absorbed and not visible in the ROSAT band.
* The “nuclear source” (X34) has an extent of 250 pc (FWHM) and is located about 100 pc above the nucleus along the minor axis on the near side of NGC 253. It is best described as having a thermal bremsstrahlung spectrum with a temperature of T = 1.2 keV (N$`{}_{\mathrm{H}}{}^{}=\mathrm{3\hspace{0.17em}10}^{21}`$ cm<sup>-2</sup>) and L$`{}_{\mathrm{X}}{}^{\mathrm{exgal}}=\mathrm{3\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup> (corrected for Galactic foreground absorption). The nuclear source most likely marks the area where the line-of-sight absorption gets low enough that soft X-ray emission from the heated walls of the funnel drilled by the super-wind, can leak through.
* The “X-ray plume” has a hollow-cone shape (opening angle of 32° and extent of $``$ 700 pc along the SE minor axis). Its spectrum is best modeled by a composite of a thermal bremsstrahlung (N$`{}_{\mathrm{H}}{}^{}=\mathrm{3\hspace{0.17em}10}^{20}`$cm<sup>-2</sup>, T = 1.2 keV, L$`{}_{\mathrm{X}}{}^{\mathrm{exgal}}=\mathrm{4.6\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup>) and a thin thermal plasma (Galactic foreground absorption, T = 0.33 keV, L$`{}_{\mathrm{X}}{}^{\mathrm{exgal}}=\mathrm{4\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup>). It traces the interaction region between the galactic super-wind and the dense interstellar medium of the disk. The soft component with just Galactic foreground absorption stems from the halo above the disk.
* Diffuse emission from the disk is heavily absorbed and follows the spiral structure. It can be described by a thin thermal plasma spectrum ( T = 0.7 keV, intrinsic luminosity L$`{}_{\mathrm{X}}{}^{\mathrm{intr}}=\mathrm{1.2\hspace{0.17em}10}^{39}`$ erg s<sup>-1</sup>), and most likely reflects a mixture of sources (X-ray binaries, supernova remnants, and emission from H ii regions) and the hot interstellar medium. The surface brightness profile, with a bright inner and a fainter outer component along the major axis with extents of $`\pm `$3.4 kpc and $`\pm `$7.5 kpc, resembles the 1.46 GHz radio profile.
* The “coronal emission” originates from the halo immediately above the disk (scale height $`1`$ kpc). It is only detected from the near side of the disk (in the SE, T = 0.2 keV, L$`{}_{\mathrm{X}}{}^{\mathrm{intr}}=\mathrm{7.8\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup>), emission from the back (in the NW) is shadowed by the intervening interstellar medium causing a pronounced gap in the soft X-ray emission along the NW edge of the optical disk, unambiguously determining the orientation of NGC 253 in space. The X-ray corona is primarily due to the nuclear superwind that is also responsible for the emission in the outer halo. An additional component may be due to hot gas ejected by galactic fountains originating within the boiling star-forming disk.
* The emission in the outer halo can be traced to projected distances from the disk of 9 kpc, and shows a horn-like structure. Luminosities are higher (10 and $`\mathrm{5\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup>, respectively) and spectra harder in the NW halo than in the SE. The emission most likely originates from a strong galactic wind emanating from the starburst nucleus. A two temperature thermal plasma model with temperatures of 0.13 and 0.62 keV or a thin thermal plasma model with temperature of 0.15 keV and Gaussian components above $``$0.7 keV and Galactic foreground absorption are needed to arrive at acceptable fits for the NW halo. This may be explained by starburst-driven super-winds or by effects of a non-equilibrium cooling function in a plasma expanding in a fountain or wind.
While NGC 253 does not beat other nearby edge-on starburst galaxies with its X-ray halo parameters, it stands out due to the low foreground N<sub>H</sub>, distance, and the favorable inclination, which all together allowed the detection of so many details with ROSAT. Of specific interest are three questions that only could be touched but not finally answered by ROSAT and ASCA/BeppoSAX:
1. Is the X-ray emission in the outer halo caused by the nuclear super-wind or by hot gas transported into the halo via fountains that are fed from the star forming, boiling disk?
2. Can the X-ray spectra of the halo really be described by a delayed recombination model?
3. Is the nuclear starburst region responsible for the nuclear and X-ray plume emission or is there a contribution from jet-like activity of a hidden AGN?
These and more questions should be solved with the next generation of X-ray instruments on board XMM-Newton and Chandra, that will have a broader energy coverage, higher collecting area, and much better spatial and energy resolution. Therefore, NGC 253 promises to stay a key object to understand the formation and fueling of hot galaxy halos with future X-ray observations.
###### Acknowledgements.
Hartmut Schulz kindly provided us with the H$`\alpha `$ images for the HRI overlays of the central field. We are very grateful to Dieter Breitschwerdt for discussions of the effects of non-equilibrium cooling, and Andy Read and Ginevra Trinchieri for a careful reading of the manuscript. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, CALTECH, under contract with the National Aeronautics and Space Administration. To overlay the X-ray data we used an image based on photographic data of the National Geographic Society – Palomar Observatory Sky Survey (NGS-POSS), obtained using the Oschin Telescope on Palomar Mountain. The NGS-POSS was funded by a grant from the National Geographic Society to the California Institute of Technology. The plates were processed to the present compressed digital form with their permission. The Digitized Sky Survey was produced at the Space Telescope Science Institute under US Government grant NAG W-2166. The ROSAT project is supported by the German Bundesministerium für Bildung und Forschung (BMBF/DLR) and the Max-Planck-Gesellschaft (MPG).
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# PRODUCTION OF LITHIUM IN THE GALACTIC DISK.
## 1 Introduction and observational base
A few seconds after the Big Bang, four light isotopes were produced: D,<sup>3</sup>He,<sup>4</sup>He and <sup>7</sup>Li (see eg. Walker et al. 1991,Copi et al. 1995,Shramm and Turner 1998); all of these warrant careful study, and here we are focusing on <sup>7</sup>Li.The importance of understanding the evolution of the Galactic abundance of Li was highlighted in the key discovery by Spite $`\&`$ Spite (SS) (1982) that the observed abundance of Li in Galactic stars does not continue to fall uniformily with decreasing iron abundance below \[Fe/H\]$``$-1, but levels off to a ”plateau” at a level of logN(<sup>7</sup>Li)$``$2, which SS interpreted as corresponding to the abundance produced by big bang nucleosynthesis (SBBN).Spite and Spite measured the <sup>7</sup>Li abundance as a function of metallicity (iron abundance) and surface temperature.They found that the <sup>7</sup>Li abundance is flat for surface temperatures greater than about 5600K, and further, it is also flat for the stars with the lowest iron abundance.The first plateau suggests that the stars with the highest surface temperatures are not destroying their <sup>7</sup>Li by convection (the depth of the convective zone depends on surface temperature and is shallowest for stars with the highest surface temperatures).The second plateau indicates that any post-big-bang production must be insignificant for the most metal-poor stars because the <sup>7</sup>Li abundance does not increase with iron abundance.The case against major depletion (and hence for a plateau abundance that reflects the primeval abundance) was strengthened by the observation of <sup>6</sup>Li in certain population II stars (Smith et al. 1993, Hobbs and Thorburn 1994).Big-bang production of <sup>6</sup>Li is negligible; the <sup>6</sup>Li seen was probably produced by cosmic-ray processes (along with beryllium and boron).Because <sup>6</sup>Li is much more fragile than <sup>7</sup>Li and yet still survived with the abundance relative to Be and B expected from cosmic-ray production, depletion of of <sup>7</sup>Li cannot have been very significant (Steigman et al. 1993). Using this interpretation the primordial abundance is given by logN<sub>P</sub>(<sup>7</sup>Li)=2.2($`\pm `$0.2), a value confirmed in detailed work by a succession of authors (Rebolo,Beckman $`\&`$ Molaro 1987,Hobbs $`\&`$ Thorburn 1991,Spite 1991,Thorburn 1994) which can be combined with the SBBN produced abundance of <sup>4</sup>He (see e.g. Pagel et al. 1992), to infer basic cosmological parameters: the universal baryon density $`\mathrm{\Omega }`$<sub>b</sub>, and the number of massless two-component neutrino types N<sub>ν</sub>.To be sure that the population II abundance of <sup>7</sup>Li is a largely undepleted SBBN abundance, entails two essential steps: showing that population II stars (with T<sub>eff</sub>$``$5500 K) have not depleted or barely depleted their <sup>7</sup>Li, and showing that most of the <sup>7</sup>Li in population I stars is of Galactic origin.The first step has already been accomplished via the theoretical work of the Yale group, who showed (Pinsonneault, Deliyannis $`\&`$ Demarque 1992) that sub-surface convective transport, and hence <sup>7</sup>Li depletion is strongly suppressed at low metallicities.As a result of this, and of steadily accumulating observations, opinion (see Spite $`\&`$ Spite 1993) has swung strongly behind the view that logN(<sup>7</sup>Li)$``$2.2 is the SBBN value.Thorburn’s (1994) refined work on the ”plateau”, has brought out a scatter in the <sup>7</sup>Li v. \[Fe/H\] plot below \[Fe/H\]=-1.5, which is incompatible with zero production of <sup>7</sup>Li in the halo, but in practice strongly supports a primordial value for <sup>7</sup>Li not far above logN<sub>P</sub>(<sup>7</sup>Li)$``$2. The second step is quite complicated, because during the disk lifetime there may have been a number of significant production processes for Li, and also a number of destruction, or depletion processes.The initial primordial abundance masks any Li evolution in the halo, so we concentrate our attention in the present paper on production in the disk and its interpretation. The evolution of the lithium abundance in the Galactic disk can be followed via observations which define the upper envelope of the lithium abundance in stars over the range of iron metallicity, -1.5$``$\[Fe/H\]$``$0.1, which characterizes the disk population.The underlying assumption is that while lithium is in general depleted within stars, for a given value of metallicity the highest observed abundance value for a set of stars will correspond to minimum depletion, and hence to an optimum approximation to the Galactic interstellar lithium abundance at the epoch when the stars were formed.Following the evolution of lithium should give similar insight into the processes which form it to that which we can obtain by following the evolution of any other element.The rise in the ratio O/Fe with decreasing Fe, for example, gives the key to understanding the origin of a major fraction of Galactic oxygen in supernovae of type II, whereas iron is formed in all stars with masses greater than or equal to 1 solar mass.
Although the general trend in Galactic Li evolution can be followed via the Li-Fe envelope, the effect of depletion imposes the need for the greatest care when interpreting the observed abundance in any single object.Depletion is well-known to occur in cool stars (see e.g. Pinsonneault et al. 1992,Deliyannis et al. 1990): those with T<sub>eff</sub> less than the solar value, and occurs also in the ”lithium gap” (Cayrel et al. (1984), Boesgaard (1987)) in the middle-F range of spectral classes.Stellar depletion renders lithium particularly interesting as a probe of stellar structure (Steigman et al. 1993),but makes it more difficult to interpret measured abundances in terms of production processes.
However, if we are not able to account adequately for the present-day and post-solar system abundance values: logN(Li)$``$3, there must remain some room for doubt about the BBN value.For this reason, as well as for its intrinsic importance as a test of Galactic evolution models, the source(s) of Galactic lithium continue to be of considerable research interest.A number of production mechanisms have been proposed: cosmic ray spallation of CNO (Reeves,Fowler & Hoyle (1970), Meneguzzi,Audouze & Reeves (1971), Walther,Mathews & Viola (1989)), nucleosynthesis in novae (Arnould & Norgaard (1979), Starrfield et al. (1978)), in the atmospheres of red giants (Cameron & Fowler (1971)), in supernovae (Dearborn et al. 1989,Woosley et al. 1990), and in AGB stars and carbon stars (D’Antona and Matteucci (1991)) and in black hole binaries (Martin et al. (1994)).It is well accepted that processes in the interior of normal stars not only fail to yield lithium, but tend to deplete it.In spite of the detection of individual lithium rich objects which might be characteristic sources, models which incorporate such sources into a Galactic evolution scheme (Audouze et al. (1983), Abia & Canal (1988), D’Antona & Matteucci (1991)) do not give good agreement with the observed lithium-iron envelope.Further, the spatial homogeneity of the Fe-Li curve points against sparse sets of point sources, even distributed sources, and in favour of a more diffuse origin for the lithium. Recent detailed models of Li production assumed in carbon stars, massive AGB stars,SNII and novae (Romano 1999), do not give fair fits to the very high slope of the Li abundance vs. \[Fe/H\] near \[Fe/H\]$``$-0.3 (see Fig. 1 of Romano 1999).
The light nuclide <sup>6</sup>Li is not produced dignificantly in SBBN and is expected to be produced over the lifetime of the Galaxy in Galactic cosmic ray spallation as well as $`\alpha `$+$`\alpha `$ fusion reactions.Its high fragility to stellar processing makes it a less useful tool than <sup>7</sup>Li to constrain big bang nucleosynthesis, but many authors have modelled <sup>6</sup>Li time evolution due to the assumed conexion with the <sup>9</sup>Be and B abundances (Yoshii et al. 1997,Lemoine et al. 1997,Vangioni-Flam et al. 1999,Fields and Olive 1999, Ryan et al. 1999).
It has been suggested (see,e.g.,Steigman et al. (1993)) that since (due to dust grain depletion, and ionization equilibrium uncertainly) isotope ratios can be determined more reliably in the interstellar medium than absolute abundances or ratios of different elements, the interstellar isotope ratio <sup>6</sup>Li/<sup>7</sup>Li might offer a better parameter to test source models than the absolute lithium abundance estimated directly in the ISM.However, measurements of the local interstellar <sup>7</sup>Li/<sup>6</sup>Li ratio (e.g. Lemoine et al. (1993),Meyer,Hawkins & Wright (1993)) show major variations, with differences of up to an order of magnitude from one IS cloud to another.Further, given the extreme difficulty of the ratio measurement in a stellar atmosphere and the consequent extreme paucity of such data as a function of metallicity together with the difficult interpretation of these data in terms of differential stellar depletion as a function of stellar surface temperature, it would be especially risky to attempt to draw conclusions at this stage by using a chemical evolution model to predict the evolution of the isotope ratio against, say, iron abundance.Because of the apparent spatial inhomogeneity it is not even safe to place too much emphasis on the well measured solar system <sup>7</sup>Li/<sup>6</sup>Li ratio of 12.5 (Mason (1971)).These considerations have led us to the present approach of concentrating on the overall Li abundance envelope as a key model constraint.
In this paper we adopt the technique we were the first to use in Rebolo et al. (1988) of assuming that the upper envelope of lithium vs. iron abundance plot is at least a close approximation to the undepleted curve.We do in fact examine the alternative hypothesis: that this envelope represents a depletion curve, and show that this interpretation is quantitatively improbable, so leaving the way clear for the use of the lithium vs. iron envelope as a test of lithium production processes.The purpose of the paper is to show, using this envelope, which types of production processes are excluded, and which permitted.Without going into any numerical detail it is evident from inspection (see Fig. 1) that the rise in the lithium abundance towards the values found in objects close to solar (iron) metallicity occurs relatively late in the Galactic disk evolution time scale; the lithium rise lags the rise in iron, precisely the opposite case to that of oxygen (see Fig. 3).A direct implication is that processes associated with type II supernovae could, but with difficulty, yield the observed lithium production.This consideration not only covers hypothetical processes within the supernovae, but interactions of the energetic particles which they produce in processes occurring in the interstellar medium (ISM).This is just an example of how we can hope to constrain the Galactic lithium production process using the available observational data.Below we will use quantitative modelling (both analytical and numerical) with the aim of reproducing the Li-Fe envelope, thereby eliminating processes which predict significantly different envelopes.What remains will be candidate material for the process (or processes) which gave rise to some 90$`\%`$ of the lithium we can observe today.
In section 2 we show, using simplified analytical models, how the overall shape of the lithium-iron curve for the Galactic disk can be reproduced on the assumptions of delayed production of lithium and increasing infall of gas to the disk.In section 3 we describe briefly a numerical chemical evolution model used to handle the detailed evolution of lithium.In section 4 we examine quantitatively the problem of the galactic cosmic ray (GCR) flux as a candidate source for lithium.In section 5 we compare some of the suggested production mechanisms for lithium.Finally we draw some conclusions about the primordial abundance of lithium.
## 2 Analytical and semi-analytical models for the temporal evolution of Li: comparison with data and with previous modelling.
Numerical modelling of Galactic chemical evolution, which we consider in section 3, offers the advantage of being able in principle to match with realism the physical variables which have driven it, in other words given correct assumptions to give exact fits to the relevant observations.However, without descriptions of inordinate length and including full listing of complete codes, such numerical models are not as transparent as desired.While analytical models are inevitably too simple to reproduce most data sets, their use is more didactic, enabling the underlying physics to be better demonstrated.For this reason we have chosen first to show analytical models which illustrate the physical requirements of any scheme that predicts the observed lithium-iron relation, before presenting our numerical models.The analytical models are designed according to standard methodology, which is based on the paradigmatic work of Tinsley (1980). As the observations give us directly the evolution of the abundance of one element vs. that of another, and because the analytical treatment predicts the evolution of abundances vs. time, we adopt hereafter our well tested (see Figs. 2a,2b) numerical results for the translation from the metallicity (taken as \[Fe/H\] or \[O/H\]) plane to the time plane and vice versa. Firstly, in the volume under consideration, we set the star formation rate, SFR, proportional to the gas fraction $`\sigma _g`$, following Schmidt (1959), so that
$$SFR(t)=\gamma \sigma {}_{g}{}^{}(t)$$
(1)
and the time evolution of the star formation rate is given by
$$\frac{dSFR(t)}{dt}=\gamma SFR(t)+\gamma E(t)$$
(2)
where E(t) is the gas acquired by the zone per unit time due to expulsion by stars plus any net infall of gas to the volume , and SFR(t) is the rate of conversion of gaseous mass into stellar mass.Integrating equation (2) gives
$$SFR(t)=\gamma e^{\gamma (t+{\scriptscriptstyle G(t)𝑑t})}$$
(3)
where
$$G(t)=\frac{E(t)}{SFR(t)}$$
(4)
To simplify the treatment we will first approximate E(t) to SFR(t); this is the case, for example, where the star forming process has efficiency close to 100$`\%`$ and rapidly consumes all the available gas in the volume so that in each time interval new star formation uses only gas expelled from existing stars.This is akin to the assumption of instantaneous recycling.In this case G(t)$``$1 and SFR(t)$``$$`\gamma `$.A similar expression would be obtained if the infalling mass of gas added to the gas expulsion by stars, at time t, were comparable to the gas consumption by star formation at the same time t.The first assumption to test is that lithium is produced either in SNe of type II or by processes in the interstellar medium caused by energetic particles expelled from type II supernovae.In this case the rate of lithium production will be proportional to the star formation rate, which gives
$$\frac{dLi(t)}{dt}SFR(t)\gamma $$
(5)
which integrating and translating from the time plane to the metallicity plane through simple parabolic fit (of the form \[Fe/H\]$``$-(1-$`\frac{t}{15}`$)<sup>2</sup> with t in Gyr.) to the data of Fig.2a, gives
$$Li(t)Li(0)+\gamma (1([Fe/H]){}_{}{}^{1/2})$$
(6)
where Li(0) is the initial lithium abundance.From this we can deduce that models in which the bulk of Galactic Li is produced in processes involving type II SNe could satisfy the observational requirements, but with some degree of difficulty (see Fig.3).
One can obtain an approximation to the effects of infall in this type of models in a less direct way, but which serves to illustrate the principle.Here we must refer ahead to a numerical model developed in Section 3, from which we take an approximate time-dependence of the rate of SNeII.It turns out to be approximately parabolic (see Fig. 5) centered at t=80 in unit model steps of 100 Myr, and we fit
$$\frac{dLi(t)}{dt}1.510{}_{}{}^{8}(t80){}_{}{}^{2}+510^5$$
(7)
which gives on integration
$$Li(t)0.510{}_{}{}^{8}(t80){}_{}{}^{3}+510{}_{}{}^{5}t+Li(0)$$
(8)
The numerical coefficients show that for all times of interest (t$``$5Gyr) up to the present age of the disk,the cubic term can in fact be neglected, and the lithium abundance grows essentially as
$$Li(t)510{}_{}{}^{5}t+Li(0)$$
(9)
which shows the same behaviour than that of eq. (6) in the metallicity plane.
Thus processes whose rate is proportional to the number of SNII in the disk (either in stars or in the ISM) are not ruled out by the lithium-iron envelope test, in all scenarios, with or without infall of gas.
Processes which depend on the SNII rate included among current explanations for lithium production in the ISM (see e.g. Ramaty et al. 1997): spallation of CNO by highly energetic alphas, has been adduced to account for the major fraction of the lithium produced.An alternative non-stellar mechanism is the interaction of moderate energy alphas with the existing abundant He nuclei in the ISM, producing lithium via the $`\alpha `$+$`\alpha `$ fusion reaction.The sources of these low energy alphas can be the winds of normal stars (see section 4).In this case an additional IS acceleration mechanism is required to give the alphas sufficient energy, at least a few MeV, required for the $`\alpha `$+$`\alpha `$ fusion reaction.One could assume that the effects of the presence of supernovae on the ISM can produce the required acceleration (but, we have shown, using simple calculations, that the same results would follow using the acceleration in wind termination shocks of stars of all masses as proposed by Rosner $`\&`$ Bodo (1996)).The lithium production rate at a given epoch will then be proportional to the mass outflow from stellar winds multiplied by the supernova rate.We can approximate the latter as constant (cf. above) and the mass outflow rate from stars of a given mass will be proportional to the number of stars of that mass; for low mass stars (those with masses less than 1M, and hence supplying gas only via winds because their lifetimes are greater than the life of the disk) this number is fully cumulative, and we have approximately (integrating in time SFR(t)=constant)
$$N{}_{stars}{}^{}(t)t$$
(10)
and so
$$\frac{dLi(t)}{dt}t$$
(11)
hence, integrating, and translating from the t-plane to the metallicity plane through the same fit to data than was taken for eq. (6), one has:
$$Li(t)Li(0)+\gamma ^{}(1([Fe/H]){}_{}{}^{1/2})^2$$
(12)
with $`\gamma `$’ a constant; see Fig. 3 to compare the predictions of eq. (12) with data. Now, considering intermediate mass stars (those with masses 1M$``$m$``$3M, which supply gas mainly via processes in their late evolutionary stages) as the main producers of Li, because of their fairly long lifetimes, their expulsion of gas accumulates over times long after their birth, and for the lower part of this mass range, keeps accumulating until the present epoch.We may approximate analytically the collective gas expulsion rate by an exponential, while SN(t) can be taken as approximately constant, as above.In this case one has
$$\frac{dLi(t)}{dt}e^{kt}$$
(13)
and integrating and translating to the metallicity plane as above, we have
$$Li(t)\gamma ^{\prime \prime }(1+e{}_{}{}^{k^{}(1([Fe/H]){}_{}{}^{1/2})})$$
(14)
with $`\gamma `$” and k’ constants.
In Fig. 3 we can see the fit of this expression to the observations. One can go further and analyze in more detail the time dependence of the Li abundance on the properties of stellar mass ranges as follows: Assuming an SFR approximately constant, and assuming that the produccion of Li is associated with the gas expulsion by stars of all masses one has
$$\frac{d(Li(t))}{dt}_{m_t}^{m_u}m{}_{}{}^{2.35}R(m)𝑑m$$
(15)
Taking an analytical approximation for R(m) in the form m<sup>0.2</sup>-0.58 (based on the numerical values of Renzini and Voli (1981)) and integrating, one has
$$\frac{d(Li(t))}{dt}0.005+\frac{m{}_{t}{}^{}^{1.15}}{1.15}\frac{0.58}{1.35}m{}_{t}{}^{}^{1.35}$$
(16)
Using the approximation relating the mass of a star and its lifetime as m$`{}_{t}{}^{}(11.7){}_{}{}^{1/2}t^{1/2}`$, translating to the metallicity plane as below, and integrating one has, neglecting the first term which has a very low value compared with the other two:
$$Li(t)Li(0)+k^{\prime \prime }(0.005(1([Fe/H]){}_{}{}^{1/2})0.049(1([Fe/H]){}_{}{}^{1/2}){}_{}{}^{1.675}+0.13(1([Fe/H]){}_{}{}^{1/2}){}_{}{}^{1.575})$$
(17)
where k” is a constant; see Fig. 3.If we place an upper limit on the mass range of stars contributing to the gas which yields Li, this analytical approximation leads to a sharp increase from zero Li production on short timescales to a high value when times reach the scale of lifetime of the stars whose masses are those of the upper limit.For example for an upper mass limit of 1M the Li production will be zero until the time is near 12Gyr or \[Fe/H\]$``$0.0 (see Fig. 3).Thus we can see, in general terms, that using a gas expulsion timescale and modulating the upper mass limit, we can find solutions leading to late-time production of Li, as apparently required by the observations.
For comparison one can set out a similar formulation using the assumption that Li production is proportional to the cumulative number of stars of all masses at a given time, which implies production within the stars or in their envelopes, rather than in the ISM via expelled gas:
$$\frac{d(Li(t))}{dt})_{m_L}^{m_t}m^{2.35}dm$$
(18)
Performing these integrals with the same approximation for the mass-time and time metallicity dependences as before, and taking m<sub>L</sub>=0.1M, one has
$$Li(t)Li(0)+k^{\prime \prime \prime }(0.084(1([Fe/H]){}_{}{}^{1/2}){}_{}{}^{1.675}+16.58(1([Fe/H]){}_{}{}^{1/2}))$$
(19)
where k”’ is a constant. A comparison of these stylized models: with Li production proportional to the cumulative numbers of moderate mass stars or alternatively to the cumulative flux of expelled gas, is given in Fig. 3.
In all the approximations based on proportionality of Li production rate to the gas expulsion rate by stars of low or intermediate masses (equations (12), (14)) one can see considerable resemblance to the observed lithium growth profile.The reason for the importance of intermediate mass stars, rather than high mass stars, as producers of alpha particles leading to Li production is that the former turn out to be particularly efficient in expelling He. For stars with masses greater than 3M the central temperature becomes high enough for the ignition of the triple-alpha reaction (which transforms He to C) before the giant stage is reached.At the moderate densities of the central regions, this nuclear process gains importance in a gradual manner.On the other hand for stars with masses between 0.5M and $``$3M the central regions become degenerate and the triple-alpha reaction ignites via the violent helium flash.
Although we would certainly not wish at this stage to exclude the SNe as key producers of Li in the Galactic Disk, to be coherent with our general scenario of chemical evolution for the Galaxy, we have used the $`\alpha `$+$`\alpha `$ process in the ISM as illustrative of processes which follow the behaviour of stars of intermediate and low masses, those whose lifetimes are long, and whose numbers in the disk have therefore grown cumulatively with time.Any process with equivalent time-dependent characteristics might, in general terms, satisfy this global observational constraint.Production of Li in novae (Arnould $`\&`$ Norgaard 1979,Starrfield et al. 1978), or in SN type I, might also satisfy the criterion of delayed production because there the rate of production would be proportional to the number of stars with intermediate masses at each time (cumulative as their lifetimes are long).But, as one can see in Fig. 4, where we plot the results from different numerical models in the case of GCR flux proportional to the cumulative number of stars with masses less than or equal to 1M (for upper mass limits greater than 1M the increase in the number of stars with time is proportionally less and less), the prediction, although not bad, is clearly inferior in fit to that obtained using gas expulsion by stars of masses less than or equal to 3M (and is in fact also worse than that obtained using gas expulsion of stars with masses less than or equal to 2M).
Another possibility which has been discussed is the production of Li in compact objects such as neutron stars and black holes, but the time evolution of the numbers of compact objects would be proportional to the rate of gas expulsion by all stars (which is mainly that of SNI+SNII).This production is shown in Fig. 1 from our numerical modelling.In fact, as one would expect,the curve is similar to that for the model in which Li production is proportional to the SFR, also shown in Fig. 1.
Further possibilities for the production of Li, such as those in AGB stars or carbon stars (Matteucci et al. (1995)) are unable to satisfy the detailed observational constraints, as will be seen in Fig. 5.
The aim of this section has been to show those categories of models, and hence of lithium sources, which can best account for the lithium-iron envelope observations.However for a valid test of any process which is a candidate to have produced the observed disk lithium we have no choice but to use quantitative, numerical, modelling methods.
## 3 Numerical evolutionary models: the basic formalism.
The model we have used for the evolution of the disk in the solar neighborhood empoys the formalism already explained in Casuso & Beckman (1997),which embodies a numerical rather than an analytical approach,in order to take all the relevant physics adequately into account.The model allows us to follow the evolution,within a fixed volume of space,of the gaseous mass fraction $`\sigma _g`$ and the abundances X<sub>i</sub> of 6 nuclides:<sup>4</sup>He,<sup>12</sup>C,<sup>13</sup>C,<sup>14</sup>N,<sup>16</sup>O, and <sup>56</sup>Fe,selected because observations of their evolutionary abundance behaviour are available.The set of basic equations employed,in which the units are mass fraction per unit time interval,are:
$$d\sigma _g=SFR(t)+E(t)$$
(20)
$$d(\sigma _gX{}_{i}{}^{})=_{m_t}^{m_u}SFR(tt{}_{m}{}^{})\varphi (m)(Q{}_{i}{}^{}(m)+X{}_{i}{}^{}(tt{}_{m}{}^{})(R(m)Q{}_{i}{}^{}(m))R(m)X{}_{i}{}^{}(t))dm+J(t)$$
(21)
with
$$J(t)=P(t)(X{}_{i}{}^{}{}_{}{}^{}(t)X{}_{i}{}^{}(t))$$
(22)
with
$$E(t)=_{m_t}^{m_u}SFR(tt{}_{m}{}^{})\varphi (m)R(m)dm+P(t)$$
(23)
in which t<sub>m</sub> is the lifetime of a star of mass m,X<sub>i</sub>’(t) is the halo abundance,P(t) is a term which represents net inflow of material to the volume under study, SFR(t) is the star formation rate,and $`\varphi `$(m) the initial mass function (IMF) of the stars.Within this volume there is a population of stars whose lifetimes t<sub>m</sub>,for a mass m$``$m<sub>t</sub> (where m<sub>t</sub> is the mass of a star which has a lifetime t),are less than or equal to the value of the time variable t,and which eject the products of their internal nucleosynthesis at a rate proportional to SFR(t-t<sub>m</sub>),the star formation rate at their birth.We term the fraction of its mass which a star ejects during its lifetime R(m),and Q<sub>i</sub>(m) is the yield of nuclide i from a star of mass m.The total mass which has been added to the ISM,either by stellar evolution or by net inflow into the volume under consideration,is called E(t).All relevant stages of stellar evolution have been taken into account, including post-main-sequence phases (e.g. AGB-stars, planetary nebulae and other gas ejection stages) and explosive processes.
We have used a simple proportionality law for the dependence of the star formation rate on the gas fraction,viz. SFR(t)=$`\gamma \sigma _g`$<sup>k</sup>(t) where,following the classical approach of Schmidt (1959) we have used k=1,and the value of $`\gamma `$ is 0.11 Gyr<sup>-1</sup>.The observed parameters of chemical evolution for the solar neighborhood are reasonably reproduced.As a suitable approximation to the IMF we have used the Salpeter (1955) law,i.e. $`\varphi `$(m) $``$ m<sup>-(1+x)</sup> with x=1.35,between 73 M and 0.5 M,and have approximated the flattening observed at low masses (see e.g. Scalo (1986),Kroupa et al. (1993)) with a plateau of value $`\varphi `$(0.5) between 0.5M and 0.1M.This approximation is convenient for computations,and represents a reasonable fit to the observations.Slightly better but more complicated fits to the observations would not affect any of the conclusions reached here.Finally we approximated the stellar lifetime t<sub>m</sub> as a function of mass m,following Arimoto $`\&`$ Yoshii (1986) by t<sub>m</sub>=11700/m<sup>2</sup> in units of Myrs for t,and solar masses for m.
We have shown that our model can reproduce the well established observed chemical abundance parameters of the galactic disk in the solar neighborhood, the <sup>9</sup>Be/H and <sup>10+11</sup>B/H temporal evolution (see Casuso & Beckman (1997)), as well as that of D/H, <sup>7</sup>Li/<sup>6</sup>Li and <sup>11</sup>B/<sup>10</sup>B (see Casuso& Beckman 1999).
It is well known that closed box models (with P(t)=0),of galactic chemical evolution fail to reproduce several of the disk constraints,most notably the metallicity distribution in the disk,characterized by the low numbers of G dwarfs with low metallicities.They also fail to reproduce the disk evolution of Be and B vs. Fe, and models which give adequate fits to these observed data sets require increasing infall to the disk of metal-free or metal-poor gas (Casuso & Beckman (1997)).We have therefore adopted the same representation of the infall as was used in that paper:
$$P(t)=\frac{e^{\lambda t}}{M(t)}$$
(24)
where M(t) is the total mass of the zone at time t,and $`\lambda ^1`$ is a time constant which must be in the range of a few Gyr.It is notable that recent observations of abundances in the local interstellar medium (Fitzpatrick (1996)) showing undepleted elements with abundances significantly below solar, give support to the idea of steady dilution by infall of non-enriched gas to the disk. Including the global effect of depletion as in the model of Casuso and Beckman (1999) does not in practice yield significant improvements in the data fit compared to non-depleted models, within the limits of error.
## 4 Application of numerical modelling: incorporation of $`\alpha `$+$`\alpha `$ by GCR in the ISM.
A seminal early paper describing the physics of light element production by spallation and fusion reactions,authored by Meneguzzi,Audouze and Reeves (1971) has been the basis of much of the intervening work in the field.These authors showed that nuclides of light elements can be produced by spallation during collisions of galactic cosmic ray (GCR) protons and alpha particles with nuclei of C,N and O in the interstellar medium (ISM),and also by CNO in the GCR colliding with protons and alphas in the ISM.The contribution of the latter set of reactions to the production of Li,Be and B nuclides by spallation has been estimated to be some 20$`\%`$ of the former (Meneguzzi and Reeves 1975),and this should not have been very different in the past,assuming that the cosmic ray composition reflects that of the ISM.This has meant that as far as spallation is concerned we needed to calculate in detail only the former reactions,taking the latter into account by proportion.A further source of light element nuclides,whose importance has been recognized more recently,is the production of <sup>6</sup>Li and <sup>7</sup>Li by fusion reactions between GCR alpha particles and those of the ISM.The relative importance of this mechanism may have declined somewhat,as the abundaces of CNO in the ISM have grown relative to that of <sup>4</sup>He (although the abundance of the latter is still overwhelmingly greater), but in the early phases of the disk it was certainly an important mechanism (Montmerle,1977,Steigman $`\&`$ Walker,1992),and as we will see,it must still play a major role today.
In this work we have used the standard expression for the production of light element nuclides by GCR protons and alphas in the ISM:
$$\frac{dY_k}{dt}=Y{}_{j}{}^{ISM}(t)F{}_{i}{}^{GCR}(E,t)\sigma _{ij}^k(E)dE$$
(25)
where Y<sub>j</sub>(t) are the abundances,by number,of the various species,and j refers to <sup>12</sup>C,<sup>13</sup>C,<sup>14</sup>N,<sup>16</sup>O or <sup>4</sup>He,k refers to <sup>6,7</sup>Li and <sup>9</sup>Be,<sup>10</sup>B or <sup>11</sup>B,and the variable i refers to GCR protons or alphas.F$`{}_{}{}^{GCR}{}_{i}{}^{}`$(E,t) is the interstellar GCR flux spectrum,and $`\sigma _{ij}^k`$(E) is the cross-section for each reaction i+j$``$k,which has a corresponding energy threshold E<sub>T</sub>.The quantities Y$`{}_{}{}^{ISM}{}_{j}{}^{}`$(t) are computed from the galactic chemical evolution models described in section 3.They are thus observationally constrained,and we have reasonably good estimates of their evolution with t during the disk lifetime.The spallation and fusion cross sections $`\sigma _{ij}^k`$(E) are also well-known (see Read $`\&`$ Viola 1984, Mercer,Austin and Glagola 1997) within narrow limits of error.These cross sections show rather similar global behaviour,starting from thresholds E<sub>T</sub> close to 10-20 MeV/nucleon,peaking somewhere between 20 and 70 MeV/nucleon,and declining rapidly to a plateau above $``$100MeV/nucleon.There is,however,a key difference between the $`\alpha `$+$`\alpha `$<sup>6,7</sup>Li reaction cross-sections and the remaining cross-sections:while the peak value for the former is some 500 times that of the plateau value,for the latter the corresponding ratio is only 5 to 10.This implies that while for the processes that can give rise to <sup>9</sup>Be,<sup>10</sup>B and <sup>11</sup>B the whole of the energy spectrum of the GCR,up to the GeV range,comes into play,for the $`\alpha `$+$`\alpha `$ process only the lowest energy particles,those with less than 100 MeV/nucleon,yield significant <sup>6,7</sup>Li.(Of course a fraction of <sup>6,7</sup>Li is indeed formed by spallation in the higher energy range,but none of the <sup>9</sup>Be or <sup>10,11</sup>B can be formed via the $`\alpha `$+$`\alpha `$ process).This dichotomy has important consequences for the observationally very different time dependence,and hence the metallicity dependence,of <sup>6,7</sup>Li on the one hand,and <sup>9</sup>Be/<sup>10,11</sup>B on the other.
In order to try to reproduce the evolution of the light elements it is clear that we need estimates of the current energy spectrum of the GCR component,of its flux,and of how these parameters have varied with time.For the present epoch the magnitude and spectral shape of the flux of GCR particles reaching the earth are fairly well determined by direct experiment for particles with energies higher than a few hundred MeV/nucleon:those which are directly observed at the earth’s orbit. For these particles a spectrum of form F(E,t<sub>0</sub>)$``$E<sup>-2.2</sup> is found,up to a few GeV/nucleon and F(E,t<sub>0</sub>)$``$E<sup>-2.6</sup> at higher energies (Ip and Axford 1985).However at lower energies the GCR spectrum must be demodulated to take into account the blocking effects of the heliosphere.This has been a well-known cause of difficulties for light nuclide production theory(Meneguzzi et al. 1971;Meneguzzi $`\&`$ Reeves 1979;Reeves $`\&`$ Meyer 1978),and the problem of determining the spectral dependence and the amplitude of the unmodulated interstellar component of the GCR spectrum below 100 MeV/nucleon still lacks an entirely acceptable solution.In the present modelling exercise we follow Reeves $`\&`$ Meyer (1978) in taking a most probable value for solar demodulation of 5,for the whole GCR spectrum,and supplement this with a further mean factor of 15 for the particles with energies having energies E$``$100MeV/nucleon,which is consistent with the estimates made by McDonald et al. (1990) and McKibben (1991) from observations of the low energy component of $`\alpha `$ particles out to a heliocentric distance of 43 a.u. with Pioneer 10 and 11.More recent studies do not claim major improvements here,because no direct measurement beyond the heliopause has yet been made.To summarize,we take the spectral dependence of the flux to be proportional to (E+E<sub>0</sub>),where E<sub>0</sub> is the rest energy of the proton,and $`\lambda `$ takes values of 2.6 below 0.1GeV/nucleon,2.2 between 0.1 GeV/nucleon and 1GeV/nucleon,and again 2.6 at energies higher than this,and in these we follow previous studies on light nuclide production by Walker et al. (1985) and by Steigman $`\&`$ Walker (1992).Finally we normalized the flux to match the constraint imposed by the measured GCR proton flux for energies greater than 0.1 GeV/nucleon:12.5 cm<sup>-2</sup>s<sup>-1</sup>GeV<sup>-1</sup>nucleon<sup>-1</sup> (Gloeckler $`\&`$ Jokipii 1967),together with a ratio of the $`\alpha `$/p fluxes of 0.15,consistent with the observations of Gloeckler $`\&`$ Jokipii (1967) as re-examined by Webber $`\&`$ Lezniak (1974).Of course, one must invoke (as must all models invoking GCR reactions to produce the light elements) effective magnetic confinement of GCR’s in the Galaxy in order to obtain the required high absolute fluxes of GCR protons and alphas.
In the present work we are concerned with Li production, but in the models we have included depletion for the gas which has been processed into stars (we have assumed here that stars which expel their gas into the ISM have completely depleted their Li so that this expelled gas has zero Li abundance), and the ”impoverishment” due to the infall of gas to the disk from the halo: we assume that this gas has the initial,i.e. the primordial, Li abundance.However the inclusion, or not, of these effects, influences the model results, i.e. the effective production curve, significantly only in the range -0.2$``$\[Fe/H\]$``$+0.2, where the effect of infall has been to cause a slight fall in the observable Li abundance.
One of the arguments of the present paper rests on a correct understanding of how the low energy GCR particle flux has developed with time.Specialists in cosmic ray physics have previously proposed that, since the measured abundances of GCR nuclides show a dependence on the first ionization potentials of the parent atoms,the principal sources of these particles must be the atmospheres of relatively low mass,relatively cool stars (see Cass$`\stackrel{´}{e}`$ $`\&`$ Goret (1978),Meyer (1985,1993)).The consequences of this for the time dependence of the galactic GCR flux in the range of energies required to produce the light element nuclides,and specifically for those below 100 MeV/nucleon which participate in the $`\alpha `$+$`\alpha `$ fusion reaction,is an increase of the flux at later times due to the accumulation of stars with lifetimes comparable to that of the disk.In the seminal study of Meneguzzi,Audouze $`\&`$ Reeves (1971) and also in the careful re-examination of light element abundance production by Walker et al. (1985),the zero-order assumption was made of a GCR flux constant with time.A number of other workers in the field (Reeves $`\&`$ Meyer 1978,Mathews et al.1990) used a time-varying scheme in which the GCR flux F(E,t) is proportional to the supernova rate,SNR(t),which in turn was set proportional to the star formation rate SFR(t),in their models.
In their study Prantzos,Cass$`\stackrel{´}{e}`$ $`\&`$ Vangioni-Flam (1993) were well aware of the importance of the time-dependence of the GCR flux,and also assumed that it followed the supernova rate.However,in the present study we use the assumption that this flux follows the expulsion rate of gas from stars,and we have in fact varied the upper limit of the mass range from which expulsion is considered.The delay entailed allows the model predictions to avoid one of the main difficulties encountered by Prantzos et al.:that without this delay there would have to have been an early sharp increment in the disk of Li (in particular) against Fe,in the range of \[Fe/H\] between -2.0 and -1.0,an increment which is not observed.This delay is due to the fact that although in the early disk there would have been a high SN rate,required to accelerate the GCR particles,the cumulative number of low and intermediate mass stars required to inject major quantities of He nuclei (Meyer 1985) was still low.Both injection and secondary acceleration are required to yield MeV range GCR,and these conditions have been fulfilled simultaneously with increasing effect in the later disk,which explains the observed delay in the onset of disk Li production, even with respect to Fe (and a fortiori with respect to O).
In Fig 1 we contrast the evolution of the Li abundance in a typical model in which the GCR flux is proportional to the SFR,with that in models chosen from those we have applied in the present paper,in which the flux is proportional to the gas expulsion rate from the stellar population at a given epoch.The qualitative difference is evident,and the relative reduction of the GCR flux in the early disk,compared with more recent epochs,is clear.There are two further points about the acceleration and propagation of the GCR which we should make here.As a result of many studies over the past 30 years it is now a widely accepted possibility that the majority of the observed particles in the GCR flux have been accelerated in collisionless mode by shock waves which originate in supernova explosions and propagate through the dilute interstellar plasma (Lagage $`\&`$ Cesarsky 1983,Blandford $`\&`$ Eichler 1987).In the model of GCR propagation by Prantzos et al. (1993) the high energy part of the GCR flux spectrum is modulated according to the escape length of the particles as a function of their energy;this effect has changed with epoch,in such a way that the current spectrum has a greater slope than the spectrum at early Galactic epochs.The low energy fraction of the GCR flux has remained,however,virtually unaffected by this change,suggesting that the evolution of the Li production rate has been due rather to the time variation of alpha-particle density than to the variation of the spectral index.Secondly we must emphasize that provided there has been sufficient volume occupied by SN-affected ISM,the flux of low energy $`\alpha `$-particles will depend principally on the population density of the injectors:low-mass stars,rather than the SN remnants which accelerate them.A final point here is that low energy GCR’s may in fact be accelerated by wind termination shocks due to stars of the full mass range.This process has been invoked by Rosner and Bodo (1996) to explain the diffuse non-thermal Galactic radio emission.Clearly acceleration via this process is not proportional to the SN rate but to the cumulative gas expelled by all stars present at a given epoch.Although taken alone it does not lead to sufficient delay to explain the abrupt rise in the Li-Fe envelope it is, a promising mechanism for $`\alpha `$ acceleration in the context of the Li abundance observations.
Here we should allude to the as yet not fully resolved question of the origin of cosmic rays.In early models of SN shock theory, the thermal gas in the ISM was regarded as the reservoir of seed particles which can became cosmic-ray nuclei.But this clashes with the source composition of the GCR (Meyer 1985).To solve this problem one needs to invoke an injection of suprathermal ions.There are two main types of scenario here: one assumes that the observed local flux of GCRs has its origin in supernovae accelerating their own ejecta, and the other assumes an origin in the atmospheres of intermediate and low mass stars (for discussions see Meyer et al. 1997,Ellison et al. 1997,Ramaty et al. 1997,Ramaty et al. 1998, Higdon et al. 1999).While the early Galactic beryllium data suggest production by cosmic rays originating from SN accelerating their own ejecta, the observed composition of the cosmic-ray source material reflect a correlation with first ionization potentials, leading to the suggestion that cosmic-ray source material originates in the atmospheres of stars.As evidence for this, we know that the abundances of elements with low first-ionization potentials are enhanced in the solar corona and in solar energetic particles, suggesting that similar shock acceleration in low-mass, cool stars could provide a particle injection source for acceleration by supernova shocks in the ISM.Both origins (SN or low-mass stars) have many problems as complete explanations of the origin of GCRs with energies greater than 1 GeV per nucleon (see Ramaty et al. 1998). We cannot consider as coincidental the similarity between the GCRS composition and that of the solar corona which is biased according to first ionization potential, and we must take very seriously the asertion of Ellison et al. (1997) that ”in the outer solar atmosphere the solar coronal gas, the solar wind, and the $``$MeV solar energetic particles have undoubtedly a composition biased according to FIP”, together with the fact that the hydrogen and precisely helium are not well fitted by the alternative model of Meyer et al. (1997) and Ellison et al. (1997) based on volatility and mass to charge to explain the GCRS.Also, we must note that the cosmic-ray electrons have very different spectra from that of the nuclear species at GeV energies, and may, in fact, have entirely different origins (Berezinskii et al. 1990).In addition, atomic collisions of low-energy ions (corresponding to a distinct low-energy cosmic-ray component) produce characteristic nonthermal X-ray emission.On this point Tatischeff et al. (1999) have shown that a distinct Galaxy-wide low energy cosmic-ray component could account for the hard component of the Galactic ridge X-ray emission in the 0.5-10 keV energy domain.Also, one must note the different behaviour of helium and hydrogen data with respect to the other GCR nuclei when the energy of these particles is increased , as inferred from the observations in the solar corona, in the solar wind, in the solar energetic particles and in the GCR (see Meyer 1985, Meyer 1993).One can see how He and H abundances decrease systematically as the energy increases, while the abundances of the other nuclei remain invariant.All these considerations point to an origin for the low-energy $`\alpha `$ particles which optimize Li production, which could be quite different from that for the GCR nuclei at higher energies.In our coherent scenario of chemical evolution for the Galaxy, we point to the origin in low-mass stars of the low-energy (those below 0.1 GeV per nucleon) $`\alpha `$-particles of GCRs, consistent with the fact that the Li production cross section for the $`\alpha `$-$`\alpha `$ fusion reaction falls very steeply outside the energy range between 0.01GeV/nucleon and 0.1 GeV/nucleon (see e.g. Ramaty et al. 1997).
In order to quantify our model, we will account for the energy needed for the $`\alpha `$+$`\alpha `$ fusion reaction be coherent with the energy supplied by the intermediate mass stars in a range when the Li production is efficient.Because our concern is essentially with alphas, and because the $`\alpha `$+$`\alpha `$ fusion reaction which yields Li has a distribution of measured cross sections which is very low outside the 10-100 MeV/nucleon range, we calculate the energy budget by integrating the flux in that range.So, the energy per SNII will be the total GCR energy in this range during e.g. 10<sup>8</sup>yr., taking 0.1 cm<sup>-3</sup> as an average current He abundance in GCRs , and a section corresponding to the 500pc radius taken here as the size of the selected circumsolar volume, all divided by the number of SNII needed by our numerical model ($``$13000 during each 10<sup>8</sup> yr.):
$$Energy/SNII=515_{0.01}^{0.1}E(E+E{}_{0}{}^{}){}_{}{}^{2.6}kdE10^80.1\pi 500{}_{}{}^{2}/13000$$
(26)
with E<sub>0</sub>=0.938GeV, 5 is the solar demodulation factor and 15 a factor required to give a good fit to the data, and k is the normalization constant which is obtained fitting the observed GCR flux of alphas at energies above 0.1GeV:
$$0.1512.5cm{}_{}{}^{2}s{}_{}{}^{1}GeV{}_{}{}^{1}n{}_{}{}^{1}=_{0.1}^1k(E+E{}_{0}{}^{}){}_{}{}^{2.2}dE+_1^{inf}k(E+E{}_{0}{}^{}){}_{}{}^{2.6}dE$$
(27)
From this one obtains the energy per SNII needed for our model to produce the Li observed in the disk by alphas, of 2$``$10<sup>50</sup> ergs. which is in very reasonable agreement for energy available from SN model estimates.Taking the IMF used here, we also obtain that the energy per star of intermediate mass (1-3 M) needed for the alphas produced in this stars, is of 10<sup>49</sup> ergs, very reasonable for the helium flash or coronal mass injections.
The net accelerating power of the OB star enviroment in the local spiral arm has not,in this model, varied by more than a moderate fraction during the disk lifetime. The fact that the model produces some twice as much local GCR flux at the present epoch as at the begining of the disk is due entirely to the accumulation of particles (H,He,C,N,O) emitted at low energies from low mass stars,and not to any substantial change in the net efficiency of the SN mechanism which subsequently accelerates them,and which has been present constantly throughout the disk lifetime.
This scenario is consistent with other parameters such as the lifetimes of $`\alpha `$-particles in the ISM.First we must note that due to the very high temperatures needed to deplete <sup>4</sup>He, efficient destruction occurs in stellar interiors.However, the $`\alpha `$-particles can also disappear in principle due to fusion reactions with other $`\alpha `$s,<sup>12</sup>C,<sup>14</sup>N and <sup>16</sup>O in the ISM, leading to <sup>6</sup>He,<sup>6</sup>Li,<sup>7</sup>Li,<sup>7</sup>Be,<sup>9</sup>Be,<sup>10</sup>Be,<sup>10</sup>B,<sup>11</sup>B,<sup>10</sup>C and <sup>11</sup>C.These latter reactions have cross sections below 100mb, and so, taking densities of the ISM below 10<sup>5</sup>cm<sup>-3</sup>, the lifetimes are greater than 10<sup>12</sup> yr, i.e. greater than the age of the Universe.
On the other hand, the flux of $`\alpha `$-particles in the range of energies concerned here, i.e. between 10MeV/n and 100 MeV/n, could comes from the $`\alpha `$-particles produced in later type stars at energies below 2.5MeV/n. In the mass range 0.5M$``$M$``$3M the central regions become degenerate and the triple-alpha reaction ignites via the violent helium flash.The central energy-generation rate at the peak of this helium flash exceeds 10<sup>13</sup> times that in the center of the Sun causing the well attested expansion to the giant phase before the bulk of the He has been consumed.For these stars the expulsion of He nuclei into the ISM occurs with much greater efficiency than for the less accelerated transformation to the giant phase accompanied by the burning of the He which occurs in stars with higher masses.Another realistic possibility to explain the required low energy alpha flux is that the ions can be injected (at MeV energies) via coronal mass injections (mainly from the coronae of dMe and dKe dwarfs, by far the most numerous stars in the Galaxy) (Shapiro 1999).
One way of seeing the problem is via a two-stage scenario similar to that of Meyer (1985), which assumes the OB associations as the best sites of production of $`\alpha `$-particles of GCRs; there, a large number of later type stars are being formed together with a few short-lived massive stars; the former have a very high surface activity owing to their youth and should emit lots of suprathermal particles, while the latter provide stellar wind and SN shock waves within their few 10<sup>6</sup> yr lifetime; so injectors and high energy accelerators are closely linked in space and time.Energies as low as 0.01-0.1MeV/n are sufficient for suprathermal particles to be accelerated much more efficiently than the thermal gas.Particles with energies below 2.5MeV/n undergo significant coulomb energy losses, which brake and thermalize the particles, impeding $`\alpha `$+$`\alpha `$ fusion production of Li.However, the time for thermalization is inversely proportional to the density of the medium in which suprathermal particles propagate.The $`\alpha `$-particles propagating in dense clouds, thermalize within 10<sup>4</sup> yr, but in the diffuse hot interstellar medium (HIM) where n<sub>H</sub>$``$3.10<sup>-3</sup>cm<sup>-3</sup>, the thermalization time is several times 10<sup>6</sup> yr.We can then assume that the mechanisms already proposed by Meyer (1985) can operate.Ionized media confine energetic particles, so that suprathermal particles emitted in the HIM will not in general traverse any neutral, dense medium.In dense cloud complexes later type stars can form continuously while OB star formation will disperse the complex rapidly; and Meyer (1985) brings out the possibility of reacceleration of suprathermal particles emitted by young late type stars having migrated into the HIM just nearby the cloud complex.
But one might not in fact need a two-stage scenario if the intermediate stars produce and also accelerate alphas to energies in the range 10-100MeV/n, where the coulomb energy losses are not so efficient.
Another possibility is that the $`\alpha `$-particles needed come from the so called ”anomalous” component of He nuclei which is observed precisely at the range of low energies here considered (below 100 MeV/n).These He nuclei had reported as an unusually flat helium spectrum, apparently unrelated with GCR spectrum (Webber 1989).In fact, the solar modulation effects on this helium anomalous component, as observed on the Pioneer 10 spacecraft at $``$40 AU from the sun, between 1985 and 1987 (when the solar modulation reached its minimum) show a change of a factor $``$100 in the intensity of these particles between 10 and 20 MeV/n as they rapidly emerge from the background of low energy GCRs (Webber 1989), in good agreement with this work where we need a total modulation of 5$``$15=75 over the ”normal” GCRs.
We can now summarize the reasons why the present family of models gives an adequate prediction of the observed Li vs. Fe evolution curve.In those older models where a constant GCR flux was used (e.g. Walker et al. 1985) which is in fact not too bad a first-order approximation to the time-delayed evolution for the low energy flux which we obtain here,the importance of the $`\alpha `$+$`\alpha `$ fusion reaction was not realized.In more recent models,on the other hand,where $`\alpha `$+$`\alpha `$ reactions have been well included (Steigman $`\&`$ Walker (1992),Steigman (1993),Prantzos et al. (1993)), sophisticated time-evolution schemes for the GCR have been used, which unfortunately do not explore the delayed contribution of lower mass stars to the low energy GCR spectrum.In the former models normalized to give the correct Li abundance at, say, \[Fe/H\]$``$-1.5, there was simply not enough contemporary Li produced and in the latter,while it would be feasible to attain the contemporary abundance value: log N (Li) $``$ 3,this would entail abundances of Li at \[Fe/H\] $``$ -1 which are too high by more than half an order of magnitude.In the next section we show that our models yield results in much better agreement with the observations.
## 5 Predictions of the family of numerical models: comparison with observations.
In the present section we present the results of our modelling exercises.We have already shown in Fig. 1 the observations that we have set out to model.We began with data which we ourselves reported (Rebolo,Molaro $`\&`$ Beckman 1988) because of the ready availability of the complete data set, to which we have added a newer and extensive set of results from Spite (1996).In Fig 1 we show observations of Li vs. Fe abundances over a wide range of surface temperatures and metallicities.The assumption we will make in interpreting these data is that the upper envelope shows,essentially,the evolution of Li with Fe,while the points which fall below the envelope refer to Li depleted in the individual objects observed.
In Fig 1 we see that the Li abundance remained essentially constant during the halo period (in which the Fe abundance was evolving from its lowest values to around -1.5) and then began to rise.The approximate plateau at low \[Fe/H\],the ”Spite” plateau,corresponds roughly to primordial Li,while the later rise represents the presence of galactic Li production modulated by any averaged depletion which may take place.A key observational result is the rather abrupt rise of the Li envelope at \[Fe/H\]$``$-0.2, to an essentially constant value between -0.2 and +0.2.This second ”plateau”, although it is not as clear as the ”Spite plateau” appears to be consistent with our general model predictions as one can see in Fig. 1.This plateau is a natural consequence of the ”loop back” in abundance of Li already shown in Casuso & Beckman (1997) to occur for the Be and B disk abundances (where it appears more distinctly because stellar depletion is much less important than for Li), and is due to the increase with time in the infall of gas to the disk, which dilutes the Li abundance and more so the Fe abundance, reducing the latter in recent epochs from a broad peak attained several Gyr ago.
The model shown in Fig 1,in which the GCR flux was held proportional to the gas expulsion rate from the whole stellar population represents a first approach using delayed $`\alpha `$+$`\alpha `$ as the principal source of Li,and its relative success is encouraging,but in using gas expelled from the whole stellar mass range it does not try to take into account the fact that the low energy component of GCR,responsible for the $`\alpha `$+$`\alpha `$ process which we are postulating as the principal source of Li, may well originate mainly in lower mass stars.The most direct way to do this is to place upper stellar mass limits on the expulsion rate of gas for which the GCR flux,at each epoch,is deemed proportional.Curves (ii) and (iii) show the results of allowing the GCR flux which enters the $`\alpha `$+$`\alpha `$ reaction to be proportional to the gas expelled per unit time from stars of all masses and up to mass limit of 3M respectively.It is clear that the model with the upper limit of 3M gives a much better fit to the observed envelope shown here, and goes far in demonstrating the need to include the implied time delay in the build-up of the GCR flux in the relevant energy range.The curve for the 2M upper limit shows a good fit to the form of the observations but yields rather low Li production.
Even given the observational uncertainties in the Li-Fe dependence we can use these data and these models to constrain broadly the stellar mass range which serves as a significant source of low energy GCR $`\alpha `$ particles.The fact that models with an upper mass restriction give better fits to the data is itself an argument in favour of intermediate and lower mass stars as the principal sources of the GCR flux at low energies.
One further note should be added.Our evolutionary models were designed to account for the G-dwarf metallicity distribution, and B and Be evolution in the disk.In fact the latter data begin to be very weak statistically for \[Fe/H\] greater than +0.1, which is where the previous model predictions terminate.In Fig. 6 we show the result of a modified model where the upper limiting Fe abundance is +0.3.This is an $`\alpha `$+$`\alpha `$ model with the relevant GCR flux proportional to the gas expelled from stars with masses $``$3M, and it also accounts well for the observations.We must note, however, that the systematic uncertainties admitted by the observers for the stars with high Li abundances and high Fe abundance, are in the sense of requiring lower \[Fe/H\] (0.1 instead of 0.3 or 0.4) (see e.g. Boesgaard & Tripicco (1986)).
## 6 Discussion: alternative sources of disk lithium.
We have used evolutionary models for the galactic disk in the solar neighborhood to re-examine the theme of Li production.In the first part of this paper we presented analytical and semi-analytical models, with their simplifying assumptions and limitations, and showed that models incorporating a delay in Li production compared with that of Fe give a fair description of the observed evolutionary history of Li,using the Fe abundance as a reference parameter.In the subsequent sections we showed that numerical models with infall of non-enriched gas during the disk lifetime give good fits to the observations, the best fits coming from models where the infall has shown a tendency to increase (see Fig.5).One consequence of this has been the non-monotonic evolution of the metallicity with time:\[Fe/H\] has grown to somewhat greater than solar values,and then fallen back slowly.This circumstance yields metal-metal plots with a characteristic fold-back close to solar metallicities.The Li production depends in fine detail,but not in principle or in broad trend on this form of the infall model.
Applying this general evolutionary scheme to Li production we have introduced one novel assumption which proves capable of resolving the hitherto difficult to reproduce Li-Fe curve.This assumption is that the part of the GCR flux responsible for the production of <sup>7</sup>Li and <sup>6</sup>Li by the $`\alpha `$+$`\alpha `$ fusion process:the low energy flux (at least) is emitted principally by intermediate mass stars,an assumption well supported in the literature on GCR production (see e.g. Meyer 1985,Meyer 1993 and references therein).The novelty which this introduces into the models is that the production rate of <sup>7</sup>Li and <sup>6</sup>Li is then constrained to follow not the SFR,but the rate of net expulsion of gas from stars within a mass range whose upper limit becomes an independent variable in the modeling scheme.This assumption leads to the delayed production of Li in the models,in good agreement with the observations.
The concept of late-time production of Li is not,of course,newly introduced within the present model,even though its use with a GCR source here is original.Other mechanisms for producing Li associated predominantly with lower mass objects might in principle be able to satisfy the observational constraints of the Li-Fe envelope of Fig. 1,and a number of these have been suggested.Sites which have been put forward as serious candidates for a major fraction of galactic Li production include novae,neutron stars,stellar mass black holes, red giants, AGB stars, and carbon stars,all of which might satisfy the constraint of delaying the Li production with respect to that of Fe and also processes in supernovae,which satisfy this condition with greater or lesser degree of difficulty.It is not possible to dismiss these sources but the fact that they each fare some difficulty:theoretical or observational,tends to reinforce our view that GCR rather than stellar production of Li has in fact predominated.
Results of earlier work by Arnould $`\&`$ Norgaard (1975) on novae,followed by the more detailed study of Starfield et al. (1978) have been more recently called into question by Boffin,Paulus $`\&`$ Arnould (1993) using new reaction cross-sections.The latter authors conclude that it is much more difficult to produce a significant quantity of Li in novae than previously predicted.A search for lithium in late-type companions of several dwarf and classical novae has not yielded detections (Martin et al. (1995))
Measurements of Li in the youngest stars,in some of whose atmospheres production has been postulated to occur,tend to yield abundances close to the ”canonical” value for moderately young stars,of log N (Li) $``$ 3.2 or 3.3,and do not seem to show sufficient Li to be strong candidates for major galactic production.Martin et al. (1992) showed that in some cool companions of hot (i.e. young) stars,where simple atmospheric model analysis could appear to show Li abundances of up to 3.7,a more careful NLTE study yields upper values of 3.4,with most objects falling below this.Similar results have been obtained by Duncan (1991),by Magazzu et al. (1992) and by Martin et al. (1992) for T Tauri stars.Here again careful analysis reduced apparently very high values of the Li abundance in some objects to values well within the range of the normal young stellar population.
There do appear to be giants especially over-abundant in Li,notably a sub-set of the C-stars (Abia et al. (1991)),and a fraction of normal K-giants (De la Reza & Da Silva (1992)).Observations here are still few,and conclusions made somewhat more difficult by the convective tendency of giants to deplete Li (this of course strengthens any argument in favor of such stars being Li sources if strong Li absorption lines are seen in their spectra).In the most extensive sets of observations of field giants,however,very few indeed have particularly strong Li abundances (Brown et al. 1989).Thus if certain types of giants are important Li producers,since they are few,and therefore need to be strong sources,while sparsely distributed,one might expect more scatter in population I Li than is observed.Nevertheless production in giants remain a possible source,and the main argument we offer against its being the main source can only be that the GCR model presented is able to account for the relevant observations without a major extra Li contribution.
Similar consideration may be given to the postulated importance of Li production in late-type companions to neutron stars and black-hole candidates.In a paper on these objects by Martin et al. (1994),Li abundances ranging up to 3.3 for Cen X-4 are detected which might be increased if there is substantial overionization of LiI due to UV and X-ray flux coming from the compact object.These authors claim that since Li undergoes depletion by convection in late-type stars,the presence of relatively high Li abundances in these objects marks them as Li producers,and therefore candidates for major enrichment of the galactic disk.Here again while we see no immediate argument which can rule out this possibility one may doubt that there are sufficient such sources.Martin et al. (1994) put forward the idea that there may well have been more X-ray binaries,especially high-mass binaries,in the past but the curve of Li as a function of Fe implies that the production mechanism should not be associated with high mass objects.Nevertheless we are not in a position here to claim that processes in X-ray binaries cannot be responsible for a significant part of galactic Li production,only that these appears to be no requirement for this as a major source.
In all the analytical approximations based on proportionality of Li production rate to the gas expulsion rate by stars of low or intermediate masses (equations (12) (14)) one can see considerable resemblance to the observed lithium growth profile in the zone of interest, i.e., with metallicities \[Fe/H\] between -1.0 and 0.0 (see Fig. 1 and Fig. 3).
We have used the $`\alpha `$+$`\alpha `$ process in the ISM as illustrative of processes which follow the behaviour of stars of intermediate and low masses, those whose lifetimes are long, and whose numbers in the disk have therefore grown cumulatively with time.Any process with equivalent time-dependent characteristics could, in general terms, satisfy this global observational constraint.Production of Li in flares of red giants (Cameron $`\&`$ Fowler 1971), or in novae (Arnould $`\&`$ Norgaard 1979,Starrfield et al. 1978), could also satisfy, in principle, the criterion of delayed production because there the rate of production would be proportional to the number of stars with intermediate masses at each time (cumulative because of their long lifetimes) and not to the gas expulsion rate.But, as one can see in Fig. 4, where we plot the results from numerical model in the case of GCR flux proportional to the cumulative number of stars with masses less than or equal to 1M (increasing this limit yields a decreasing rate of accumulation of Li producers), the prediction, although not too bad, is by no means as good as that obtained using gas expulsion by stars of masses less than or equal to 3M, and is also in fact worse than that obtained using gas expulsion of stars with masses less than or equal to 2M (see Fig. 4).
As far as mechanisms which depend on the presence of SNI are concerned, two arguments appear to weaken their claims.One is that the Li-Fe envelope in Fig. 1 is not linear, which would be the dependence if the Li were either produced directly by the impact of SNIe on their immediate surroundings, or by processes involving GCR in a wider volume of space, produced by SNIe.The other is that the locally measured GCR abundances favour an origin, for the lower energy particles at least, in the thermal equilibria pertaining in the atmospheres of stars of moderate mass, rather than in the extreme conditions of a supernova.However here again we argue in terms of probabilities rather than claming that this mechanism is excluded.
Another possibility is the production of Li in compact objects such as neutron stars and black holes, but the time evolution of compact objects would be proportional to the gas expulsion rate from all stars (which is mainly that of SNI+SNII).This production is shown in Fig. 1 from numerical modelling, and one can see how this implies overabundance of Li with respect to the observational constraints; in fact its tendency is not too different from that of those GCR models where the flux is proportional to the SFR.
Other possibilities for production of Li, such as in AGB stars or carbon stars (Matteucci et al. (1995), Romano (1999)) are shown not to reproduce the observations really as well as the delayed models, as one can see, for example, in Fig. 4.
## 7 Conclusions.
We have surveyed mechanisms of Li production in the disk, and confronted them with the upper envelope of the Li-Fe observations, which we have taken to represent the Li-Fe evolution curves in the absence of stellar depletion for the individual objects observed.As a result we can conclude that:
1)Mechanisms relying on SNIIe to produce Li cannot be excluded in at least an approximate explanation of the observations.This is true for production within the SNe themselves, but also holds for GCR fluxes originating in SNIIe.
2)Mechanisms whose time dependence is that of the SFR give either too much Li in the early disk or too little in the later disk.
3)Mechanisms which rely on SNIe to produce the Li (again either in the immediate surroundings of the SNa or via a more generally dissipated GCR flux originating in SNIe) predict that the disk Li should grow proportionally to Fe, which does not appear to fit the observations.
4)An attempt to reproduce the results on the assumption that the contemporary maximum abundance, logN(Li)$``$3.4 is the true primordial abundance, and that the Li-Fe envelope is a pure depletion curve also fails by a wide margin.
5)Mechanisms which produce an increase in disk Li significantly delayed with respect to that of Fe can explain the observations very well.
We have in this article explored one such mechanism: the production of Li via $`\alpha `$+$`\alpha `$ fusion reaction in the ISM due to low energy cosmic rays whose source of origin is the atmospheres of low and intermediate mass stars.This mechanism has the virtue that these stars have lifetimes comparable with that of the disk, so that their collective gas expulsion rate has accumulated progressively throughout the disk lifetime, leading automatically to a delay with respect to Fe in the Li production curve.We have explained that even if the acceleration of the GCR is due to SNe envelopes, the product of injection rates and acceleration rates retains the delay implied by the observations (further work on acceleration mechanisms such as that due to stellar wind termination shocks is, however, well worth exploring in this context).Support for the possibility of this mechanism is provided by the observed similarity between the GCRS composition and that of the solar corona which is biased according to the first ionization potential, and we note in this context the statement of Ellison et al. (1997) that ”in the outer solar atmosphere the solar coronal gas, the solar wind, and the $``$MeV solar energetic particles have undoubtedly a composition biased according to FIP”, together with the fact that the hydrogen and precisely helium are not well fitted by the alternative model of Meyer et al. (1997) and Ellison et al. (1997) based on volatility and mass to charge to explain the GCRS.These considerations permit an origin in an environment close to thermal equilibrium, i.e. typical of stars of moderate mass.We have incorporated the mechanism in an evolutionary model of the disk previously demonstrated to be capable of accounting well for the Be and B vs. Fe observations (Casuso & Beckman (1997)), and which gives a particularly good account of the G-dwarf metallicity distribution in the solar neighborhood.The resulting Li-Fe plots include very fair fits to the observed Li-Fe envelope.
We have included in this scenario a natural mechanism of differential depletion (Casuso & Beckman 1999) operating within red supergiant envelopes, which can account for the observed D/H v. time and isotopic ratios of <sup>7</sup>Li/<sup>6</sup>Li and <sup>11</sup>B/<sup>10</sup>B v. time.
However we would not at this stage wish to rule out the possibility of other mechanism or mechanisms for disk lithium production.The observational weight of the stellar Li abundances, as we have shown, does place some strong constraints on Li-production models.One of the clearest conclusions we can draw is that the ”high” value log N(Li)$``$3.4 for the primordial Li abundance can be quantitatively rejected using the Li-Fe observational constraint.The assignation of a value close to the ”Spite plateau” (Spite & Spite (1982)) value: log N(Li)$``$2.2 as primordial is thereby strengthened.In this context the comprehensive study by Thorburn (1994) of Li in halo stars, in which a contribution to the plateau produced by the $`\alpha `$+$`\alpha `$ reaction due to the halo GCR flux is shown to account well for the observed scatter and slight rise in the Li abundance below \[Fe/H\]=-1.5, makes a suggestive link with the disk model tested in the present paper.The importance of the $`\alpha `$+$`\alpha `$ process has almost certainly been previously underestimated in the disk, and the powerful constraint on evolutionary processes and models implied by the Li vs. Fe observations has not been adequately taken into account; it is these aspects of the lithium puzzle which the present paper has been designed to expose.
Note added in Proof: Newly observations of Li and <sup>7</sup>Li/<sup>6</sup>Li in ISM (toward o Per and $`\zeta `$ Per) by Knauth, Federman, Lambert, Crane (Nature in press), give a variation in <sup>7</sup>Li/<sup>6</sup>Li ratio (from near 2 which is the expected for Li production from spallation or alpha-alpha fusion reactions purely, to near 11 which is very similar than that of solar value), together with very similar reported values for Li/H abundance (near 11x10<sup>-10</sup>) for the two clouds in contrast with the solar value of 20x10<sup>-10</sup>).Also, the two clouds are near the star forming region IC 348. All of these data agree very well with our picture of production of light elements in the ISM via GCRs (Be,B) (Casuso and Beckman 1997) and via alphas of low-energy (Li).We explained this variation (in fact a fall off) via a model in which the envelopes of red-supergiant stars (so, star forming region) deplete differentially <sup>6</sup>Li and <sup>7</sup>Li , and the increasing infall of non-depleted gas with time (Casuso and Beckman 1999).And also, we explained in the present article the decay on Li/H abundance from solar to actual ISM due precisely to the depletion in star forming regions in addition with the infall of non-enriched gas (see Fig. 4).So, we can explain these data without the problem inh erent to the explanation by Knauth et al., which point to the differential production of Li in the o Per direction and in the $`\zeta `$ Per direction because of the higher flux of cosmic rays in the o Per direction, while observations point to almost the same total Li/H abundance.
Acknowledgments: We are happy to thank F. Spite for supplying his lithium abundance data compilation, and for helpful suggestions, and E.L. Martin for useful discussions.The anonymous referee made a number of valuable suggestions which led to significant improvements in the paper.This research was supported in part by grant PB97-0219 of the Spanish DGICYT.
Figure Captions
Fig. 1 a) Abundance of Li v. iron metallicity \[Fe/H\] in the solar neighborhood.Compilation from the work of the group of the present authors (Rebolo,Beckman,Molaro,1987,Rebolo,Molaro and Beckman,1988), and from Spite (1996).A conservative error bar is shown in the left upper corner.The upper envelope characterizes the stars least depleted in Li at the epochs implied by each value of metallicity.Stars below the upper envelope have suffered notable internal Li depletion. In the present article we model only the envelope for the disk,i.e. for \[Fe/H\] $``$ -1.5. b) Disc production curves of Li for three models in which the low energy GCR flux responsible for the Li,via $`\alpha `$+$`\alpha `$,is proportional to:(i)The star formation rate (dotted line),(ii)The gas expulsion rate for the population of stars with all masses (dashed line),(iii)The gas expulsion rate for the population of stars with masses $``$ 3M (full line).Each curve is normalized to give the observed Li abundance of 2.4 at \[Fe/H\]=-1.3.The curves are shown in comparison with the data.The axes in this figure, and in Figs. 3, 4 and 6, are: for \[Fe/H\] the logarithmic abundance with the solar value \[Fe/H\]=0, for Li the logarithmic abundance where logN(H)=12.
Fig. 2 a) Iron metallicity in the solar neighborhood as a function of age.The curve show the prediction of our model.Observation comparison points with error bars are from: Meusinger et al. (1991) (crosses) and Twarog (1986) (crossed circles).b) The ratio oxygen/iron as a function of the iron abundance.\[O/Fe\] v. \[Fe/H\] data are from: Rebolo et al. (1994),Nissen et al. (1994),Israelian et al. (1998).The full line corresponds to our chemical evolution model using the yields of Fe theoretically calculated for stars with very low metallicities, and the dotted line is for yields from stars with intermediate metallicities.However, in the present paper we are critically concerned with iron metallicities greater than -1.5 where there are no major problems with the dispersion of data.
Fig. 3 Normalized plots based on models embodying analytical schematic approximations to Li production rates and their evolution with Fe, to illustrate how some selected scenarios for the evolution of disk Li are much less in agreement with the form of the observations (the upper envelope of the points in Fig. 1a), than are others.Points are the data as in Fig. 1.All curves are from analytic approximations by mathematical functions with two free parameters which are constrained to match the data envelope at \[Fe/H\]$``$-1.0 (where the full disk initiates) and at \[Fe/H\]$``$0.0.Long dashed line: from eq. (17).Full line: from eqs. (6) and (19) which give a very similar result.Dashed line: from eq. (14).Dotted line: from eq. (12).
Fig. 4 a) through d).Observations of Li vs. Fe compared with the predictions of models taken from the literature, with parameters described in the text, and also with models developed in the present paper.Graphs common to all panels: Model due to Prantzos et al. (1993): dashed plus dotted line.Model due to Matteucci et al. (1995): dotted line.Halo model due to Casuso & Beckman (1997): solid line for \[Fe/H\]$``$-1.3 only.Differential curves: a) $`\alpha `$+$`\alpha `$ model with GCR flux $``$ gas expulsion rate by stars of masses $``$3M (full line); model with GCR flux $``$ gas expulsion rate by stars of masses $``$2M (dashed line). b) $`\alpha `$+$`\alpha `$ model with GCR flux $``$ gas expulsion rate by stars with masses $``$3M, with exponentially increasing infall (full line), with no infall (dashed line). c) $`\alpha `$+$`\alpha `$ model with GCR flux $``$ gas expulsion rate by stars with masses $``$3M (full line); with GCR flux $``$ cumulative number of stars with masses $``$1M (long dashed + dotted line). d) $`\alpha `$+$`\alpha `$ model with GCR flux $``$ gas expulsion rate by stars with masses $``$3M (full line); model with no disk production but with linearly time-dependent stellar depletion from a ”primordial” value of log N(Li)$``$3.5 (long and short dashed line).
Fig. 5 Example showing the essential difference in GCR flux as a function of time for two key models: proportional to SFR(t) (full line) and proportional to the gas expulsion rate by stars with masses less than or equal to 3M (dotted line) against time, from our numerical model with increasing infall of gas to the solar neighborhood.The steps in the curves are unsmoothed constructs of the model due to the finite time intervals employed.
Fig. 6 Extrapolated model curve which can account for Li abundances observed with \[Fe/H\] greater than 0.1.
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# On the Asymptotics of Takeuchi Numbers
## 1 Introduction
In a paper entitled “Textbook Examples of Recursion,” Donald E. Knuth discusses recurrence equations related to the properties of recursive programs \[Knuth 1991\], among them Takeuchi’s function \[Takeuchi 1978, Takeuchi 1979\]
$$t(x,y,z)=\text{if }xy\text{ then }y\text{ else }t(t(x1,y,z),t(y1,z,x),t(z1,x,y)).$$
(1)
Let $`T(x,y,z)`$ denote the number of times the else clause is invoked when $`t(x,y,z)`$ is evaluated recursively.<sup>1</sup><sup>1</sup>1Note that it is the recursive evaluation of $`t(x,y,z)`$ rather than the actual value of $`t(x,y,z)`$ that is of interest. See Knuth’s paper for an explicit expression of $`t(x,y,z)`$. For non-negative integers $`n`$, the Takeuchi numbers $`T_n`$ are defined as $`T_n=T(n,0,n+1)`$. The first few values of $`T_n`$ for $`n=0,1,2,\mathrm{}`$ are
$$0,\mathrm{\hspace{0.33em}1},\mathrm{\hspace{0.33em}4},\mathrm{\hspace{0.33em}14},\mathrm{\hspace{0.33em}53},\mathrm{\hspace{0.33em}223},\mathrm{\hspace{0.33em}1034},\mathrm{\hspace{0.33em}5221},\mathrm{\hspace{0.33em}28437},\mathrm{\hspace{0.33em}165859},\mathrm{}.$$
(2)
Knuth gives the recurrence
$$T_{n+1}=\underset{k=0}{\overset{n}{}}\left\{\left(\genfrac{}{}{0pt}{}{n+k}{n}\right)\left(\genfrac{}{}{0pt}{}{n+k}{n+1}\right)\right\}T_{nk}+\underset{k=1}{\overset{n+1}{}}\left(\genfrac{}{}{0pt}{}{2k}{k}\right)\frac{1}{k+1}\text{}n0\text{,}$$
(3)
and deduces a functional equation for the generating function $`T(z)=_{n=0}^{\mathrm{}}T_nz^n`$:
$$T(z)=\frac{C(z)1}{1z}+\frac{z(2C(z))}{\sqrt{14z}}T(zC(z)),$$
(4)
where
$$C(z)=\frac{1}{2z}(1\sqrt{14z})=\underset{n=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\frac{1}{n+1}$$
(5)
is the generating function for the Catalan numbers $`C_n=\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\frac{1}{n+1}`$. Lastly, he gives asymptotically valid bounds for $`T_n`$,
$$e^{n\mathrm{log}nn\mathrm{log}\mathrm{log}nn}<T_n<e^{n\mathrm{log}nn+\mathrm{log}n}\text{for all sufficiently large }n\text{,}$$
(6)
and poses obtaining further information about the asymptotic properties of $`T_n`$ as an open problem.
In this paper I give arguments leading to two conjectures about the asymptotic behaviour of the Takeuchi numbers $`T_n`$. In Section 2, I present an explicit asymptotic formula for $`T_n`$ which improves upon the bounds (6) based on numerical evidence and a heuristic argument. For this, I briefly discuss the related asymptotic behavior of the Bell numbers and give an argument based on a numerical observation which leads directly to an explicit asymptotic formula for $`T_n`$ as $`n`$ tends to infinity. The formula, as described in Conjecture 1, is exact up to $`O((\mathrm{log}n/n)^2)`$ and contains a constant $`C_T`$ which is numerically determined to 25 significant digits. Section 3 presents a heuristic analytic argument which gives the asymptotic behavior up to $`o(1)`$ and enables one to identify the constant $`C_T`$ in terms of an explicit expression, as stated in Conjecture 2. In the final section I show that the method developed in Section 3 can give insight into the asymptotic behavior of a larger class of problems.
I conclude this introduction by briefly discussing the structure of the recurrence (3) and the related functional equation (4). It is clear from the asymptotic bounds (6) for $`T_n`$, that the generating function $`T(z)`$, defined as a formal power series, does not converge, and is therefore at best only an asymptotic expansion to an actual solution of the functional equation. This is also evident from the structure of the functional equation. This structure becomes clearer upon a change of variables,
$$T(z)=\frac{1}{z}T(zz^2)\frac{1}{(1z)(1z+z^2)},$$
(7)
where one sees directly that the transformation involved is $`g(z)=zz^2`$, which is only marginally contracting at its fixed point $`z=0`$. While functional equations with a transformation $`g(z)`$ that has an expanding or contracting fixed point ($`|g^{}(0)|1`$) are very well understood \[Kuczma 1990\], it is precisely the fact that $`|g^{}(0)|=1`$ which is at the root of the underlying difficulty of the problem discussed in this paper.
## 2 Numerical Observations
Our starting point is Knuth’s observation \[Knuth 1991\] that for $`n>0`$ the Bell numbers $`B_n`$ are a lower bound to $`T_n`$. Here, $`B_n`$ is defined as
$$B_{n+1}=\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)B_{nk},B_0=1.$$
(8)
The asymptotics of $`B_n`$ is discussed in great detail by de Bruijn \[DeBruijn 1961\]. In \[Moser 1955\] one finds a systematic way of generating higher order terms in the asymptotic expansion by means of a contour integral representation using the well-known fact that
$$\underset{n=0}{\overset{\mathrm{}}{}}B_n\frac{z^n}{n!}=\mathrm{exp}(e^z1).$$
(9)
Alternatively, one can also expand the right-hand side of (9) in $`z`$ to get
$$B_n=\frac{1}{e}\underset{m=0}{\overset{\mathrm{}}{}}\frac{m^n}{m!},$$
(10)
which can then be evaluated asymptotically by use of the Euler-MacLaurin formula. Either way, in the course of the computation of the asymptotics of $`B_n`$ it turns out that a convenient asymptotic scale is given in terms of $`W(n)`$ rather than in terms of $`n`$, where $`W(x)`$ is Lambert’s $`W`$ function, which is defined as the real solution of
$$W(x)\mathrm{exp}W(x)=x.$$
(11)
The sum (10) is dominated by terms around $`m=e^{W(n)}`$, and one can easily calculate that, written in terms of $`w=W(n)`$, the Bell numbers $`B_n`$ behave asymptotically as
$`\mathrm{log}B_n`$ $`=`$ $`e^w(w^2w+1){\displaystyle \frac{1}{2}}\mathrm{log}(1+w)1{\displaystyle \frac{w(2w^2+7w+10)}{24(1+w)^3}}e^w`$
$`{\displaystyle \frac{w(2w^4+12w^3+29w^2+40w+36)}{48(1+w)^6}}e^{2w}+O(e^{3w}),`$
and it is straightforward to calculate additional terms. It is more conventional to state this formula with the exponentials $`e^{kw}`$ replaced by $`(n/w)^k`$, but this obscures the fact that the asymptotic expansion is obtained in terms of $`w`$ rather than $`n`$.
Having rather explicit control over this lower bound, it is natural to now try to compare $`B_n`$ and $`T_n`$ more closely. There is a principal difficulty coming from the fact that the asymptotic scale presumably also involves $`w=W(n)`$, which grows more slowly than $`\mathrm{log}n`$. (As an example, $`W(1000)5.2496`$ and $`W(10000)7.2318`$.) Thus, one would expect that a direct numerical investigation of $`T_n/B_n`$ is not very insightful, due to the presence of slowly varying correction terms of unknown form.
However, as Figure 1 shows, if one compares the growth rates $`T_n/T_{n1}`$ and $`B_n/B_{n1}`$ instead, one is led to observe the surprisingly simple relationship
$$\underset{n\mathrm{}}{lim}\left(\frac{T_{n+1}}{T_n}\frac{B_n}{B_{n1}}\right)=1.$$
(13)
In fact, the left hand side approaches $`1`$ rather quickly,
$$\frac{B_n}{B_{n1}}+1\frac{T_{n+1}}{T_n}\frac{B_n}{B_{n1}}+1+O(e^w).$$
(14)
This (unproven) numerical observation leads to a straightforward derivation of an asymptotic formula. From (2) it follows easily that $`B_{n1}/B_n=e^w+O(e^{2w})`$. Now one takes logarithms and sums up successively, from whence it follows that
$$\mathrm{log}T_{n+1}=\mathrm{log}B_n+\frac{1}{2}w^2+w+O(1).$$
(15)
Now that I have guessed the leading asymptotic form, I can again resort to numerical work to try to improve upon it. In fact, numerically it appears that the convergence is even better than expected due to a chance cancellation of higher order correction terms. Figure 2 shows that
$$T_{n+1}=C_TB_n\mathrm{exp}(\frac{1}{2}w^2+w+O(e^{2w})),$$
(16)
and from the first $`1000`$ series terms I am able to deduce by iterative application of standard series extrapolation methods that
$$C_T=\mathrm{2.23943\hspace{0.33em}31040\hspace{0.33em}05260\hspace{0.33em}73175\hspace{0.33em}4785}(1).$$
(17)
Using the known asymptotic form of the Bell numbers, I can now give an explicit asymptotic expression for $`T_n`$ in terms of $`w=W(n)`$ alone, as stated in the following conjecture.
###### Conjecture 1
As $`n`$ tends to infinity, one has
$`\mathrm{log}T_n`$ $`=`$ $`e^w(w^2w+1)+{\displaystyle \frac{1}{2}}w^2{\displaystyle \frac{1}{2}}\mathrm{log}(1+w)+`$
$`+\mathrm{log}C_T1{\displaystyle \frac{w(26w^2+67w+46)}{24(1+w)^3}}e^w+O(e^{2w}).`$
Here $`w=W(n)`$ and the constant $`C_T`$ is some positive real number.
Numerically, $`\mathrm{log}C_T1=\mathrm{0.19377\hspace{0.33em}72447\hspace{0.33em}31916\hspace{0.33em}75890\hspace{0.33em}1157}(1)`$. Of course it would be desirable to find an analytic expression for this number. In the next section I shall present a heuristic argument giving such an expression.
Dropping one correction term and comparing with the asymptotic expression (2) for the Bell numbers, Conjecture 1 implies the nice formula given in the abstract,
$$T_nC_TB_n\mathrm{exp}\frac{1}{2}W(n)^2.$$
(19)
## 3 Analytic Results
In view of the previous section it seems promising to exploit the apparent affinity between Takeuchi numbers $`T_n`$ and Bell numbers $`B_n`$. Given a recurrence of the general form
$$a_n=\underset{k=1}{\overset{n}{}}c_{n,k}a_{nk}+b_n,$$
(20)
I choose to write
$$a_n=\frac{1}{e}\underset{m=0}{\overset{\mathrm{}}{}}\frac{1}{m!}f_{m,n}\text{and}b_n=\frac{1}{e}\underset{m=0}{\overset{\mathrm{}}{}}\frac{1}{m!}b_n.$$
(21)
Inserting (21) into the recurrence (20) and shifting the summation index by one, I next equate terms to get
$$f_{m,n}=m\underset{k=1}{\overset{n}{}}c_{n,k}f_{m1,nk}+b_n,$$
(22)
This Ansatz might seem less arbitrary when considering that in the case of Bell numbers it reduces to $`f_{m,n}=m^n`$. In general, one observes that $`f_{m,n}`$ must be a polynomial in $`m`$ of at most degree $`n`$, which I write as
$$f_{m,n}=m^n\underset{k=0}{\overset{n}{}}d_{n,k}m^k.$$
(23)
If one further requires the coefficients $`c_{n,k}`$ in the recurrence to be polynomials of degree $`k`$ in $`n`$, it follows that $`d_{n,k}`$ are polynomials of degree $`k`$ in $`n`$, which I write as
$$d_{n,k}=n^k\underset{l=0}{\overset{k}{}}r_{k,l}n^l.$$
(24)
Combining equations (23) and (24) gives
$$m^nf_{m,n}=\underset{k=0}{\overset{n}{}}\underset{l=0}{\overset{k}{}}r_{k,l}n^{kl}m^k=\underset{l=0}{\overset{n}{}}m^l\underset{k=0}{\overset{nl}{}}r_{l+k,l}(n/m)^k.$$
(25)
In order to get an idea about the asymptotic behavior of this double sum, I now replace the quotient $`n/m`$ by a new variable $`v`$ and consider the formal limit of taking the summation bounds to infinity, leading to
$$s_m(v)=\underset{l=0}{\overset{\mathrm{}}{}}m^lr_l(v)\text{with}r_l(v)=\underset{k=0}{\overset{\mathrm{}}{}}r_{l+k,l}v^k.$$
(26)
Applying this method to Takeuchi numbers, one now inserts $`c_{n,k}=\left\{\left(\genfrac{}{}{0pt}{}{n+k2}{n1}\right)\left(\genfrac{}{}{0pt}{}{n+k2}{n}\right)\right\}`$ and $`b_n=_{k=1}^n\left(\genfrac{}{}{0pt}{}{2k}{k}\right)\frac{1}{k+1}`$. With this choice, $`r_0(v)`$ is trivially zero as $`c_{n,k}`$ are polynomials in $`n`$ of degree $`k1`$, and one gets a rather interesting result for $`l1`$. In fact,
$`r_1(v)`$ $`=`$ $`e^{\frac{1}{2}v^2+v}`$ (27)
$`r_2(v)`$ $`=`$ $`e^{\frac{1}{2}v^2}\left(2e^{2v}{\displaystyle \frac{1}{2}}(v^3+v^2+4v+2)e^v\right)`$ (28)
$`r_3(v)`$ $`=`$ $`e^{\frac{1}{2}v^2}({\displaystyle \frac{1}{8}}e^{3v}(v^3+3v^2+7v+6)e^{2v}+`$
$`+{\displaystyle \frac{1}{24}}(3v^6+6v^5+47v^4+52v^3+144v^2+74v+51)e^v)`$
$`r_4(v)`$ $`=`$ $`e^{\frac{1}{2}v^2}({\displaystyle \frac{347}{108}}e^{4v}+{\displaystyle \frac{1}{16}}(v^3+5v^2+12v+12)e^{3v}+`$
$`+{\displaystyle \frac{1}{12}}(3v^6+18v^5+89v^4+226v^3+411v^2+406v+195)e^{2v}`$
$`{\displaystyle \frac{1}{432}}(9v^9+27v^8+315v^7+603v^6+3024v^5+`$
$`+3384v^4+8757v^3+4707v^2+5484v+772)e^v).`$
From this, I conjecture
$$r_l(v)=e^{\frac{1}{2}v^2}\left(p_{l,0}(v)e^{lv}+p_{l,1}(v)e^{(l1)v}+\mathrm{}+p_{l,l1}(v)e^v\right),$$
(31)
where $`p_{l,k}(v)`$ are polynomials in $`v`$ of degree $`3k`$. (This pattern has been verified for $`l8`$.) For $`v`$ large, this conjecture implies
$$r_l(v)\lambda _le^{\frac{1}{2}v^2+lv},$$
(32)
and with a little effort one can compute the next values of $`\lambda _l=p_{l,0}`$. One gets
$$\lambda _0=0,\lambda _1=1,\lambda _2=2,\lambda _3=\frac{1}{8},\lambda _4=\frac{347}{108},\lambda _5=\frac{28201}{3456},$$
$$\lambda _6=\frac{3172987}{216000},\lambda _7=\frac{822813607}{93312000},\lambda _8=\frac{2183235065857}{16003008000},\mathrm{},$$
from whence I am led to conjecture that one can write $`\lambda _n=\mu _n/[(n1)!]^3`$ where $`\mu _n`$ is integer. (Unfortunately, I did not find a way to compute the coefficients $`\lambda _l`$ in a closed form!) One would expect from this behavior that
$$s_m(v)e^{\frac{1}{2}v^2}h(e^v/m),$$
(33)
where $`h(x)`$ is given in some sense by
$$h(x)=\underset{l=0}{\overset{\mathrm{}}{}}\lambda _lx^l.$$
(34)
I caution here that the series may be divergent and just valid as an asymptotic expansion.
Keeping in mind that the evidence for the existence of $`h(x)`$ is rather sketchy, I nevertheless proceed under the assumption that for $`n>>m>>1`$ one can write
$$f_{m,n}m^ne^{\frac{1}{2}(n/m)^2}h(e^{n/m}/m).$$
(35)
This now enables a heuristic computation of the Takeuchi numbers $`T_n`$. I approximate
$`T_n`$ $`=`$ $`{\displaystyle \frac{1}{e}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}f_{m,n}{\displaystyle \frac{1}{e}}{\displaystyle \underset{mm_{max}(n)}{}}{\displaystyle \frac{1}{m!}}f_{m,n}`$ (36)
$``$ $`{\displaystyle \frac{1}{e}}{\displaystyle \underset{mm_{max}(n)}{}}{\displaystyle \frac{m^n}{m!}}e^{\frac{1}{2}(n/m)^2}h(e^{n/m}/m),`$ (37)
where in the last step it is assumed that $`n>>m_{max}(n)>>1`$. This sum is indeed dominated around $`m_{max}(n)e^w`$, where the argument of $`h`$ simplifies to $`1`$. A careful asymptotic analysis of
$$\widehat{T}_n=\frac{1}{e}\underset{mm_{max}(n)}{}\frac{m^n}{m!}e^{\frac{1}{2}(n/m)^2}h(e^{n/m}/m)$$
(38)
gives
$`\mathrm{log}\widehat{T}_n`$ $`=`$ $`e^w(w^2w+1)+{\displaystyle \frac{1}{2}}w^2{\displaystyle \frac{1}{2}}\mathrm{log}(1+w)+h_01`$
$`+{\displaystyle \frac{w(12w^5+24w^4+36w^3+58w^2+29w10)}{24(w+1)^3}}e^w`$
$`+{\displaystyle \frac{(w+1)(h_1^2+h_2)+(2w^2+w+2)h_1}{2}}e^w+O(e^{2w}),`$
where one has expanded $`h(x)`$ around $`x=1`$ as $`\mathrm{log}h(x)=h_0+h_1(x1)+h_2(x1)^2/2+O((x1)^3)`$.
As long as the corrections made on passing from $`T_n`$ to $`\widehat{T}_n`$ are small enough, it follows easily from this that asymptotically
$$T_nB_ne^{\frac{1}{2}w^2}h(1),$$
(40)
and one can identify the constant $`C_T`$ from equation (19) with $`h(1)`$. Provided the series expansion of $`h(x)=_{k=0}^{\mathrm{}}\lambda _kx^k`$ converges at $`x=1`$, I can thus conjecture an explicit expression for the constant $`C_T`$, which in principle is computable.
###### Conjecture 2
The constant $`C_T`$ in Conjecture $`1`$ is given by
$$C_T=h(1)=\underset{k=0}{\overset{\mathrm{}}{}}\lambda _k.$$
(41)
While the approximation of $`T_n`$ by $`\widehat{T}_n`$ may be correct up to $`O(e^w)`$, no choice of $`h(x)`$ can match the next term in (3) with the expansion of $`T_n`$. Thus, one also gets an indication of the size of the error made.
It seems that a careful asymptotic evaluation of the $`f_{m,n}`$ promises to be a suitable way of providing rigorous proof for the asymptotics of the Takeuchi numbers. Of course one could also try to find a direct proof of our numerically observed equation (13).
## 4 A Generalization
In the derivation of the functional equation (4) for the Takeuchi numbers $`T_n`$, it is crucial that
$$\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{n+2k}{k}\right)z^k=C(z)^k/\sqrt{14z},$$
(42)
as this identity allows the explicit summation of the terms in the recurrence (3). The identity used is a special case of the following nice identity
$$\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{n+(\lambda +1)k}{k}\right)z^k=\left\{\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{(\lambda +1)k}{k}\right)z^k\right\}\left\{\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{(\lambda +1)k}{k}\right)\frac{z^k}{1+\lambda k}\right\}^n.$$
(43)
This identity can be proved by inserting $`z=y/(1+y)^{\lambda +1}`$, which after expanding leads to
$$\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{(\lambda +1)k}{k}\right)\frac{z^k}{1+\lambda k}=1+y$$
(44)
and
$$\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{n+(\lambda +1)k}{k}\right)z^k=\frac{(1+y)^{n+1}}{1\lambda y}.$$
(45)
I use this now as a motivation for the study of the family of recursions (with parameter $`\lambda `$)
$$A_{n+1}=\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n+\lambda k}{k}\right)A_{nk},A_0=1.$$
(46)
Due to equation (43) one is able to derive a functional equation for the corresponding generating function $`A(z)=_{n=0}^{\mathrm{}}A_nz^n`$:
$$A(z)=1+z\frac{1+y}{1\lambda y}A(z(1+y)),z=y/(1+y)^{\lambda +1}.$$
(47)
For $`\lambda =0`$ one recovers the recursion for the Bell numbers, and for $`\lambda =1`$ one has something which is at least “morally” related to the Takeuchi numbers.
Inserting the Ansatz (21) into (46), one can easily repeat the analysis of the previous section. The result is now
$$A_nB_n\mathrm{exp}\lambda \left\{\frac{1}{2}W(n)^2+W(n)+d(\lambda )\right\}$$
(48)
for any fixed value of $`\lambda `$. Again, one have an identification of the kind $`d(\lambda )=h_\lambda (1)`$, where the first terms in the series expansion of $`h_\lambda (x)`$ are
$`h_\lambda (x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\lambda 1)x{\displaystyle \frac{1}{24}}(2\lambda ^2+18\lambda 5)x^2{\displaystyle \frac{1}{216}}(33\lambda ^3+90\lambda ^2329\lambda +54)x^3`$ (49)
$``$ $`{\displaystyle \frac{1}{960}}(52\lambda ^4520\lambda ^3+4240\lambda 502)x^4+O(x^5),`$
and one sees that the $`k`$th coefficient is a polynomial in $`\lambda `$ of degree $`k`$ (this has been verified up to $`k=7`$). I caution again that convergence of this series expansion is an open question.
Finally, one can establish numerically the next term in the asymptotic expansion of $`A_n`$. For any fixed value of $`\lambda `$, one finds
$$\mathrm{log}A_n=\mathrm{log}B_n+\lambda \left(\frac{w^2}{2}+w+d(\lambda )\frac{\lambda +1}{2}e^w\right)+O(e^{2w}).$$
(50)
Indeed, this result even seems to hold for complex values of $`\lambda `$.
I conclude with remarking that even though Takeuchi’s function has been labelled a “Textbook Example,” it provides an exciting open question for asympotic analysis.
## Acknowledgements
I thank Philippe Flajolet for bringing this problem to my attention.
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# Use of harmonic inversion techniques in the periodic orbit quantization of integrable systems
## 1 Introduction
A question of fundamental interest for systems with both regular and chaotic dynamics is how quantum mechanical eigenvalues can be obtained by quantization of classical orbits. The EBK torus quantization method of Einstein, Brillouin, and Keller Ein17 ; Bri26 ; Kel58 is restricted to integrable systems, i.e., the method cannot be generalized to systems with an underlying chaotic or mixed regular-chaotic dynamics Ein17 . Furthermore, EBK quantization requires the knowledge of all the constants of motion, which are not normally given in explicit form, and therefore practical EBK quantization based on the direct or indirect numerical construction of the constants of motion turns out to be a formidable task Per77 . As an alternative, EBK quantization was recast as a sum over all periodic orbits of a given topology on respective tori by Berry and Tabor Ber76 . In contrast to EBK-quantization, periodic orbit theory can be applied to systems with more general classical dynamics: Gutzwiller’s trace formula Gut67 ; Gut90 for chaotic systems and the corresponding Berry-Tabor formula for regular systems Ber76 give the semiclassical approximation for the density of states as a sum over the periodic orbits of the underlying classical system. However, a fundamental problem of these periodic orbit sums is that they usually do not converge, or if they do, the convergence is extremely slow. During recent years, various techniques have been developed to overcome this problem. Most of them are especially designed for chaotic systems Cvi89 ; Aur92 ; Ber90 and cannot be applied to systems with regular or mixed regular-chaotic dynamics, or they depend on special properties of the system such as the existence of a symbolic dynamics. They are therefore restricted to a relatively small number of physical systems. It would be desirable to have a method at hand which is universal in the sense that it is applicable for all types of underlying classical dynamics.
Recently, a method for periodic orbit quantization, based on harmonic inversion of a semiclassical signal has been developed and successfully applied to classically chaotic systems Mai97b ; Mai98 ; Mai99a . The aim of the present paper is to demonstrate that this technique is equally powerful in reproducing the spectra of regular systems. The semiclassical quantization of integrable and chaotic systems on an equal footing will be the basis for applications to systems with even more general, i.e., mixed regular-chaotic dynamics Mai99d . Furthermore, the harmonic inversion technique is generalized in two directions: Firstly, the periodic orbit quantization will be extended to include higher order $`\mathrm{}`$ corrections Mai98c , and, secondly, the use of cross-correlated periodic orbit sums Mai99c ; Mai99b ; Hor00 allows us to significantly reduce the required number of orbits for semiclassical quantization, i.e., to improve the efficiency of the semiclassical method. As a representative of regular systems, we choose the circle billiard whose periodic orbits and quantum eigenvalues can easily be obtained. The paper is organized as follows:
In Section 2 we give a brief overview over the periodic orbit theory for integrable systems, especially the Berry-Tabor formula, which is the analogue for integrable systems to Gutzwiller’s trace formula for chaotic systems. We then calculate the explicit expression for the density of states of the circle billiard from the Berry-Tabor formula. The equations are generalized in two directions, firstly, to the density of states weighted with the diagonal matrix elements of one or more given operators Mai99c , and, secondly, to include higher order $`\mathrm{}`$ corrections in the periodic orbit sum Mai98c .
The high precision analysis of quantum spectra and the method for the analytic continuation of non-convergent periodic orbit sums applied in this paper are based on the harmonic inversion of time signals. In Section 3 we briefly introduce harmonic inversion by filter-diagonalization Wal95 ; Man97 . We also discuss an extension of the filter-diagonalization method to cross-correlation functions Wal95 ; Nar97 ; Man98 , which can be used to extract semiclassical eigenvalues and matrix elements from cross-correlated periodic orbit sums with a significantly reduced set of periodic orbits Mai99b .
Harmonic inversion circumvents the uncertainty principle of the conventional Fourier transform and can be used for the high precision analysis of quantum spectra Mai99a ; Mai97a . In Section 4 the method will be applied to the quantum spectra of the circle billiard. The analysis will verify the validity of the Berry-Tabor formula and its generalization to spectra weighted with diagonal matrix elements discussed in Section 2.3. Furthermore, harmonic inversion will be applied to determine the higher order $`\mathrm{}`$ contributions to the periodic orbit sum. The Gutzwiller and the Berry-Tabor formula are only the leading order contributions of an expansion of the density of states in terms of $`\mathrm{}`$ and therefore only yield semiclassical approximations to the eigenvalues. By analyzing the difference spectrum between exact and semiclassical eigenvalues, first order $`\mathrm{}`$ corrections to the periodic orbit sum can be determined, as we will demonstrate in Section 4.2. The results are compared with the analytic expressions for the $`\mathrm{}`$ expansion of the periodic orbit sum given in Section 2.4.
In Section 5, we turn to the periodic orbit quantization of integrable systems. Firstly, we show how in general the problem of extracting semiclassical eigenvalues from periodic orbit sums can be reformulated as a harmonic inversion problem: A semiclassical signal is constructed from the periodic orbit sum, the analysis of which yields the semiclassical eigenvalues of the system. The general formulae are then applied to the circle billiard and the results are compared to the exact quantum and the EBK eigenvalues. In Section 5.3 it is demonstrated how the accuracy of the semiclassical eigenvalues can be significantly improved with the help of higher order $`\mathrm{}`$ corrections to the periodic orbit sum.
In Section 5.4 we address the question of how to improve the efficiency of the semiclassical quantization method, i.e., how to extract the same number of eigenvalues with a reduced set of periodic orbits, which is important especially when the orbits must be searched numerically. This is achieved by constructing cross-correlated periodic orbit sums as introduced in Section 2.3 which are then harmonically inverted with the generalized filter-diagonalization method of Section 3.2. The efficiency of the method will be discussed for various sets of operators and various sizes of the cross-correlation matrix. It is also possible to include higher order $`\mathrm{}`$ corrections in the cross-correlation signal which will allow us to calculate $`\mathrm{}`$ corrections even for nearly degenerate states.
Some concluding remarks are given in Section 6.
## 2 Periodic orbit theory for integrable systems
### 2.1 EBK quantization and Berry-Tabor formula
Integrable systems are characterized by the property that their dynamics can be expressed in action-angle variables. The action variables, which are defined on certain “irreducible” paths, are constants of motion. In the $`2n`$-dimensional phase space, the motion of an integrable system is restricted to $`n`$-dimensional tori, which are given by the values of the action variables.
A well-established method for the semiclassical quantization of integrable systems is the EBK torus quantization scheme, which was developed by Einstein, Brillouin and Keller Ein17 ; Bri26 ; Kel58 . In the EBK theory, the energy eigenvalues of the system are directly associated with certain classical tori. These tori are defined by the EBK conditions, which select special sets from all possible values of the action variables of the system. Each such set corresponds to a quantum mechanical eigenstate of the system. The tori selected by the EBK conditions are usually not rational, i.e., the orbits on these tori are usually not periodic.
For many physical systems the application of the EBK quantization scheme is a nontrivial task. Especially for non-separable or near-integrable systems the irreducible paths are difficult to find. Most importantly, as already discussed by Einstein Ein17 the torus quantization scheme cannot be extended to chaotic systems. For chaotic systems, Gutzwiller derived a semiclassical expression for the density of states in terms of the periodic orbits of the corresponding classical system: The semiclassical density of states consists of a smooth background and an oscillating part
$$\rho (E)=\rho _0(E)+\rho ^{\mathrm{osc}}(E)$$
(1)
with
$$\rho ^{\mathrm{osc}}(E)=\frac{1}{\pi \mathrm{}}\underset{\mathrm{po}}{}\frac{T_{\mathrm{po}}}{r|det(M_{\mathrm{po}}\mathrm{𝟏})|^{1/2}}\mathrm{cos}\left(\frac{S_{\mathrm{po}}}{\mathrm{}}\mu _{\mathrm{po}}\frac{\pi }{2}\right).$$
(2)
The sum runs over all periodic orbits (po) of the system, including multiple traversals. Here, $`T`$ and $`S`$ are the period and the action of the orbit, $`M`$ and $`\mu `$ are the Monodromy matrix and the Maslov index, and the repetition number $`r`$ counts the traversals of the underlying primitive orbit.
For integrable systems an analogous formula for the density of states in terms of a smooth part and oscillating periodic orbit contributions was derived by Berry and Tabor Ber76 . While in chaotic systems the periodic orbits are isolated, the periodic orbits of integrable systems are all those orbits lying on rational tori – i.e., tori on which the frequencies of the motion are commensurable – and thus are non-isolated. The Berry-Tabor formula gives the density of states in terms of the rational tori:
$$\rho (E)=\rho _0(E)+\rho ^{\mathrm{osc}}(E)$$
(3)
with
$$\rho ^{\mathrm{osc}}(E)=\frac{2}{\mathrm{}^{\frac{1}{2}(n+1)}}\underset{𝐌}{}\frac{\mathrm{cos}(S_𝐌/\mathrm{}\frac{1}{2}\pi 𝜶𝐌+\frac{1}{4}\pi \beta _𝐌)}{|𝐌|^{\frac{1}{2}(n1)}|𝝎_𝐌||K(𝐈_𝐌)|^{\frac{1}{2}}}.$$
(4)
The sum runs over all rational tori at energy $`E`$, characterized by the frequency ratios given by the ray of integer numbers $`𝐌`$. The sum includes cases where the $`M_i`$ are not relatively prime, $`𝐌=r𝝁`$, which corresponds to multiple traversals of the primitive periodic orbits on the torus characterized by $`𝝁`$. Here, $`n`$ is the dimension of the system, $`𝐈_𝐌`$ and $`𝝎`$<sub>M</sub> are the values of the action variables and the frequencies on the torus, $`S_𝐌`$ is the action of the periodic orbits on the torus, and $`K`$ is the scalar curvature of the energy contour. The components of $`𝜶`$ are the Maslov indices of the irreducible paths on which the action variables are defined, and the phase $`\beta `$ is obtained from the second derivative matrix of the action variables in terms of the coordinates. In contrast to the EBK torus quantization, there is no direct relation between the eigenvalues of the system and the tori which enter the Berry-Tabor formula.
Both the EBK torus quantization and the Berry-Tabor formula are semiclassical theories delivering lowest order $`\mathrm{}`$ approximations to the exact quantum eigenvalues. In general, the results of the two approaches can only be expected to be the same in lowest order of $`\mathrm{}`$ but not necessarily beyond. However, it was shown in Ref. Rei96 that for the circle billiard, which will be discussed in the following sections, the two approaches are in fact equivalent and should yield exactly the same results.
Eq. (4) can be simplified for the special case $`n=2`$, i.e., for two-dimensional regular systems Ull96 :
$$\rho ^{\mathrm{osc}}(E)=\frac{1}{\pi \mathrm{}^{3/2}}\underset{𝐌}{}\frac{T_𝐌}{M_2^{3/2}|g_E^{\prime \prime }|^{1/2}}\mathrm{cos}\left(\frac{S_𝐌}{\mathrm{}}\eta _𝐌\frac{\pi }{2}\frac{\pi }{4}\right).$$
(5)
The sum runs over all rational tori at energy $`E`$, characterized by the frequency ratio given by the integer numbers $`𝐌=(M_1,M_2)`$, including multiple traversals (i.e., cases where $`M_1,M_2`$ are not relatively prime). Here, $`T_𝐌`$ is the traversal time, and $`g_E`$ is the function describing the energy surface: $`H(I_1,I_2=g_E(I_1))=E`$, where $`I_1`$ and $`I_2`$ are the action variables. The Maslov index $`\eta _𝐌`$ is obtained from the Maslov indices $`\alpha _1`$, $`\alpha _2`$ of the paths on which the action variables are defined:
$$\eta _𝐌=(M_1\alpha _1+M_2\alpha _2)\mathrm{\Theta }(g_E^{\prime \prime }),$$
(6)
where $`\mathrm{\Theta }`$ is the Heaviside step function.
The semiclassical density of states can be expressed in terms of the semiclassical response function $`g(E)`$:
$$\rho (E)=\frac{1}{\pi }\mathrm{Im}g(E).$$
(7)
For both chaotic and regular systems the response function is of the form
$$g(E)=g_0(E)+\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}e^{\frac{i}{\mathrm{}}S_{\mathrm{po}}},$$
(8)
where the amplitudes are given by the Gutzwiller or the Berry-Tabor formula, respectively.
In practical applications both the Gutzwiller formula (2) and the Berry-Tabor formula (4) suffer from the property that the periodic orbit sums usually do not converge. Depending on the system in question, this problem may be overcome, e.g., by convolution of the periodic orbit sum with a suitable averaging function Rei96 . But even then the convergence will usually be slow, and a large number of orbits has to be included in order to obtain the semiclassical eigenvalues. In the following sections, we will demonstrate how the convergence problem can be circumvented by the harmonic inversion method and the eigenvalues can be calculated from a relatively small number of periodic orbits.
### 2.2 Application to the circle billiard
We now apply the Berry-Tabor formula to the circle billiard as a specific separable system with two degrees of freedom.
For completeness and comparisons with the results from periodic orbit theory we first briefly review the quantum mechanical expressions and the EBK quantization condition. The exact quantum mechanical energy eigenvalues of the circle billiard with radius $`R`$ are given by the condition
$$J_m(kR)=0,m,E=\frac{\mathrm{}^2k^2}{2M},$$
(9)
where $`J_m`$ are the Bessel functions of integer order. Here, $`M`$ denotes the mass of the particle, $`E`$ is the energy, and $`k=\sqrt{2ME}/\mathrm{}`$ is the wavenumber. The corresponding wave functions are given by
$$\psi (r,\phi )=J_m(kr)e^{im\phi }.$$
(10)
As $`J_m(x)=(1)^mJ_m(x)`$, all energy eigenvalues belonging to nonzero angular momentum quantum number $`(m0)`$ are twofold degenerate. In the following the exact quantum mechanical results for the circle billiard are used as a benchmark for the development and application of semiclassical quantization methods for integrable systems.
The circle billiard problem in two dimensions is separable in polar coordinates. The semiclassical expressions for both EBK torus quantization and the Berry-Tabor formula for the density of states are based on the action-angle variables associated with the angular $`\phi `$-motion and the radial $`r`$-motion Rei96 ; Meh99 . The action variables are given by
$`I_\phi `$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle p_\phi 𝑑\phi }=L`$ (11)
$`I_r`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle p_r𝑑r}`$ (12)
$`=`$ $`{\displaystyle \frac{1}{\pi }}\left(\sqrt{2MER^2L^2}|L|\mathrm{arccos}{\displaystyle \frac{|L|}{\sqrt{2ME}R}}\right),`$
where $`E`$ and $`L`$ are the energy and the angular momentum, respectively. Quantization of the action variables
$`I_\phi `$ $`=`$ $`\left(m+{\displaystyle \frac{\alpha _\phi }{4}}\right)\mathrm{},m`$ (13)
$`I_r`$ $`=`$ $`\left(n+{\displaystyle \frac{\alpha _r}{4}}\right)\mathrm{},n=0,1,2,\mathrm{}`$ (14)
with $`\alpha _\phi =0`$ and $`\alpha _r=3`$ for the circle billiard yields the EBK quantization condition
$$\sqrt{(kR)^2m^2}|m|\mathrm{arccos}\frac{|m|}{kR}=\left(n+\frac{3}{4}\right)\pi ,$$
(15)
where $`L=m\mathrm{}`$ are the angular momentum eigenvalues.
The frequencies of the classical motion on the two-dimensional tori are given by
$`\omega _\phi `$ $`=`$ $`{\displaystyle \frac{E}{I_\phi }}={\displaystyle \frac{2E}{\sqrt{2MER^2L^2}}}\mathrm{arccos}{\displaystyle \frac{|L|}{\sqrt{2ME}R}}`$ (16)
$`\omega _r`$ $`=`$ $`{\displaystyle \frac{E}{I_r}}={\displaystyle \frac{2\pi E}{\sqrt{2MER^2L^2}}}.`$ (17)
The Berry-Tabor formula includes all tori with a rational frequency ratio, i.e., tori on which the orbits are periodic. In the case of the circle billiard, the rational tori are given by the condition
$$\frac{\omega _\phi }{\omega _r}=\frac{M_\phi }{M_r}$$
(18)
with positive integers $`M_r`$, $`M_\phi `$ and the restriction
$$M_r2M_\phi .$$
(19)
The periodic orbits of the circle billiard have the form of regular polygons. The numbers $`M_r`$ and $`M_\phi `$ can be shown to be identical with the number of sides of the corresponding polygon and its number of turns around the center of the circle, respectively Bal72 . Some examples are given in Figure 1.
A pair of numbers $`(M_r,M_\phi )`$ which are not relatively prime corresponds to multiple traversals of a primitive periodic orbit.
The classical action of the periodic orbits is given by
$`S_𝐌`$ $`=`$ $`2\pi M_\phi I_\phi ^{(𝐌)}+2\pi M_rI_r^{(𝐌)}`$ (20)
$`=`$ $`\sqrt{2ME}R2M_r\mathrm{sin}\left({\displaystyle \frac{M_\phi }{M_r}}\pi \right).`$
As in all billiard systems, the action scales like
$$S/\mathrm{}=ws,$$
(21)
here with the scaling parameter
$$w\sqrt{2ME}R/\mathrm{}=kR$$
(22)
and the scaled action
$$s2M_r\mathrm{sin}\left(\frac{M_\phi }{M_r}\pi \right).$$
(23)
The form of the corresponding classical trajectory is independent of $`w`$. For the circle billiard with unit radius $`R=1`$, the scaling parameter $`w`$ is identical with the wavenumber $`k`$, and the scaled action is the length of the orbit.
For the semiclassical density of states, we start from the special version of the Berry-Tabor formula presented in Eq. (5). Using the relation
$$\rho (w)=\frac{dE}{dw}\rho (E),$$
(24)
valid for billiard systems, we introduce the density of states depending on the scaling parameter $`w`$. Evaluating the different expressions in (5) for the circle billiard then finally leads to
$$\rho ^{\mathrm{osc}}(w)=\frac{1}{\pi }\mathrm{Im}g^{\mathrm{osc}}(w),$$
(25)
with
$$g^{\mathrm{osc}}(w)=\sqrt{\frac{\pi }{2}}\sqrt{w}\underset{𝐌}{}m_𝐌\frac{s_𝐌^{3/2}}{M_r^2}e^{i(ws_𝐌\frac{3}{2}M_r\pi \frac{\pi }{4})},$$
(26)
where we have used the relations $`\alpha _\phi =0`$ and $`\alpha _r=3`$ for the Maslov indices. The sum runs over all pairs of positive integers $`𝐌=(M_r,M_\phi )`$ with $`M_r2M_\phi `$. The degeneracy factor
$$m_𝐌=\{\begin{array}{cc}1;\hfill & M_r=2M_\phi \hfill \\ 2;\hfill & M_r>2M_\phi \hfill \end{array},$$
(27)
accounts for the fact that all trajectories with $`M_r2M_\phi `$ can be traversed in two directions.
Due to the rapid increase of the number of periodic orbits with growing action, the sum (26) does not converge. In our case, the problem is even more complicated by the fact that there exist accumulation points of periodic orbits at scaled actions of multiples of $`2\pi `$ (see Fig. 2), which means that we cannot even include all periodic orbits up to a given finite action.
In Ref. Rei96 the convergence problem of the sum (26) was solved by averaging it with a Gaussian function. The semiclassical eigenvalues of the circle billiard were calculated from the periodic orbit sum by including a very large number of periodic orbits. Our aim is to demonstrate that by using the harmonic inversion scheme, we can obtain eigenvalues of $`w=kR`$ from a relatively small number of periodic orbits. We will return to this problem in Section 5.
### 2.3 Semiclassical matrix elements
The semiclassical trace formula for both regular and chaotic systems can be extended to include diagonal matrix elements. The calculation of individual semiclassical matrix elements is an objective in its own right. Furthermore, the extended trace formulae will allow us to construct cross-correlated periodic orbit signals and thus to significantly reduce the required number of orbits for periodic orbit quantization, as we will demonstrate in Section 5.4. Here, we will briefly recapitulate the basic ideas and equations.
Both Gutzwiller’s and Berry and Tabor’s formula give the semiclassical response function as a sum over contributions from periodic orbits (see Eq. (8)). The quantum mechanical response function is the trace over the Green function $`G_E^+`$,
$$g_{\mathrm{qm}}(E)=\underset{n}{}\frac{1}{EE_n+i0}=\mathrm{tr}G_E^+.$$
(28)
As a generalization, one can consider the quantum mechanical response function weighted with the diagonal matrix elements of some operator $`\widehat{A}`$, i.e.,
$$g_{A,\mathrm{qm}}(E)=\underset{n}{}\frac{n|\widehat{A}|n}{EE_n+i0}=\mathrm{tr}(G_E^+\widehat{A}).$$
(29)
The semiclassical approximation to the extended response function (29) is obtained by weighting the contributions of the periodic orbits in (8) with the average $`\overline{A}_p`$ of the corresponding classical quantity $`A(𝐪,𝐩)`$ over the periodic orbits:
$$g_A(E)=g_{A,0}(E)+\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}\overline{A}_pe^{\frac{i}{\mathrm{}}S_{\mathrm{po}}}.$$
(30)
For chaotic systems, the average is taken over one period $`T_p`$ of the isolated periodic orbit Wil87 ; Wil88 ; Eck92 :
$$\overline{A}_p=\frac{1}{T_p}_0^{T_p}A(𝐪(t),𝐩(t))𝑑t.$$
(31)
For an $`N`$-dimensional integrable system, the quantity $`A`$ has to be expressed in action-angle variables ($`𝐈`$, $`𝜽`$) and averaged over the rational torus Meh99 :
$$\overline{A}_p=\frac{1}{(2\pi )^N}A(𝐈,𝜽)d^N\theta .$$
(32)
Eq. (29) can even be further generalized by introducing a second operator $`\widehat{B}`$ and considering the quantity Mai99c
$$g_{AB,\mathrm{qm}}(E)=\underset{n}{}\frac{n|\widehat{A}|nn|\widehat{B}|n}{EE_n+i0}.$$
(33)
If either $`\widehat{A}`$ or $`\widehat{B}`$ commutes with the Hamiltonian, Eq. (33) can be written as a trace formula and a calculation similar to that in Eck92 yields the semiclassical approximation
$$g_{AB}(E)=g_{AB,0}(E)+\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}\overline{A}_p\overline{B}_pe^{\frac{i}{\mathrm{}}S_{\mathrm{po}}}.$$
(34)
Note that for general operators $`\widehat{A}`$ and $`\widehat{B}`$, Eq. (33) cannot be written as a trace any more. However, strong numerical evidence was provided (for both regular and chaotic systems) that Eq. (34) is correct in general, i.e., even if neither operator $`\widehat{A}`$ nor $`\widehat{B}`$ commutes with the Hamiltonian Mai99c . For chaotic systems, a mathematical proof of Eq. (34) is given in Ref. Hor00 . An analogous rigorous derivation for integrable systems is, to our knowledge, still lacking.
In Refs. Mai99c ; Hor00 , the relations (33) and (34) were generalized to products of diagonal matrix elements of more than two operators. As a further extension, we can also introduce functions of diagonal matrix elements in the response function:
$$g_{f(A),\mathrm{qm}}(E)=\underset{n}{}\frac{f(n|\widehat{A}|n)}{EE_n+i0}.$$
(35)
By a Taylor expansion of the (sufficiently smooth) function $`f`$ and using the results of Refs. Mai99c ; Hor00 for multiple products of matrix elements, we obtain the semiclassical approximation
$$g_{f(A)}(E)=g_{f(A),0}(E)+\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}f(\overline{A}_p)e^{\frac{i}{\mathrm{}}S_{\mathrm{po}}}.$$
(36)
We will use the extended trace formulae in combination with an extension of the harmonic inversion procedure to cross-correlated signals in order to significantly reduce the number of orbits which have to be included in the periodic orbit sum.
The diagonal matrix elements obtained from the extended trace formulae are semiclassical approximations to the exact quantum matrix elements. For the circle billiard, we can compare these values to those given by EBK theory. According to EBK theory, the diagonal matrix element of an operator $`\widehat{A}`$ with respect to an eigenstate $`|n`$ is obtained by averaging the corresponding classical quantity $`A(𝐈,𝜽)`$ over the quantized torus related to this eigenstate:
$$n|\widehat{A}|n=\frac{1}{(2\pi )^N}A(𝐈_n,𝜽_n)d^N\theta _n$$
(37)
Note the difference to Eq. (32), where the average is taken over the rational tori.
### 2.4 Higher order $`\mathrm{}`$ corrections
The Berry-Tabor formula for integrable systems and Gutzwiller’s trace formula for chaotic systems are only the leading order terms of an expansion of the density of states in terms of $`\mathrm{}`$. In billiard systems, the scaling parameter $`w`$ of the classical action (cf. Eq. (21)) is proportional to $`\mathrm{}^1`$ and thus plays the role of an inverse effective Planck constant,
$$w=\mathrm{}_{\mathrm{eff}}^1.$$
(38)
The $`\mathrm{}`$ expansion of the response function can therefore be written as a power series in terms of $`w^1`$ Mai98c :
$$g^{\mathrm{osc}}(w)=\underset{n=0}{\overset{\mathrm{}}{}}g_n(w)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{w^n}\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}^{(n)}e^{is_{\mathrm{po}}w}.$$
(39)
The zeroth order amplitudes $`𝒜_{\mathrm{po}}^{(0)}`$ are those of the Berry-Tabor or Gutzwiller formula, respectively, whereas for $`n>0`$, the amplitudes $`𝒜_{\mathrm{po}}^{(n)}`$ give the $`n^{\mathrm{th}}`$ order corrections $`g_n(w)`$ to the response function. Explicit expressions for the first order correction terms for chaotic systems were developed by Gaspard and Alonso Alo93 ; Gas93 and by Vattay and Rosenqvist Vat94 ; Vat96 ; Ros94 , following two different approaches. Vattay and Rosenqvist compute the corrections by solving the local Schrödinger equation in the neighborhood of periodic orbits. They introduce a quantum generalization of the Gutzwiller formula which contains these local eigenvalues. The results of Refs. Vat94 ; Vat96 ; Ros94 cannot directly be applied to integrable systems, as the derivations are valid only for isolated periodic orbits. To our knowledge, a general theory for $`\mathrm{}`$ corrections to the Berry-Tabor formula does not yet exist.
Nevertheless, for the circle billiard we have succeeded in obtaining an explicit expression for the first order $`\mathrm{}`$ corrections to the Berry-Tabor formula. The calculations are quite lengthy and are therefore deferred to Appendix A. Our final result for the first order $`\mathrm{}`$ amplitude of the circle billiard in (39) reads:
$$𝒜_{\mathrm{po}}^{(1)}=𝒜_{\mathrm{po}}^{(0)}\frac{i}{2}M_r\left(\frac{1}{3\mathrm{sin}\gamma }\frac{5}{6\mathrm{sin}^3\gamma }\right),$$
(40)
with $`\gamma \pi M_\phi /M_r`$ and $`𝒜_{\mathrm{po}}^{(0)}`$ the zeroth order amplitudes given by the Berry-Tabor formula. Using the zeroth order amplitudes from Eq. (26), we finally obtain
$$𝒜_{\mathrm{po}}^{(1)}=\sqrt{w}\sqrt{\pi M_r}\frac{2\mathrm{sin}^2\gamma 5}{6\mathrm{sin}^{3/2}\gamma }e^{i(\frac{3}{2}M_r\pi \frac{\pi }{4})}.$$
(41)
As explained above our derivation of Eq. (41) cannot be applied to general integrable systems. It will be an interesting task for the future to develop a general theory for the higher order $`\mathrm{}`$ corrections to the Berry-Tabor formula.
## 3 Harmonic inversion by filter-diagonalization
The quantization of the periodic orbit sum as well as the analysis of quantum spectra in terms of the periodic orbits can be reformulated as a harmonic inversion problem of formulae which have been introduced in the previous Section 2. Before discussing these applications in Sections 4 and 5 we will now briefly recapitulate the basic ideas and the technical tools of harmonic inversion by filter-diagonalization. In Section 3.1 we will start with the harmonic inversion of a single function. The equations will be generalized in Section 3.2 to the harmonic inversion of cross-correlated signals.
### 3.1 Harmonic inversion of a single function
The harmonic inversion problem can be formulated as a nonlinear fit of a signal $`C(t)`$ to the form
$$C(t)=\underset{k}{}d_ke^{i\omega _kt},$$
(42)
where $`d_k`$ and $`\omega _k`$ are generally complex variational parameters. Other than, e.g., in a simple Fourier transformation of the signal, there is no restriction to the closeness of the frequencies $`\omega _k`$. Solving $`(\text{42})`$ will therefore yield a high resolution analysis of the signal $`C(t)`$. The signal length required for resolving the frequencies $`\omega _k`$ by harmonic inversion can be estimated to be
$$t_{\mathrm{max}}4\pi \overline{\rho }(\omega ),$$
(43)
where $`\overline{\rho }(\omega )`$ is the mean density of frequencies in the range of interest.
A method which has proven very useful for solving the harmonic inversion problem is the filter-diagonalization procedure Wal95 ; Man97 . This procedure allows us to compute the frequencies $`\omega _k`$ in any small interval $`[\omega _{\mathrm{min}},\omega _{\mathrm{max}}]`$ given. The idea is to consider the signal $`C(t)`$ on an equidistant grid
$$c_n=C(n\tau );n=0,1,2,\mathrm{}$$
(44)
and to associate $`c_n`$ with an autocorrelation function of a suitable fictitious dynamical system, described by a complex symmetric effective Hamiltonian $`H_{\mathrm{eff}}`$:
$$c_n=\left(\mathrm{\Phi }_0|e^{in\tau H_{\mathrm{eff}}}\mathrm{\Phi }_0\right).$$
(45)
Here, the brackets denote a complex symmetric inner product $`(a|b)=(b|a)`$, i.e., no complex conjugation of either $`a`$ or $`b`$. The harmonic inversion problem can then be reformulated as solving the eigenvalue problem for the effective Hamiltonian $`H_{\mathrm{eff}}`$. The frequencies $`\omega _k`$ are the eigenvalues of the Hamiltonian
$$H_{\mathrm{eff}}|\mathrm{{\rm Y}}_k)=\omega _k|\mathrm{{\rm Y}}_k),$$
(46)
and the amplitudes are obtained from the eigenvectors $`\mathrm{{\rm Y}}_k`$:
$$d_k=\left(\mathrm{\Phi }_0|\mathrm{{\rm Y}}_k\right)^2.$$
(47)
The filter-diagonalization method solves this eigenvalue problem in a small set of basis vectors $`\mathrm{\Psi }_j`$. The Hamiltonian and the initial state $`\mathrm{\Phi }_0`$ do not have to be known explicitly but are given implicitly by the quantities $`c_n`$. In detail, the procedure works as follows:
A small set of values $`\phi _j`$ in the frequency interval of interest is chosen. The set must be larger than the number of frequencies in this interval. The values $`\phi _j`$ are used to construct the small Fourier-type basis
$$\mathrm{\Psi }_j=\underset{n=0}{\overset{M}{}}e^{in(\phi _j\tau H_{\mathrm{eff}})}\mathrm{\Phi }_0.$$
(48)
The matrix elements of the evolution operator at a given time $`p\tau `$ in this basis can be expressed in terms of the quantities $`c_n`$:
$$U_{jj^{}}^{(p)}\left(\mathrm{\Psi }_j|e^{ip\tau H_{\mathrm{eff}}}\mathrm{\Psi }_j^{}\right)=\underset{n=0}{\overset{M}{}}\underset{n^{}=0}{\overset{M}{}}e^{i(n\phi _j+n^{}\phi _j^{})}c_{n+n^{}+p}.$$
(49)
The frequencies $`\omega _k`$ are then obtained by solving the generalized eigenvalue problem
$$𝐔^{(p)}𝐁_k=e^{ip\tau \omega _k}𝐁_k.$$
(50)
The amplitudes $`d_k`$ can be calculated from the eigenvectors and are given by
$$d_k=\left(\underset{j}{}B_{jk}\underset{n=0}{\overset{M}{}}c_ne^{in\phi _j}\right)^2.$$
(51)
The values of $`\omega _k`$ and $`d_k`$ obtained by the above procedure should be independent of $`p`$. This condition can be used to identify non-converged frequencies by comparing the results for different values of $`p`$. The difference between the frequency values obtained for different $`p`$ can be used as a simple error estimate.
### 3.2 Harmonic inversion of cross-correlated signals
An important extension of the filter-diagonalization method for harmonic inversion is the generalization to cross-correlation functions Wal95 ; Nar97 ; Man98 ; Mai99b . This extended method allows us to significantly reduce the signal length required to resolve the frequencies contained in the signal. The idea is not to consider a single signal $`C(t)`$ as given in Eq. (42) but a set of cross-correlated signals
$$C_{\alpha \alpha ^{}}(t)=\underset{k}{}d_{\alpha \alpha ^{},k}e^{i\omega _kt};\alpha ,\alpha ^{}=1,\mathrm{},N$$
(52)
with the restriction
$$d_{\alpha \alpha ^{},k}=b_{\alpha ,k}b_{\alpha ^{},k}.$$
(53)
This set of signals considered on an equidistant grid
$$c_{n\alpha \alpha ^{}}=C_{\alpha \alpha ^{}}(n\tau );n=0,1,2,\mathrm{}$$
(54)
is now associated with a time cross-correlation function between an initial state $`\mathrm{\Phi }_\alpha `$ and a final state $`\mathrm{\Phi }_\alpha ^{}`$:
$$c_{n\alpha \alpha ^{}}=\left(\mathrm{\Phi }_\alpha ^{}|e^{in\tau H_{\mathrm{eff}}}\mathrm{\Phi }_\alpha \right).$$
(55)
Again, the frequencies $`\omega _k`$ are obtained as the eigenvalues of the effective Hamiltonian $`H_{\mathrm{eff}}`$. The amplitudes are now given by the relation
$$b_{\alpha ,k}=(\mathrm{\Phi }_\alpha |\mathrm{{\rm Y}}_k).$$
(56)
In analogy to the procedure described in Section 3.1, this eigenvalue problem is solved in a small set of basis vectors $`\mathrm{\Psi }_{j\alpha }`$ in order to obtain the frequencies in a given interval $`[\omega _{\mathrm{min}},\omega _{\mathrm{max}}]`$.
The advantage of the above procedure becomes evident if one considers the information content of the set of signals. Due to the restriction (53), the $`N\times N`$ set of signals $`C_{\alpha \alpha ^{}}(t)`$ may contain $`N`$ independent signals, which are known to possess the same frequencies $`\omega _k`$. This means that, at constant signal length, the matrix may contain $`N`$ times as much information about the frequencies as a single signal, provided that the whole set is inverted simultaneously. On the other hand, the information content is proportional to the signal length. This means that the signal length required to resolve the frequencies in a given interval is reduced by a factor of $`N`$. This statement clearly holds only approximately and for small matrix dimensions $`N`$. However, a significant reduction of the required signal length can be achieved.
## 4 High resolution analysis of quantum spectra
Harmonic inversion is a powerful tool to calculate the classical periodic orbit contributions to the density of states from the quantum mechanical eigenvalues or from experimental spectra, thus delivering a high resolution analysis of the spectra in terms of the classical orbits. The method allows us, e.g., to resolve clusters of orbits or to discover hidden ghost orbit contributions in the spectra, which would not be resolved by conventional Fourier analysis of the spectra Mai99a ; Mai97a . Here, we will analyze the quantum spectra of the circle billiard as a representative of integrable systems. The analysis will verify the validity of the Berry-Tabor formula and its extensions to semiclassical matrix elements and higher order $`\mathrm{}`$ corrections discussed in Section 2.
### 4.1 Leading order periodic orbit contributions to the trace formula
#### 4.1.1 General procedure
In this section we develop the general procedure for the analysis of quantum spectra in terms of periodic orbits by harmonic inversion. This procedure is universal in the sense that it can be applied to both regular and chaotic systems. We will apply it to the circle billiard as a representative of regular systems in the next section.
We start from the semiclassical density of states given by the Berry-Tabor or the Gutzwiller formula. As in Section 2, we consider scaling systems where the density of states depends on the scaling parameter $`w`$ \[$`w=kR`$ for the circle billiard\], i.e., $`\rho (w)=(1/\pi )\mathrm{Im}g(w)`$ with $`g(w)`$ the semiclassical response function. Both the Berry-Tabor and the Gutzwiller formula give the oscillating part of the response function in the form
$$g^{\mathrm{osc}}(w)=\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}e^{iws_{\mathrm{po}}},$$
(57)
where the sum runs over all rational tori (regular systems) or all periodic orbits (chaotic systems) of the underlying classical system, respectively. Here, $`S_{\mathrm{po}}`$ is the action of the periodic orbit. The form of the amplitude $`𝒜_{\mathrm{po}}`$ depends on whether the system is chaotic or regular and also contains phase information.
The exact quantum mechanical density of states is given by
$$\rho _{\mathrm{qm}}(w)=\underset{k}{}m_k\delta (ww_k),$$
(58)
where the $`w_k`$ are the exact quantum eigenvalues of the scaling parameter and the $`m_k`$ are their multiplicities. The analysis of the quantum spectrum in terms of periodic orbit contributions can now be reformulated as adjusting the exact quantum mechanical density of states (58) to the semiclassical form
$`\rho ^{\mathrm{osc}}(w)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}g^{\mathrm{osc}}(w)`$ (59)
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{\mathrm{po}}{}}\left(𝒜_{\mathrm{po}}e^{iws_{\mathrm{po}}}𝒜_{\mathrm{po}}^{}e^{iws_{\mathrm{po}}}\right).`$
If the amplitudes $`𝒜_{\mathrm{po}}`$ do not depend on $`w`$, the semiclassical density of states is of the form (42) \[here, with $`w`$ playing the role of $`t`$ and $`s_{\mathrm{po}}`$ playing the role of $`\omega _k`$\]. That means, we have reformulated the problem of extracting the periodic contributions from the quantum spectrum as a harmonic inversion problem. In the fitting procedure, we ignore the non-oscillating, smooth part of the density of states. This part does not fulfill the ansatz (42) of the harmonic inversion method and therefore acts as a kind of noise, which will be separated from the oscillating part of the “signal” by the harmonic inversion procedure.
In practice, in order to regularize the $`\delta `$ functions in (58), we convolute both expressions (58) and (59) with a Gaussian function,
$$C_\sigma (w)=\frac{1}{\sqrt{2\pi }\sigma }_{\mathrm{}}^{\mathrm{}}\rho (w^{})e^{(ww^{})^2/2\sigma ^2}𝑑w^{}.$$
(60)
In our calculations, we usually took the convolution width to be about three times the step width $`\tau `$ in the signal (44). Typical values are $`\tau =\mathrm{\Delta }w=0.002`$ and $`\sigma =0.006`$. The resulting quantum mechanical signal is
$$C_{\mathrm{qm},\sigma }(w)=\frac{1}{\sqrt{2\pi }\sigma }\underset{k}{}m_ke^{(ww_k)^2/2\sigma ^2},$$
(61)
and the corresponding semiclassical quantity reads
$$C_\sigma (w)=\frac{1}{2\pi i}\underset{\mathrm{po}}{}\left(𝒜_{\mathrm{po}}e^{iws_{\mathrm{po}}}𝒜_{\mathrm{po}}^{}e^{iws_{\mathrm{po}}}\right)e^{\frac{1}{2}\sigma ^2s_{\mathrm{po}}^2}.$$
(62)
The above procedure still works if the amplitudes in (57) are not independent of $`w`$ but possess a dependency of the form
$$A_{\mathrm{po}}=w^\alpha a_{\mathrm{po}},$$
(63)
which is, for example, the case for regular billiards. This dependency can be eliminated Mai98c by replacing the semiclassical response function $`g(w)`$ with the quantity
$$g^{}(w)=w^\alpha g(w)=w^\alpha g_0(w)+\underset{\mathrm{po}}{}a_{\mathrm{po}}e^{iws_{\mathrm{po}}}.$$
(64)
When introducing the corresponding quantum mechanical response function
$$g_{\mathrm{qm}}^{}(w)=\underset{k}{}\frac{m_kw_k^\alpha }{ww_k+i0}$$
(65)
the procedure can be carried out for $`\rho ^{}(w)=(1/\pi )\mathrm{Im}g^{}(w)`$ as described above.
In addition to considering the pure density of states, the relations of Section 2.3 can be used to calculate the averages of classical quantities over the periodic orbits from the quantum diagonal matrix elements of the corresponding operators. If we start from the extended quantum response function (29), including diagonal matrix elements of some operator $`\widehat{A}`$, the analysis of the signal should again yield the actions $`s_{\mathrm{po}}`$ as frequencies but with the amplitudes weighted with the classical averages $`\overline{A}_p`$ of the corresponding classical quantities. In the same way, we can also use the extended response function (33), which includes diagonal matrix elements of two different operators.
#### 4.1.2 Application to the circle billiard
For the circle billiard, the oscillating part $`g^{\mathrm{osc}}(w)`$ of the semiclassical response function is given by Eq. (26). If one eliminates the dependency of the amplitudes on $`w`$ by defining
$$\rho ^{}(w)=\frac{1}{\sqrt{w}}\rho (w),$$
(66)
the resulting expression for the density of states
$`\rho ^{}(w)={\displaystyle \frac{1}{\sqrt{8\pi }}}{\displaystyle \underset{𝐌}{}}m_𝐌{\displaystyle \frac{s_𝐌^{3/2}}{M_r^2}}(`$ $`e^{i(\frac{3}{2}M_r\pi \frac{\pi }{4})}e^{iws_𝐌}`$
$`+`$ $`e^{i(\frac{3}{2}M_r\pi \frac{\pi }{4})}e^{iws_𝐌}`$ $`)`$ (67)
is of the form (42), here with $`S_𝐌`$ playing the role of $`w_k`$. The quantum mechanical quantity corresponding to (66) is
$$\rho _{\mathrm{qm}}^{}(w)=\underset{k}{}\frac{m_k}{\sqrt{w}_k}\delta (ww_k).$$
(68)
In addition to analyzing the pure quantum spectrum of the circle billiard, we also considered spectra weighted with diagonal matrix elements of different operators (cf. Section 2.3). We used three different operators, viz.
* the absolute value of the angular momentum $`L`$ as an example of a constant of motion,
* the distance $`r`$ from the center as an example of a quantity which is no constant of motion,
* the variance of the radius $`r^2r^2`$ as an example using the relation (34) for products of operators.
The classical angular momentum $`L`$ is proportional to $`w`$, which means that when constructing the signal for $`L`$, $`g(w)`$ now has to be multiplied by $`w^{3/2}`$ instead of $`w^{1/2}`$ (cf. Eq. (64)).
We calculated the scaled actions and classical amplitudes of the periodic orbits in the interval $`s_𝐌[15,23]`$. The signal was constructed from the exact zeros of the Bessel functions, up to a value of $`w_{\mathrm{max}}=500`$. The accuracy of results is improved if we cut off the lower part of the signal, using only zeros larger than $`w_{\mathrm{min}}=300`$. A possible explanation for this is that the low zeros are in a sense “too much quantum” for the semiclassical periodic orbit sum.
Figure 3 shows the results of our calculation. The positions of the solid lines are the scaled actions of the classical periodic orbits, their heights are the classical amplitudes $`m_𝐌s_𝐌^{3/2}/M_r^2`$ times the respective averaged classical quantity. The crosses are the results obtained by harmonic inversion of the signal constructed from the zeros of the Bessel functions. There is an excellent agreement between the spectra, illustrating the validity of the Berry-Tabor formula and its extension to semiclassical matrix elements discussed in Section 2.3. The examined interval contains an accumulation point of orbits ($`s=6\pi `$). Here, only those orbits were resolved which were still sufficiently isolated.
It might be surprising that, although the Berry-Tabor formula only gives a semiclassical approximation to the density of states and we started from the exact quantum mechanical density, our results for the periodic orbit contributions are exact and do not show any deviations due to the error of the semiclassical approximation. The reason for this will become obvious in the following Section.
### 4.2 Higher order $`\mathrm{}`$ corrections to the trace formula
An interesting application of the method described in the previous Section 4.1 is the determination of higher order $`\mathrm{}`$ contributions to the periodic orbit sum. The higher orders can be obtained by analysis of the difference spectrum between the exact quantum and semiclassical eigenvalues, as we will show below.
As explained in Section 2.4, the Berry-Tabor formula for integrable systems as well as the Gutzwiller formula for chaotic systems are the leading order terms of an expansion of the density of states in terms of $`\mathrm{}`$. For scaling systems, this expansion can be put in the form (cf. (39))
$$g^{\mathrm{osc}}(w)=\underset{n=0}{\overset{\mathrm{}}{}}g_n(w)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{w^n}\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}^{(n)}e^{is_{\mathrm{po}}w}.$$
(69)
Provided that the amplitudes $`𝒜_{\mathrm{po}}^{(n)}`$ in (69) do not depend on $`w`$, only the zeroth order term fulfills the ansatz (42) for the harmonic inversion procedure with constant amplitudes and frequencies. In systems like regular billiards, where the amplitudes possess a $`w`$ dependence of the form $`𝒜_{\mathrm{po}}^{(n)}=w^\alpha a_{\mathrm{po}}^{(n)}`$, the same argumentation holds if we consider $`g^{}(w)=w^\alpha g(w)`$ instead of $`g(w)`$ (cf. Section 4.1.1). This is the reason why the analysis of the quantum spectrum yields exactly the lowest order amplitudes $`𝒜_{\mathrm{po}}^{(0)}`$, without any deviations due to the semiclassical error: As the higher order terms do not fulfill the ansatz, the $`𝒜_{\mathrm{po}}^{(0)}`$ are the best fit for the amplitudes. The higher oder terms have similar properties as a weak noise and are separated from the “true” signal by the harmonic inversion procedure.
Although the direct analysis of the quantum spectrum only yields the lowest order amplitudes, higher order corrections can still be extracted from the spectrum by harmonic inversion. Assume that the exact eigenvalues $`w_k`$ and their $`(n1)^{\mathrm{st}}`$ order approximations $`w_{k,n1}`$ are given. We can then calculate the difference between the exact quantum mechanical and the $`(n1)^{\mathrm{st}}`$ order response function
$$g^{\mathrm{qm}}(w)\underset{j=0}{\overset{n1}{}}g_j(w)=\underset{j=n}{\overset{\mathrm{}}{}}g_j(w)=\underset{j=n}{\overset{\mathrm{}}{}}\frac{1}{w^j}\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}^{(j)}e^{is_{\mathrm{po}}w}.$$
(70)
The leading order terms in (70) are $`w^n`$, i.e., multiplication with $`w^n`$ yields
$$w^n\left[g^{\mathrm{qm}}(w)\underset{j=0}{\overset{n1}{}}g_j(w)\right]=\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}^{(n)}e^{is_{\mathrm{po}}w}+𝒪\left(\frac{1}{w}\right).$$
(71)
In (71) we have restored the functional form (42). The harmonic inversion of the function (71) will now provide the periods $`s_{\mathrm{po}}`$ and the $`n^{\mathrm{th}}`$ order amplitudes $`𝒜_{\mathrm{po}}^{(n)}`$ of the $`\mathrm{}`$ expansion (69).
In practice, we will follow the procedure outlined in Section 4.1.1 to construct a smooth signal, i.e., we consider the densities of states $`\rho (w)=(1/\pi )\mathrm{Im}g(w)`$ rather than the response functions $`g(w)`$, and smoothen the signal by convoluting it with a Gaussian function.
For the circle billiard, the exact quantum eigenvalues are given by the condition (9), while the zeroth order eigenvalues are equal to the EBK eigenvalues given by (15) (cf. Section 2.1). From the difference between the exact and the semiclassical density of states, we can calculate the amplitudes $`𝒜_{\mathrm{po}}^{(1)}`$ of the first order correction to the trace formula.
We analyzed the difference spectrum between exact and EBK eigenvalues of the circle billiard in the range $`100<w<500`$. Figure 4 shows a small part of this difference spectrum. The results of the harmonic inversion of the spectrum are presented in Figure 5.
The crosses mark the values
$$f(\gamma )\frac{2}{\sqrt{\pi M_r}}\frac{1}{\sqrt{w}}|𝒜_{\mathrm{po}}^{(1)}|,$$
(72)
with $`\gamma =\pi M_\phi /M_r`$ which we obtained for the periodic orbits by harmonic inversion of the difference spectrum. The crosses are labeled with the numbers $`(M_r,M_\phi )`$ of the orbits. The solid line in Fig. 5 is the theoretical curve
$$f(\gamma )=\frac{52\mathrm{sin}^2\gamma }{3\mathrm{sin}^{3/2}\gamma },$$
(73)
which results from our analytical expression (41) for the first order amplitudes discussed in Section 2.4. The results obtained by harmonic inversion are in excellent agreement with the theoretical curve, which clearly illustrates the validity of Eq. (41).
## 5 Periodic orbit quantization
### 5.1 General procedure
We now turn to the problem of extracting eigenvalues from the periodic orbit sum. We will demonstrate that the harmonic inversion procedure, which has already been successfully applied to extract the eigenvalues of chaotic systems Mai97b ; Mai98 , can be used for integrable systems as well when starting from the Berry-Tabor formula.
As previously (see Section 2), we consider scaling systems and start from the response function
$$g(w)=g_0(w)+\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}e^{iws_{\mathrm{po}}},$$
(74)
depending on the scaling parameter $`w`$. The amplitudes $`𝒜_{\mathrm{po}}`$ are those of the Berry-Tabor or the Gutzwiller formula for regular and chaotic systems, respectively. The periodic orbit sum in (74) usually does not converge, or, at least, the convergence will be very slow. In practice, especially for chaotic systems, only the orbits with small scaled actions will be known. Nevertheless, the eigenvalues of the scaling parameter can still be extracted from the periodic orbit sum. The central idea is to adjust Eq. (74), with the sum including periodic orbits up to a finite action $`s_{\mathrm{max}}`$, to the functional form of the corresponding quantum mechanical response function
$$g_{\mathrm{qm}}(w)=\underset{k}{}\frac{m_k}{ww_k+i0}.$$
(75)
This fitting problem cannot be solved directly, but can be reformulated as a harmonic inversion problem Mai97b ; Mai98 . The first step of the reformulation is a Fourier transformation of the response functions with respect to $`w`$:
$$C(s)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}g(w)e^{isw}𝑑w.$$
(76)
In the semiclassical response function, we only consider the oscillating part of $`g(w)`$. The smooth part, which does not possess a suitable form for the harmonic inversion method, would only give a contribution to the signal for very small $`s`$. Assuming that the amplitudes in (74) do not depend on $`w`$, the result of the Fourier transformation is
$`C(s)`$ $`=`$ $`{\displaystyle \underset{\mathrm{po}}{}}𝒜_{\mathrm{po}}\delta (ss_{\mathrm{po}}),`$ (77)
$`C_{\mathrm{qm}}(s)`$ $`=`$ $`i{\displaystyle \underset{k}{}}m_ke^{isw_k}.`$ (78)
Like in Section 4.1.1 we convolute the signals (77) and (78) with a Gaussian function with width $`\sigma `$, resulting in
$`C_\sigma (s)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }\sigma }}{\displaystyle \underset{\mathrm{po}}{}}𝒜_{\mathrm{po}}e^{(ss_{\mathrm{po}})^2/2\sigma ^2},`$ (79)
$`C_{\mathrm{qm},\sigma }(s)`$ $`=`$ $`i{\displaystyle \underset{k}{}}m_ke^{\frac{1}{2}\sigma ^2w_k^2}e^{isw_k}.`$ (80)
Typical values of the convolution width are $`\sigma =0.006`$ for signals with step width $`\mathrm{\Delta }s=0.002`$. The eigenvalues of the scaling parameter are now obtained by adjusting the signal $`C_\sigma (s)`$ to (80), which is of the functional form (42). The frequencies $`w_k`$ obtained by harmonic inversion of the signal (79) are the eigenvalues of the scaling parameter $`w`$; from the amplitudes $`d_k`$, the multiplicities $`m_k`$ can be calculated.
Like the general procedure for analyzing quantum spectra (see Section 4.1.1), the above procedure still works if the amplitudes in (74) are not independent of $`w`$ but possess a dependency of the form
$$𝒜_{\mathrm{po}}=w^\alpha a_{\mathrm{po}}.$$
(81)
We can again eliminate this dependency by replacing $`g(w)`$ with the quantity
$$g^{}(w)=w^\alpha g(w).$$
(82)
The semiclassical signal is then given by
$$C_\sigma (s)=\frac{1}{\sqrt{2\pi }\sigma }\underset{\mathrm{po}}{}a_{\mathrm{po}}e^{(ss_{\mathrm{po}})^2/2\sigma ^2},$$
(83)
and the corresponding quantum mechanical signal reads
$$C_{\mathrm{qm},\sigma }(s)=i\underset{k}{}m_kw_k^\alpha e^{\frac{1}{2}\sigma ^2w_k^2}e^{isw_k}.$$
(84)
### 5.2 Semiclassical eigenvalues of the circle billiard
#### 5.2.1 Construction of the periodic orbit signal
The semiclassical response function of the circle billiard is given by Eq. (26). The amplitudes in (26) are proportional to $`w^{1/2}`$. As described in Section 5.1, we eliminate this dependency on $`w`$ by introducing the quantity
$$g^{}(w)=w^{1/2}g(w).$$
(85)
Applying Eqs. (83) and (84), we now obtain the semiclassical and the corresponding exact quantum signal for the circle billiard:
$`C_\sigma (s)`$ $`=`$ $`{\displaystyle \frac{e^{i\frac{\pi }{4}}}{2\sigma }}{\displaystyle \underset{𝐌}{}}m_𝐌{\displaystyle \frac{s_𝐌^{3/2}}{M_r^2}}e^{i\frac{3}{2}M_r\pi }e^{(ss_𝐌)^2/2\sigma ^2},`$
$`C_{\mathrm{qm},\sigma }(s)`$ $`=`$ $`i{\displaystyle \underset{k}{}}{\displaystyle \frac{m_k}{\sqrt{w}_k}}e^{\frac{1}{2}\sigma ^2w_k^2}e^{isw_k}.`$ (87)
Eq. (87) possesses the functional form (42) with
$$d_k=i\frac{m_k}{\sqrt{w}_k}e^{\frac{1}{2}\sigma ^2w_k^2}.$$
(88)
Applying the harmonic inversion method to the signal (LABEL:csigmas) should yield the eigenvalues of $`w`$ as frequencies, with the amplitudes given by Eq. (88).
#### 5.2.2 Results for the lowest eigenvalues
We calculated the eigenvalues of the scaling parameter $`w=kR`$ for the lowest states of the circle billiard from a signal of length $`s_{\mathrm{max}}=150`$. For the construction of the signal, we chose a minimum length for the sides of the periodic orbits as cut-off criterion at the accumulation points (cf. Fig. 2). We observed that the results were nearly independent of the choice of this parameter, as long as the minimum length was not chosen too large.
Table 1 presents the semiclassical eigenvalues $`w_{\mathrm{hi}}`$ and multiplicities $`m_{\mathrm{hi}}`$ obtained by harmonic inversion of the periodic orbit signal (LABEL:csigmas). For comparison, the exact quantum mechanical and the EBK results are also given in Table 1. The eigenvalues obtained by harmonic inversion clearly reproduce the EBK eigenvalues within an accuracy of $`10^4`$ or better. The deviation of the $`w_{\mathrm{hi}}`$ from the EBK eigenvalues is significantly smaller than the error of the semiclassical approximation. The improvement of the semiclassical quantization by including higher order $`\mathrm{}`$ corrections to the periodic orbit sum will be discussed in Section 5.3.
Note that for calculating the eigenvalues of the circle billiard by a direct evaluation of the periodic orbit sum, a huge number of periodic orbit terms is required, e.g., orbits with maximum length $`s_{\mathrm{max}}=\mathrm{30\hspace{0.17em}000}`$ were included in Ref. Rei96 . We obtained similar results using only orbits up to length $`s_{\mathrm{max}}=150`$. This demonstrates the high efficiency of the harmonic inversion method in extracting eigenvalues from the periodic orbit sum. The efficiency can even be further increased with the help of cross-correlated periodic orbit sums as will be demonstrated in Section 5.4.
In Table 1 the exact multiplicities $`m_{\mathrm{ex}}`$ of eigenvalues and the multiplicities $`m_{\mathrm{hi}}`$ obtained by harmonic inversion also agree to very high precision. The deviations are one or two orders of magnitude larger than those in the frequencies, which reflects the fact that, with the harmonic inversion method using filter-diagonalization, the amplitudes usually converge more slowly than the frequencies.
In the frequency interval shown, there are two cases of nearly degenerate frequencies which have not been resolved by harmonic inversion of the periodic orbit signal with $`s_{\mathrm{max}}=150`$. The harmonic inversion yielded only one frequency, which is the average of the two nearly degenerate ones, with the amplitudes of the two added. These nearly degenerate states can be resolved when the signal length is increased to about $`s_{\mathrm{max}}=500`$ or with the help of cross-correlated periodic orbit sums (see Section 5.4).
#### 5.2.3 Semiclassical matrix elements
Using the extended periodic orbit sums discussed in Section 2.3, we can now also calculate semiclassical diagonal matrix elements for the circle billiard.
Following the procedure described in Section 5.1, a semiclassical signal can be constructed from the extended response function, the analysis of which should again yield the eigenvalues $`w_k`$ as frequencies but with the amplitudes weighted with the diagonal matrix elements. As examples, we used the same operators as in Section 4.1.2 to calculate the diagonal matrix elements $`L`$, $`r`$, and the variance of the radius, $`r^2r^2`$.
Figure 6 shows the results in the range $`25w30`$. For comparison, Fig. 6a presents the spectrum for the identity operator. The positions of the solid lines are the EBK eigenvalues, their heights are the semiclassical matrix elements obtained from EBK theory times the multiplicities. The crosses are the results of the harmonic inversion of a signal of length $`s_{\mathrm{max}}=300`$. The diagrams show excellent agreement between the results obtained by harmonic inversion and EBK torus quantization. This is even the case for the variance of $`r`$, which is a very small quantity.
For the states shown in Fig. 6 we have also compared the semiclassical to the exact quantum mechanical matrix elements. The agreement is also excellent. The deviations between the semiclassical and quantum matrix elements are typically of the order of $`10^3`$, which can well be understood by the semiclassical approximation.
### 5.3 Higher order $`\mathrm{}`$ corrections
The eigenvalues of the circle billiard obtained in the previous Section 5.2 are not the exact quantum mechanical eigenvalues but semiclassical approximations for the reason that the Berry-Tabor and the Gutzwiller formula are only the leading order terms of an expansion of the density of states in terms of $`\mathrm{}`$ (see Section 2.4). We will now demonstrate how to obtain corrections to the semiclassical eigenvalues from the $`\mathrm{}`$ expansion (39) of the periodic orbit sum
$$g^{\mathrm{osc}}(w)=\underset{n=0}{\overset{\mathrm{}}{}}g_n(w)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{w^n}\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}^{(n)}e^{is_{\mathrm{po}}w},$$
with $`w=\mathrm{}_{\mathrm{eff}}^1`$ an effective inverse Planck constant (see Eq. (38)). The amplitudes $`𝒜_{\mathrm{po}}^{(0)}`$ are those of the Berry-Tabor or Gutzwiller formula. For $`n>0`$, the amplitudes $`𝒜_{\mathrm{po}}^{(n)}`$ (including also phase information) give the $`n^{\mathrm{th}}`$ order corrections $`g_n(w)`$ to the response function $`g^{\mathrm{osc}}(w)`$. For simplicity, we will assume in the following that the amplitudes $`𝒜_{\mathrm{po}}^{(n)}`$ in (39) do not depend on $`w`$. Again, in systems where the amplitudes possess a $`w`$ dependence of the form $`𝒜_{\mathrm{po}}^{(n)}=w^\alpha a_{\mathrm{po}}^{(n)}`$, the same line of arguments holds if we consider $`g^{}(w)=w^\alpha g(w)`$ instead of $`g(w)`$ (cf. Section 5.1).
For periodic orbit quantization the zeroth order contributions $`𝒜_{\mathrm{po}}^{(0)}`$ are usually considered only. As demonstrated in Section 5.1 (see Eqs. (77) and (78)), the Fourier transform of the principal periodic orbit sum
$$C_0(s)=\underset{\mathrm{po}}{}𝒜_{\mathrm{po}}^{(0)}\delta (ss_{\mathrm{po}})$$
is adjusted by application of the harmonic inversion technique to the functional form of the exact quantum expression
$$C_{\mathrm{qm}}(s)=i\underset{k}{}m_ke^{iw_ks},$$
with $`\{w_k,m_k\}`$ the eigenvalues and multiplicities.
For $`n1`$, the asymptotic expansion (39) of the semiclassical response function suffers from the singularities at $`w=0`$. It is therefore not appropriate to harmonically invert the Fourier transform of (39) as a whole, although the Fourier transform formally exists. This means that the method of periodic orbit quantization by harmonic inversion cannot straightforwardly be extended to the $`\mathrm{}`$ expansion of the periodic orbit sum. Instead, we will calculate the correction terms to the semiclassical eigenvalues separately, order by order Mai98c .
Let us assume that the $`(n1)^{\mathrm{st}}`$ order approximations $`w_{k,n1}`$ to the semiclassical eigenvalues have already been obtained and the $`w_{k,n}`$ are to be calculated. The difference between the two subsequent approximations to the quantum mechanical response function reads
$`g_n(w)`$ $`=`$ $`{\displaystyle \underset{k}{}}\left({\displaystyle \frac{m_k}{ww_{k,n}+i0}}{\displaystyle \frac{m_k}{ww_{k,n1}+i0}}\right)`$ (89)
$``$ $`{\displaystyle \underset{k}{}}{\displaystyle \frac{m_k\mathrm{\Delta }w_{k,n}}{(w\overline{w}_{k,n}+i0)^2}},`$
with $`\overline{w}_{k,n}=\frac{1}{2}(w_{k,n}+w_{k,n1})`$ and $`\mathrm{\Delta }w_{k,n}=w_{k,n}w_{k,n1}`$. Integration of (89) and multiplication by $`w^n`$ yields
$$𝒢_n(w)=w^ng_n(w)𝑑w=\underset{k}{}\frac{m_kw^n\mathrm{\Delta }w_{k,n}}{w\overline{w}_{k,n}+i0},$$
(90)
which has the functional form of a quantum mechanical response function but with residues proportional to the $`n^{\mathrm{th}}`$ order corrections $`\mathrm{\Delta }w_{k,n}`$ to the semiclassical eigenvalues. The semiclassical approximation to (90) is obtained from the term $`g_n(w)`$ in the periodic orbit sum (39) by integration and multiplication by $`w^n`$, i.e.,
$`𝒢_n(w)`$ $`=`$ $`w^n{\displaystyle g_n(w)𝑑w}`$ (91)
$`=`$ $`i{\displaystyle \underset{\mathrm{po}}{}}{\displaystyle \frac{1}{s_{\mathrm{po}}}}𝒜_{\mathrm{po}}^{(n)}e^{iws_{\mathrm{po}}}+𝒪\left({\displaystyle \frac{1}{w}}\right).`$
We can now Fourier transform both (90) and (91), and obtain ($`n1`$)
$`C_n(s)`$ $``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝒢_n(w)e^{iws}𝑑w`$ (92)
$`=`$ $`i{\displaystyle \underset{k}{}}m_k(\overline{w}_k)^n\mathrm{\Delta }w_{k,n}e^{i\overline{w}_ks}`$
$`\stackrel{\mathrm{h}.\mathrm{i}.}{=}`$ $`i{\displaystyle \underset{\mathrm{po}}{}}{\displaystyle \frac{1}{s_{\mathrm{po}}}}𝒜_{\mathrm{po}}^{(n)}\delta (ss_{\mathrm{po}}).`$ (93)
Equations (92) and (93) imply that the $`\mathrm{}`$ expansion of the semiclassical eigenvalues can be obtained, order by order, by harmonic inversion (h.i.) of the periodic orbit signal in (93) to the functional form of (92). \[In practice, we will again convolute both expressions with a Gaussian function (cf. Section 5.1) in order to regularize the $`\delta `$ functions in (93).\] The frequencies $`\overline{w}_k`$ of the periodic orbit signal (93) are the semiclassical eigenvalues, averaged over different orders of $`\mathrm{}`$. Note that the accuracy of these values does not necessarily increase with increasing order $`n`$. We indicate this in (92) by omitting the index $`n`$ at the eigenvalues $`\overline{w}_k`$. Our numerical calculations for the first order $`\mathrm{}`$ corrections show that, in practice, the frequencies $`\overline{w}_k`$ we obtain are approximately equal to the zeroth order $`\mathrm{}`$ eigenvalues rather than the exact average between zeroth and first oder eigenvalues. The corrections $`\mathrm{\Delta }w_{k,n}`$ to the eigenvalues are not obtained from the frequencies, but from the amplitudes, $`m_k(\overline{w}_k)^n\mathrm{\Delta }w_{k,n}`$, of the periodic orbit signal.
We will now apply the above technique to the circle billiard to obtain the first order corrections to the semiclassical eigenvalues obtained in Section 5.2.2. In Section 2.2, we derived the zeroth order amplitudes of the circle billiard (cf. Eq. (26)):
$$\frac{1}{\sqrt{w}}𝒜_{\mathrm{po}}^{(0)}=\sqrt{\frac{\pi }{2}}m_𝐌\frac{s_𝐌^{3/2}}{M_r^2}e^{i(\frac{3}{2}M_r\pi +\frac{\pi }{4})},$$
(94)
with $`s_𝐌`$ and $`m_𝐌`$ the action and multiplicity of the orbit, respectively. The first order amplitudes are given by (cf. Sections 2.4 and 4.2):
$$\frac{1}{\sqrt{w}}𝒜_{\mathrm{po}}^{(1)}=\sqrt{\pi M_r}\frac{2\mathrm{sin}^2\gamma 5}{6\mathrm{sin}^{3/2}\gamma }e^{i(\frac{3}{2}M_r\pi \frac{\pi }{4})},$$
(95)
with $`\gamma \pi M_\phi /M_r`$.
Using these expressions, we have calculated the first order corrections $`\mathrm{\Delta }w_{k,1}`$ to the lowest eigenvalues of the circle billiard, by harmonic inversion of periodic orbit signals with $`s_{\mathrm{max}}=200`$. Part of the spectrum is presented in Figure 7. The peak heights (squares) are the corrections $`\mathrm{\Delta }w_{k,1}=w_{k,1}w_{k,0}`$ times the multiplicities. For comparison, the differences between the exact and the EBK eigenvalues at the positions of the EBK eigenvalues are also plotted (see the crosses in Fig. 7). Both spectra are in excellent agreement. The small deviations of the peak heights arise from second or higher order corrections to the eigenvalues.
An appropriate measure for the accuracy of semiclassical eigenvalues is the deviation from the exact quantum eigenvalues in units of the average level spacings, $`\mathrm{\Delta }w_{\mathrm{av}}=1/\overline{\rho }(w)`$. Figure 8 presents the semiclassical error in units of the average level spacings $`\mathrm{\Delta }w_{\mathrm{av}}4/w`$ for the zeroth order (diamonds) and first order (crosses) approximations to the eigenvalues. In zeroth order approximation the semiclassical error for the low lying states is about 3 to 10 percent of the mean level spacing. This error is reduced in the first order approximation by at least one order of magnitude for the least semiclassical states with radial quantum number $`n=0`$. The accuracy of states with $`n1`$ is improved by two or more orders of magnitude.
### 5.4 Reduction of required signal length via harmonic inversion of cross-correlated periodic orbit sums
As described in the sections above, the harmonic inversion method is able to extract quantum mechanical eigenvalues from the semiclassical periodic orbit sum including periodic orbits up to a finite action $`s_{\mathrm{max}}`$. This means that in practice, although the periodic orbit sum does not converge, the eigenvalues can be obtained from a finite set of periodic orbits. The required signal length for harmonic inversion depends on the mean density of states, i.e., $`s_{\mathrm{max}}4\pi \overline{\rho }(w)`$ (cf. (43)). Depending on the mean density of states, the action $`s_{\mathrm{max}}`$ up to which the periodic orbits have to be known may therefore be large. Due to the rapid proliferation of the number of periodic orbits with growing action, the efficiency and practicability of the procedure depends sensitively on the signal length required. This is especially the case when the periodic orbits have to be found numerically.
The quantization method can be improved with the help of cross-correlated periodic orbit sums. The extended response functions weighted with products of diagonal matrix elements discussed in Section 2.3, in combination with the method for harmonic inversion of cross-correlation functions presented in Section 3.2, can be used to significantly reduce the signal length required for the periodic orbit quantization Mai99c ; Mai99b . This technique is particularly helpful for chaotic systems, where the periodic orbits must be found numerically and where the number of periodic orbits grows exponentially with their action. However, for regular systems the number of orbits which have to be included can also be significantly reduced as will be demonstrated for the circle billiard.
The basic idea is to construct a set of signals where each individual signal contains the same frequencies (i.e., semiclassical eigenvalues) and the amplitudes are correlated by obeying the restriction (53). This can be achieved with the help of the generalized periodic orbit sum (34) introduced in Section 2.3:
A set of operators $`\widehat{A}_\alpha `$, $`\alpha =1,\mathrm{},N`$ is chosen. Following the procedure described in Section 5.1, the signals $`C_{\alpha \alpha ^{}}(s)`$ are obtained as Fourier transform of the generalized response functions $`g_{\alpha \alpha ^{}}^{\mathrm{osc}}(w)`$, i.e.,
$`g_{\alpha \alpha ^{}}^{\mathrm{osc}}(w)`$ $`=`$ $`{\displaystyle \underset{\mathrm{po}}{}}𝒜_{\mathrm{po}}\overline{A}_{\alpha ,p}\overline{A}_{\alpha ^{},p}e^{is_{\mathrm{po}}w},`$ (96)
$`C_{\alpha \alpha ^{}}(s)`$ $`=`$ $`{\displaystyle \underset{\mathrm{po}}{}}𝒜_{\mathrm{po}}\overline{A}_{\alpha ,p}\overline{A}_{\alpha ^{},p}\delta (ss_{\mathrm{po}}),`$ (97)
where for integrable systems the means $`\overline{A}_{\alpha ,p}`$ are defined by (32). According to Sections 2.3 and 5.1, the corresponding quantum mechanical signal is given by
$$C_{\mathrm{qm},\alpha \alpha ^{}}(s)=i\underset{k}{}m_kk|\widehat{A}_\alpha |kk|\widehat{A}_\alpha ^{}|ke^{isw_k},$$
(98)
where the amplitudes have the required form (53). As in Section 5.1 the semiclassical eigenvalues $`w_k`$ are obtained by adjusting the periodic orbit signal (97) (after convolution with a Gaussian function) to the functional form of the cross-correlated quantum signal (98) with the important difference that we now apply the extension of harmonic inversion to cross-correlation functions (see Section 3.2).
For the circle billiard, the mean density of states – with all multiplicities taken as one – is given by $`\overline{\rho }(w)=w/4`$. According to (43), the signal length required for a single signal to resolve the frequencies in a given interval around $`w`$ is therefore approximately given by
$$s_{\mathrm{max}}4\pi \overline{\rho }(w)=\pi w=2S_H,$$
(99)
where $`S_H=2\pi \overline{\rho }`$ is the Heisenberg period (which is action instead of time for scaling systems). By using an $`N\times N`$ set of signals, it should be possible to extract about the same number of semiclassical eigenvalues from a reduced signal length $`s_{\mathrm{max}}2S_H`$, or, vice versa, if the signal length is held constant, the resolution and therefore the number of converged eigenvalues should significantly increase.
To demonstrate the power of the cross-correlation technique, we first analyze a $`2\times 2`$ cross-correlated periodic orbit signal of the circle billiard with $`\widehat{A}_1=\mathrm{𝟏}`$ the identity operator and $`\widehat{A}_2=r`$. For comparison with the results in Section 5.2.2 we choose the same signal length $`s_{\mathrm{max}}=150`$. By contrast to the eigenvalues in Table 1 obtained from the one-dimensional signal the nearly degenerate states around $`w11.05`$ and $`w13.3`$ are now resolved as can be seen in Table 2. Note that a signal length $`s_{\mathrm{max}}500`$ is required to resolve these states without application of the cross-correlation technique.
As in all other calculations concerning cross-correlated signals, the results were improved by not making a sharp cut at the accumulation points but by damping the amplitudes of the orbits near these points. With the same cut-off criterion at the accumulation points, the total number of orbits in the signal with $`s_{\mathrm{max}}=150`$ was about 10 times smaller than in the signal with $`s_{\mathrm{max}}=500`$ . This means, we could reduce the required number of orbits by one order of magnitude. For chaotic systems, where the number of orbits grows more rapidly (exponentially) with the maximum action, the improvement in the required number of orbits may even be better.
We now investigate the number of eigenvalues which do converge for fixed signal length but different sets and dimension of the cross-correlation matrix. Indeed, the highest eigenvalue $`w_{\mathrm{max}}`$ which can be resolved increases significantly when the cross-correlation technique is applied. However, the detailed results depend on the operators chosen. Furthermore, with increasing dimension of the matrix, the range in which the transition from resolved to unresolved eigenvalues takes place becomes broader, and the amplitudes in this region become less well converged. Rough estimates of $`w_{\mathrm{max}}`$ for various sets of operators and fixed signal length $`s_{\mathrm{max}}=150`$ are given in Table 3. For some of the signals, we used the extension of the trace formula to functions of matrix elements, Eqs. (35) and (36). The improvement achieved by increasing the dimension of the matrix by one is most distinct for very small $`N`$; for $`N5`$, the improvement is only small. This suggests that, at given signal length and frequency range, the matrix dimension should not be chosen too large, i.e., there exists an optimal matrix dimension, which at constant signal length becomes larger with increasing eigenvalues $`w`$.
With a $`5\times 5`$ signal of length $`s_{\mathrm{max}}=150`$, eigenvalues up to the region $`w130`$ can be resolved. The results are presented in Figure 9. There are two points which should be emphasized: The first point is that, even in this dense part of the spectrum, the error of the method is still by about one order of magnitude smaller than the semiclassical error, which is illustrated in Fig. 9 by the peak heights. The squares and crosses mark the semiclassical error $`|w_{\mathrm{EBK}}w_{\mathrm{ex}}|`$ and the numerical error $`|w_{\mathrm{hi}}w_{\mathrm{EBK}}|`$ of the harmonic inversion procedure in units of the mean level spacing $`\mathrm{\Delta }w_{\mathrm{av}}4/w`$. The second point concerns the signal length compared to the Heisenberg action $`S_H=2\pi \overline{\rho }`$. For $`w=130`$, one obtains $`S_H204.2`$. A one-dimensional signal would have required a signal length $`s_{\mathrm{max}}2S_H`$. With the cross-correlation technique, we calculated the eigenvalues from a signal of length $`s_{\mathrm{max}}=1500.735S_H`$. This is about the same signal length as required by the semiclassical quantization method of Berry and Keating Ber90 , which, however, only works for ergodic systems.
In summary, our results demonstrate that by analyzing cross-correlated signals instead of a single signal, the required signal length can indeed be significantly reduced. Clearly, the signal length cannot be made arbitrarily small, and the method is restricted to small dimensions of the cross-correlation matrix. However, the number of orbits which have to be included can be very much reduced. Another advantage of the method is that not only the frequencies and the multiplicities but also the diagonal matrix elements of the chosen operators are obtained by one single calculation.
#### 5.4.1 Including higher order $`\mathrm{}`$ corrections
In the cases discussed so far, we have constructed the cross-correlated signal by including different operators and making use of Eq. (34). By this procedure, we could obtain the semiclassical eigenvalues from a signal of reduced length or improve the resolution of the spectrum at constant signal length, while simultaneously obtaining the diagonal matrix elements of the operators. We can now even go one step further and include higher $`\mathrm{}`$ corrections in the signal. Here we make use of the results of Section 5.3. The first order correction term, which in Section 5.3 was harmonically inverted as a single signal, is now included as part of a cross-correlated signal. This procedure combines all the techniques developed in the previous sections.
Formally, the frequencies in the zeroth and first order $`\mathrm{}`$ parts of the cross-correlated signal are not exactly the same \[see the denominators in Eqs. (75) and (90)\], however, as already mentioned in Section 5.3, the numerically obtained values for the frequencies $`\overline{w}_k`$ in (92) are equal to the lowest order $`\mathrm{}`$ eigenvalues rather than the exact average of the zeroth and first order eigenvalues. In practice, the cross-correlated signal is therefore in fact of the form (52). We can now, on the one hand, improve the resolution of the spectrum, and, on the other hand, obtain semiclassical matrix elements and the first order corrections to the eigenvalues with the same high resolution by one single harmonic inversion of a cross-correlated signal.
As an example, we built a $`3\times 3`$ signal for the circle billiard from the first order correction term given by (95) and the operators $`\widehat{A}_1=\mathrm{𝟏}`$ (identity) and $`\widehat{A}_2=\widehat{r}`$. Again, we chose a signal length of $`s_{\mathrm{max}}=150`$. By harmonic inversion of the cross-correlated signal, we obtained simultaneously the semiclassical eigenvalues, their first order order corrections, and the semiclassical matrix elements of the operator $`\widehat{r}`$. The results for the zeroth order approximations $`w_{k,0}`$ to the eigenvalues and the first order approximations $`w_{k,1}=w_{k,0}+\mathrm{\Delta }w_{k,1}`$ are presented in Table 4. For comparison the exact and the EBK eigenvalues are also given. As for the results presented in Table 2, we were able to resolve the nearly degenerate states in the zeroth order approximation, which for a single signal would have required a signal length of $`s_{\mathrm{max}}500`$. Moreover, in contrast to the results of Section 5.3, we could now also resolve the first order approximations to the nearly degenerate states.
## 6 Conclusion
The harmonic inversion method has been introduced as a powerful tool for the calculation of quantum eigenvalues from periodic orbit sums as well as for the high resolution analysis of quantum spectra in terms of classical periodic orbits. We have demonstrated that this method, which has already successfully been applied to classically chaotic systems, yields excellent results for regular systems as well. Harmonic inversion has thus been shown to be a universal method, which, in contrast to other high resolution methods, does not depend on special properties of the system such as ergodicity or the existence of a symbolic code.
With the harmonic inversion method, we are able to calculate the contributions of the classical periodic orbits to the trace formula from the quantum eigenvalues with high precision and high resolution. By analyzing the difference spectrum between exact and semiclassical eigenvalues, we could determine higher order $`\mathrm{}`$ corrections to the periodic orbit sum of the circle billiard.
Up to now, no theory for the $`\mathrm{}`$ corrections to the Berry-Tabor formula for regular systems has been developed. We have numerically found an expression for the first order $`\mathrm{}`$ corrections to the Berry-Tabor formula by harmonic inversion of the difference spectrum. The same expression can be derived analytically by using Vattay’s and Rosenqvist’s method for chaotic systems and introducing some reasonable ad-hoc assumption for the circle billiard. As this is clearly not a strict derivation, it is an interesting task for the future to develop a general theory for the higher order $`\mathrm{}`$ corrections to the trace formula for regular systems.
In addition to calculating semiclassical eigenvalues from the usual periodic orbit sum, we have demonstrated how further information can be extracted from the parameters of the classical orbits by applying the harmonic inversion technique to different extensions of the trace formula. Using a generalized trace formula including an arbitrary operator, we have shown that the method also allows the calculation of semiclassical diagonal matrix elements from the parameters of the periodic orbits. Furthermore we have demonstrated how higher order $`\mathrm{}`$ corrections to the semiclassical eigenvalues can be obtained by harmonic inversion of correction terms to the periodic orbit sums. For the case of the circle billiard, we found that, including the first order correction, the accuracy of the semiclassical eigenvalues compared to the exact quantum eigenvalues could be improved by one or more orders of magnitude.
Although by harmonic inversion the quantum eigenvalues can be calculated from a semiclassical signal of finite length, i.e., from a finite set of periodic orbits, the number of orbits which have to be included may still be large. We have demonstrated that by a generalization of the harmonic inversion method to cross-correlation functions the required signal length may be significantly reduced, even below the Heisenberg time. Because of the rapid proliferation of periodic orbits with growing period, this means that the number of orbits which have to be included may be reduced by about one to several orders of magnitude.
## Appendix A Calculation of the first order $`\mathrm{}`$ correction terms to the semiclassical trace formula
The calculation of higher order $`\mathrm{}`$ corrections to the semiclassical eigenvalues introduced in Section 2.4 requires the knowledge of the $`n`$th order amplitudes $`𝒜_{\mathrm{po}}^{(n)}`$ in the periodic orbit sum (39). In this Appendix, we briefly outline the derivation of the first order amplitudes $`𝒜_{\mathrm{po}}^{(1)}`$ and the application to the circle billiard given in Eq. (40).
Two different methods for the calculation of higher order $`\mathrm{}`$ correction terms in chaotic systems have been derived by Gaspard and Alonso Alo93 ; Gas93 and Vattay and Rosenqvist Vat94 ; Vat96 ; Ros94 . The latter method has been specialized to two-dimensional chaotic billiards in Ros94 . Here, we follow the approach of Vattay and Rosenqvist. However, it is important to note that both methods cannot straightforwardly be applied to integrable systems and additional assumptios will be necessary to derive Eq. (40) for the circle billiard. A general theory for the calculation of higher order $`\mathrm{}`$ corrections to the Berry-Tabor formula (4) for integrable systems is, to the best of our knowledge, still lacking.
Vattay and Rosenqvist give a quantum generalization of Gutzwiller’s trace formula based on the path integral representation of the quantum propagator. The basic idea of their method is to express the global eigenvalue spectrum in terms of local eigenvalues computed in the neighbourhood of periodic orbits. The energy domain Green function $`G(q,q^{},E)`$ is connected to the spectral determinant $`\mathrm{\Delta }(E)=\mathrm{\Pi }_n(EE_n)`$, with $`E_n`$ the energy eigenvalues or resonances, by
$$\mathrm{Tr}G(E)=𝑑qG(q,q,E)=\frac{d}{dE}\mathrm{ln}\mathrm{\Delta }(E).$$
(100)
The trace of the Green function can be expressed in terms of contributions from periodic orbits
$$\mathrm{Tr}G(E)=\underset{\mathrm{p}.\mathrm{o}.}{}\mathrm{Tr}G_p(E),$$
(101)
with the local traces connected to the local spectral determinants by
$$\mathrm{Tr}G_p(E)=\frac{d}{dE}\mathrm{ln}\mathrm{\Delta }_p(E).$$
(102)
The trace of the Green function can therefore be calculated by solving the local Schrödinger equation around each periodic orbit, which yields the local eigenspectra.
To obtain the local eigenspectra, the ansatz
$$\psi =\mathrm{\Phi }e^{iS/\mathrm{}}$$
(103)
is inserted into the Schrödinger equation, yielding the following differential equations for $`\mathrm{\Phi }`$ and $`S`$.
$`_tS+{\displaystyle \frac{1}{2}}(S)^2+U=0`$ (104)
$`_t\mathrm{\Phi }+\mathrm{\Phi }S+{\displaystyle \frac{1}{2}}\mathrm{\Phi }\mathrm{\Delta }S{\displaystyle \frac{i\mathrm{}}{2}}\mathrm{\Delta }\mathrm{\Phi }=0,`$ (105)
where $`U`$ is the potential entering the Schrödinger equation.
The spectral determinant can be calculated from the local eigenvalues of the amplitudes $`\mathrm{\Phi }`$. For arbitrary energy E, the amplitudes $`\mathrm{\Phi }_p^l`$ of the local eigenfunctions fulfill the equation
$$\mathrm{\Phi }_p^l(t+T_p)=R_p^l(E)\mathrm{\Phi }_p^l(t),$$
(106)
where $`T_p`$ is the period of the classical orbit. Using Eq. (102), the trace formula can be expressed in terms of the eigenvalues $`R_p^l(E)`$:
$`\mathrm{Tr}G(E)={\displaystyle \frac{1}{i\mathrm{}}}`$ $`{\displaystyle \underset{p}{}}{\displaystyle \underset{l}{}}\left(T_p(E)i\mathrm{}{\displaystyle \frac{d\mathrm{ln}R_p^l(E)}{dE}}\right)`$ (107)
$`\times `$ $`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}(R_p^l(E))^re^{\frac{i}{\mathrm{}}rS_p(E)}.`$
This is the quantum generalization of Gutzwiller’s trace formula and holds exactly.
The amplitudes and their eigenvalues are now expanded in powers of $`\mathrm{}`$:
$`\mathrm{\Phi }^l`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{i\mathrm{}}{2}}\right)^m\mathrm{\Phi }^{l(m)}`$ (108)
$`R^l(E)`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{i\mathrm{}}{2}}\right)^mC_l^{(m)}\right\}`$ (109)
$``$ $`\mathrm{exp}(C_l^{(0)})\left(1+{\displaystyle \frac{i\mathrm{}}{2}}C_l^{(1)}+\mathrm{}\right)`$ (110)
The terms $`C_l^{(0)}`$ yield the Gutzwiller trace formula as zeroth order approximation, while the terms with $`m>0`$ give $`\mathrm{}`$ corrections.
To solve Eqs. (104) and (105) in different order of $`\mathrm{}`$, the Schrödinger equation and the functions $`\mathrm{\Phi }^{l(m)}`$ and $`S`$ are Taylor expanded around the periodic orbit,
$`S(𝐪,t)`$ $`=`$ $`{\displaystyle \frac{1}{𝐧!}s_𝐧(t)(𝐪𝐪_p(t))^𝐧}`$ (111)
$`\mathrm{\Phi }^{l(m)}(𝐪,t)`$ $`=`$ $`{\displaystyle \frac{1}{𝐧!}\varphi _𝐧^{l(m)}(t)(𝐪𝐪_p(t))^𝐧},`$ (112)
resulting in a set of differential equations for the different orders of the Taylor expansions and different orders in $`\mathrm{}`$. In one dimension, these equations read explicitly:
$$\dot{s}_ns_{n+1}\dot{q}+\frac{1}{2}\underset{k=0}{\overset{n}{}}\frac{n!}{(nk)!k!}s_{nk+1}s_{k+1}+u_n=0,$$
(113)
where $`u_n`$ are the coefficients of the Taylor expanded potential, and
$`\dot{\varphi }_n^{(m+1)}`$ $``$ $`\varphi _{n+1}^{(m+1)}\dot{q}`$ (114)
$`+`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{n!}{(nk)!k!}}\times `$
$`(\varphi _{nk+1}^{(m+1)}s_{k+1}+{\displaystyle \frac{1}{2}}\varphi _{nk}^{(m+1)}s_{k+2})`$
$`\varphi _{n+2}^{(m)}=0`$
This set of differential equations can be solved iteratively. The $`l`$-th eigenfunction is characterized by the condition $`\varphi _n^{(m)}0`$ for $`n<l`$. The different orders of $`\mathrm{}`$ are connected by the last term in (114). Starting from zeroth order $`\mathrm{}`$ and the lowest nonvanishing order of the Taylor expansion, the functions can be determined order by order. For higher dimensional systems, the coefficient matrices obey similar equations, and the structure of the set of equations remains the same.
To obtain the first order $`\mathrm{}`$ correction to the Gutzwiller trace formula, one has to calculate the quantities $`C_l^{(1)}`$. To obtain these quantities one has to solve the set of equations (114) up to the lowest nonvanishing first order $`\mathrm{}`$ coefficient function $`\varphi _l^{l(1)}`$, respectively. As $`\varphi _n^{l(m)}0`$ for $`n<l`$, this involves solving the equations for $`s_2`$, $`s_3`$ and $`s_4`$, and for the zeroth order $`\mathrm{}`$ coefficient functions $`\varphi _l^{l(0)}`$, $`\varphi _{l+1}^{l(0)}`$ and $`\varphi _{l+2}^{l(0)}`$. If the initial conditions are set to be $`\mathrm{\Phi }_l^{l(0)}(0)=1`$ and $`\mathrm{\Phi }_l^{l(m)}(0)=0`$ for $`m>0`$, the correction term $`C_l^{(1)}`$ is then given by the relation
$$C_l^{(1)}=\frac{\varphi _l^{l(1)}(T_p)}{\mathrm{exp}(C_l^{(0)})},$$
(115)
which follows from the $`\mathrm{}`$ expansion of the eigenequation (106).
An explicit recipe for the calculation of the first $`\mathrm{}`$ correction for two-dimensional chaotic billiards is given in Ros94 . For billiards, the potential $`U`$ in the Schrödinger equation equals zero between two bounces at the hard wall. The functions $`S`$ and $`\mathrm{\Phi }`$ now have to be Taylor expanded in two dimensions:
$`S(x,y,t)`$ $`=`$ $`S_0+S_x\mathrm{\Delta }x+S_y\mathrm{\Delta }y`$ (116)
$`+`$ $`{\displaystyle \frac{1}{2}}(S_{x^2}(\mathrm{\Delta }x)^2+2S_{xy}\mathrm{\Delta }x\mathrm{\Delta }y+S_{y^2}(\mathrm{\Delta }y)^2)`$
$`+`$ $`\mathrm{}`$
and similarly for $`\mathrm{\Phi }`$. If the coordinate system is chosen in such a way that $`x`$ is along the periodic orbit and $`y`$ is perpendicular to the orbit, derivatives with respect to $`x`$ can be expressed in terms of the derivatives with respect to $`y`$ using the stationarity conditions
$$S_{x^{n+1}y^m}=\frac{\dot{S}_{x^ny^m}}{S_x},\varphi _{x^{n+1}y^m}=\frac{\dot{\varphi }_{x^ny^m}}{S_x}.$$
(117)
The quantity $`S_x`$ is equal to the classical momentum of the particle.
For the free motion between the collisions with the wall, the set of differential equations correspoding to (113) and (114) then reduces to a set of equations involving only derivatives with respect to $`y`$. These equations can be solved analytically, with the general solution still containing free parameters. Setting $`S_x=1`$, the first coefficient functions of the Taylor expanded phase are given by:
$`S_{yy}(t)`$ $`=`$ $`{\displaystyle \frac{1}{t+t_0}}`$ (118)
$`S_{yyy}(t)`$ $`=`$ $`{\displaystyle \frac{A}{(t+t_0)^3}}`$ (119)
$`S_{yyyy}(t)`$ $`=`$ $`{\displaystyle \frac{3}{(t+t_0)^3}}+{\displaystyle \frac{B}{(t+t_0)^4}}+{\displaystyle \frac{3A^2}{(t+t_0)^5}}`$ (120)
where $`t_0`$, $`A`$ and $`B`$ are free parameters. For given $`l`$, the first nonvanishing coefficients of the amplitude read
$`\varphi _{y^l}^{(0)}(t)`$ $`=`$ $`E\left({\displaystyle \frac{t_0}{t+t_0}}\right)^{l+1/2}`$ (121)
$`\varphi _{y^{l+1}}^{(0)}(t)`$ $`=`$ $`{\displaystyle \frac{E}{(t+t_0)^{l+3/2}}}\left[C+(l+1)^2{\displaystyle \frac{A}{2}}{\displaystyle \frac{t_0^{l+1/2}}{(t+t_0)}}\right]`$ (122)
$`\varphi _{y^{l+2}}^{(0)}(t)`$ $`=`$ $`{\displaystyle \frac{E}{(t+t_0)^{l+5/2}}}\{D+{\displaystyle \frac{1}{t+t_0}}\times `$
$`\left[(l+2)^2{\displaystyle \frac{AC}{2}}+(l+2)(l+1)({\displaystyle \frac{l}{3}}+{\displaystyle \frac{1}{2}}){\displaystyle \frac{B}{2}}t_0^{l+1/2}\right]`$
$`+`$ $`{\displaystyle \frac{A^2t_0^{l+1/2}}{2(t+t_0)^2}}\times `$
$`[{\displaystyle \frac{1}{4}}(l+2)^2(l+1)^2+{\displaystyle \frac{3}{2}}(l+2)(l+1)({\displaystyle \frac{l}{3}}+{\displaystyle \frac{1}{2}})]\}`$
Again, $`D`$ and $`E`$ are free parameters.
At the collisions with the hard wall, the phase and amplitude have to obey the bouncing conditions
$`S(x,y,t_0)`$ $`=`$ $`S(x,y,t_{+0})+i\pi `$ (124)
$`\mathrm{\Phi }(x,y,t_0)`$ $`=`$ $`\mathrm{\Phi }(x,y,t_{+0}),`$ (125)
from which the bouncing conditions for the coefficients of the Taylor expanded functions $`S`$ and $`\mathrm{\Phi }`$ can be derived. While the general solutions between the bounces are valid for all billiards, the bouncing conditions in their Taylor expanded form depend explicitely on the shape of the hard wall.
An additional condition which the solutions $`S`$ and $`\mathrm{\Phi }`$ have to obey is periodicity along the orbit. With every traversal the phase gains the same constant contributions at the collisions with the wall. The derivatives of the phase are periodic. The amplitude collects the same factor with each traversal, which means that all Taylor coefficients of the amplitude are periodic apart from a constant factor. These conditions together with the bouncing conditions determine the values of the free constants in the general solutions between the collisions.
The solutions can in general be found numerically, by choosing suitable initial conditions and following the evolution of the phase and amplitude functions along the orbit. After several iterations around the orbit, the parameters should converge against their periodic solution. The correction terms $`C_l^{(1)}`$ are then given by the integral
$$C_l^{(1)}=\frac{\varphi _{y^{l+2}}^{l(0)}+\varphi _{y^lx^2}^{l(0)}}{\varphi _{y^l}^{l(0)}}𝑑t,$$
(126)
which can be computed explicitly from the solutions found above.
As already explained, this method is designed for chaotic systems, as its derivation is based on the assumption that the periodic orbits are isolated. Nevertheless, we obtain reasonable results when applying the method to the circle billiard, taking one periodic orbit from each rational torus and introducing some additional assumptions.
Because of the symmetry of the orbits, we can assume that every side of the orbit contributes in the same manner to the $`\mathrm{}`$ correction term for the whole orbit. This means, if we reset $`t=0`$ at the start of each side, the free parameters in the general solutions (118) to (LABEL:gen\_sol\_phi3) must be the same for each side, apart from the parameter $`E`$, which collects the same factor during every collison with the wall. With these assumptions, the differential equations can be solved analytically.
However, it turns out that the bouncing conditions resulting from (124) and (125) are not sufficient to determine all free parameters, as some of the conditions are automatically fulfilled. We need additional conditions for the parameters. These can be obtained from the rotational symmetry of the system: Because of this symmetry, we can assume that the amplitude of the wave function does not depend on the polar angle $`\phi `$. The same holds for all derivatives of the amplitude with respect to the radius $`r`$. For the zeroth order $`\mathrm{}`$ amplitudes, expressed in polar coordinates $`(r,\phi )`$, this gives us the additional conditions we need:
$$\frac{}{\phi }\frac{^n\mathrm{\Phi }^{l(0)}}{r^n}=0.$$
(127)
If we further assume that the phase separates in polar coordinates
$$S(r,\phi )=S_r(r)+S_\phi (\phi ),$$
(128)
which implies that all mixed derivatives vanish, it turns out that we do not need the bouncing conditions at all. All parameters can be determined from the symmetry of the system, and the bouncing conditions are automatically fulfilled. We considered only the case $`l=0`$, for which we used Eq. (128) together with the conditions
$$\frac{}{\phi }\mathrm{\Phi }^{0(0)}=0,\frac{}{\phi }\frac{\mathrm{\Phi }^{0(0)}}{r}=0.$$
(129)
Our final results for the constant parameters in Eqs. (118) to (LABEL:gen\_sol\_phi3) are
$`t_0`$ $`=`$ $`\mathrm{sin}\gamma `$ (130)
$`A`$ $`=`$ $`\mathrm{cos}\gamma `$ (131)
$`B`$ $`=`$ $`0`$ (132)
$`C`$ $`=`$ $`0`$ (133)
$`D`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{sin}^{1/2}\gamma `$ (134)
with $`\gamma `$ as defined in Section 2.4 (see Fig. 1). The radius of the billiard was taken to be $`R=1`$. Inserting these solutions in (126) finally leads to
$$C_0^{(1)}=M_r\left(\frac{1}{3\mathrm{sin}\gamma }\frac{5}{6\mathrm{sin}^3\gamma }\right)$$
(135)
where $`M_r`$ is the number of sides of the orbit.
The first order amplitudes $`𝒜^{(1)}`$ are obtained by inserting the $`\mathrm{}`$ expansion (110) in the trace formula (107) and comparing the result with the $`\mathrm{}`$ expansion (39). In the units we have used here (radius $`R=1`$ and momentum $`\mathrm{}k=1`$), the scaling parameter $`w`$ is equal to $`\mathrm{}`$. If we use only the $`l=0`$ contributions and assume that the terms $`\mathrm{exp}(C_0^{(0)})`$ are equal to the amplitudes given by the Berry-Tabor formula, we finally end up with the expression (40). Although we cannot strictly justify the last step, our analysis of the quantum spectrum in Section 4.2 provides strong numerical evidence that Eq. (40) is correct. It will be an interesting task for the future to develop a general theory for the $`\mathrm{}`$ correction terms of integrable systems and thus to provide a more rigorous mathematical proof of Eq. (40).
## Acknowledgements
This work was supported by the Deutsche Forschungsgemeinschaft.
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# Objective Classification of Galaxy Spectra using the Information Bottleneck Method
## 1 Introduction
Very large numbers of galaxy spectra are being generated by modern redshift surveys: for example, the Anglo-Australian Observatory 2-degree-Field (2dF) Galaxy Survey (e.g. Colless 1998; Folkes et al. 1999) aims to collect 250,000 spectra and has already gathered over $`100,000`$ redshifts (as of May 2000). The Sloan Digital Sky Survey (e.g. Gunn & Weinberg 1995; Fukugita 1998) will observe the spectra of millions of galaxies. Such large data sets will provide a wealth of information pertaining to the distribution and properties for a vast variety of galaxy types. A major goal of such surveys is to determine the relative numbers of these different galaxy populations, and eventually to gain clues as to their physical origin. Traditional methods of classifying galaxies “by eye” are clearly impractical in this context. The analysis and full exploitation of such data sets require well justified, automated and objective techniques to extract as much information as possible. In this paper we present a new approach to objective classification of galaxy spectra, by utilising a recently proposed method based on information theory.
Broadly speaking, inference from galaxy spectra can be considered in two ways. One approach is to compare each galaxy spectrum to the most likely one of a library of model spectra (e.g. based on age, metallicity and star formation history). The other model independent approach is to consider an ensemble of observed spectra and to look for patterns in analogy with the stellar HR-diagram or the Hubble sequence for galaxy morphology.
The concept of spectral classification goes back to Humason (1936) and Morgan & Mayall (1957). Recent attempts to analyse galaxy properties from spectra in a model independent way have been made using Principal Component Analysis (e.g. Connolly et al. 1995; Folkes, Lahav & Maddox 1996). The PCA is effective for data compression, but if one wishes to break the ensemble into classes it requires a further step based on a training set (e.g. Bromley et al. 1998, Folkes et al. 1999). Unlike PCA, the information bottleneck (IB) method presented here is non-linear, and it naturally yields a principled partitioning of the galaxies into classes. These classes are obtained such that they maximally preserve the original information between the galaxies and their spectra.
The end goal of galaxy classification is a better understanding of the physical origin of different populations and how they relate to one another. In order to interpret the results of any objective classification algorithm, we must relate the derived classes to the physical and observable galaxy properties that are intuitively familiar to astronomers. For example, important properties in determining the spectral characteristics of a galaxy are its mean stellar age and metallicity, or more generally its full star formation history. This is in turn presumably connected with *morphological* properties of the galaxy; eg. its Hubble or T-type. An assumed star formation history can be translated into a synthetic spectrum using models of stellar evolution (e.g., Bruzual & Charlot 1993, 1996; Fioc & Rocca-Volmerange 1997). Spectral features are also affected by dust reddening and nebular emission lines.
In one example, typical of such an approach (Ronen, Aragaon-Salamance & Lahav 1999), the star formation history was parameterized as a simple single burst or an exponentially decreasing star formation rate. However, the construction of the ensemble of galaxy spectra was done in an ad-hoc manner. Here we try a similar exercise using a *cosmologically motivated ensemble* of synthetic galaxies, with realistic star formation histories. These histories are determined by the physical processes of galaxy formation in the context of hierarchical structure formation. We construct such an ensemble using semi-analytic models of galaxy formation set within the Cold Dark Matter (CDM) framework. While this approach has the disadvantage of relying on numerous assumptions about poorly-understood physics, it has the advantage of a certain self-consistency, and of producing ensembles of galaxies with global properties that agree well with observations (eg. local luminosity functions, colours, Tully-Fisher relations, metallicity-luminosity relation, etc).
In this paper, we shall analyze a subset of about $`6000`$ spectra from the ongoing 2dF survey. Using the semi-analytic models we construct a “mock-2dF” catalogue with the same magnitude limit, spectral resolution, wavelength coverage and noise as in the 2dF data. We analyze the synthetic and observed spectra in the same way using the IB method. Using the mock catalogue, we then interpret the resulting classes in terms of familiar, intuitive physical properties of galaxies. We also draw a connection between the IB approach and the results of a PCA analysis for both the 2dF and mock data.
The organization of the rest of this paper is as follows. In section 2 we present the IB method and the classification algorithm. In section 3 we describe the sample of observed spectra and the semi-analytic models used to produce the mock galaxy catalogues. In section 4 we present our results, and in section 5 we relate the IB classification to PCA. We summarize our conclusions in section 6.
## 2 The Classification Algorithm: the Information Bottleneck Method
Consider a galaxy spectrum as an array of photon counts in different wavelength bins. <sup>1</sup><sup>1</sup>1 We note that other representations of the spectra are possible (e.g. flux instead of counts). Thus each galaxy is represented in a high-dimensional space, where each component corresponds to the counts in a given spectral bin. We can also view the ensemble of such spectra as the joint distribution of the galaxy-wavelength variables. By normalizing the total photon counts in each spectrum to unity, we can consider it as a conditional probability, the probability of observing a photon at a specific wavelength from a given galaxy. This view of the ensemble of spectra as a conditional probability distribution function enables us to undertake the information theory-based approach that we describe in this section.
Our goal is to group the galaxies into classes that preserve some objectively defined spectral properties. Ideally, we would like to make the number of classes as small as possible (i.e. to find the ‘least complex’ representation) with minimal loss of the ‘important’ or ‘relevant’ information. In order to do this objectively, we need to define formal measures of ‘complexity’ and ‘relevant information’.
Some classification methods are based on a training set of labeled data (e.g. morphological types of galaxies defined by a human expert; e.g. Naim et al. 1995a,b). Such prior labels introduce a bias towards existing classification schemes. On the other hand, our goal here is to develop an unsupervised classification method, which is free of this bias, and thus to provide objective, ‘meaningful’, categorization. This problem, however, is ill defined without a better definition of ‘meaningful’. Almost all existing algorithms begin with some pairwise ‘distances’ (e.g. Euclidean) between the points in the high-dimensional representation space, or with a distortion measure between the data points and candidate group representative or model. The ‘meaning’ is then dictated through this, sometimes arbitrary, choice of the distance or distortion measure. In addition, it is difficult in such cases to objectively evaluate the quality of the obtained classes.
Recently, Tishby, Pereira & Bialek (1999) proposed an information theoretical approach to this problem which avoids the arbitrary choice of the distance or distortion measures. It also provides a natural quality measure for the resulting classification. Their algorithm is extremely general and may be applied to different problems in analogous ways. This method has been successfully applied to the analysis of neural codes (Bialek, Nemenman & Tishby 2000), linguistic data (word sense-disambiguation, Pereira, Tishby & Lee, 1993) and for classification of text documents (Slonim & Tishby 2000). In the latter case for example, one may see an analogy between an ensemble of galaxy spectra and a set of text documents. The words in a document play a similar role to the wavelengths of photons in a galaxy spectrum, i.e. the frequency of occurrence of a given word in a given document is equivalent to the number of photon counts at a given wavelength in a given galaxy spectrum. In both cases, the specific patterns of these occurances may be used in order to classify the galaxies or documents.
### 2.1 The concept of mutual information
In the following we denote the set of galaxies by $`G`$ and the array of wavelength bins by $`\mathrm{\Lambda }`$. As already mentioned, we view the ensemble of spectra as a joint distribution $`p(g,\lambda )`$, which is the joint probability of observing a photon from galaxy $`gG`$ at a wavelength $`\lambda \mathrm{\Lambda }`$. We normalize the total photon counts in each spectrum (galaxy) to unity, i.e. we take the prior probability $`p(g)`$ of observing a galaxy $`g`$ to be uniform: $`p(g)=\frac{1}{N_G}`$, where $`N_G`$ is the number of galaxies in this sample $`G`$. This is a standard statistical procedure since these galaxies are sample points from the underlying (unknown) galaxy distribution. Note, however, that the prior probability $`p(\lambda )`$ of observing a photon in a wavelength $`\lambda `$ is not considered uniform.
Given two random variables, a fundamental question is: to what extent can one variable be predicted from knowledge of the other variable? Clearly, when the two variables are statistically independent, no information about one variable can be obtained through knowledge of the other one. This question is quantitatively answered through the notion of mutual information between the variables. Moreover, this is the only possible measure, up to a multiplicative constant, that captures our intuition about the information between variables, as was shown by Shannon (1948a,b) in his seminal work (see also Acze‘l and Darotzky, 1975). In our case, the two variables are the galaxy identity and the photon counts as a function of wavelength (spectral densities). These variables are clearly not independent, as the galaxy identity determines its spectrum through its physical properties.
The mutual information between two variables can be shown (see e.g. Cover & Thomas 1991) to be given by the amount of uncertainty in one variable that is removed by the knowledge of the other one. In our case this is the reduction of the uncertainty in the galaxy identity through the knowledge of its spectrum. The uncertainty of a random variable is measured by its entropy, which for the case of the galaxy variable $`G`$ is given by
$$H(G)=\underset{g}{}p(g)\mathrm{log}p(g).$$
(1)
Since in our case $`p(g)`$ is uniform (over the sample) we get $`H(G)=\mathrm{log}N_G`$. <sup>2</sup><sup>2</sup>2 When we take the logarithm to base $`2`$ the information is measured in bits. This means that in the absence of other knowledge, the amount of information needed to specify a galaxy $`g`$ out of the sample $`G`$ is exactly $`\mathrm{log}N_G`$ bits.
The amount of uncertainty in the galaxy identity, given its spectral density, is given by the conditional entropy of the galaxies on their spectra. More formally stated, the conditional entropy of $`G`$ given $`\mathrm{\Lambda }`$ is defined by
$$H(G|\mathrm{\Lambda })=\underset{\lambda }{}p(\lambda )\underset{g}{}p(g|\lambda )\mathrm{log}p(g|\lambda ).$$
(2)
Obviously, knowledge about $`\mathrm{\Lambda }`$ can only reduce the uncertainty in $`G`$, i.e. $`H(G|\mathrm{\Lambda })H(G)`$. The amount of reduction in the uncertainty is thus the mutual information, which is now given by
$$I(G;\mathrm{\Lambda })=H(G)H(G|\mathrm{\Lambda })=\underset{g,\lambda }{}p(g,\lambda )\mathrm{log}\frac{p(g,\lambda )}{p(g)p(\lambda )},$$
(3)
or, using $`p(g,\lambda )=p(g)p(\lambda |g)`$,
$$I(G;\mathrm{\Lambda })=\underset{g,\lambda }{}p(g)p(\lambda |g)\mathrm{log}\frac{p(\lambda |g)}{p(\lambda )}.$$
(4)
It is easy to see that $`I(G;\mathrm{\Lambda })`$ is symmetric and non-negative, and is equal to zero if and only if $`g`$ and $`\lambda `$ are independent. As $`I(G;\mathrm{\Lambda })`$ measures the reduction of uncertainty in $`G`$ for known $`\mathrm{\Lambda }`$, it is a measure of the amount of information about the galaxy identity contained in the spectrum.
### 2.2 The Bottleneck Variational Principle
Our goal is to find a mapping of the galaxies $`gG`$ into classes $`cC`$ such that the class $`c(g)`$ provides essentially the same prediction, or information, about the spectrum as the specific knowledge of the galaxy. The partitioning may be “soft”, i.e. each galaxy is associated with each class through the conditional probability $`p(c|g)`$.
The prior probability for a specific class $`c`$ is then given by
$$p(c)=\underset{g}{}p(g)p(c|g).$$
(5)
Using the fact that the only statistical dependence of $`\mathrm{\Lambda }`$ on $`C`$ is through the original statistical dependence of $`\mathrm{\Lambda }`$ on $`G`$ (since the distribution of $`C`$ is determined completely by $`p(c|g)`$, a Markov condition) we get,
$$p(\lambda |c)=\underset{g}{}p(\lambda |g)p(g|c),$$
(6)
where $`p(\lambda |c)`$ can be clearly interpreted as the spectral density associated with the class $`c`$. Using these equations, we can now calculate the mutual information between a set of galaxy classes $`C`$ and the spectral wavelengths $`\mathrm{\Lambda }`$. Specifically, this information is given by
$$I(C;\mathrm{\Lambda })=\underset{c,\lambda }{}p(c)p(\lambda |c)\mathrm{log}\frac{p(\lambda |c)}{p(\lambda )}.$$
(7)
A basic theorem in information theory, known as data processing inequality, states that no manipulation of the data can increase the amount of (mutual) information given in that data. Specifically this means that by grouping the galaxies into classes one can only lose information about the spectra, i.e. $`I(C;\mathrm{\Lambda })I(G;\mathrm{\Lambda })`$. Our goal is then to find a non-trivial classification of the galaxies that preserves, as much as possible, the original information about the spectra. In other words, we wish to maximize $`I(C;\mathrm{\Lambda })`$. However, based on the above inequality, maximizing this information is trivial: each galaxy $`g`$ is a class $`c`$ of its own, which formally means $`CG`$. To avoid this trivial solution, one must introduce a formal constraint that will force the classification into a more compact representation.
It turns out that the compactness of a classification is directly governed by the mutual information between the classes and the galaxies. This mutual information is given by
$$I(C;G)=\underset{c,g}{}p(g)p(c|g)\mathrm{log}\frac{p(c|g)}{p(c)}.$$
(8)
To understand this expression, it is useful to consider its behaviour at two extremes. One extreme is when the new representation is the most compact one possible, i.e. there is only one class and all $`gG`$ are assigned to it with probability $`1`$. In this case there is no dependence between $`G`$ and $`C`$, thus $`I(C;G)`$ is trivially minimized to zero. This agrees with our intuition that a single global class carries no information at all about the original identity of a galaxy (i.e. its unique spectrum). In the other extreme, the classification is maximally complex when every $`gG`$ is assigned to a class of its own, i.e. $`CG`$. In this case $`I(C;G)`$ is maximized. Accordingly, the class of a specific galaxy provides the full information about its identity. The interesting cases are of course in between, where the number of classes is relatively small (but larger than one). In fact, in general the mutual information $`I(C;G)`$ gives a well justified measure for the complexity of the classification (Tishby, Pereira & Bialek 1999). Moreover, the maximal amount of information that the class can provide about the spectrum, $`I(C;\mathrm{\Lambda })`$, for a given amount of information preserved about the galaxies, $`I(C;G)`$, is a characteristic function of the data which does not depend on any specific classification algorithm.
We are now ready to give a full formulation of the problem: how do we find classes of galaxies that maximize $`I(C;\mathrm{\Lambda })`$, under a constraint on their complexity, $`I(C;G)`$? This constrained information optimization problem was first presented in Tishby et al. (1999) and their solution was termed the information bottleneck method. In effect we pass the information that $`G`$ provides about $`\mathrm{\Lambda }`$ through a “bottleneck” formed by the classes in $`C`$. The classes $`C`$ are forced to extract the relevant information between $`G`$ and $`\mathrm{\Lambda }`$.
Under this formulation, the optimal classification is given by maximizing the functional
$$[p(c|g)]=I(C;\mathrm{\Lambda })\beta ^1I(C;G),$$
(9)
where $`\beta ^1`$ is the Lagrange multiplier attached to the complexity constraint. For $`\beta 0`$ our classification is as non-informative (and compact) as possible — all galaxies are assigned to a single class. On the other hand, as $`\beta \mathrm{}`$ the representation becomes arbitrarily detailed. By varying the single parameter $`\beta `$, one can explore the tradeoff between the preserved meaningful information, $`I(C;\mathrm{\Lambda })`$, and the compression level, $`I(C;G)`$, at various resolutions.
Perhaps surprisingly, this general problem of extracting the relevant information — formulated in Eq. (9) — can be given an exact formal solution. In particular, the optimal assignment that maximizes Eq. (9) satisfies the equation
$$p(c|g)=\frac{p(c)}{Z(g,\beta )}\mathrm{exp}\left[\beta \underset{\lambda }{}p(\lambda |g)\mathrm{log}\frac{p(\lambda |g)}{p(\lambda |c)}\right],$$
(10)
where $`Z(g,\beta )`$ is the common normalisation (partition) function. <sup>3</sup><sup>3</sup>3 We note that $`\beta `$ here is analogous to the inverse temperature in the Boltzmann’s distribution function. The value in the exponent can be considered the relevant “distortion function” between the class and the galaxy spectrum. It turns out to be the familiar cross-entropy (also known as the ‘Kullback-Leibler divergence’, e.g. Cover & Thomas 1991), defined by
$$D_{KL}[p(\lambda |g)p(\lambda |c)]=\underset{\lambda }{}p(\lambda |g)\mathrm{log}\frac{p(\lambda |g)}{p(\lambda |c)}.$$
(11)
We emphasise that this effective distortion measure emerges here from first principles of information preserving and was not imposed as an ad-hoc measure of spectral similarity. Note that Eqs. (5, 6, 10) must be solved together in a self-consistent manner.
### 2.3 Relations to conventional classification approaches
We may gain some intuition into this method by contrasting Eq. 10 with more standard clustering algorithms. Suppose we start from Bayes’ theorem, where the probability for a class $`c`$ for a given galaxy $`g`$ is
$$p(c|g)p(c)p(g|c),$$
(12)
and $`p(c)`$ is the prior probability for class $`c`$. As a simple ad-hoc example, we can take the conditional probability $`p(g|c)`$ to be a Gaussian distribution with variance $`\sigma ^2`$
$$p(g|c)=\frac{1}{\sqrt{2\pi }\sigma }\mathrm{exp}(\frac{1}{2\sigma ^2}D_E^2),$$
(13)
with
$$D_E=\sqrt{\underset{\lambda }{}[p(\lambda |c)p(\lambda |g)]^2}.$$
(14)
The Euclidean distance $`D_E`$ is analogous to the cross-entropy $`D_{KL}`$ of Eq. 11. However, unlike our earlier formulation, neither the distance $`D_E`$ nor the obtained classes have good theoretical justifications. The variance $`\sigma ^2`$ (which may be due cosmic scatter as well as noise) plays a somewhat analogous role to the Lagrange multiplier $`\beta `$, but unlike $`\beta `$ it forces a fixed ‘size’ for all the clusters. Hence $`\sigma `$ can be viewed as the ‘resolution’ or the effective ‘size’ of the class in the high-dimensional representation space. We note that the Euclidean distance is commonly used in supervised spectral classification using ‘template matching’ (e.g. Connolly et al. 1995; Benitez 1999), in which galaxies are classified by matching the observed spectrum with a template obtained either from a model or from an observed standard galaxy. By comparing our method with an “Euclidean algorithm”, we find that our approach yields better class boundaries and preserves more information for a given number of classes.
### 2.4 The agglomerative information bottleneck algorithm
The initial approach to the solution of the three self-consistent Eqs. (5, 6, 10), applied already in Pereira, Tishby & Lee (1993), was similar to the “deterministic annealing” method (see e.g. Rose 1998). This is a top-down hierarchical algorithm that starts from a single class and undergoes a cascade of class splits (through second order phase transitions) into a “soft” tree of classes. Here we use an alternative algorithm, first introduced in Slonim & Tishby (1999), based on a bottom-up merging process. This algorithm generates “hard” classifications, i.e. every galaxy $`gG`$ is assigned to exactly one class $`cC`$. Therefore, the membership probabilities $`p(c|g)`$ may only have values of $`0`$ or $`1`$. Thus, a specific class $`c`$ is defined by the following equations, which are actually the “hard” limit $`\beta \mathrm{}`$ of the general self-consistent Eqs. (5, 6, 10) presented previously,
$$\{\begin{array}{c}p(c)=_{gc}p(g)\hfill \\ \\ p(\lambda |c)=\frac{1}{p(c)}_{gc}p(\lambda |g)p(g)\hfill \\ \\ p(c|g)=\{\begin{array}{cc}1\hfill & \text{if }gc\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}\hfill \end{array}$$
(15)
where for the second equation we used Bayes’ theorem, $`p(g|c)=\frac{1}{p(c)}p(c|g)p(g)`$.
The algorithm starts with the trivial solution, where $`CG`$ and every galaxy is in a class of its own. In every step two classes are merged such that the mutual information $`I(C;\mathrm{\Lambda })`$ is maximally preserved. The merging process is formally described as follows. Assume that we merge the two classes $`c_1,c_2`$ into a new class $`c^{}`$. Then the equations characterizing the new class are naturally defined by
$$\{\begin{array}{c}p(c^{})=p(c_1)+p(c_2)\hfill \\ \\ p(\lambda |c^{})=\frac{p(c_1)}{p(c^{})}p(\lambda |c_1)+\frac{p(c_2)}{p(c^{})}p(\lambda |c_2)\hfill \\ \\ p(c^{}|g)=\{\begin{array}{cc}1\hfill & \text{if }g\{c_1\}\{c_2\}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}\hfill \end{array}$$
(16)
The information loss with respect to $`\mathrm{\Lambda }`$ due to this merger is given by
$$\delta I(c_1,c_2)I(C_{before};\mathrm{\Lambda })I(C_{after};\mathrm{\Lambda })0,$$
(17)
where $`I(C_{before};\mathrm{\Lambda })`$ and $`I(C_{after};\mathrm{\Lambda })`$ are the mutual information between the galaxy classes and the wavelengths before and after the merge, respectively. Using Eqs.(15,16,17) it can be shown after some algebra (Slonim & Tishby, in preparation) that
$$\delta I(c_1,c_2)=(p(c_1)+p(c_2))D_{JS}(c_1,c_2),$$
(18)
where $`D_{JS}`$ is the Jensen-Shannon divergence (Lin 1991; El-Yaniv, Fine & Tishby 1997), defined by
$$D_{JS}(c_1,c_2)=\underset{i=1}{\overset{2}{}}p(c_i)D_{KL}[p(\lambda |c_i)\underset{i=1}{\overset{2}{}}p(c_i)p(\lambda |c_i)].$$
(19)
An intuitive interpretation is that the “merging cost” (in information terms) is equal to the “distance” $`D_{JS}(c_1,c_2)`$ between the classes before merging <sup>4</sup><sup>4</sup>4This distance has an interesting statistical interpretation as the distance to the most likely joint source of the two classes. Alternately, it can be viewed as analogous to the physical mixing entropy of two pure gases (see El-Yaniv et al. 1997; Bialek, Nemenman and Tishby 2000). multiplied by their “weight”, $`p(c_1)+p(c_2)`$. The algorithm is now straightforward — in each step we perform “the best possible merger”, i.e. we merge the two classes which minimize $`\delta I(c_i,c_j)`$. In this way, we maximize $`I(C;\mathrm{\Lambda })`$ in every step (but note that this does not necessarily guarantees a global maximum at the endpoint). In figure 1, we give the pseudo-code for this procedure. Note that this algorithm naturally finds a classification for any desired number of classes with no need to take into account the theoretical constraint via $`\beta `$ (Eq. 9). This is due to the fact that the agglomerative procedure contains an inherent algorithmic compression constraint, i.e. the merging process. A more general version of this algorithm which directly implements Eq. (9) is described elsewhere (Slonim & Tishby, in preparation).
For comparison with some conventional grouping algorithms, we also implemented an algorithm which uses the Euclidean metric instead of the Jensen-Shannon divergence used in Eq. (18). In this case, in each step we merge the pair that minimizes $`(p(c_i)+p(c_j))\sqrt{_\lambda (p(\lambda |c_i)p(\lambda |c_j))^2}`$, while ignoring the statistical meaning of the distributions. We refer to this procedure as the Euclidean algorithm.
## 3 Spectral Ensembles
### 3.1 Observed Spectra from the 2dF Survey
The 2dF Galaxy Redshift Survey (2dFGRS; Colless 1998, Folkes et al. 1999) is a major new redshift survey utilising the 2dF multi-fibre spectrograph on the Anglo-Australian Telescope (AAT). The observational goal of the survey is to obtain high quality spectra and redshifts for 250,000 galaxies to an extinction-corrected limit of $`b_J`$=19.45. The survey will eventually cover approximately 2000 sq deg, made up of two continuous declination strips plus 100 random $`2^{}`$-diameter fields. Over 100,000 galaxy spectra have been obtained as of May 2000. The spectral scale is 4.3Åper pixel and the FWHM resolution is about 2 pixels. Galaxies at the survey limit of $`b_J`$=19.45 have a median S/N of $`14`$, which is more than adequate for measuring redshifts and permits reliable spectral types to be determined.
Here we use a subset of 2dF galaxy spectra, previously used in the analysis of Folkes et al. (1999). We emphasize that the spectra are left in terms of photon counts (as opposed to energy flux). The spectra were de-redshifted to their rest frame and re-sampled to a uniform spectral scale with 4Å bins. Since the galaxies cover a range in redshift, the rest-frame spectra cover different wavelength ranges. To overcome this problem, only objects with redshifts in the range $`0.01z0.2`$ were included in the analysis. All the objects meeting this criterion then have rest-frame spectra covering the range 3700Å to 6650Å (the lower limit was chosen to exclude the bluest end of the spectrum where the response function is poor). Limiting the analysis to this common wavelength range means that all the major optical spectral features between O<sub>II</sub> (3727Å) and H$`\alpha `$ (6563Å) are included. In order to make the spectral classifications as robust as possible, objects with low S/N were eliminated by imposing a minimum mean flux of 50 counts per bin. The spectra were then normalised so that the mean flux over the whole spectral range was unity. The final sample contains 5869 galaxies, each described by 738 spectral bins. Throughout this paper, we refer to this ensemble as the “2dF catalogue”.
We corrected each spectrum by dividing it by a global system response function (Folkes et al. 1999). However, it is known that due to various problems related to the telescope optics, the seeing, the fibre aperture etc. the above correction is not perfect. In fact, each spectrum should be corrected by an individual response function. Unfortunately this incomplete correction mainly affects the continuum of the spectrum, i.e. the galaxy ‘colour’. As we shall see below, the IB analysis on the mock data (which obviously is free of the above problems) shows that the colour is a significant indicator of the underlying astrophysics. This highlights the need to correct properly each individual spectrum (work in progress).
### 3.2 Model Spectra from Semi-Analytic Hierarchical Merger Models
Our goal is to produce an ensemble of synthetic spectra with a representative admixture of different types of galaxies with realistic star formation histories. One way to accomplish this is to use semi-analytic modelling techniques (cf. Kauffmann, White, & Guiderdoni 1993; Cole et al. 1994; Somerville & Primack 1999 (SP) and references therein). Semi-analytic models have the advantage of being computationally efficient, while being set within the fashionable hierarchical framework of the Cold Dark Matter (CDM) scenario of structure formation. In addition to model spectra, this approach provides many physical properties of the galaxies, such as the mean stellar age and metallicity, size, mass, bulge-to-disk ratio, etc. This allows us to determine how effectively a given method can extract this type of information from the spectra, which are determined in a self-consistent way. We have used the code developed by Somerville (1997), which has been shown to produce good agreement with many properties of local (SP) and high-redshift (Somerville, Primack & Faber 2000; SPF) galaxies. Below we briefly summarize the models.
The formation and merging of dark matter halos as a function of time is represented by a “merger tree”, which we construct using the method of Somerville & Kolatt (1999). The number density of halos of various masses is determined by an improved version of the Press-Schechter model (Sheth & Tormen 1999), which mostly cures the usual discrepancy with N-body simulations. The cooling of gas, formation of stars, and reheating and ejection of gas by supernovae within these halos is modelled by simple recipes. Chemical evolution is traced assuming a constant yield of metals per unit mass of new stars formed. Metals are cycled through the cold and hot gas phase by cooling and feedback, and the stellar metallicity of each generation of stars is determined by the metal content of the cold gas at the time of its formation. All cold gas is assumed to initially cool into, and form stars within, a rotationally supported disk; major mergers between galaxies destroy the disks and create spheroids. New disks may then be formed by subsequent cooling and star formation, producing galaxies with a range of bulge-to-disk ratios. Galaxy mergers also produce bursts of star formation, according to the prescription described in SPF. Thus the star formation history of a single galaxy is typically quite complex and is a direct consequence of its environment and gas accretion and merger history.
These star formation histories are convolved with stellar population models to calculate magnitudes and colors and produce model spectra. We have used the multi-metallicity GISSEL models (Bruzual & Charlot, in preparation) with a Salpeter IMF to calculate the stellar part of the spectra. Emission lines from ionized $`H_{\mathrm{II}}`$ regions have been added using the empirical library included in the PEGASE models (Fioc & Rocca-Volmerange 1997). This library provides the total luminosity of all of the major emission lines as a function of the age of the stellar population. Any dependence of the line strengths on metallicity, ionization state, geometry, etc. is neglected. We adjusted the resolution of the mock spectra to be comparable to the 2dF spectra. The width of the lines is then determined by the resolution of the grating, and we have modelled them as Gaussians with a width of 4 Å.
Dust extinction is included using an approach similar to that of Guiderdoni & Rocca-Volmerange (1987). Here, the mass of dust is assumed to be a function of the gas fraction times the metallicity of the cold gas. We then use a standard Galactic extinction curve and a “slab” model to compute the extinction as a function of wavelength and inclination (see Somerville 1997 or Somerville et al. 2000b for details). The extinction correction is applied indiscriminantly to the stellar and line emission. This is probably unphysical as it is likely that the star-forming regions that produce the emission lines are more heavily extinguished than the underlying old stellar population, but at our current level of modelling, we ignore this effect.
As described in SP, we set the free parameters of the models by reference to a subset of local galaxy data; in particular, we require a typical $`L_{}`$ galaxy to obey the observed I-band Tully-Fisher relation and to have a gas fraction of $`0.1`$ to $`0.2`$, consistent with observed gas contents of local spiral galaxies. If we assume that mergers with mass ratios greater than $``$ 1:3 form spheroids, we find that the models produce the correct morphological mix of spirals, S0s and ellipticals at the present day (we use the mapping between bulge-to-disk ratio and morphological type from Simien & de Vaucouleurs 1986). This critical value for spheroid formation is what is predicted by N-body simulations of disk collisions (cf. Barnes & Hernquist 1992).
In SP, we found that a cosmology with $`\mathrm{\Omega }_0=\mathrm{\Omega }_\mathrm{\Lambda }=0.5`$ and $`H_0=60`$ km/s/Mpc produced very good agreement with the 2dF $`b_J`$-band luminosity function (for all types combined), as well as the observed K-band luminosity function, Tully-Fisher relation, metallicity-luminosity relation, and colours of local galaxies. We use the same fiducial model here, with a few minor modifications: we incorporate self-consistently the modelled metallicity of the hot gas in the cooling function, and use the multi-metallicity SED’s (instead of solar metallicity) with a Salpeter (instead of Scalo) IMF. Another minor detail is that ejected material is eventually returned to the halos as described in the updated models of SPF. We find that these minor modifications do not significantly change our previous results for local galaxies.
We construct a “mock 2dF catalogue” of $`2611`$ model galaxies with the same magnitude limit, wavelength coverage and spectral resolution, and redshift range as the 2dF survey (described above). The synthetic spectra are expressed in terms of photon counts and the total number of counts in each spectrum is normalized to unity, as in the prepared observed spectra. We shall present elsewhere comparison of the mock catalogues with preliminary data from the 2dF survey
To study the effect of noise on the classification, we added Poisson noise to the simulated spectra. This was done using an empirical relation between the mean photon counts $`\overline{N}_{ph}`$ per observed 2dF spectrum and the corresponding APM $`b_J`$ magnitude, $`\overline{N}_{ph}=9638.0757.8b_J+13.9b_J^2`$ (D. Madgwick, private communication). Each simulated galaxy was assigned a mean number of counts based on its $`b_J`$ magnitude as given by the models, and Poisson deviants per spectral bin were drawn at random. <sup>5</sup><sup>5</sup>5 We note that an alternative approach could be to filter out the noise of the observed spectra, e.g. by using PCA (Folkes et al. 1996). We have ignored the effects of the response function of the fibres, aperture effects, and systematic errors related to the placement of fibres in the holding plate (see above).
Figure 2 shows the mean spectrum for the 2dF and mock+noise catalogues, obtained by simply averaging the photon counts in each wavelength bin for all the galaxies in the ensemble. The mean spectra for the observed and mock catalogues are seen to be similar. The magnitude limit that we have chosen is such that our ensembles are dominated by fairly bright, moderately star-forming spiral galaxies, and the mean spectra show familiar features such as the 4000 Å break, the Balmer series, and metal lines such as O<sub>II</sub> and O<sub>III</sub>. One can see that the 2dF spectrum appears to bend downwards relative to the models towards both ends of the wavelength range. This may be due to an inaccurate correction for the response function, as discussed above.
## 4 Results
We now apply the IB algorithm to both the 2dF and the mock data. Recall that our algorithm begins with one class per galaxy, and groups galaxies so as to minimize the loss of information at each stage. Figure 3 shows how the information content of the ensemble of galaxy spectra decreases as the galaxies are grouped together and the number of classes decreases. In the left panel, we show the ‘normalized’ information content $`I(C;\mathrm{\Lambda })/I(G;\mathrm{\Lambda })`$ as a function of the reduced complexity $`N_C/N_G`$, where $`N_G`$ is the number of galaxies in the ensemble and $`N_C`$ is the number of classes. One may think of this as the fraction of information about the original ensemble that would be preserved if we threw away all the individual galaxies and kept only the representative spectra of the classes. Remarkably, we find that if we keep about five classes, about $`85`$ and $`75`$ percent of the information is preserved for the mock and mock+noise simulations, respectively. This indicates that the wavelength bins in the model galaxy spectra are highly correlated. This may not seem very surprising, since many of the spectral features arise from the same stellar physics. For example, emission lines will be stronger in a galaxy with significant recent star formation, and the whole sequence of metal absorption lines will be deeper for a galaxy with an old stellar population or a high metallicity. However, if we keep in mind that each of our model galaxies has a very complex star formation history and is therefore comprised of stars with a distribution of ages and metalicities, and is affected differently by dust extinction, this result may seem somewhat more surprising.
In contrast with the mock samples, for the 2dF catalogue, only about $`50`$ percent of the information is preserved by five classes. Thus the wavelength bins in the real spectra are also highly correlated, but not to the same degree as the model spectra. This is unlikely to be solely due to the effect of noise. Adding noise to the mock catalogues does change the curve in Figure 3 (left panel), bringing it closer to the 2dF curve, but the effect is not large enough to explain the whole discrepancy, and moreover the shapes of the mock+noise and 2dF information curves are still different. We also see that the ‘absolute information’ for 2dF (right panel of Figure 3) is much higher than for the mock samples. This discrepancy may be partially due to the influence on the real spectra of more complicated physics than what is included in our simple models. It could also be due to systematic observational errors (see Section 3.1). We are in the process of attempting to model these systematic errors in detail to better understand this result.
The information curves are not sensitive to the number of galaxies used in the analysis. When we take a random subset of 2611 galaxies out of the 2dF sample (to make it identical to the number of galaxies in the mock sample), the differences in the information curve w.r.t the full 2dF set are minor. We also experimented with the ‘Euclidean’ algorithm (see Section 2.3) for both the mock and 2dF data. We find that for the mock data, the Euclidean algorithm produces nearly indistinguishable results from our fiducial algorithm, however, for the 2dF data, the difference between the two is more significant. The full IB algorithm preserves more information for any given number of classes and is therefore superior. This suggests that the Euclidean-type algorithms may be sufficient for certain problems, but inadequate for more complex data.
### 4.1 The IB Classes
For the remainder of this paper, we present the results obtained for five classes. <sup>6</sup><sup>6</sup>6We note that galaxy images can be reliably classified by morphology into no more than 7 or so classes (e.g. Lahav et al. 1995; Naim et al. 1995a) Figure 4 shows the representative spectra for these five classes for both the 2dF and mock+noise catalogues. The corresponding five spectra for the noise-free mock data were very similar to the mock+noise spectra shown. We ‘matched’ each of the classes obtained for the 2dF data with one from the mock+noise data by minimizing the average $`D_{JS}`$ ‘distance’ between the pairs. The classes are then ordered by their mean $`BV`$ colour. Note that the five classes produced by the algorithm appear similar for both catalogues — there was certainly no guarantee that this would be the case.
It is also interesting to examine the relative fractions of galaxies in each class, $`p(c)`$, for the observed and mock catalogues. These values are given on the appropriate panels of Figure 4. There is a partial agreement between the weights of the matched classes for the 2dF and mock+noise catalogues (although particularly for $`c_3`$ and $`c_5`$, the agreement is not very good). We might hope that this could provide a way to improve the physics included in the models by constraining the relative composition of different types of galaxies in more detail than was previously possible. However, we find that adding noise to the mock catalogues causes some galaxies to “jump” to different classes, thus changing the relative fractions. When we add noise at the level of the 2dF data as described previously, $`23`$ percent of the galaxies are assigned to a different class than in the noise-free case. Of these, the great majority ($`22`$ percent) are assigned to adjacent classes (i.e. $`\mathrm{\Delta }c=\pm 1`$). We note these results are for the ‘hard’ version of the algorithm. We expect the ‘soft’ version of the algorithm to be less sensitive to noise.
We are also concerned that some of the noticeable discrepancy in the shape of the mean spectra of the 2dF classes and those of the synthetic galaxies (also noted in the comparison of the mean spectra, Figure 2) may be due to systematic observational effects such as inaccurate modelling of the response function. In the future we hope to be able to model and correct for these effects.
More generally, we can see that the algorithm is sensitive to the overall slope (or colour) of the spectrum, and also to the strength of the emission lines. The classes clearly preserve the familiar physical correlation of colour and emission line strength; the five classes form a sequence from $`c_1`$, which has a blue continuum with strong emission lines, to $`c_5`$, with red continuum with no emission lines <sup>7</sup><sup>7</sup>7Recall that the order of the classes as produced by the algorithm is arbitrary, and we have placed them in this sequence by hand, but the correlation of colour and emission line strength within the classes is produced by the algorithm with no help from us. Already, we may form the impression that the algorithm has classified the galaxies in a way that is reminiscent of conventional spectral classes. It may be interesting to compare the mean spectrum of $`c_1`$ with the spectrum of the Sm/Irr galaxy NGC449, and $`c_5`$ with the Sa galaxy NGC775 from Figure 2a of Kennicutt (1992). Apparently, the $`c_1`$ class corresponds to late type galaxies (Sm/Irr) and $`c_5`$ to early types (Sa-E).
### 4.2 Correlation of the IB Classes with Physical Properties
In order to gain a better understanding of the IB classes, we now use the noise-free mock catalogue and investigate the physical properties of the galaxies in each class as given by the same models that we use to produce the spectra. The parameters that we chose to investigate are: $`BV`$ colour, the ratio of the present to past-averaged star formation rate (Kennicutt b-parameter, $`b_{KENN}SFR/SFR`$; Kennicutt 1983), the ratio of the mass of the bulge to the total stellar mass of the galaxy ($`B/T`$), the fraction of the total baryonic mass in cold gas ($`f_{\mathrm{gas}}m_{\mathrm{cold}}/(m_{\mathrm{cold}}+m_{\mathrm{star}})`$), the mean mass-weighted stellar age and metallicity, the internal velocity dispersion $`\sigma _v`$ or circular velocity $`V_c`$, and the B-band absolute magnitude $`M_B`$.
Figure 5 shows the trends of these physical parameters with class number. The strongest dependence is of $`BV`$ colour and present-to-past-averaged star formation rate $`b_{\mathrm{KENN}}`$. Class $`c_1`$ contains blue galaxies that are forming stars at rates that are one to two orders of magnitude higher than the average over their past history. As one moves towards $`c_5`$, galaxies are redder and formed a larger fraction of their stars in the past. This is consistent with the observed strong correlation of the Kennicutt b-parameter with $`BV`$ colour and morphological type in nearby galaxies (Kennicutt et al. 1994). For reference, note that the mean colour for Sm/Irr galaxies in the local Universe is $`BV=0.42`$, for Sbc-Sc $`BV=0.55`$, for Sab-Sb $`BV=0.64`$, for S0a-Sa $`BV=0.78`$, and for E-S0 $`BV=0.90`$ (Roberts & Haynes 1994). From the colours alone, we might guess that galaxies in $`c_1`$ and $`c_2`$ are starburst galaxies, galaxies in $`c_2`$ are Sm-Irr, $`c_4`$ corresponds to Sb/Sbc and $`c_5`$ to Sa/S0/E. This bolsters our initial impression based on the visual appearance of the spectra. Weaker trends are visible in other properties: the sequence $`c_1`$-$`c_5`$ shows an increasing B/T ratio, as expected, but with a large scatter within each class. There is quite a large scatter in the observed correspondence between morphological T-type and bulge-to-total ratio, but values of $`B/T0.40.5`$ are typical of very early-type spirals (Sa) or lenticular galaxies (cf. Simien & de Vaucouleurs 1986). Weak trends are also visible in the mean stellar age and metallicity and the internal velocity dispersion: as we move from $`c_1`$$`c_5`$, galaxies tend to be older, more metal rich, and more massive. This is in accord with the observed trends of these quantities with morphological type (cf. Roberts & Haynes 1994) and the trends we have noted previously.
We find no trend of dust extinction with the IB classes, probably because we have assumed that dust extinction attenuates the emission lines and the continuum by the same factor. Therefore the dust only changes the overall colour of a galaxy, but does not affect the ratio of emission or absorption lines to continuum.
We investigate the composition of the classes in more detail by examining the distributions of some of the physical parameters for the different classes. In each of figures 6-10, in the upper left-most panel we show the distribution of the relevant physical parameter for the whole ensemble of model (noise-free) galaxies. In the other five panels we show the distributions for the galaxies in each class separately. Figure 6 shows the distribution of Kennicutt’s $`b_{KENN}`$ parameter (ratio of present-to-past-averaged SFR). This parameter shows a particularly strong correlation with the IB class assignment, although the distributions show significant overlap. Figure 7 shows the distribution of B-band absolute magnitudes. Class $`c_1`$ contains an excess of bright galaxies, whereas $`c_2`$ and to a lesser extent $`c_5`$ contain excesses of faint galaxies compared to the overall distribution. Figure 8 shows the distribution of the internal velocity dispersion $`\sigma _v`$ of the galaxies, a measure of the dynamical mass of the galaxy. Here we note the curious fact that the distribution is skewed towards small $`\sigma _v`$ galaxies in $`c_1`$-$`c_3`$, and towards large $`\sigma _v`$ galaxies in $`c_4`$-$`c_5`$. This suggests that galaxies in $`c_1`$, which are preferentially bright and with small mass are starburst galaxies. We return to this point later.
In Figure 9, we see that $`c_1`$ is entirely composed of disk-dominated galaxies ($`B/T\stackrel{<}{}0.4`$). The classes become progressively more skewed towards larger $`B/T`$, bulge-dominated galaxies as we move towards $`c_5`$. However, despite the fact that we noted that the spectral appearance, colors and star formation rates of galaxies in $`c_4`$-$`c_5`$ are typical of observed galaxies with bulges (Sb-E), in the models these classes contain a significant fraction of nearly pure disks ($`B/T\stackrel{<}{}0.2`$). This may be an indication that the connection between star formation and morphology is not being modelled properly.
The distribution of gas fraction is shown in Figure 10. In $`c_1`$, there are quite a lot of galaxies with fairly low gas fractions. This is somewhat in conflict with our previous association of this class with late-type galaxies, but perhaps not if most of these galaxies have experienced a very recent starburst which would have tended to consume the gas supply very quickly. The class $`c_2`$ is highly skewed towards high gas fractions, but with some objects with low gas fractions — perhaps this class is composed of a combination of post-starburst galaxies and quiescent, gas rich galaxies. Classes $`c_4`$-$`c_5`$ are composed solely of rather gas-poor galaxies, confirming that these galaxies have probably been evolving passively, with little new star formation.
It is also useful to see where the classes are located in the two-dimensional space of pairs of the physical parameters. In the following figures, we once again show the entire ensemble in the upper left panel, and the breakdown by class in the other five panels. Figure 11 shows the mean stellar metallicity of galaxies as a function of their mean stellar age. We see that there is a weak trend, with a large scatter, between these two quantities in our models. This trend is stronger for classes $`c_1`$-$`c_3`$, and becomes mostly washed out for $`c_4`$-$`c_5`$. Galaxies in $`c_1`$-$`c_2`$ also tend to have higher metallicities for their age than galaxies in $`c_4`$-$`c_5`$. As applied here, the algorithm suffers from the familiar age-metallicity degeneracy. However, an alternative way of applying the algorithm (by asking it to preserve the maximum information about a particular physical parameter, e.g. age or metallicity) may prove to be effective for this problem. We intend to pursue this in future work.
We also examine how the classes inhabit a classical colour-colour diagram. Figure 12 shows the far-UV/optical (1500 Å-B) and optical/near-IR (B-K) colours. Recall that the spectra that we provided to the IB algorithm did not contain any information about the far-UV (1500 Å) or near-IR (K-band) magnitudes. Galaxies in $`c_1`$ have a narrow range of very blue UV-B colours but a broad range of B-K colours. The shape and orientation of this clump changes as we move from $`c_1`$-$`c_5`$. The lines on the top right panel show the tracks for instantaneous bursts of fixed age and metallicity from the Bruzual & Charlot models (see figure caption). This illustrates the complex manner in which both age and metallicity determine the location of galaxies in this diagram.
Figure 13 shows the absolute B-magnitude as a function of circular velocity. This is essentially what is usually known as the Tully-Fisher (TF) relation, although it should be noted that our model magnitudes contain the effects of dust extinction, unlike observed TF samples where at least the effects of inclination are generally removed. Nor have we made any cut on morphology or gas fraction in the models. This is why the slope and scatter of the relation plotted look so different from the usual TF relation. It was shown in SP that when the above effects are accounted for, we obtain reasonable agreement with the observed zero-point, slope, and scatter of the I-band TF relation in these models. The interesting thing to note is the way the classes cut the two-dimensional space of this diagram. Galaxies in $`c_1`$ lie at preferentially bright magnitudes for their velocity/mass (i.e., they are starbursts), whereas galaxies in $`c_5`$ lie at preferentially faint magnitudes for their velocity/mass. An increasing curvature of the relation is also seen from $`c_1`$ to $`c_5`$. The progressive offset of the TF relation with varying Hubble type is recognized (Burstein et al. 1997) but not very well understood. This result offers a hint as to its origin, and also suggests that the familiar observed TF relation and its small scatter may be a special feature of the particular type of galaxies that are generally selected for these samples.
## 5 Comparison with PCA
PCA has previously been applied to data compression and classification of spectral data of stars (e.g. Murtagh & Heck 1987; Bailer-Jones et al. 1997), QSO (e.g. Francis et al. 1992) and galaxies (e.g. Connolly et al. 1995a; Folkes, Lahav & Maddox 1996; Sodre & Cuevas 1997; Galaz & de Lapparent 1997; Bromley et al. 1998; Glazebrook, Offer & Deeley 1998; Ronen et al. 1999; Folkes et al. 1999). As we noted earlier, while PCA operates as an efficient data compression algorithm, it is purely linear, based only on the variance of the distribution. PCA on its own does not provide a rule for how to divide the galaxies into classes. <sup>8</sup><sup>8</sup>8 For example, in Folkes et al. (1999) the classification was done by drawing lines in the $`PC1PC2`$ plane using training sets. One training set was based on visual inspection of the spectra by a human expert. This led to classification which is more sensitive to emission and absorption lines, rather than to the continuum (which is affected by observational problems). It is therefore interesting to see where the IB classes reside in the space of the PC components. The PC eigenvectors are defined in the usual way (see the above references for more details), and the projections of the (noiseless) model spectra onto the first two components of this basis (PC1 and PC2) are shown in Figure 14. The overall pattern reminds the one seen in the models of Ronen et al. (1999) and in the 2dF data (Folkes et al. 1999), but note that the PC eigenvectors are different for every data set. We note that our 5 IB classes form fairly well-separated “clumps” in PC1-PC2 space, and that to a first approximation, the IB classification is along PC1. This highlights the need to correct properly the artifacts due to optics so that the continuum (colour) and hence PC1 can be determined accurately. The PC-space of the IB clumps looks quite different from the partitioning (based on training sets) given in Folkes et al. (1999). It has been shown (Ronen et al. 1999) that PC1 and PC2 are correlated with colour and emission line strength, and the sequence from $`c_1`$-$`c_5`$ is again sensible in this context.
## 6 Discussion
Unsupervised classification methods are generally used to obtain efficient representation of complex data. One can identify two general classes of techniques for achieving this goal, geometrical and statistical. Geometrical methods begin by an embedding of the data in a high dimensional, usually Euclidean, space, and then searching for a low dimensional manifold that captures the essential variation of the data. The simplest of such methods is PCA, which provide a linear projection of the data. PCA can be generalized to more powerful linear projections, e.g. projection pursuit (Friedman and Tukey 1974) or to nonlinear projections that maximize statistical independence, such as Independent Component Analysis (ICA; Bell and Sejnowski 1995). These methods provide a low dimensional representation, or compression, in which one might hope to identify the relevant structure more easily. Another approach of identifying classes of objects in a parameter-space (based on a training set) is by utilising Artificial Neural Networks (e.g. used for morphological classification of galaxies; Naim et al. 1995b; Lahav et al. 1996).
Statistical methods assume that the data is sampled from an underlying distribution with some knowledge of its parametric structure. Finding the structure of data amounts then to the estimation of the unknown parameters of the distribution. Such methods include the familiar Gaussian mixture estimation and similar vector quantization techniques. Modern formulation of many of the statistical as well as geometrical methods is based on information maximization principles for finding either efficient compression or statistical independence among the features.
Other compression methods have also been proposed for our problem, e.g. by starting with standard Maximum Likelihood desired parameters (e.g. age) and utilising the Fisher information matrix to define a set of optimal axes (Heavens, Jimenez & Lahav 2000).
The information bottleneck method presented here provides a new statistical approach to structure extraction. Unlike all other techniques it aims directly at the problem of the extraction of the relevant structure or features, where the relevance is determined through the information on another, carefully chosen, variable. The goal of the method is well defined and objective, with natural information theoretical figures of merit. It is superior to both geometric and statistical methods since it makes no model-dependent assumptions on the data origin, nor about the similarity or metric among data points.
An important issue, common to most unsupervised classification methods, is model order estimation: what is the correct number of classes? This question is closely related to the sampling noise issue — the obtained classes should not be sensitive to the specific sample, thus should be robust to sampling noise. This criterion can be checked, using techniques such as cross-validation, in most clustering algorithms including ours. Yet it is important to emphasize that the “true” or “correct” number of classes may be an ill-defined quantity for real data sets. The number should be determined by the desired resolution, or preserved information in our case.
We have shown how the IB algorithm can be used to classify galaxy spectra in a principled and objective way. The number of distinct classes is specified by an acceptable degree of information loss. We have applied the algorithm to a subset of spectra from the ongoing 2dF redshift survey, and to a mock-2dF catalogue of synthetic spectra obtained from semi-analytic hierarchical (CDM) models of galaxy formation.
We find that five classes preserve about 50 percent of the information about the ensemble of 2dF spectra. The same number of classes preserves 85 percent of the information about the model spectra in the absence of Poisson noise. When noise is added to the models with S/N comparable to the 2dF data, five classes preserve 75 percent of the information.
Examining the mean spectra of the five classes produced by our algorithm, we first see that there is a good matching between the five average spectra obtained for the mock data and the five average spectra obtained for 2dF. It is also apparent that these spectra form a sequence from blue galaxies with strong emission lines to red galaxies with strong absorption lines and no emission lines. This corresponds well with the general approach usually followed in more subjective spectral classification methods (i.e. “by eye”).
For the model galaxies, we also show that the classes form sequences in several physical quantities, such as the present-to-past-averaged star formation rate (Kennicutt $`b`$ parameter), morphology (as represented by the ratio of bulge-to-total stellar mass), and stellar mass/velocity dispersion. Since the spectra obviously do not contain any of this information directly, the existence of these correlations (which are in accord with known observational correlations among the physical parameters studied) seems to indicate two things. First, that the physics used to create the model galaxies is fairly sensible. Second, and more novel, grouping the galaxies in a way that formally preserves the information with respect to the spectra (as our method does), discovers interesting physical correlations.
Again for the models, we find that the classes occupy different parts of bivariate diagrams in pairs of the physical parameters, such as age vs. metallicity, color-color, and luminosity vs. circular velocity (Tully-Fisher). These results may hint at important clues as to how to constrain the star formation histories of different types of galaxies and the physical origin of these sorts of relationships.
We compare our results with those of a Principal Component Analysis. We find that the classes produced by the IB algorithm form fairly well-defined clumps in the PC1-PC2 space (the projections onto the first two principle component eigenvectors).
We conclude that this method provides a way to classify galaxies that is fully automated and objective, yet is related to the physical properties of galaxies and the intuition that astronomers have built up over the years using more subjective methods. A further advantage is that this method can be applied in exactly the same way to observations and models such as the ones investigated here, allowing comparisons between theory and observations to be made on the same footing. We intend to apply the algorithm to the full set of galaxy spectra obtained from the 2dF redshift survey. Obvious applications are then to study the spatial clustering of different classes and relative biasing among them, and luminosity functions divided by class. If spectra with adequate signal-to-noise can be obtained, the same method could be applied to high redshift spectra to study the evolution with redshift of different types of galaxies.
## 7 acknowledgment
We thank S. Folkes, D. Madgwick, A. Naim, S. Ronen, and the 2dFGRS team for their contribution to the work presented here. This research was supported by a grant from the Ministry of Science, Israel.
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# The structure of Lie algebras and the classification problem for partial differential equations.
## Introduction
Modeling phenomena in nature with partial differential equations is one of the central problems of mathematical physics and applied mathematics. One can even say that mathematical physics in its classical form was created in order to provide a rigorous mathematical foundation for describing different phenomena in physics, chemistry and biology by partial differential equations. However, when one has to decide which differential equation fits in the best way as a model for the process under study, one has to select from a broad class of possible partial differential equations. Even if one has taken into account all the peculiarities of the process under study (which is hardly possible!), there is still great freedom in choosing possible models. One of the principal criteria for choosing the partial differential equations modeling real processes is the symmetry selection principle. By this we mean that from the whole set of admissible models, those models which have the highest symmetry should be selected. This point of view is supported by the fact that the most successful mathematical models in theoretical and applied science have a rich symmetry structure. Indeed, the basic equations of modern physics, the wave, Schrödinger, Dirac and Maxwell equations are distinguished from the whole set of partial differential equations by their Lie and non-Lie (hidden) symmetry (see for more details on symmetry properties of these equations).
The effectiveness of the symmetry (group-theoretical) approach to the classification of admissible partial differential equations relies heavily upon the availability of a constructive way of describing transformation groups leaving invariant the form of a given partial differential equation. This is done via the well-known infinitesimal method developed by Sophus Lie (see, e.g., ). Given a partial differential equation, the problem of investigating its maximal (in some sense) Lie invariance group reduces to solving an over-determined system of linear partial differential equations, called the determining equations. However, if the equation under study contains arbitrary elements (functions), then one has to solve an intermediate classification problem. Namely, it is necessary to describe all the possible forms of the functions involved such that this equation admits a non-trivial invariance group.
In principle, the classification problem is solved with the help of the Lie algebra approach. However, since the determining equations involve some arbitrary functions, there is an evident need for a modification of the basic Lie technique in order to obtain an efficient and systematic way of classifying these arbitrary elements. The idea of this modification was suggested by Sophus Lie himself. Indeed, his way of obtaining all ordinary differential equations in one variable admitting non-trivial symmetry algebras tells us what is to be done in the case at hand. We should first construct all the possible inequivalent realizations of symmetry algebras within some class of Lie vector fields. If we succeed in doing this, then the symmetry algebras will be specified, so that we can apply directly Lie’s infinitesimal algorithm, thus getting inequivalent classes of invariant equations. In this way, Sophus Lie obtained his famous classification of realizations of all inequivalent complex Lie algebras in the plane . Recently, Lie’s classification was exploited by Olver and Heredero in order to classify nonlinear wave equations in two independent variables that are invariant with respect to transformation groups not changing the temporal variable.
A systematic implementation of these ideas for partial differential equations has been worked out by Ovsjannikov . His approach is based on the concept of an equivalence group, which is a Lie transformation group acting in the extended space of independent variables, functions and their derivatives, and preserving the class of partial differential equations under study. It is possible to modify Lie’s algorithm in order to make it applicable for the computation of this group . Next, one constructs the optimal system of subgroups of the equivalence group. The last step uses Lie’s algorithm for obtaining specific partial differential equations, that (a) belong to the class under study, and (b) are invariant with respect to the subgroups mentioned above. This approach has been applied to a number of equations of nonlinear gas dynamics and diffusion equations (Akhatov, Gazizov and Ibragimov ). Ovsjannikov’s ideas have also been exploited by Torrisi and co-workers in order to perform a preliminary group classification of some nonlinear diffusion and heat conductivity equations . Ibragimov and Torrisi have obtained a number of important results on the group classification of nonlinear detonation equations and nonlinear hyperbolic type equations . There is a number of papers (see, e.g., and the references therein) devoted to a direct computation of equivalence groups of some PDEs. Since the transformations of the equivalence group are used in their finite form, this approach has the merit of giving the possibility of finding discrete equivalence groups or even non-local ones.
However, the possibility of implementing Ovsjannikov’s approach in its full generality presupposes that we are able to construct the optimal system of subgroups of the equivalence group. So that, even for the case when the equivalence group has a finite number of parameters, there arise major algebraic difficulties, since for a number of known finite-parameter Lie groups the classification problem has not yet been solved (to say nothing about infinite-parameter Lie groups, where this problem is completely open). Consequently, there is an evident need for Ovsjannikov’s approach to be modified so that it can be applied to the case of infinite-parameter equivalence groups.
In the paper we have developed a new approach that enables us to solve efficiently the symmetry classification problem for partial differential equations even for the case of infinite-dimensional equivalence groups. It is mainly based on the following facts:
* If the partial differential equation possesses non-trivial symmetry, then it is invariant under some finite-dimensional Lie algebra of differential operators which is completely determined by its structural constants. In the event that the maximal algebra of invariance is infinite-dimensional, then it contains, as a rule, some finite-dimensional Lie algebra.
* If there are local non-singular changes of variables which transform a given differential equation into another, then the finite-dimensional Lie algebra of invariance of these equations are isomorphic, and in the group-theoretic analysis of differential equations such equations are considered to be equivalent.
What we have suggested in is a preliminary classification of inequivalent realizations of low-dimensional Lie algebras within some specific class of first-order linear differential operators. This class is determined by the structure of the equation under study. Its elements form a representation space for realizations of Lie algebras of symmetry groups admitted by the equations belonging to the class of partial differential equations under study. A natural equivalence relation is introduced on the set of all possible realizations. Namely, two realizations are called equivalent if they are transformed into each other by the action of the equivalence group. In other words, solving the problem of symmetry classification of partial differential equations having some prescribed form, is equivalent to constructing a representation theory of Lie transformation groups (or Lie algebras of first-order differential operators) realized as symmetry groups (algebras) of the equations in question.
The first aim of the present paper is to give a detailed exposition of our approach. A full understanding of the techniques applied requires some basic facts from the general theory of Lie groups and algebras, some of which are dispersed in the literature and are not available in English (this is the case for the papers of Mubarakzyanov and Morozov). So, in addition to the exposition of the classification results, we give a survey of results on the structure of Lie algebras (with special emphasis on low-dimensional Lie algebras), which are of vital importance for the effective implementation of our approach.
The second aim is obtaining a complete description of the nonlinear heat conductivity equations of the form
$$u_t=F(t,x,u,u_x)u_{xx}+G(t,x,u,u_x)$$
(0.1)
that admit non-trivial symmetry group. Hereafter $`u=u(t,x)`$$`F`$$`G`$ are sufficiently smooth functions of the corresponding arguments, $`u_t={\displaystyle \frac{u}{t}}`$$`u_x={\displaystyle \frac{u}{x}}`$, $`u_{xx}={\displaystyle \frac{^2u}{x^2}}`$, $`F0`$.
Note that the above equation is, in some sense, the most general evolution equation in one dimension. Indeed, any equation of the most general form
$$u_t=H(t,x,u,u_x,u_{xx})$$
(0.2)
which admits at least one-parameter symmetry group, not changing the temporal variable, can be reduced to the form (0.1) by a non-point transformation. So that our group classification of equations (0.1) will also cover invariant equations of the form (0.2) excepting for the small subclass of equations whose symmetry algebras are spanned by operators with non-vanishing coefficients by $`\frac{}{t}`$.
The principal scheme of the paper is as follows. Section I contains a general description of our approach. In the next section we give a brief overview of the necessary facts from the general theory of Lie algebras. Section III is devoted to group classification of PDEs (0.1). We consider subsequently, the cases of semi-simple, semi-direct sum of semi-simple and solvable and solvable symmetry algebras thus getting the full solution of the classification problem for nonlinear heat conductivity equations belonging to the class (0.1). The last section contains discussion of the results obtained and some conclusions.
## I Description of the method
The approach to the classification of partial differential equations which we propound is, in fact, a synthesis of Lie’s infinitesimal method, the use of equivalence transformations and the theory of classification of abstract finite-dimensional Lie algebras. It constitutes a constructive solution of the problem the group classification of partial differential equations possessing large classes of arbitrary elements and admitting non-trivial finite-dimensional invariance algebras.
The realization of group classification in the proposed approach consists in the implementation of the following algorithm:
* The first step involves finding the form of the infinitesimal operators which generate the symmetry group of the equation under consideration, and the construction of the equivalence group of this equation. To find the form of the infinitesimal operators one uses the usual Lie algorithm. As a result we obtain a system of linear partial differential equations of first order, which connect the coefficients of the infinitesimal operators with the arbitrary term of the equation. In what follows, we call this system the characterizing system of the equation. In order to construct the equivalence group $``$ of the equation under consideration, one can use the infinitesimal as well as the direct method.
* In the second step, one carries out the group classification of those equations of the given form which admit finite-dimensional Lie algebras of invariance.
For this, one carries out a step-by-step classification of finite-dimensional Lie algebras within the specified class of infinitesimal operators, up to equivalence under transformations of the group $``$. In this, one has to see if each algebra obtained in this way can be an invariance algebra of the equation at hand before proceeding from the realization of Lie algebras of lower dimension to the realization of Lie algebras of higher dimension. This eliminates superfluous realizations of Lie algebras. Also, those realizations of Lie algebras which are invariance algebras of the equation will, as their dimension increases, correspond to greater fixing of the arbitrary term.
This procedure is continued until the arbitrary term in the equation is completely determined or until it is no longer possible to extend the realization of Lie algebras beyond a given dimension within the specified class of infinitesimal operators.
* The third step is then to exploit the characterizing system or the infinitesimal method of Lie in order to find, for each of the particular choices of the arbitrary term, the maximal invariance algebra of the equation under consideration. Furthermore, the equivalence of the equations obtained in this manner is determined. We note that, in as much as equivalent equations have isomorphic invariance algebras, we may test the realizations of the invariance algebras for equivalence rather than the equations themselves.
Note that similar ideas have been used by Gangon and Winternitz in order to classify symmetries of nonlinear Schroödinger equations having variable coefficients.
## II Lie-algebraic structures involved in the classification algorithm
Let us take a more detailed look at the second step of the algorithm. As is clear from what has been said above, carrying out this step assumes that there is a classification of non-isomorphic finite-dimensional Lie algebras (in particular, we are interested in a classification of Lie algebras over the real numbers).
One of the central theorems which deals with the structure of Lie algebras is the Levi-Mal’cev theorem:
###### Theorem 2.1
Let $`L`$ be a finite-dimensional Lie algebra over $`𝐑`$ or $`𝐂,`$ and let $`N`$ denote its radical (the largest solvable ideal in $`L`$). Then there exists a semi-simple Lie subalgebra $`S`$ of $`L`$ such that
$$L=S+N$$
(2.1)
Equation (2.1) is called the Levi decomposition of the Lie algebra $`L,`$ and the semi-simple subalgebra $`S`$ is called the Levi factor.
The Levi-Mal’cev decomposition gives us
$$[N,N]N,[S,S]S,[N,S]N,$$
so that any Lie algebra $`L`$ is the semi-direct sum $`L=S+N`$ of its maximal solvable ideal $`N`$ and the semi-simple subalgebra $`S.`$ We see then that this result reduces the task of classifying all Lie algebras to the following problems:
* the classification of all semi-simple Lie algebras;
* the classification of all solvable Lie algebras;
* the classification of all algebras which are semi-direct sums of semi-simple Lie algebras and solvable Lie algebras.
### II.1 Semi-simple Lie algebras.
Of the problems listed above, only that of classifying all semi-simple Lie algebras is completely solved. We have the well-known theorem due to Cartan:
###### Theorem 2.2
(Cartan’s theorem) Any semi-simple complex or real semi-simple Lie algebra can be decomposed into a direct (Lie algebra) sum of ideals which are mutually orthogonal simple subalgebras. Here, orthogonality is with respect to the Cartan-Killing form $`(X,Y)=Tr(adX,adY).`$
Let $`L`$ be a semi-simple Lie algebra. Then, by Cartan’s theorem, we have
$$L=S_1S_2\mathrm{}S_m,$$
where $`S_1,\mathrm{},S_m`$ are simple Lie algebras. Thus, the problem of classifying semi-simple Lie algebras is equivalent to that of classifying all non-isomorphic simple Lie algebras. This classification is known (see, for instance, ).
There are four sequences of classical Lie algebras $`A_n(n1),B_n(n1),C_n(n1),D_n(n1)`$ and five exceptional Lie algebras $`G_2,F_4,E_6,E_7,E_8`$ which together exhaust all the simple complex Lie algebras. There are some isomorphisms between some of these algebras. Indeed there are the following isomorphisms:
$$A_1B_1C_1,B_2C_2,A_3D_3,D_2A_1A_1$$
and there are no other isomorphisms between the series.
The dimensions of the classical complex Lie algebras $`A_n,B_n,C_n`$ and $`D_n`$ are given in the following table:
| Algebra | $`A_n`$ | $`B_n`$ | $`C_n`$ | $`D_n`$ |
| --- | --- | --- | --- | --- |
| Dimension | $`n(n+2)`$ | $`n(2n+1)`$ | $`n(2n+1)`$ | $`n(2n1)`$ |
The dimensions of the exceptional Lie algebras are all even: $`dimG_2=14,dimF_4=52,dimE_6=78,dimE_7=133,dimE_8=248.`$
To describe the real simple Lie algebras one uses the fact that every simple Lie algebra over the reals $`𝐑`$ is either a simple algebra over the complex field $`𝐂`$ (considered as an algebra over $`𝐑`$), or it is the real form of a simple Lie algebra over $`𝐂.`$
The real classical Lie algebras play an important role in the group analysis of differential equations. Below, we give a more detailed description of these Lie algebras. The symbol $`L_k`$ denotes a compact simple Lie algebra.
* Real forms of the algebras $`sl(n,𝐂)(A_{n1},n2)`$
* $`L_k=su(n),`$ the Lie algebra of all skew-symmetric matrices $`Z`$ of order $`n`$ with $`TrZ=0`$ of order $`n`$ with $`TrZ=0.`$
* $`sl(n,𝐑),`$ the Lie algebra of all real matrices $`X`$ of order $`n`$ with $`TrX=0.`$
* $`su(p,q),p+q=n,pq,`$ the Lie algebra of all matrices of the form
$$\left[\begin{array}{cc}Z_1& Z_2\\ Z_2^{}& Z_3\end{array}\right]$$
where $`Z_1,Z_3`$ are skew-symmetric matrices of order $`p`$ and $`q`$ respectively, $`Tr(Z_1+Z_3)=0,`$ and $`Z_2`$ is an arbitrary matrix of order $`q.`$
* $`su^{}(2n),`$ the Lie algebra of all complex matrices of order $`2n`$ of the form
$$\left[\begin{array}{cc}Z_1& Z_2\\ \overline{Z}_2& \overline{Z}_1\end{array}\right]$$
where $`Z_1,Z_2`$ are complex matrices of order $`n`$ with $`Tr(Z_1+\overline{Z}_1)=0.`$
* Real forms of the algebras $`so(2n,𝐂)(D_n,n1)`$
* $`L_k=so(2n),`$ the Lie algebra of all real skew-symmetric matrices of order $`2n.`$
* $`so(p,q),p+q=2n,pq,`$ the Lie algebra of all real matrices of order $`2n`$ of the form
$$\left[\begin{array}{cc}X_1& X_2\\ X_2^T& X_3\end{array}\right]$$
where all the $`X_i`$ are real matrices, and $`X_1,X_3`$ are skew-symmetric matrices of order $`p`$ and $`q`$ respectively, and $`X_2`$ is an arbitrary matrix of order $`q.`$
* $`so^{}(2n),`$ the Lie algebra of all complex matrices of order $`2n`$ of the form
$$\left[\begin{array}{cc}Z_1& Z_2\\ \overline{Z}_2& \overline{Z}_1\end{array}\right]$$
$`Z_1`$ skew-symmetric and $`Z_2`$ Hermitian.
* Real forms of the algebras $`so(2n+1,𝐂)w(B_n,n1)`$
* $`L_k=so(2n+1),`$ the Lie algebra of all real skew-symmetric matrices of order $`2n+1.`$
* $`so(p,q),p+q=2n+1,pq,`$ the Lie algebra of all real matrices of order $`2n+1`$ of the form
$$\left[\begin{array}{cc}X_1& X_2\\ X_2^T& X_3\end{array}\right]$$
where all the $`X_i`$ are real matrices, and $`X_1,X_3`$ are skew-symmetric matrices of order $`p`$ and $`q`$ respectively, and $`X_2`$ is an arbitrary matrix of order $`q.`$
* Real forms of the algebras $`sp(n,𝐂)(C_n,n1)`$
* $`L_k=sp(n),`$ the Lie algebra of all matrices of order $`2n`$ of the form
$$\left[\begin{array}{cc}Z_1& Z_2\\ Z_3& Z_1^T\end{array}\right]$$
where all the $`Z_i`$ are complex matrices of order $`n`$ and $`Z_2,Z_3`$ are symmetric.
* $`sp(n,𝐑),`$ the Lie algebra of all real matrices of order $`2n`$ of the form
$$\left[\begin{array}{cc}X_1& X_2\\ X_3& X_1^T\end{array}\right]$$
where $`X_1,X_2,X_3`$ are all real matrices of order $`n,`$ and $`X_2,X_3`$ are symmetric.
* $`sp(p,q),p+q=n,pq,`$ the Lie algebra of all complex matrices of order $`2n`$ of the form
$$\left[\begin{array}{cccc}Z_{11}& Z_{12}& Z_{13}& Z_{14}\\ Z_{12}^{}& Z_{22}& Z_{14}^T& Z_{24}\\ \overline{Z}_{13}& \overline{Z}_{14}& \overline{Z}_{11}& \overline{Z}_{12}\\ Z_{14}^{}& \overline{Z}_{24}& \overline{Z}_{12}^T& \overline{Z}_{22}\end{array}\right]$$
where the $`Z_{ij}`$ are complex matrices, $`Z_{11}`$ and $`Z_{13}`$ are of order $`p`$, $`Z_{12}`$ and $`Z_{14}`$ are $`p\times q`$ matrices, $`Z_{11}`$ and $`Z_{22}`$ are skew-Hermitian, and $`Z_{13}`$ and $`Z_{24}`$ are symmetric.
The structure of the above real, simple classical Lie algebras is such that every algebra of a higher dimension contains, as a subalgebra, an algebra of the same class but of lower dimension. This allows us to proceed step-by-step when we study the realizations of these algebras as vector fields, at each stage extending the realizations of lower dimension to realizations of higher dimension. If at some stage in this procedure the chain stops, then this implies that there are no realizations within the given type of vector fields of Lie algebras of higher dimension.
In searching for realizations of the classical simple Lie algebras over $`𝐑,`$ it is important to take into account the isomorphisms for the lower-dimensional classical Lie algebras:
$`su(2)so(3)sp(1);`$
$`sl(2,𝐑)su(1,1)so(2,1)sp(1,𝐑);`$
$`so(5)sp(2);`$
$`so(3,2)sp(2,𝐑);`$
$`so(4,1)sp(1,1);`$
$`so(4)so(3)so(3)sp(1)sp(1);`$
$`so(5)sp(2);`$
$`so(2,2)sl(2,𝐑)sl(2,𝐑);`$
$`sl(2,C)so(3,1);`$
$`su(4)so(6);`$
$`sl(4,𝐑)so(3,3);`$
$`su(2,2)so(4,2);`$
$`su(3,1)so^{}(6);`$
$`su^{}(4)so(5,1);`$
$`so^{}(8)so(6,2);`$
$`so^{}(4)su(2)sl(2,𝐑).`$
It is not difficult to see that the Lie algebras of the first two rows have the lowest dimension $`n=3.`$ Thus, in constructing realizations of the classical simple Lie algebras over $`𝐑`$ one may begin with the algebras
$$so(3)=e_1,e_2,e_3,[e_1,e_2]=e_3,[e_2,e_3]=e_1,[e_3,e_1]=e_2;$$
$$sl(2,𝐑)=e_1,e_2,e_3,[e_1,e_3]=2e_2,[e_1,e_2]=e_1,[e_2,e_3]=e_3.$$
Maximal compact subalgebras play an important role for the structure of the simple (and semi-simple) Lie algebras. We have:
###### Theorem 2.3
(Cartan’s Theorem) A semi-simple real Lie algebra has a decomposition of the form
$$L=K\dot{+}P,$$
(2.2)
where
$$[K,K]K,[K,P]P,[P,P]K,$$
(2.3)
and
$`(X,X)<0\mathrm{for}X0\mathrm{in}K,`$ (2.4)
$`(Y,Y)>0\mathrm{for}Y0\mathrm{in}P.`$
If the conditions (2.3), (2.4) are satisfied, then $`K`$ is a maximal compact subalgebra of $`L.`$
The decomposition (2.2) for a real semi-simple Lie algebra is called the Cartan decomposition.
Consider as an example $`so(3,1),`$which is the Lie algebra of the Lorentz group. Denoting by $`K_i(i=1,2,3)`$ the generators of the compact algebra $`so(3)`$, and by $`N_i(i=1,2,3)`$ the generators of the Lorentz boosts, we obtain the commutation relations
$`[K_i,K_j]`$ $`=`$ $`\epsilon _{ijl}K_l,`$
$`[K_i,N_j]`$ $`=`$ $`\epsilon _{ijl}N_l,`$
$`[N_i,N_j]`$ $`=`$ $`\epsilon _{ijl}K_l.`$
Thus, when looking for realizations of the Lie algebra $`so(3,1),`$ one may use the realizations obtained for the Lie algebra $`so(3).`$
In the table below we give the maximal compact subalgebras of the real classical Lie algebras which are non-compact:
| No/o | $`L`$ | $`K`$ | No/o | $`L`$ | $`K`$ |
| --- | --- | --- | --- | --- | --- |
| 1 | $`sl(n,𝐑)`$ | $`so(n)`$ | 5 | $`so^{}(2n)`$ | $`u(n)`$ |
| 2 | $`su(p,q)`$ | $`s(u(p)u(q))`$ | 6 | $`sp(n,𝐑)`$ | $`u(n)`$ |
| 3 | $`su^{}(2n)`$ | $`sp(n)`$ | 7 | $`sp(p,q)`$ | $`sp(p)sp(q)`$ |
| 4 | $`so(p,q)`$ | $`so(p)so(q)`$ | | | |
Here, $`u(n)`$ is the Lie algebra of the unitary group $`U(n),`$ and $`s(u(p)u(q))`$ is the set of all elements $`xu(p)u(q)`$ such that $`Trx=0.`$ Note that the matrix $`e_{ij},`$ defined as a matrix of order $`n`$ with $`1`$ in the $`(i,j)`$ position and zeroes in all other entries, is an element of $`u(n).`$
Because of their large dimension, the exceptional Lie algebras do not play as important a role as the classical simple Lie algebras do in the group analysis of differential equations. For this reason, we only mention briefly the real forms of the algebras of the type $`G_2,F_4,E_6,E_7,E_8,`$ and we consider those subalgebras whose realizations one may use for the construction of realizations of the real exceptional simple Lie algebras. More details about these algebras can be found in .
The algebra $`G_2`$ has real compact form $`g_2`$ and one real non-compact form $`g_2^{^{}}`$. Moreover, $`g_2g_2^{^{}}su(2)su(2)`$.
The algebra $`F_4`$ has real compact form $`f_4`$ and two real non-compact forms $`f_4^{^{}},f_4^{^{\prime \prime }}`$. We also have $`f_4^{^{}}f_4sp(3)su(2),f_4^{^{\prime \prime }}f_4so(9).`$
The algebra $`E_6`$ has real compact form $`e_6`$ and four real non-compact forms $`e_6^{^{}}`$, $`e_6^{^{\prime \prime }}`$, $`e_6^{^{\prime \prime \prime }}`$, $`e_6^{IV}`$. Moreover, $`e_6^{^{}}e_6sp(4),`$ $`e_6^{^{\prime \prime }}e_6su(6)su(2)`$, $`e_6^{^{\prime \prime \prime }}e_6so(10)𝐑`$, $`e_6^{IV}e_6f_4`$.
The algebra $`E_7`$ has real compact form $`e_7`$ and four real non-compact forms $`e_7^{^{}}`$, $`e_7^{^{\prime \prime }}`$, $`e_7^{^{\prime \prime \prime }}`$. We also have $`e_7^{^{}}e_7su(8)`$, $`e_7^{^{\prime \prime }}e_7so(12)su(2)`$, $`e_7^{^{\prime \prime \prime }}e_7e_6𝐑`$.
The algebra $`E_8`$ has real compact form $`e_8`$ and two real non-compact forms $`e_8^{^{}},e_8^{^{\prime \prime }}`$. Also, $`e_8^{^{}}e_8e_7su(2)`$, $`e_8^{^{\prime \prime }}e_8so(16)`$.
### II.2 Solvable Lie algebras.
The problem of classifying solvable Lie algebras up to isomorphism is, as far as we know, completely solved only for real Lie algebras of dimension up to and including six (see for example ). The difficulty in the classification of these algebras is, above all, connected with the fact that the number of non-isomorphic Lie algebras increases considerably with increasing dimension, beginning with dimension five. Thus, according to , there are 66 classes of non-isomorphic real, solvable Lie algebras of dimension five. Furthermore, for dimension six, there are 99 classes of non-isomorphic algebras just amongst the real solvable algebras containing a nilpotent element .
Let us consider in more detail at the structure of solvable Lie algebras over the field $`𝐑`$ with dimension no greater than five. We give a method of searching for their realizations in the class of differential operators.
Let $`L_n`$ denote a solvable Lie algebra of dimension $`n,`$ over a field of characteristic zero. It is known () that there exists a series of subalgebras
$$L_nL_2\mathrm{}L_1L_0=\{0\}$$
such that each subalgebra $`L_i(i=1,\mathrm{},n1)`$ is an ideal of the algebra $`L_{i+1}.`$ This series is called the composition series of the algebra $`L_n.`$
The existence of the composition series for a real solvable Lie algebra allows us to make the following important conclusion: if, in the given class of differential operators, there is a realization of the solvable Lie algebras with $`dimLm,`$ and there is no realization for algebras with $`dimL=m+1,`$ then those realizations which appear give a complete description of the realizations of solvable algebras in the given class of vector fields.
Further, we shall use the following notation: $`A_{k.i}=e_1,\mathrm{},e_k`$ denotes a Lie algebra of dimension $`k,`$ $`e_j(j=1,\mathrm{},k)`$ is its basis, and the index $`i`$ denotes the number of the class to which the given Lie algebra belongs.
Fixing the type of the algebra $`A_{k.i},`$ we shall give only the non-zero commutation relations between the basis elements. Among the solvable Lie algebras over $`𝐑`$ of lowest dimension, we have only one algebra which is one-dimensional $`A_1=e_1,`$ and two algebras which have dimension two:
$`A_{2.1}`$ $`=`$ $`e_1,e_2=A_1A_1=2A_1;`$
$`A_{2.2}`$ $`=`$ $`e_1,e_2,[e_1,e_2]=e_2.`$
Further, we shall call decomposable Lie algebras those algebras which can be decomposed as a direct sum of solvable algebras of lower dimension. We give a list of all solvable Lie algebras, up to and including dimension five, in Appendix 1.
It is clear that the search for realizations of solvable Lie algebras over $`𝐑`$ must be begun with the description of the inequivalent forms of the general infinitesimal operator, up to equivalence under the transformations of $`.`$ Each of the operators obtained will be a basis for the inequivalent realizations of one-dimensional Lie algebras. Further, the completion of the basis operators of each of the one-dimensional Lie algebras, by an infinitesimal operator of the most general form, is done by extension of the realizations of the one-dimensional Lie algebras to realizations of two-dimensional Lie algebras. In doing this, in order to simplify the form of the second basis operator one uses those transformations from $``$ which leave invariant the form of the first basis operator.
Analogously, the realizations of the two-dimensional Lie algebras which one obtains, are extended to realizations of three-dimensional solvable Lie algebras, and then the realizations of the three-dimensional Lie algebras are extended in the same way to realizations of the four-dimensional algebras, and so on.
In the extension of the realizations of Lie algebras of lower dimension to realizations of decomposable solvable Lie algebras is done simply by adding to each realization a basis operator which commutes with all the other basis elements.
For the construction of the realizations of non-decomposable solvable Lie algebras, as is shown by an analysis of their structure above, one may also carry out the extension of the realizations of Lie algebras of lower dimension to realizations of non-decomposable solvable Lie algebras of higher dimension.
### II.3 Semi-direct sums of semi-simple and solvable Lie algebras.
Lie algebras which are semi-direct sums of semi-simple and solvable Lie algebras can be divided into two classes:
* those Lie algebras which are direct sums of semi-simple and solvable Lie algebras (decomposable algebras);
* algebras which cannot be written as a direct sum of semi-simple and solvable Lie algebras (indecomposable algebras).
Since decomposable algebras have the structure
$$L=SN$$
where $`S`$ is the Levi factor and $`N`$ is the radical (maximal solvable ideal of $`L`$), then a complete description of these algebras is easily obtained by combining the semi-simple and solvable Lie algebras. However, since the classification of solvable Lie algebras has only been done partially, there is a corresponding incompleteness in the classification of decomposable Lie algebras. The classification of indecomposable Lie algebras has been done only as far as for Lie algebras of dimension eight (). These are Lie algebras whose Levi factor is $`sl(2,𝐑)`$ or $`so(3).`$
We give a complete list of these algebras in Appendix 2. We use the following notation:
$`sl(2,𝐑)=e_1,e_2,e_3;[e_1,e_2]=2e_2,[e_1,e_3]=2e_3,[e_2,e_3]=e_1`$
$`so(3)=e_1,e_2,e_3;[e_1,e_2]=e_3,[e_1,e_3]=e_2,[e_2,e_3]=e_1.`$
We note that the basis of $`sl(2,𝐑)`$ given here differs from that given previously, but it is not difficult to see that they are isomorphic. Indeed, if we make the transformations
$$e_12e_2,e_2e_3,e_3e_1$$
then we have an isomorphism from the basis given here to the basis given previously.
In denoting the radicals $`N=e_4,\mathrm{},e_m(m=dimN3),`$ we keep to the notation used for the classification of solvable algebras given above. Moreover, the corresponding commutation relations for the basis operators of $`N`$ can easily be obtained from those given by replacing the index $`i`$ by the index $`i+3.`$ Thus, in listing the algebras which are semi-direct sums of semi-simple and solvable Lie algebras, we give only those commutators $`[e_i,e_j]=c_{ij}^ke_k,i=1,2,3;j,k=4,\mathrm{},m`$ which are non-zero.
Taking into account the above classification of finite-dimensional real Lie algebras, we take a closer look at the second step of the algorithm for the group classification of differential equations.
After having completed the first step of the algorithm, we have the general form of the infinitesimal symmetry operator (together with a defining system) for the given equation, and we have group $``$ of equivalence transformations of this equation.
At the beginning of the second step we have to bring the symmetry operator to the simplest form, using transformations from $`.`$ We note that it is well-known () that linearization of the vector field is not possible since the group $``$ is a subgroup of all the local transformations of the manifold $`V`$ of dependent and independent variables which enter into the differential equation. Thus, we will obtain, up to equivalence, some finite set of simplest forms for the symmetry operator.
Further, using the determining equations, we find from the equation at hand, equations which admit the operators we have obtained as symmetry operators. With this, we will obtain the group classification of the differential equations of the given form which admit one-dimensional Lie algebras of invariance.
The list of simplest forms of the symmetry operator which we find, allows us to take one of the symmetry operators in its simplest form, when we consider the realizations of Lie algebras of higher dimension. Moreover, we may first consider realizations of semi-simple and solvable Lie algebras.
When we consider realizations of semi-simple Lie algebras, we must, as well as looking at those semi-simple Lie algebras which appear in our list, also take into account semi-simple Lie algebras which are direct sums of semi-simple Lie algebras which do not appear in the list given above.
We take a closer look at low-dimensional semi-simple Lie algebras. As we noted above, the semi-simple Lie algebras of lowest dimension are the algebras $`sl(2,𝐑)`$ and $`so(3),`$ both having dimension 3. Then we have Lie algebras of dimension 6 (the algebras $`so(4)so(3)so(3),so(3,1),so(2,2)sl(2,𝐑)sl(2,𝐑),so^{}(4)so(3)sl(2,𝐑)`$); dimension 8 (the algebras $`su(3),sl(3,𝐑),su(2,1)`$).These algebras can be found in our list. However, the semi-simple Lie algebras of dimension 9 (the algebras $`so(4)so(3),so(4)sl(2,𝐑),so(3,1)so(3),so(3,1)sl(3,𝐑),`$ $`so(2,2)so(3),so(2,2)sl(2,𝐑)`$) are not to be found in our list.
For solvable Lie algebras, one may extend realizations of lower dimensional algebras to realizations of higher-dimensional algebras according to the scheme given above. Moreover, for solvable Lie algebras of higher dimension, the composition series play an important role: if there is a realization in terms of operators of a given class of solvable Lie algebras of dimension $`m1,`$ but no realization of a solvable Lie algebra of dimension $`m,`$ the there will, a priori, be no realizations of solvable Lie algebras of dimension greater than $`m.`$
When looking at the realization of Lie algebras which are semi-direct sums of semi-simple and solvable Lie algebras, one must extend a given realization of semi-simple Lie algebras by operators which will be basis operators of the corresponding radicals. Moreover, it is only necessary to take into account those radicals which are isomorphic to solvable Lie algebras which have realizations in the given class of vector fields.
Finally, we note that since the classification of solvable Lie algebras, and those algebras which are semi-direct sums of semi-simple and solvable Lie algebras, is incomplete, then it is not possible to give a complete group classification of differential equations within the present framework. However, the problem of the complete group classification of, for instance, scalar equations in two-dimensional space-time which are invariant under finite-dimensional Lie algebras, is constructive in this approach.
## III Classification results
First of all let us mention some papers in which group classification of particular equations of the form (0.1) has been carried out.
| Ovsjannikov (1959) | $`F=F(u)`$, | $`G={\displaystyle \frac{dF}{du}}u_x^2`$ ; |
| --- | --- | --- |
| Akhatov et al (1987) | $`F=F(u_x)`$, | $`G=0`$ ; |
| Dorodnitsyn (1982) | $`F=F(u)`$, | $`G={\displaystyle \frac{dF}{du}}u_x^2+g(u)`$ ; |
| Oron & Rosenau (1986), | | |
| Edwards (1994) | $`F=F(u)`$, | $`G={\displaystyle \frac{dF}{du}}u_x^2+f(u)u_x`$ ; |
| Gandarias (1996) | $`F=u^n`$, | $`G={\displaystyle \frac{dF}{du}}u_x^2+g(x)u^mu_x+f(x)u^s`$ ; |
| Cherniha & Serov (1998) | $`F=F(u)`$, | $`G={\displaystyle \frac{dF}{du}}u_x^2+f(u)u_x+g(u)`$ ; |
| Zhdanov & Lahno (1999) | $`F=1`$, | $`G=G(t,x,u,u_x)`$ . |
We shall apply the algorithm described above in order to perform an exhaustive group classification of invariant equations of the general form (0.1). That is, we shall describe all inequivalent forms of functions $`F,G`$ such that the corresponding equation admits a non-trivial symmetry group.
### III.1 Computation of the equivalence group admitted by equation (0.1)
The first step of the algorithm is the determination of the most general form of the infinitesimal symmetry operator admitted by the PDE (0.1). To this end, we use Lie’s method and look for a symmetry generator in the form
$`Q=\tau (t,x,u)_t+\xi (t,x,u)_x+\eta (t,x,u)_u,`$ (3.1)
where $`\tau `$$`\xi `$$`\eta `$ are arbitrary, real-valued smooth functions defined in some subspace of the space $`V=XR^1`$ of the independent variables $`X=t,x`$ and the dependent variable $`R^1=u`$.
As a result, we find that the operator (3.1) generates a one-parameter symmetry group of equation (0.1) iff
$$\phi ^t[\tau F_t+\xi F_x+\eta F_u+\phi ^xF_{u_x}]u_{xx}\phi ^{xx}F\tau G_t\xi G_x\eta G_u\phi ^xG_{u_x}|_{u_t=Fu_{xx}+G}=0,$$
(3.2)
where
$`\phi ^t`$ $`=`$ $`D_t(\eta )u_tD_t(\tau )u_xD_t(\xi ),`$
$`\phi ^x`$ $`=`$ $`D_x(\eta )u_tD_x(\tau )u_xD_x(\xi ),`$
$`\phi ^{xx}`$ $`=`$ $`D_x(\phi ^x)u_{tx}D_x(\tau )u_{xx}D_x(\xi )`$
and $`D_t`$$`D_x`$ are operators of total differentiation in $`t`$ and $`x`$ respectively.
After simplifying (3.2) we arrive at the following assertion.
###### Lemma 3.1
The symmetry group of the nonlinear heat equation PDE (0.1) is generated by the infinitesimal operators of the form
$$Q=a(t)_t+b(t,x,u)_x+c(t,x,u)_u,$$
(3.3)
where $`a,b,c`$ are real-valued functions that satisfy the system of PDEs
$`(2b_x+2u_xb_u\dot{a})F=aF_t+bF_x+cF_u+(c_x+u_xc_uu_xb_xu_x^2b_u)F_{u_x},`$
$`c_tu_xb_t+(c_u\dot{a}u_xb_u)G+(u_xb_{xx}c_{xx}2u_xc_{ux}u_x^2c_{uu}+`$ (3.4)
$`+2u_x^2b_{xu}+u_x^3b_{uu})F=aG_t+bG_x+cG_u+(c_x+u_xc_uu_xb_xu_x^2b_u)G_{u_x}.`$
In the rest of this paper we use the notation $`\dot{a}={\displaystyle \frac{da}{dt}}`$, $`\ddot{a}={\displaystyle \frac{d^2a}{dt^2}}`$.
In order to construct the equivalence group $``$ of the class of PDEs (0.1) one has to select from the set of invertible changes of variables of the space $`V`$
$`\overline{t}=\alpha (t,x,u),\overline{x}=\beta (t,x,u),v=\gamma (t,x,u),{\displaystyle \frac{D(\alpha ,\beta ,\gamma )}{D(t,x,u)}}0,`$ (3.5)
those changes of variable which do not alter the form of the class of PDEs (0.1).
###### Lemma 3.2
The maximal equivalence group $``$ of the class of PDEs (0.1) reads as
$`\overline{t}=T(t),\overline{x}=X(t,x,u),v=U(t,x,u),`$ (3.6)
where $`\dot{T}0,`$ $`{\displaystyle \frac{D(X,U)}{D(x,u)}}0`$.
Proof. Let (3.5) be an invertible change of variables that transforms equation (0.1) into another equation of the same form (0.1), namely,
$`v_{\overline{t}}=\stackrel{~}{F}(\overline{t},\overline{x},v,v_{\overline{x}})v_{\overline{x}\overline{x}}+\stackrel{~}{G}(\overline{t},\overline{x},\overline{v},\overline{v}_x).`$ (3.7)
Computing $`u_x`$ according to (3.5) we get
$`u_x={\displaystyle \frac{v_{\overline{t}}\alpha _x+v_{\overline{x}}\beta _x\gamma _x}{\gamma _uv_{\overline{t}}\alpha _uv_{\overline{x}}\beta _u}}.`$ (3.8)
As the functions $`F`$, $`G`$ in (0.1) and $`\stackrel{~}{F}`$, $`\stackrel{~}{G}`$ in (3.7) are arbitrary functions of the corresponding arguments, we must have that
$$u_xg(\overline{t},\overline{x},v,v_{\overline{x}})$$
for some function $`g.`$ This implies that $`\alpha _x=\alpha _u=0`$ in (3.8). Consequently, $`\alpha =T(t),\dot{T}0.`$
Next, making the change of variables (3.5), where $`\alpha =T(t)`$, we arrive at the relations
$`u_tv_{\overline{t}}\dot{T}(\gamma _uv_{\overline{x}}\beta _u)^1+\theta _1(\overline{t},\overline{x},v,v_{\overline{x}}),`$
$`u_{xx}v_{\overline{x}\overline{x}}\theta _2(\overline{t},\overline{x},v,v_{\overline{x}})+\theta _3(\overline{t},\overline{x},v,v_{\overline{x}}),`$ (3.9)
where $`\theta _1,\theta _20,\theta _3`$ are some functions of $`\alpha ,\beta ,\gamma `$ and of their derivatives. Then, inserting $`u_t,u_{xx}`$ from (3.9) into (0.1), we arrive at a PDE of the form (3.7). The lemma is proved.
### III.2 Classification of equations (0.1) invariant under semi-simple Lie algebras
Now we proceed to solving the classification problem for the nonlinear heat-conductivity equation (0.1). Our first step is to construct realizations of finite-dimensional real Lie algebras whose representation space is spanned by operators of the form (3.3). It should be noted that the realizations are constructed up to equivalence as determined by the transformations (3.6). In the second step, we choose those realizations which are invariance algebras of PDE (0.1) and thus specify the form of the functions $`F,G`$. Finally, in the third step, we find the maximal symmetry groups of the equations we obtain, and thus complete the group classification of PDE (0.1).
By the Levi-Mal’cev theorem, we need only consider the cases of semi-simple, solvable and semi-direct sums of semi-simple and solvable symmetry algebras. This will yield an exhaustive description of invariant equations of the form (0.1). In this subsection, we analyze the case of semi-simple symmetry algebras. As semi-simple Lie algebras can always be decomposed into a direct sum of simple Lie algebras (which is the content of Cartan’s theorem), we begin with the lowest dimensional simple Lie algebras $`sl(2,𝐑)`$ and $`so(3)`$.
We begin by proving the following useful lemma:
###### Lemma 3.3
There are changes of variables (3.6), that reduce an operator (3.3) to one of the operators below:
$`Q`$ $`=`$ $`_t,`$ (3.10)
$`Q`$ $`=`$ $`_x.`$ (3.11)
Proof. Making the change of variables (3.6) transforms operator (3.3) to the following one:
$`QQ^{}=a\dot{T}_{\overline{t}}+(aX_t+bX_x+cX_u)_{\overline{x}}+(aU_t+bU_x+cU_u)_v.`$ (3.12)
Suppose $`a0.`$ Then, choosing in (3.6) the function $`T`$ to be a solution of the equation $`\dot{T}=a^1`$ and the functions $`X`$ and $`U`$ to be independent fundamental solutions of the first-order PDE
$$aY_t+bY_x+cY_u=0,Y=Y(t,x,u),$$
we find that the operator (3.12) takes the form $`Q^{}=_{\overline{t}}`$.
Now suppose $`a=0`$. Then $`b^2+c^20`$. If $`b0`$, then choosing in (3.6) a particular solution of PDE $`bX_x+cX_u=1`$ as the function $`X`$ and a fundamental solution of PDE $`bU_x+cU_u=0`$ as the function $`U`$, we transform (3.12) to become $`Q^{}=_{\overline{x}}`$.
If $`b=0`$$`c0`$, then making the change of variables (3.6) with $`\overline{t}=t,\overline{x}=u,v=x,`$ we again get the case $`b0`$.
By the direct calculation we can verify that there is no transformation from $``$, that reduce operator (3.10) to the form (3.11).
The lemma is proved.
###### Theorem 3.1
Within the equivalence relation $`,`$ there exists only one realization of the algebra $`so(3)`$ by operators of the form (3.3):
$`_x,\mathrm{tan}u\mathrm{sin}x_x+\mathrm{cos}x_u,\mathrm{tan}u\mathrm{cos}x_x\mathrm{sin}x_u,`$ (3.13)
It is the invariance algebra of an equation from the class (0.1). Furthermore, the most general form of the functions $`F,G`$ allowing for PDE (0.1) to be invariant under the above realization is given by
$$F=\frac{\mathrm{sec}^2u}{1+\omega ^2},G=\frac{2\omega ^2+1}{1+\omega ^2}\mathrm{tan}u+\sqrt{1+\omega ^2}\stackrel{~}{G}(t),\omega =u_x\mathrm{sec}u.$$
(3.14)
Provided the function $`\stackrel{~}{G}`$ is arbitrary, the realization (3.13) is the maximal symmetry algebra of the corresponding equation (0.1).
Proof. The Lie algebra $`so(3)=Q_1,Q_2,Q_3`$ is defined by the following commutation relations:
$`[Q_1,Q_2]=Q_3,[Q_1,Q_3]=Q_2,[Q_2,Q_3]=Q_1.`$ (3.15)
To describe all inequivalent realizations of the algebra $`so(3)`$ we take operators of the form (3.3) as the basis elements $`Q_i`$ $`(i=1,2,3)`$ of $`so(3)`$ and then study the restrictions imposed on their coefficients by relations (3.15). We also use transformations (3.6) in order to simplify the final forms of the basis elements.
In view of Lemma 3.3, we can take one of the basis elements of the algebra $`so(3)`$ (say, $`Q_1`$) either in the form $`_t`$ or $`_x`$.
Let $`Q_1=_t`$. Using the first two commutation relations from (3.15) yields
$`Q_2`$ $`=`$ $`\lambda \mathrm{cos}t_t+[b\mathrm{cos}t+\beta \mathrm{sin}t]_x+[c\mathrm{cos}t+\gamma \mathrm{sin}t]_u,`$
$`Q_3`$ $`=`$ $`\lambda \mathrm{sin}t_t+[b\mathrm{sin}t+\beta \mathrm{cos}t]_x+[c\mathrm{sin}t+\gamma \mathrm{cos}t]_u,`$
where $`\lambda =\mathrm{const}𝐑,b=b(x,u),c=c(x,u),\beta =\beta (x,u),\gamma =\gamma (x,u)`$ are arbitrary smooth functions.Then, using the third commutation relation, we arrive at the equation $`\lambda ^2=1`$ which has no real solutions $`\lambda `$. Consequently, in the case when the operator $`Q_1`$ is equivalent to the operator $`_t,`$ there are no realizations of the algebra $`so(3)`$.
Turn now to the case $`Q_1=_x`$. As a straightforward calculation shows, the most extensive subgroup of the equivalence group $``$ not altering the form of $`Q_1`$ is of the form
$`\overline{t}=T(t),\overline{x}=x+X(t,u),v=U(t,u),\dot{T}0,U_u0.`$ (3.16)
Using the first two commutation relations from (3.15) we get
$`Q_2`$ $`=`$ $`\alpha \mathrm{cos}(x+\gamma )_x+\beta \mathrm{cos}(x+\theta )_u,`$
$`Q_3`$ $`=`$ $`\alpha \mathrm{sin}(x+\gamma )_x\beta \mathrm{sin}(x+\theta )_u,`$ (3.17)
where $`\alpha =\alpha (t,u),\gamma =\gamma (t,u),\beta =\beta (t,u),\theta =\theta (t,u)`$ are arbitrary smooth functions. Now, either $`\beta =0`$ or $`\beta 0.`$ If $`\beta =0,`$ the third commutation relation givews $`\alpha ^2=1,`$ which has no real solutions. Consequently, $`\beta 0`$.
Choosing in (3.16$`X=\theta `$ and furthermore, taking an arbitrary solution of the equation $`U_u=\beta ^1`$, as $`U`$, we simplify the forms of the operators $`Q_2,Q_3`$ to obtain
$`Q_2`$ $`=`$ $`\alpha \mathrm{cos}(x+\gamma )_x+\mathrm{cos}x_u,`$
$`Q_3`$ $`=`$ $`\alpha \mathrm{sin}(x+\gamma )_x\mathrm{sin}x_u,`$
where $`\alpha =\alpha (t,u),\gamma =\gamma (t,u)`$ are arbitrary smooth functions (here and in the following, we keep the initial designations for the transformed operators to simplify the notation).
The third commutation relation for the operators $`Q_1,Q_2`$ which we have obtained, yields the equations $`\mathrm{cos}\gamma =0,\alpha ^2+\alpha _u\mathrm{sin}\gamma =1`$, whence
$`Q_2`$ $`=`$ $`\mathrm{tan}[u\pm \stackrel{~}{\alpha }(t)]\mathrm{sin}x_x+\mathrm{cos}x_u,`$
$`Q_3`$ $`=`$ $`\mathrm{tan}[u\pm \stackrel{~}{\alpha }(t)]\mathrm{cos}x_x\mathrm{sin}x_u,`$
where $`\stackrel{~}{\alpha }(t)`$ is an arbitrary smooth function.
Finally, putting $`T=t,X=0,U=u\pm \stackrel{~}{\alpha }(t)`$ in (3.16) we find that the above realization is equivalent to (3.13).
To complete the proof we have to verify whether there exists an equation of the form (0.1), whose symmetry algebra contains subalgebra (3.13). Invariance of (0.1) with respect to the one-parameter group having the generator $`Q_1`$ means that $`F=F(t,u,u_x),G=G(t,u,u_x)`$. Writing down condition (3.4) for the operators $`Q_2`$$`Q_3`$ we get the following system of PDEs:
$`F_uu_x\mathrm{tan}uF_{u_x}=2\mathrm{tan}uF,`$
$`(1+u_x^2\mathrm{sec}^2u)F_{u_x}=2u_x\mathrm{sec}^2uF,`$
$`u_x\mathrm{sec}^2uG+u_x\mathrm{tan}u(12u_x^2\mathrm{sec}^2u)F=(1+u_x^2\mathrm{sec}^2u)G_{u_x},`$
$`(1+2u_x^2\mathrm{sec}^2u)F=G_uu_x\mathrm{tan}uG_{u_x}.`$
Solving the first two equations of the above system gives
$$F=\frac{\mathrm{sec}^2u}{1+\omega ^2}\stackrel{~}{F}(t),$$
where $`\omega =u_x\mathrm{sec}u`$. Integrating the fourth equation yields $`G={\displaystyle \frac{2\omega ^2+1}{1+\omega ^2}}\mathrm{tan}u\stackrel{~}{F}(t)+\overline{G}(t,\omega )`$. Finally, solving the third equation we get the form of $`\overline{G}(t,\omega )`$
$$\overline{G}(t,\omega )=\sqrt{1+\omega ^2}\stackrel{~}{G}(t).$$
Thus PDE (0.1) is invariant with respect to the algebra (3.13) iff
$$F=\frac{\mathrm{sec}^2u}{1+\omega ^2}\stackrel{~}{F}(t),G=\frac{2\omega ^2+1}{1+\omega ^2}\mathrm{tan}u\stackrel{~}{F}(t)+\sqrt{1+\omega ^2}\stackrel{~}{G}(t),\omega =u_x\mathrm{sec}u.$$
(3.18)
with arbitrary smooth functions $`\stackrel{~}{F}(t)`$, $`\stackrel{~}{G}(t)`$ provided that $`\stackrel{~}{F}(t)0`$.
Evidently, the change of variables (3.16) with $`X=0,U=u`$ does not alter the forms of the operators of the realization (3.13). Choosing a solution of the equation $`\dot{T}=\stackrel{~}{F}`$ as $`T`$, we get $`\stackrel{~}{F}(t)=1`$. By direct computation one shows that if the function $`\stackrel{~}{G}(t)`$ is arbitrary, then the algebra (3.13) is the maximal invariance algebra admitted by the equation obtained. The theorem is proved.
###### Theorem 3.2
There exist five inequivalent realizations of the algebra $`sl(2,𝐑)`$ by operators (3.3), which are admitted by PDEs of the form (0.1)
$`2t_t+x_x,t^2_ttx_x+x^2_u,_t,`$ (3.19)
$`2t_t+x_x,t^2_t+x(x^2t)_x,_t,`$ (3.20)
$`2x_xu_u,x^2_x+xu_u,_x,`$ (3.21)
$`2x_xu_u,(u^4x^2)_x+xu_u,_x,`$ (3.22)
$`2x_xu_u,(u^4+x^2)_x+xu_u,_x.`$ (3.23)
The forms of the functions $`F,G`$ determining the corresponding invariant equations are given as follows:
| $`sl(2,𝐑)`$ | $`F`$ | $`G`$ |
| --- | --- | --- |
| (3.19) | $`\stackrel{~}{F}(\omega )`$ | $`x^2\left[\stackrel{~}{G}(\omega )2u\stackrel{~}{F}(\omega )+u^2u\omega \right],\omega =2uxu_x`$ |
| (3.20) | $`\omega ^3`$ | $`x^2\left[{\displaystyle \frac{1}{4}}\omega +3\omega ^2+\omega ^1\stackrel{~}{G}(u)\right],\omega =xu_x`$ |
| (3.21) | $`u^4`$ | $`2u^5u_x^2`$ |
| (3.22) | $`u^4\left(1+4\omega ^2\right)^1`$ | $`u\left[\sqrt{1+4\omega ^2}\stackrel{~}{G}(t){\displaystyle \frac{10\omega ^2+1}{8\omega ^2+2}}\right],\omega =u^3u_x`$ |
| (3.23) | $`u^4\left(14\omega ^2\right)^1`$ | $`u\left[\sqrt{|14\omega ^2|}\stackrel{~}{G}(t)+{\displaystyle \frac{10\omega ^21}{8\omega ^22}}\right],\omega =u^3u_x`$ |
If the functions $`\stackrel{~}{F},\stackrel{~}{G}`$ are arbitrary, then the corresponding realizations of the algebra $`sl(2,𝐑)`$ are maximal invariance algebras of the respective PDEs (0.1). Furthermore, the maximal symmetry group admitted by the third PDE from the above list
$$u_t=u^4u_{xx}2u^5u_x^2$$
is the five-dimensional Lie algebra $`sl(2,𝐑)L_{2.1}`$, where $`sl(2,𝐑)`$ is given in (3.21) and $`L_{2.1}=4t_t+u_u,_t`$.
Proof. The Lie algebra $`sl(2,𝐑)=Q_1,Q_2,Q_3`$ is defined by the following commutation relations:
$$[Q_1,Q_2]=2Q_2,[Q_1,Q_3]=2Q_3,[Q_2,Q_3]=Q_1.$$
(3.24)
In view of Lemma 3.3 we can choose the operator $`Q_3`$ either in the form $`_t`$ or $`_x`$.
Let $`Q_3=_t`$. Imposing the second commutation relation from (3.24) gives (up to equivalence under $``$) the operator $`Q_1`$ either equals to $`2t_t`$ or $`2t_t+x_x`$.
If $`Q_1=2t_t`$, then it follows from the remaining commutation relation that $`Q_2=t^2_t`$, so that we obtain the realization $`2t_t,t^2_t,_t`$. However, PDE (0.1) can admit this algebra only when the condition $`F=0`$ holds. This contradicts the assumption $`F0`$ and, consequently, there are no corresponding invariant equations within the class (0.1).
If now $`Q_1=2t_t+x_x`$, we get (up to equivalence under $``$) the realization $`2t_t+x_x,t^2_ttx_x,_t`$ and the realizations (3.19), (3.20) of the algebra $`sl(2,𝐑)`$. Substituting into the invariance conditions (3.4) shows that the first realization cannot be admitted by PDE (0.1). This leaves us with the realizations (3.19), (3.20).
The most general equations (0.1) invariant with respect to realizations (3.19), (3.20) are given by:
$`F=\stackrel{~}{F}(\omega ),G=x^2\left[\stackrel{~}{G}(\omega )2u\stackrel{~}{F}(\omega )+u^2u\omega \right],\omega =2uxu_x,`$
$`F=\omega ^3\stackrel{~}{F}(u),G=x^2\left[{\displaystyle \frac{1}{4}}\omega +3\omega ^2\stackrel{~}{F}(u)+\omega ^1\stackrel{~}{G}(u)\right],\omega =xu_x.`$
It is not difficult to show that the change of variables
$$t=t,x=x,u=U(v),U^{}0,v=v(t,x)$$
does not alter the form of the basis operators of the realization (3.20). So, choosing the function $`U`$ to be a solution of the equation $`(U^{})^3=\stackrel{~}{F}(U)`$, we can transform PDE (0.1) invariant under (3.20) in such a way that $`\stackrel{~}{F}1`$.
We now turn to the case $`Q_3=_x`$. Using the commutation relations (3.24) we find that the inequivalent realizations of $`sl(2,𝐑)`$ within the class of operators (3.3) are exhausted by the realization $`2x_x,x^2_x,_x`$ and by the realizations (3.21), (3.22), (3.23). The invariance conditions (3.4) show that the first realization cannot be an invariance algebra of PDE of the form (0.1). The remaining realizations are invariance algebras of PDEs (0.1) under proper specification of the functions $`F,G:`$
$`F=u^4\stackrel{~}{F}(t),G=2u^5u_x^2\stackrel{~}{F}(t)+u\stackrel{~}{G}(t),(\text{ for the realization (}\text{3.21}\text{)});`$ (3.25)
$`F={\displaystyle \frac{1}{u^4\left(1+4\omega ^2\right)}}\stackrel{~}{F}(t),G=u\left[\sqrt{1+4\omega ^2}\stackrel{~}{G}(t){\displaystyle \frac{10\omega ^2+1}{8\omega ^2+2}}\stackrel{~}{F}(t)\right],`$
$`\omega =u^3u_x,(\text{ for the realization (}\text{3.22});`$ (3.26)
$`F={\displaystyle \frac{1}{u^4\left(14\omega ^2\right)}}\stackrel{~}{F}(t),G=u\left[\sqrt{|14\omega ^2|}\stackrel{~}{G}(t)+{\displaystyle \frac{10\omega ^21}{8\omega ^22}}\stackrel{~}{F}(t)\right],`$
$`\omega =u^3u_x,(\text{ for the realization (}\text{3.23}\text{)}).`$ (3.27)
As the change of variables
$$\overline{t}=T,\overline{x}=x,v=U(t)u,T0,U0$$
does not alter the form of the basis operators of the realization (3.21), we can use it in order to simplify the forms of $`F,G`$. Choosing the functions $`T`$ and $`U`$ to be solutions of the equations $`\dot{U}=U\stackrel{~}{G}(t)`$$`U0,`$ $`\dot{T}=\stackrel{~}{F}U^4,`$ we obtain $`\stackrel{~}{F}1,`$ $`\stackrel{~}{G}0`$ in (3.25).
Similarly, using the change of variables
$$\overline{t}=T(t),\overline{x}=x,v=u$$
which preserve the form of the basis operators of the realizations (3.22), (3.23) we can choose $`\stackrel{~}{F}1`$ in (3.26), (3.27).
Computing the maximal invariance algebra of the PDE which admits the realization (3.21) we get the five-dimensional Lie algebra which is the direct sum of $`sl(2,𝐑)`$ having the basis elements (3.21) and the two-dimensional solvable Lie algebra $`L_{2.1}=4t_t+u_u,_t.`$
The remaining invariant equations contain an arbitrary function. If there are no additional constraints on this function, then the realizations (3.19), (3.20), (3.22), (3.23) of $`sl(2,𝐑)`$ are easily shown to be the maximal invariance algebras of the corresponding invariant equations.
The theorem is proved.
###### Theorem 3.3
The realizations of the algebras $`so(3)`$ and $`sl(2,𝐑)`$, given in Theorems 3.1, 3.2, exhaust the set of all possible realizations of semi-simple Lie algebras by operators (3.3) which are admitted by PDEs of the form (0.1).
Proof. The simple Lie algebras of the lowest dimension admit the following isomorphisms:
$$su(2)so(3)sp(1),sl(2,𝐑)su(1,1)so(2,1)sp(1,𝐑).$$
¿From this it follows that the realizations given in Theorems 3.1, 3.2 exhaust the set of all possible realizations of three-dimensional simple Lie algebras which are symmetry algebras of (0.1).
The next admissible dimension for simple Lie algebras is six. There are four distinct six-dimensional simple Lie algebras over the field of real numbers, namely, $`so(4),so(3,1),so(2,2),`$ and $`so^{}(4)`$.
As $`so(4)=so(3)so(3)`$, we have $`so(4)=Q_i,K_i|i=1,2,3`$, where $`Q_1,Q_2,Q_3=so(3),K_1,K_2,K_3=so(3)`$, and we have the commutation relations, $`[Q_i,K_j]=0,i,j=1,2,3`$. Making use of Theorem 3.1 we put the basis operators $`Q_i(i=1,2,3)`$ to be equal to the corresponding basis operators of realization (3.13). Next, the commutation relations $`[Q_i,K_j]=0,(i,j=1,2,3)`$ imply the following form of the operators $`K_j`$:
$`K_j=a_j(t)_t,a_j0,j=1,2,3.`$ (3.28)
Using the change of variables
$$\overline{t}=T(t),\overline{x}=x,v=u,$$
(which does not change the form of operators (3.13)), we can transform the operator $`K_1`$ to become $`K_1=_t`$. Checking the commutation relations for the algebra $`so(3)`$ yields that
$$K_2=\lambda \mathrm{cos}(t+\lambda _1)_t,K_3=\lambda \mathrm{sin}(t+\lambda _1),$$
where $`\{\lambda _1,\lambda \}𝐑`$ with $`\lambda ^2=1`$. Consequently, there are no realizations of the algebra $`so(4)`$ within the class of operators (3.3), which are symmetry algebras of (0.1).
We have the relation $`so^{}(4)so(3)sl(2,𝐑)`$. So, in order to construct realizations of $`so^{}(4)`$ we have to describe realizations of the algebra $`sl(2,𝐑)`$ by operators of the form (3.28). Now, in proving Theorem 3.2, we established, in particular, that there is a unique realization of $`sl(2,𝐑)`$ by operators (3.28), $`2t_t,t^2_t,_t`$, which, however, cannot be admitted by a PDE of the form (0.1). This eliminates $`so^{}(4)`$.
The algebra $`so(3,1)`$ admits the Cartan decomposition $`Q_1,Q_2,Q_3\dot{+}N_1,N_2,N_3`$, where $`Q_1,Q_2,Q_3=so(3),[Q_i,N_j]=N_k,[N_i,N_j]=Q_k,`$ $`i,j,k=\mathrm{cycle}(1,2,3)`$. Thus, taking as $`Q_i(i=1,2,3)`$ the corresponding basis operators of realization (3.13) and computing the forms of the operators $`N_1,N_2,N_3,`$ we get within the equivalence relation $``$ the following relations:
$$N_1=\mathrm{cos}u_u,N_2=\mathrm{sec}u\mathrm{cos}x_x+\mathrm{sin}u\mathrm{sin}x_u,N_3=\mathrm{sec}u\mathrm{sin}x_x+\mathrm{sin}u\mathrm{cos}x_u.$$
Imposing the invariance conditions (3.4) for the operator $`N_1`$ gives that $`F=0`$, contradicting our initial assumption $`F0`$.
In studying realizations of the algebra $`so(2,2)`$ we use the fact that $`so(2,2)sl(2,𝐑)sl(2,𝐑)`$. In view of this, we can choose the basis operators of this algebra so that $`so(2,2)=Q_i,K_i|i=1,2,3`$, where $`Q_1,Q_2,Q_3=sl(2,𝐑),K_1,K_2,K_3=sl(2,𝐑)`$, and , $`[Q_iK_j]=0,i,j=1,2,3`$. Now we can take as $`Q_1,Q_2,Q_3`$ the corresponding basis operators of the realizations of $`sl(2,𝐑)`$ given by (3.19)–(3.23). However, as further analysis shows, these realizations cannot be extended to a realization of $`so(2,2)`$ which could be a symmetry algebra of PDE of the form (0.1).
Thus there no are realizations of six-dimensional simple Lie algebras by operators (3.3), which are symmetry algebras of (0.1).
The same assertion holds true for the simple Lie algebras of dimension eight ($`sl(3,𝐑)`$$`su(3)`$$`su(2,1)`$), which is the next admissible dimension for real simple Lie algebras.
As $`su^{}(4)so(5,1)`$ and since the algebra $`so(5,1)`$ contains $`so(4)`$, we conclude that the algebras $`A_{n1}(n>1)`$ have no realizations by operators of the form (3.3), which generate symmetry algebras of (0.1), except for those given in Theorems 3.1 and 3.2.
There are also no realizations of the desired form for simple Lie algebras of the type $`D_n(n>1)`$, since the lowest dimensional algebras of this type ($`so(4),so(2,2)`$$`so^{}(4)`$) have no realizations within the class (3.3) which could be symmetry algebras of (0.1).
By the same reasoning, we conclude that the realizations (3.13), (3.19)–(3.23) exhaust the set of all possible realizations of the simple Lie algebras $`B_n(n>1)`$ and $`C_n(n1)`$. Indeed, taking the least possible value of $`n`$ and putting $`n=2`$ we see that the algebras of the type $`B_n`$ contain subalgebras that are isomorphic to $`so(4),so(1,3)`$. The same assertion for the simple Lie algebras of the type $`C_n(n1)`$ follows from the relations:
$$sp(2,𝐑)so(3,2),sp(1,1)so(4,1),sp(2)so(5),$$
if we take into account that the algebras $`so(3,2)`$$`so(4,1)`$, contain $`so(3,1)`$, and that the algebra $`so(5)`$ contains $`so(4)`$.
To complete the proof we have to consider the exceptional simple Lie algebras $`G_2`$$`F_4`$$`E_6`$$`E_7`$$`E_8`$.
We consider in detail the first two algebras, the remaining algebras being treated in the same way.
A Lie algebra of the type $`G_2`$ contains a compact real form $`g_2`$ and a non-compact real form $`g_2^{}`$. We also have $`g_2g_2^{}su(2)su(2)so(4)`$, from which we conclude that $`G_2`$ has no realizations by operators of the form (3.3) which are symmetry operators of (0.1).
A Lie algebra of the type $`F_4`$ contains a compact real form $`f_4`$ and two non-compact real forms $`f_4^{},f_4^{\prime \prime }`$. We also have $`f_4^{}f_4sp(3)su(2),f_4^{\prime \prime }f_4so(9)`$. Hence, it follows that the algebra $`F_4`$ has no realizations within the class of operators of the form (3.3) which are admitted by PDE of the form (0.1).
The theorem is proved.
### III.3 Equations invariant under semi-direct sums of simple and solvable Lie algebras
In order to describe equations of the form (0.1) which are invariant with respect to the Lie algebras that are semi-direct sums of simple and solvable Lie algebras, we could follow the same strategy as in the previous section. However, with Theorems 3.1–3.3 in hand, the most effective way is a direct application of the Lie infinitesimal method in order to specify the arbitrary functions of one variable, given in Theorems 3.1, 3.2, with the aim of obtaining all the possible extensions of the algebras $`so(3)`$$`sl(2,𝐑)`$ admitted by PDEs (0.1). In this way we will get all the possible equations of the form (0.1) admitting Lie algebras which are semi-direct sums of simple and solvable Lie algebras.
So we insert the corresponding forms of the functions $`F`$$`G`$ into invariance conditions (3.4) and then investigate the consistency of the system of determining equations which are obtained in this way. Substituting formulas (3.14) into the first equation of (3.4) yields the following system of PDEs:
$`(a)`$ $`2b_x\dot{a}2c\mathrm{tan}u=0,`$
$`(b)`$ $`b_u+c_x\mathrm{sec}^2u=0,`$
$`(c)`$ $`2c_u\dot{a}=0.`$
It follows from $`(c)`$ that $`c={\displaystyle \frac{1}{2}}\dot{a}u+\stackrel{~}{c}(t,x)`$. Then the compatibility requirement for equations $`(a)`$ and $`(b)`$ gives $`\dot{a}=0,\stackrel{~}{c}_{xx}+\stackrel{~}{c}=0`$, whence
$$\dot{a}=0,b=[f(t)\mathrm{sin}xg(t)\mathrm{cos}x]\mathrm{tan}u+h(t),c=f(t)\mathrm{cos}x+g(t)\mathrm{sin}x,$$
where $`f,g,h`$ are arbitrary smooth functions of $`t`$.
Next, substituting the expressions obtained for $`F`$ and $`G`$ into the second equation from (3.14) we see that $`a\dot{\stackrel{~}{G}}=0,\dot{f}=\dot{g}=\dot{h}=0`$. Hence it follows that extension of the symmetry algebra is only possible if $`\stackrel{~}{G}=\lambda ,\lambda =const`$. In this case, the maximal symmetry algebra of the corresponding PDE is the four-dimensional Lie algebra $`so(3)L_1`$, where $`so(3)`$ is given in (3.13), and $`L_1=_t`$.
We get similar results for PDEs invariant under the realizations (3.20), (3.22) and (3.23). Namely, extension of the symmetry algebra is only possible when $`\stackrel{~}{G}=\lambda ,\lambda =\mathrm{const}`$. Moreover, the maximal invariance algebras are the four-dimensional Lie algebras of the form $`sl(2,𝐑)L_1`$, where $`L_1=_u`$ when $`sl(2,𝐑)`$ given by (3.20), and $`L_1=_t`$ when $`sl(2,𝐑)`$ is given by (3.22) or (3.23).
We turn now to the remaining case of PDE (0.1) invariant with respect to realization (3.19) of the algebra $`sl(2,𝐑)`$. Inserting the corresponding expressions for $`F`$ and $`G`$ into (3.4), we find the following equations:
$`(A2B\omega )\stackrel{~}{F}=(C+D\omega +B\omega ^2)\dot{\stackrel{~}{F}},`$ (3.29)
$`(E+B\omega )\stackrel{~}{G}(C+D\omega +B\omega ^2)\dot{\stackrel{~}{G}}=`$
$`=K+L\omega +(M+N\omega +P\omega ^2+S\omega ^3)\stackrel{~}{F}2u(C+D\omega +B\omega ^2)\dot{\stackrel{~}{F}},`$ (3.30)
where
$`A`$ $`=`$ $`2xb_xx\dot{a}+4ub_u,B=b_u,`$
$`C`$ $`=`$ $`2xcx^2c_x2u(b+xc_uxb_x)+4u^2b_u,`$
$`D`$ $`=`$ $`b+xc_uxb_x4ub_u,`$
$`K`$ $`=`$ $`x^3c_t+2x^2ub_t+x^2uc_x+xu^2(c_u2b_x+\dot{a})2u^3b_u,`$
$`L`$ $`=`$ $`x^2b_txc+ub+ux(b_x\dot{a})+u^2b_u,`$
$`E`$ $`=`$ $`2b+x(c_u\dot{a})2ub_u,`$
$`M`$ $`=`$ $`2uE2xc+x^3c_{xx}2x^2u(b_{xx}2c_{xu})4xu^2(2b_{xu}c_{uu})8u^3b_{uu},`$
$`N`$ $`=`$ $`2ub_u+x^2(b_{xx}2c_{xu})+4xu(2b_{xu}c_{uu})+12u^2b_{uu},`$
$`P`$ $`=`$ $`x(2b_{xu}c_{uu})6ub_{uu},S=b_{uu},`$
$`\dot{\stackrel{~}{F}}`$ $`=`$ $`{\displaystyle \frac{d\stackrel{~}{F}}{d\omega }},\dot{\stackrel{~}{G}}={\displaystyle \frac{d\stackrel{~}{G}}{d\omega }},\omega =2uxu_x.`$
If $`\stackrel{~}{F}`$ is an arbitrary function of $`\omega `$, then we have $`A=B=C=D=0`$İt follows that $`b={\displaystyle \frac{1}{2}}x\dot{a}`$$`c=x^2\stackrel{~}{c}(t)`$. Equation (3.30) now takes the form
$$K+L\omega =0,$$
where
$$K=x^5\dot{\stackrel{~}{c}}+x^3u(\ddot{a}+2\stackrel{~}{c}),L=\frac{1}{2}x^3(\ddot{a}+2\stackrel{~}{c}),$$
so that, $`\dot{\stackrel{~}{c}}=0`$$`\ddot{a}+2\stackrel{~}{c}=0`$. Thus realization (3.19) of the algebra $`sl(2,𝐑)`$ is the maximal invariance algebra of the corresponding PDE (0.1).
Thus, extension of realization (3.19) is only possible when not all of the coefficients $`A,B,C,D`$ in (3.29) vanish as a result of (3.29), (3.30). In order to classify all these cases we note that (3.29) is equivalent to the following relation:
$`(k2m\omega )\stackrel{~}{F}=(n+p\omega +m\omega ^2)\dot{\stackrel{~}{F}},`$ (3.31)
where the coefficients $`k,m,n,p`$ are constant. Indeed, since $`\stackrel{~}{F}`$ is a function of $`\omega `$ only, relation (3.29) can be valid if and only if all its coefficients have the form $`\mathrm{const}.\times R(t,x,u)`$ with some non-vanishing function $`R`$. If all the coefficients in (3.29) are equal to zero, then we get the case of an arbitrary function $`\stackrel{~}{F}`$. Consequently, extension of the invariance algebra is only possible when the function $`\stackrel{~}{F}(\omega )`$ satisfies an equation of the form (3.31), where $`k,m,n,p`$, are constants not vanishing simultaneously.
Summing up, we conclude that the problem of the group classification of PDEs (0.1) invariant under realization (3.19) of the algebra $`sl(2,𝐑)`$, reduces to classifying all admissible forms of the function $`\stackrel{~}{F}`$. Solving this problem requires simple but very tedious computations and so we give only the result, omitting the intermediate calculations. The admissible forms of the functions $`\stackrel{~}{F}(\omega )`$ are:
$`\stackrel{~}{F}`$ $`=`$ $`1;`$
$`\stackrel{~}{F}`$ $`=`$ $`\lambda \omega ^\alpha ;`$
$`\stackrel{~}{F}`$ $`=`$ $`\lambda \mathrm{exp}\omega ;`$
$`\stackrel{~}{F}`$ $`=`$ $`\lambda (\omega ^2+\alpha )^1;`$ (3.32)
$`\stackrel{~}{F}`$ $`=`$ $`{\displaystyle \frac{\lambda }{(\omega +\alpha )^2}}\mathrm{exp}\left({\displaystyle \frac{2\alpha }{\omega +\alpha }}\right);`$
$`\stackrel{~}{F}`$ $`=`$ $`{\displaystyle \frac{\lambda }{(\omega +\alpha )^2+\beta ^2}}\mathrm{exp}\left({\displaystyle \frac{2\alpha }{\beta }}\mathrm{arctan}{\displaystyle \frac{\omega +\alpha }{\beta }}\right);`$
$`\stackrel{~}{F}`$ $`=`$ $`{\displaystyle \frac{\lambda }{(\omega +\alpha )^2\beta ^2}}\mathrm{exp}\left|{\displaystyle \frac{\omega +\alpha \beta }{\omega +\alpha +\beta }}\right|^{\frac{\alpha }{\beta }},`$
where $`\{\alpha ,\lambda ,\beta \}𝐑,\alpha \lambda \beta 0`$.
On analyzing the above cases, we conclude that the only forms of the function $`\stackrel{~}{F}`$ from the list (3.32) that provide an extension of invariance algebra of the equation under study are:
$$\stackrel{~}{F}=1,\stackrel{~}{F}=\omega .$$
We finally find that there exist five nonlinear equations of the form (0.1) invariant under four-dimensional algebras and two nonlinear PDEs admitting five-dimensional Lie algebras. Below we the give these equations together with their maximal symmetry algebras $`L_{\mathrm{max}}`$.
$$u_t=\frac{\mathrm{sec}^2u}{1+u_x^2\mathrm{sec}^2u}u_{xx}+\frac{1+2u_x^2\mathrm{sec}^2u}{1+u_x^2\mathrm{sec}^2u}\mathrm{tan}u+\lambda \sqrt{1+u_x^2\mathrm{sec}^2u},\lambda 𝐑,$$
$`L_{\mathrm{max}}=so(3)L_1`$$`so(3)`$ has the form (3.13), $`L_1=_t;`$
$$u_t=x^3u_x^3u_{xx}\frac{1}{4}x^1u_x+3x^4u_x^2+\lambda x^3u_x^1,\lambda 𝐑,$$
$`L_{\mathrm{max}}=sl(2,𝐑)L_1,`$ $`sl(2,𝐑)`$ has the form (3.20), $`L_1=_u;`$
$$u_t=\frac{u^2}{u^6+4u_x^2}u_{xx}\frac{10uu_x^2+u^7}{8u_x^2+2u^6}+\lambda u^2\sqrt{u^6+4u_x^2},\lambda 𝐑,$$
$`L_{\mathrm{max}}=sl(2,𝐑)L_1,`$ $`sl(2,𝐑)`$ has the form (3.22), $`L_1=_t;`$
$$u_t=\frac{u^2}{u^64u_x^2}u_{xx}+\frac{10uu_x^2u^7}{8u_x^22u^6}+\lambda u^2\sqrt{|u^64u_x^2|},\lambda 𝐑,$$
$`L_{\mathrm{max}}=sl(2,𝐑)L_1,`$ $`sl(2,𝐑)`$ has the form (3.23), $`L_1=_t;`$
$$u_t=\lambda [2uxu_x]u_{xx}+[4\gamma 4\lambda 1]x^2u^2+[1+2\lambda 4\gamma ]x^1uu_x+\gamma u_x^2,\lambda 0,\gamma 𝐑,$$
$`L_{\mathrm{max}}=sl(2,𝐑)L_1,`$ $`sl(2,𝐑)`$ has the form (3.19), $`L_1=x_x+2u_u;`$
$$u_t=u^4u_{xx}2u^5u_x^2,$$
$`L_{\mathrm{max}}=sl(2,𝐑)L_{2.1},`$ $`sl(2,𝐑)`$ has the form (3.21), $`L_{2.1}=4t_t+u_u,_t;`$
$$u_t=u_{xx}+x^1uu_xx^2u^22x^2u,$$
$`L_{\mathrm{max}}=sl(2,𝐑)+L_{2.2},`$ $`sl(2,𝐑)`$ has the form (3.19), $`L_{2.2}=t_x+[tx^1(u+2)x]_u,_x+x^1(u+2)_u.`$
The above formulas provide the full solution of the problem of describing all PDEs of the type (0.1) admitting symmetry Lie algebras which are semi-direct sums of semi-simple and solvable Lie algebras.
### III.4 Classification of equations (0.1) invariant with respect to solvable Lie algebras
To complete the classification of invariant PDEs of the form (0.1) we have to construct all possible inequivalent realizations of solvable Lie algebras within the class of operators (3.3) which are invariance algebras of (0.1). First, we shall perform a preliminary classification: we shall describe inequivalent PDEs (0.1) admitting one-, two- and three-dimensional solvable invariance algebras and then proceed to classifying equations invariant with respect to higher dimensional solvable Lie algebras.
#### III.4.1 Preliminary classification
Equations (0.1) invariant with respect to one-dimensional algebras have already been constructed, so that we can start by considering two-dimensional solvable Lie algebras. As mentioned in Section II.2 there are two inequivalent solvable Lie algebras
$`A_{2.1}`$ $`:`$ $`[e_1,e_2]=0;`$ (3.33)
$`A_{2.2}`$ $`:`$ $`[e_1,e_2]=e_2.`$
As each of the above algebras contains the algebra $`A_1`$, when studying realizations of two-dimensional Lie algebras we can take as one of the basis operators either $`\frac{}{t}`$ or $`\frac{}{x}`$. Consider in more detail the case of the algebra $`A_{2.1}`$.
Let $`e_1=_t`$ and $`e_2`$ be an operator of the form (3.3). Then it follows from (3.33) that within a choice of a basis of the algebra $`A_{2.1}`$ we can put
$$e_2=b(x,u)_x+c(x,u)_u.$$
(3.34)
Since operator (3.34) can be treated as the non-zero vector field acting on smooth functions of $`x,u`$, we can choose $`e_2=_u`$ thus getting the realization $`_t,_u`$.
Now take the case $`e_1=_x`$ and $`e_2`$ is an operator of the form (3.3). Using the commutation relation (3.33) yields
$$e_2=a(t)_t+b(t,u)_x+c(t,u)_u.$$
(3.35)
If $`a0`$, we make the change of variables
$$\overline{t}=T(t),\overline{x}=x+X(t,u),v=U(t,u),\dot{T}0,U_u0,$$
(3.36)
where $`\dot{T}=a^1`$$`aX_t+xX_u+b=0;aU_t+cU_u=0,U_u0.`$ This reduces operator (3.35) to the form $`e_2=_{\overline{t}}`$.
If in (3.35$`a=0,c0,`$ then we put $`T=t`$ in (3.36) and take as $`X`$ and $`U`$ solutions of the equations
$$cX_u+b=0,cU_u=1,$$
which reduces operator (3.35) to $`e_2=_v`$.
Finally, turning to the remaining case, when $`a=c=0`$ in (3.35). There are transformations of the form (3.36) which transform operator (3.35) to the form $`e_2=\overline{t}_{\overline{x}}`$ (if $`b_u=0`$) or to $`e_2=v_{\overline{x}}`$ (if $`b_u0`$).
Summing up we conclude that, up to equivalence defined by transformations of the group $``$, there are four inequivalent realizations of the algebra $`A_{2.1}:_t,_u`$$`_x,_u`$$`_x,t_x`$$`_x,u_x.`$
The conditions (3.4) imply that the third realization, $`_x,t_x`$, cannot be an invariance algebra of PDEs of the form (0.1). The equation invariant under the fourth realization is
$$u_t=u_x^2F(t,x)u_{xx}+u_xG(t,u).$$
(3.37)
and is linearizable by the change of variables
$$\overline{t}=t,\overline{x}=u,v=x.$$
Thus we see that there exist two inequivalent realizations of the algebra $`A_{2.1}`$ which are invariance algebras of nonlinear PDEs of the form (0.1):
$`A_{2.1}^1`$ $`=`$ $`_t,_u;`$
$`A_{2.1}^2`$ $`=`$ $`_x,_u.`$
The corresponding forms of the functions $`F`$, $`G`$ are given in Table 1.
The same reasoning gives all inequivalent realizations of the abstract Lie algebra $`A_{2.2}`$. The full list of these contains three realizations which are admitted by PDEs of the form (0.1):
$`A_{2.2}^1`$ $`=`$ $`t_tx_x,_t;`$
$`A_{2.2}^2`$ $`=`$ $`t_tx_x,_x;`$
$`A_{2.2}^3`$ $`=`$ $`x_xu_u,_x.`$
The corresponding forms of the functions $`F`$ and $`G`$ are given in Table 1.
Let us note that provided the functions $`\stackrel{~}{F},\stackrel{~}{G}`$ are arbitrary, the corresponding realizations of the two-dimensional Lie algebras are maximal invariance algebras of these equations.
Table 1. Invariance of (0.1) under two-dimensional solvable Lie algebras
| Algebra | $`F`$ | $`G`$ |
| --- | --- | --- |
| $`A_{2.1}^1`$ | $`\stackrel{~}{F}(x,u_x)`$ | $`\stackrel{~}{G}(x,u_x)`$ |
| $`A_{2.1}^2`$ | $`\stackrel{~}{F}(t,u_x)`$ | $`\stackrel{~}{G}(t,u_x)`$ |
| $`A_{2.2}^1`$ | $`x\stackrel{~}{F}(u,\omega )`$ | $`x^1\stackrel{~}{G}(u,\omega ),\omega =xu_x`$ |
| $`A_{2.2}^2`$ | $`t\stackrel{~}{F}(u,\omega )`$ | $`t^1\stackrel{~}{G}(u,\omega ),\omega =tu_x`$ |
| $`A_{2.2}^3`$ | $`u^2\stackrel{~}{F}(t,u_x)`$ | $`u\stackrel{~}{G}(t,u_x)`$ |
We begin the search for realizations of three-dimensional solvable Lie algebras $`A_3=e_1`$, $`e_2,e_3`$ by considering decomposable algebras $`A_{3.1},A_{3.2}`$.
Evidently, in order to get all the possible realizations of these algebras within the class of operators (3.3), it suffices to extend the realizations already known for the two-dimensional algebras $`A_{2.1}^i=e_1,e_2(i=1,2)`$ (for the algebra $`A_{3.1}`$) and $`A_{2.2}^i=e_1,e_2(i=1,2,3)`$ (for the algebra $`A_{3.2}`$). This follows from the definition of decomposable solvable Lie algebras. As a result we get one realization of the algebra $`A_{3.1}`$ and six inequivalent realizations of the algebra $`A_{3.2}`$, which are admissible as invariance algebras of PDEs (0.1):
$`A_{3.1}^1`$ $`=`$ $`_t,_u,_x;`$
$`A_{3.2}^1`$ $`=`$ $`t_tx_x,_t,_u;`$
$`A_{3.2}^2`$ $`=`$ $`t_tu_u,_t,xu_u;`$
$`A_{3.2}^3`$ $`=`$ $`t_tu_u,_u,t_t+x_x;`$
$`A_{3.2}^4`$ $`=`$ $`t_tx_x,_x,_u;`$
$`A_{3.2}^5`$ $`=`$ $`x_xu_u,_u,_t;`$
$`A_{3.2}^6`$ $`=`$ $`x_xu_u,_u,tx_x.`$
The explicit forms of the invariant equations (0.1) are determined by the forms of the functions $`F`$, $`G`$ which are given in Table 2, where $`\stackrel{~}{F}`$, $`\stackrel{~}{G}`$ are arbitrary smooth functions.
One can verify by direct computation that the realizations given above are the maximal invariance algebras of the corresponding equations, provided the functions $`\stackrel{~}{F}`$ and $`\stackrel{~}{G}`$ are arbitrary smooth functions.
As mentioned in Section II.2, there are seven abstract non-isomorphic non-decomposable Lie algebras. All of them contain the two-dimensional commutative ideal $`A_{2.1}=e_1,e_2`$. Thus, to construct their realizations within the class of operators under consideration, it suffices to describe all the possible extensions of the realizations of $`A_{2.1}`$ with the operator $`e_3`$ of the form (3.3). Moreover, we have to consider both the realizations $`A_{2.1}^i=e_1,e_2(i=1,2),`$ and $`\stackrel{~}{A}_{2.1}^i=\stackrel{~}{e}_1,\stackrel{~}{e}_2(i=1,2),`$ where $`\stackrel{~}{e}_1=e_2,\stackrel{~}{e}_2=e_1.`$ We will consider in more detail the procedure for constructing realizations of the Weyl algebra $`A_{3.3}`$, which is a nilpotent Lie algebra.
Table 2. Invariance of (0.1) under three-dimensional decomposable solvable Lie algebras
| Algebra | $`F`$ | $`G`$ |
| --- | --- | --- |
| $`A_{3.1}^1`$ | $`\stackrel{~}{F}(u_x)`$ | $`\stackrel{~}{G}(u_x)`$ |
| $`A_{3.2}^1`$ | $`x\stackrel{~}{F}(\omega )`$ | $`x^1\stackrel{~}{G}(\omega ),\omega =xu_x`$ |
| $`A_{3.2}^2`$ | $`u^1e^{x\omega }\stackrel{~}{F}(x)`$ | $`e^{x\omega }[\stackrel{~}{G}(x)\omega ^2\stackrel{~}{F}(x)],\omega =u^1u_x`$ |
| $`A_{3.2}^3`$ | $`t^1x^2\stackrel{~}{F}(\omega )`$ | $`x^1\stackrel{~}{G}(\omega ),\omega =t^1x^2u_x`$ |
| $`A_{3.2}^4`$ | $`t\stackrel{~}{F}(\omega )`$ | $`t^1\stackrel{~}{G}(\omega ),\omega =tu_x`$ |
| $`A_{3.2}^5`$ | $`x^2\stackrel{~}{F}(u_x)`$ | $`x\stackrel{~}{G}(u_x)`$ |
| $`A_{3.2}^6`$ | $`x^2\stackrel{~}{F}(t)`$ | $`xt^1u_x\mathrm{ln}|u_x|+xu_x\stackrel{~}{G}(t)`$ |
We begin with the realization $`A_{2.1}^1`$. If $`e_1=_t,e_2=_u`$, then it follows from the commutation relation $`[e_2,e_3]=e_1,`$ where $`e_3`$ is of the form (3.3), that the equation $`b_u_x+c_u_u=_t`$ holds true. Since this equation cannot be satisfied for any choice of the functions $`a,b,c`$ contained in $`e_3`$, this realization cannot be extended to that of a three-dimensional solvable Lie algebra. Next, if $`e_1=_u,e_2=_t`$, then $`e_3=\stackrel{~}{b}(x)_x+[t+\stackrel{~}{c}(x)]_u`$, and we obtain, up to equivalence under $``$, the following three realizations of the algebra $`A_{3.3}`$:
$`_u,_t,_x+t_u,`$
$`_u,_t,t_u,`$
$`_u,_t,(t+x)_u.`$
Checking conditions (3.4) we find that the second realization from the above list cannot be an invariance algebra of PDE (0.1). Moreover, the equation admitting the third realization is necessarily linear.
In studying the realization $`A_{2.1}^2`$ we have to take into account the existence of two possibilities. The first possibility is $`e_1=_x,e_2=_u`$ and the second one is $`e_1=_u,e_2=_x`$. However, the above realizations are transformed one into another by the change of variables
$$\overline{t}=t,\overline{x}=u,v=x.$$
(3.38)
So we may consider without loss of generality the second realization only. Performing the necessary computations yields a realization of the algebra $`A_{3.3}`$ which can be admitted by a nonlinear equation of the form (0.1): $`_u,_x,_t+x_u.`$
We conclude that there are two inequivalent realizations of the algebra $`A_{3.3}`$, which are invariance algebras of nonlinear PDEs from the class (0.1),
$`A_{3.3}^1`$ $`=`$ $`_u,_t,t_u+_x;`$
$`A_{3.3}^2`$ $`=`$ $`_u,_x,t_x+x_u.`$
The forms of the functions $`F`$ and $`G`$ defining the corresponding nonlinear heat conductivity equations are given in Table 3.
Table 3. Invariance of (0.1) with respect to the Weyl algebra
| Algebra | $`F`$ | $`G`$ |
| --- | --- | --- |
| $`A_{3.3}^1`$ | $`\stackrel{~}{F}(u_x)`$ | $`x+\stackrel{~}{G}(u_x)`$ |
| $`A_{3.3}^2`$ | $`\stackrel{~}{F}(t)`$ | $`\frac{1}{2}u_x^2+\stackrel{~}{G}(t)`$ |
The remaining non-decomposable solvable Lie algebras are treated in an analogous way. We present below those of their inequivalent realizations which are admitted by nonlinear PDEs of the form (0.1). In Table 4 we give the various forms of the functions $`F`$, $`G`$ defining the forms of the invariant equations.
$`A_{3.4}^1`$ $`=`$ $`_u,_t,t_t+x_x+[t+u]_u;`$
$`A_{3.4}^2`$ $`=`$ $`_u,_t,t_t+[t+u]_u;`$
$`A_{3.4}^3`$ $`=`$ $`_x,_u,2t_t+(x+u)_x+u_u;`$
$`A_{3.4}^4`$ $`=`$ $`_x,_u,(x+u)_x+u_u;`$
$`A_{3.5}^1`$ $`=`$ $`_t,_u,t_t+x_x+u_u;`$
$`A_{3.5}^2`$ $`=`$ $`_t,_u,t_t+u_u;`$
$`A_{3.5}^3`$ $`=`$ $`_x,_u,2t_t+x_x+u_u;`$
$`A_{3.6}^1`$ $`=`$ $`_t,_u,t_t+x_xu_u;`$
$`A_{3.6}^2`$ $`=`$ $`_t,_u,t_tu_u;`$
$`A_{3.6}^3`$ $`=`$ $`_x,_u,t_t+x_xu_u;`$
$`A_{3.6}^4`$ $`=`$ $`_x,_u,x_xu_u;`$
$`A_{3.7}^1`$ $`=`$ $`_u,_t,qt_t+x_x+u_u(q0,\pm 1);`$
$`A_{3.7}^2`$ $`=`$ $`_u,_t,qt_t+u_u(q0,\pm 1);`$
$`A_{3.7}^3`$ $`=`$ $`_x,_u,t_t+x_x+qu_u(0<|q|<1);`$
$`A_{3.7}^4`$ $`=`$ $`_x,_u,x_x+qu_u(0<|q|<1);`$
$`A_{3.8}^1`$ $`=`$ $`_x,_u,_t+u_xx_u;`$
$`A_{3.8}^2`$ $`=`$ $`_x,_u,u_xx_u;`$
$`A_{3.9}^1`$ $`=`$ $`_x,_u,_t+(u+qx)_x+(qux)_u(q>0);`$
$`A_{3.9}^2`$ $`=`$ $`_x,_u,(u+qx)_x+(qux)_u(q>0).`$
Let us note that if the functions $`\stackrel{~}{F}`$, $`\stackrel{~}{G}`$ from Tables 3, 4 are arbitrary, then the corresponding realizations are the maximal invariance algebras of the equations obtained.
#### III.4.2 Complete classification of nonlinear PDEs (0.1) invariant with respect to solvable Lie algebras
The next step of our approach to the group classification of nonlinear PDEs of the form (0.1) is to describe equations which are invariant under four-dimensional solvable Lie algebras.
Table 4. Invariance of (0.1) under non-decomposable three-dimensional solvable Lie algebras
| Algebra | $`F`$ | $`G`$ |
| --- | --- | --- |
| $`A_{3.4}^1`$ | $`x\stackrel{~}{F}(u_x)`$ | $`\stackrel{~}{G}(u_x)+\mathrm{ln}|x|`$ |
| $`A_{3.4}^2`$ | $`u_x^1\stackrel{~}{F}(x)`$ | $`\stackrel{~}{G}(x)+\mathrm{ln}|u_x|`$ |
| $`A_{3.4}^3`$ | $`u_x^2\stackrel{~}{F}(\omega )`$ | $`u_xe^{\frac{1}{u_x}}\stackrel{~}{G}(\omega ),\omega =2u_x^1\mathrm{ln}|t|`$ |
| $`A_{3.4}^4`$ | $`u_x^2\stackrel{~}{F}(t)\mathrm{exp}(2u_x^1)`$ | $`u_x\stackrel{~}{G}(t)\mathrm{exp}(u_x^1)`$ |
| $`A_{3.5}^1`$ | $`x\stackrel{~}{F}(u_x)`$ | $`\stackrel{~}{G}(u_x)`$ |
| $`A_{3.5}^2`$ | $`u_x^1\stackrel{~}{F}(x)`$ | $`\stackrel{~}{G}(x)`$ |
| $`A_{3.5}^3`$ | $`\stackrel{~}{F}(u_x)`$ | $`|t|^{\frac{1}{2}}\stackrel{~}{G}(u_x)`$ |
| $`A_{3.6}^1`$ | $`x\stackrel{~}{F}(\omega )`$ | $`x^2\stackrel{~}{G}(\omega ),\omega =x^2u_x`$ |
| $`A_{3.6}^2`$ | $`u_x\stackrel{~}{F}(x)`$ | $`u_x^2\stackrel{~}{G}(x)`$ |
| $`A_{3.6}^3`$ | $`t\stackrel{~}{F}(\omega )`$ | $`t^2\stackrel{~}{G}(\omega ),\omega =t^2u_x`$ |
| $`A_{3.6}^4`$ | $`u_x^1\stackrel{~}{F}(t)`$ | $`\sqrt{|u_x|}\stackrel{~}{G}(t)`$ |
| $`A_{3.7}^1`$ | $`|x|^{2q}\stackrel{~}{F}(u_x)`$ | $`|x|^{1q}\stackrel{~}{G}(u_x)(q=0,\pm 1)`$ |
| $`A_{3.7}^2`$ | $`|u_x|^q\stackrel{~}{F}(x)`$ | $`|u_x|^{1q}\stackrel{~}{G}(x),q0,\pm 1`$ |
| $`A_{3.7}^3`$ | $`t\stackrel{~}{F}(\omega )`$ | $`|t|^{q1}\stackrel{~}{G}(\omega ),\omega =|t|^{1q}u_x(0<|q|<1)`$ |
| $`A_{3.7}^4`$ | $`|u_x|^{\frac{2}{q1}}\stackrel{~}{F}(t)`$ | $`|u_x|^{\frac{q}{q1}}\stackrel{~}{G}(t)(0<|q|<1)`$ |
| $`A_{3.8}^1`$ | $`(1+u_x^2)^1\stackrel{~}{F}(\omega )`$ | $`\sqrt{1+u_x^2}\stackrel{~}{G}(\omega ),\omega =t+\mathrm{arctan}u_x`$ |
| $`A_{3.8}^2`$ | $`(1+u_x^2)^1\stackrel{~}{F}(t)`$ | $`\sqrt{1+u_x^2}\stackrel{~}{G}(t)`$ |
| $`A_{3.9}^1`$ | $`\frac{\mathrm{exp}\left(2q\mathrm{arctan}u_x\right)\stackrel{~}{F}\left(\omega \right)}{1+u_x^2}`$ | $`\sqrt{1+u_x^2}\mathrm{exp}(q\mathrm{arctan}u_x)\stackrel{~}{G}(\omega ),`$ |
| | | $`\omega =t+\mathrm{arctan}u_x(q>0)`$ |
| $`A_{3.9}^2`$ | $`\frac{\mathrm{exp}\left(2q\mathrm{arctan}u_x\right)\stackrel{~}{F}\left(t\right)}{1+u_x^2}`$ | $`\sqrt{1+u_x^2}\mathrm{exp}(q\mathrm{arctan}u_x)\stackrel{~}{G}(t),(q>0)`$ |
As we mentioned in Section II.2, there are ten decomposable and ten non-decomposable, non-isomorphic, solvable four-dimensional Lie algebras. Since nonlinear equations of the form (0.1) which admit three-dimensional solvable algebras contain arbitrary functions of one argument, it is only natural to expect that PDEs admitting four-dimensional algebras will depend on arbitrary parameters at most. In other words, the arbitrary functions in question will take specific forms dictated by the extension of symmetry group. This is, indeed, the case for all the invariant PDEs except for the equation
$$u_t=F(u_x)u_{xx}.$$
(3.39)
Group classification of PDEs of the form (3.39) has been carried out in and we give below the results obtained in the form of theorem.
###### Theorem 3.4
() Provided $`F`$ is an arbitrary smooth function, the maximal invariance algebra admitted by (3.39) is the four-dimensional Lie algebra
$$A_{3.1}^1+2t_t+x_x+u_u.$$
An extension of the symmetry algebra of PDE (3.39) is only possible for the three cases given below:
$`F=\mathrm{exp}u_x`$ $`:`$ $`e_5=t_tx_u;`$
$`F=u_x^n`$ $`:`$ $`e_5=nt_tu_u,n1,n0;`$
$`F={\displaystyle \frac{\mathrm{exp}(n\mathrm{arctan}u_x)}{1+u_x^2}}`$ $`:`$ $`e_5=nt_tu_xx_u,n0.`$
In view of this result, we will exclude from further consideration equations which are equivalent to an equation of the form (3.39).
Consider first the decomposable solvable four-dimensional Lie algebras
$`4A_1=A_{3.1}A_1,A_{3.2}A_1,`$
$`2A_{2.2}=A_{2.2}A_{2.2},A_{3.i}A_1(i=3,4,\mathrm{},9).`$
On analyzing extensions of the realization $`A_{3.1}^1`$ for the algebra $`4A_1`$ and of the realizations $`A_{3.2}^i(i=1,\mathrm{},6)`$ by an operator $`e_4`$ of the form (3.3), we conclude that there are no realizations of the algebras $`4A_1`$ and $`A_{3.2}A_1`$ which could be invariance algebras of PDEs of the form (0.1).
Studying realizations of the algebra $`2A_{2.2}`$ yields four inequivalent realizations admitted by PDEs from the class (0.1). We give these realizations below, as well as the corresponding forms of the functions $`F`$$`G`$.
$`2A_{2.2}^1`$ $`=`$ $`A_{3.2}^1+u_u+kx_x(k0):`$
$`F=\lambda x|\omega |^k,G=\beta x^1|\omega |^{1k},\lambda 0,\beta R,\omega =xu_x;`$
$`2A_{2.2}^2`$ $`=`$ $`A_{3.2}^2+x_x:`$
$`F=\lambda x^2u^1\mathrm{exp}\omega ,G=(\beta \lambda \omega ^2)\mathrm{exp}\omega ,\lambda 0,\beta R,\omega =xu^1u_x;`$
$`2A_{2.2}^3`$ $`=`$ $`A_{3.2}^4+u_u+kt_t(k0,1):`$
$`F=\lambda t|\omega |^{\frac{2k}{1k}},G=\beta t^1|\omega |^{\frac{1}{1k}},\omega =tu_x,\lambda 0,\beta R;`$
$`2A_{2.2}^4`$ $`=`$ $`A_{3.2}^4+u_u+t_x:`$
$`F=\lambda t,G=u_x\mathrm{ln}|tu_x|+\beta u_x,\lambda 0,\beta R.`$
Now, in order to complete group classification of the PDEs given above, one has to compute their maximal invariance algebras. To this end, it is necessary to solve the determining equations (3.4) for each choice of the functions $`F`$, $`G`$. Note, that we have simplified the forms of the functions $`F,G`$ with the use of transformations from the corresponding equivalence groups which are subgroups of $``$. As a result, we get the following simplified forms of the above invariant equations:
$`2A_{2.2}^1:u_t=|x|^{1k}|u_x|^ku_{xx}+\beta |x|^k|u_x|^{1k},\beta R,k0;`$ (3.40)
$`2A_{2.2}^2:u_t=x^2u^1\mathrm{exp}(\omega )u_{xx}+(\beta \omega ^2)\mathrm{exp}\omega ,\omega =xu^1u_x,\beta R;`$ (3.41)
$`2A_{2.2}^3:u_t=\pm |t|^{\frac{k+1}{1k}}|u_x|^{\frac{2k}{1k}}u_{xx}+ϵ|t|^{\frac{k}{1k}}|u_x|^{\frac{1}{1k}},ϵ=0,1,k0,1;`$ (3.42)
$`2A_{2.2}^4:u_t=\lambda tu_{xx}+u_x\mathrm{ln}|tu_x|,\lambda 0.`$ (3.43)
Inserting the functions $`F`$, $`G`$ defining the above PDEs (3.40)–(3.43) into the determining equations (3.4) and analyzing the equations obtained, we arrive at the following conclusions.
1. If $`k0,2,\beta \frac{k1}{k2}`$ or $`k=2,\beta \frac{5}{4}`$ in (3.40), the realization $`2A_{2.2}^1`$ is the maximal invariance algebra of the nonlinear heat conductivity equation (3.40). If $`k=2,\beta =\frac{5}{4}`$, then the maximal invariance algebra of the equation in question is five-dimensional. Its basis if formed by the operators of the realization $`2A_{2.2}^1(k=2)`$ and the operator $`4xu_xu^2_u`$. However, this algebra is isomorphic to the Lie algebra $`sl(2,R)A_{2.2}`$ and the change of variables
$$\overline{t}=t,\overline{x}=u,v=\alpha |x|^{\frac{1}{4}},\alpha 0$$
transforms its basis operators to become basis operators of the realization of $`sl(2,R)L_{2.1}`$, where $`sl(2,R)`$ is the realization (3.21) and $`L_{2.1}=4t_t+u_u,_t.`$ Thus equation (3.40) with $`k=2,\beta =\frac{5}{4}`$ is equivalent to an invariant PDE obtained in the previous section.
Finally, if $`k0,2`$ and $`\beta =\frac{k1}{k2}`$ in equation (3.40), then the latter is transformed by transformations from the group $``$ to the form (3.39).
2. If $`\beta 2`$ in equation (3.40), the realization $`2A_{2.2}^2`$ is the maximal invariance algebra of this equation. Given the condition $`\beta =2`$, the maximal invariance algebra is the five-dimensional Lie algebra spanned by the operators
$$_t,xu_u,x^2_x+\mathrm{ln}|x^2u|xu_u,2t_t+2u_ux_x,t_t+u_u.$$
The change of variables
$$\overline{t}=t,\overline{x}=x^1,v=x^1\mathrm{ln}|u|+2x^1(1+\mathrm{ln}|x|)$$
reduce equation (3.41) with $`\beta =2`$ to the equation
$$v_{\overline{t}}=\mathrm{exp}(v_{\overline{x}})v_{\overline{x}\overline{x}},$$
which is contained in the class of PDEs (3.39).
3. The realization $`2A_{2.2}^3(k0,1)`$ is the maximal invariance algebra of PDE (3.42), provided $`ϵ=1`$. If $`ϵ=0`$, then its maximal invariance algebra is the five-dimensional Lie algebra spanned by the operators
$$2A_{2.2}^3(k0,1)+|t|^{\frac{1+k}{k1}}_t.$$
However, with this choice of $`\epsilon `$, equation (3.42) is reduced through the change of variables
$$\overline{t}=\frac{1}{2}(1k)|t|^{\frac{2}{1k}},\overline{x}=x,v=u,$$
to equation
$$v_{\overline{t}}=\pm |v_{\overline{x}}|^{\frac{2k}{1k}}v_{\overline{x}\overline{x}},$$
which belongs to the class of PDEs (3.39).
4. The realization $`2A_{2.2}^4`$ is the maximal invariance algebra of PDE (3.43). Analyzing the algebra $`A_{3.3}A_1`$ we find that it has no realizations which are admissible for PDEs (0.1). Next, we get a realization $`A_{3.5}^2_x`$ of the algebra $`A_{3.5}A_1`$ but the corresponding invariant equation
$$u_t=u_x^1u_{xx}$$
belong to the class of PDEs (3.39). Studying realizations of the algebra $`A_{3.7}A_1`$ we get the nonlinear heat conductivity equation
$$u_t=u_x^2u_{xx}+u_x^1,$$
whose maximal invariance algebra is the five-dimensional algebra having the following basis elements:
$$A_5^1=A_{3.7}^2(q=2)_x,e^x_x.$$
A similar analysis of the remaining decomposable solvable four-dimensional algebras yields eight inequivalent realizations that are maximal invariance algebras of nonlinear PDEs of the form (0.1). We list the nonlinear PDEs (0.1) whose maximal invariance algebras are four-dimensional decomposable solvable Lie algebras in Table 5.
We turn now to non-decomposable algebras. There are ten non-isomorphic non-decomposable solvable four-dimensional Lie algebras $`A_4=e_i|i=1,2,3,4`$ (the full list is given in Section II.2). Their structure implies that studying realizations of these algebras can be carried out by extension of the (already known) realizations of three-dimensional solvable algebras $`A_3=e_1,e_2,e_3`$ by the operator $`e_4`$ of the form (3.3). Moreover, one must use the following extension scheme: $`A_{4.i}=A_{3.1}+e_4`$ $`(i=1,\mathrm{},6),`$ $`A_{4.i}=A_{3.3}+e_4`$ $`(i=7,8,9)`$, $`A_{4.10}=A_{3.5}+e_4.`$
There exists only one realization of the algebra $`A_{3.1}`$, and it is the maximal invariance algebra of the PDE
$$u_t=F(u_x)u_{xx}+G(u_x),$$
(3.44)
so that nonlinear PDEs invariant under realizations of the algebras $`A_{4.i}(i=1,\mathrm{},6)`$ must belong to the class of equations (3.44).
Direct computation shows that the algebra $`A_{4.1}`$ has no realizations that are admitted by PDEs (0.1). For the remaining abstract Lie algebras from the class under study we get seven realizations which are invariance algebras of nonlinear PDEs of the form (0.1).
$`A_{4.2}^1=A_{3.1}^1+qt_t+x_x+(u+x)_u(q0,1);`$
$`A_{4.2}^2=A_{3.1}^1+t_t+(t+x)_x+qu_u(q0,1);`$
$`A_{4.3}^1=A_{3.1}^1+t_t+x_u;`$
$`A_{4.3}^2=A_{3.1}^1+t_x+u_u;`$
$`A_{4.4}^1=A_{3.1}^1+t_t+(t+x)_x+(x+u)_u;`$
$`A_{4.5}^1=A_{3.1}^1+t_t+px_x+qu_u(p<q,pq0;p,q,1);`$
$`A_{4.6}^1=A_{3.1}^1+qt_t+(px+u)_x+(pux)_u(q0;p0).`$
The corresponding forms of the functions $`F`$, $`G`$ defining invariant equations (0.1) are given in Table 6. Note that, for the sake of completeness, we give in Table 6 equation (3.39), whose maximal invariance algebra for arbitrary $`F`$ is
$$A_{4.5}^2=A_{3.1}^1+t_t+\frac{1}{2}x_x+\frac{1}{2}u_u.$$
Table 5. Invariance of (0.1) under decomposable four-dimensional solvable algebras
| Algebra | $`F`$ | $`G`$ |
| --- | --- | --- |
| $`2A_{2.2}^1,(k0,2)`$ | $`|x|^{1k}|u_x|^k`$ | $`\beta |x|^k|u_x|^{1k},\beta \frac{k1}{k2}`$ |
| $`2A_{2.2}^1(k=2)`$ | $`x^1u_x^2`$ | $`\beta x^2u_x^1,\beta \frac{5}{4}`$ |
| $`2A_{2.2}^2`$ | $`x^2u^1\mathrm{exp}\omega `$ | $`(\beta \omega ^2)\mathrm{exp}\omega ,\omega =xu^1u_x,\beta 2`$ |
| $`2A_{2.2}^3(k0,1)`$ | $`\pm |t|^{\frac{k+1}{1k}}|u_x|^{\frac{2k}{1k}}`$ | $`|t|^{\frac{k}{1k}}|u_x|^{\frac{1}{1k}}`$ |
| $`2A_{2.2}^4`$ | $`\lambda t,\lambda 0`$ | $`u_x\mathrm{ln}|tu_x|`$ |
| $`A_{3.4}^2_x`$ | $`u_x^1`$ | $`\mathrm{ln}|u_x|`$ |
| $`A_{3.4}^4_t`$ | $`u_x^2\mathrm{exp}(2u_x^1)`$ | $`u_x\mathrm{exp}(u_x^1)`$ |
| $`A_{3.6}^2_x`$ | $`u_x`$ | $`u_x^2`$ |
| $`A_{3.6}^4_t`$ | $`u_x^1`$ | $`\sqrt{|u_x|}`$ |
| $`A_{3.7}^2_x(q0,\pm 1,2)`$ | $`|u_x|^q`$ | $`|u_x|^{1q}`$ |
| $`A_{3.7}^4_t(0<|q|<1)`$ | $`|u_x|^{\frac{2}{q1}}`$ | $`|u_x|^{\frac{q}{q1}}`$ |
| $`A_{3.8}^2_t`$ | $`(1+u_x^2)^1`$ | $`\sqrt{1+u_x^2}`$ |
| $`A_{3.9}^2_t(q>0)`$ | $`\frac{\mathrm{exp}\left(2q\mathrm{arctan}u_x\right)}{1+u_x^2}`$ | $`\sqrt{1+u_x^2}\mathrm{exp}(q\mathrm{arctan}u_x)`$ |
As we have already mentioned, the realizations of the algebras $`A_{4.i}(i=7,8,9)`$ are constructed by extension of the realizations of the algebra $`A_{3.3}`$ by an operator $`e_4`$ of the type (3.3). Also, while considering the realizations of $`A_{3.3}=e_1,e_2,e_3`$, we have taken into account the isomorphism of this algebra given by $`e_1e_1,e_2e_3,e_3e_2.`$
In this way, we get three inequivalent realizations of the algebras $`A_{4.7}`$ and $`A_{4.9}`$
$`A_{4.7}^1=A_{3.3}^1+t_t+(xt)_x+(2u{\displaystyle \frac{1}{2}}t^2)_u,`$
$`A_{4.7}^2=A_{3.3}^2+_t+x_x+2u_u,`$
$`A_{4.9}^1=A_{3.3}^2+(1+t^2)_t+(qt)x_x+(2qu{\displaystyle \frac{1}{2}}x^2)_u(q>0),`$
which are maximal invariance algebras of nonlinear PDEs (0.1). The corresponding forms of the functions $`F`$, $`G`$ are given in Table 6.
Next, we have constructed four inequivalent realizations of the algebra $`A_{4.8}`$ that are admitted by nonlinear PDEs from the class (0.1):
$`A_{4.8}^1=A_{3.3}^1+t_t+qx_x+(1+q)u_u(qR),`$
$`A_{4.8}^2=A_{3.3}^1+t_t+k_x+u_u(k0),`$
$`A_{4.8}^3=A_{3.3}^1+x_x+u_u+k^1(_t+x_u)(k0),`$
$`A_{4.8}^4=A_{3.3}^2+(1q)t_t+x_x+(1+q)u_u(|q|1).`$
The realizations $`A_{4.8}^2,A_{4.8}^4`$ are the maximal invariance algebras of nonlinear heat conductivity equations belonging to the class of PDEs (0.1), and the corresponding forms of the functions $`F`$, $`G`$ are given in Table 6.
The PDE invariant with respect to the realization $`A_{4.8}^1`$ reduces to the form
$$u_t=\lambda |u_x|^{2q1}u_{xx}+x+ϵ|u_x|^q,$$
(3.45)
where $`ϵ=0,`$ $`\lambda =\pm 1,`$ provided $`q=0,1`$ and $`ϵ=0,\lambda =\pm 1`$ or $`ϵ=1,\lambda 0`$ if $`q0,1`$.
Investigating the maximal symmetry admitted by (3.45)we find that for $`q\frac{1}{2}`$ the realization $`A_{4.8}^1`$ is its maximal invariance algebra. If $`q=\frac{1}{2}`$, the change of variables (3.38) reduces PDE (3.45) to the Burgers equation
$$v_{\overline{t}}=\lambda v_{\overline{x}\overline{x}}vv_{\overline{x}}.$$
The maximal invariance algebra of the Burgers equation is the semi-direct sum of the algebra $`sl(2,R)`$ and a two-dimensional solvable radical.
The equation invariant under the realization $`A_{4.8}^3`$ is
$$u_t=\pm \mathrm{exp}(2ku_x)u_{xx}+x+ϵ\mathrm{exp}(ku_x),k0,ϵ=0,1.$$
If $`ϵ=1`$, then the realization $`A_{4.8}^3`$ is the maximal invariance algebra of this equation. For $`ϵ=0`$, the change of variables
$$\overline{t}=\frac{1}{2k}e^{2kt},\overline{x}=x,v=2ku+2ktx,k0,$$
reduces the equation in question to the PDE
$$v_{\overline{t}}=\pm \mathrm{exp}(v_{\overline{x}})v_{\overline{x}\overline{x}},$$
which belongs to the class of equations (3.39).
Finally, after extending the realizations of the Lie algebra $`A_{3.5}=e_1,e_2,e_3`$ by an operator $`e_4`$ of the form (3.3), we obtain a realization of the algebra $`A_{4.10}`$ of the form
$$A_{4.10}^1=A_{3.5}^3+2kt_t+u_xx_u,k0,$$
which is the maximal invariance algebra of the equation
$$u_t=\frac{\mathrm{exp}(2k\mathrm{arctan}u_x)}{1+u_x^2}u_{xx}+\beta |t|^{\frac{1}{2}}\sqrt{1+u_x^2}\mathrm{exp}(k\mathrm{arctan}u_x),k0,\beta 0.$$
We give in Table 6 a complete list of inequivalent PDEs of the form (0.1), whose maximal invariance algebras are non-decomposable four-dimensional solvable Lie algebras.
Table 6. Invariance of (0.1) under non-decomposable four-dimensional solvable Lie algebras
| Algebra | $`F`$ | $`G`$ |
| --- | --- | --- |
| $`A_{4.2}^1`$ | $`\mathrm{exp}(2q)u_x`$ | $`\mathrm{exp}(1q)u_x,q0,1`$ |
| $`A_{4.2}^2`$ | $`|u_x|^{\frac{1}{q1}}`$ | $`(1q)^1u_x\mathrm{ln}|u_x|,q0,1`$ |
| $`A_{4.3}^1`$ | $`\mathrm{exp}(u_x)`$ | $`\mathrm{exp}(u_x)`$ |
| $`A_{4.3}^2`$ | $`1`$ | $`u_x\mathrm{ln}|u_x|`$ |
| $`A_{4.4}^1`$ | $`\mathrm{exp}u_x`$ | $`\frac{1}{2}u_x^2`$ |
| $`A_{4.5}^1`$ | $`|u_x|^{\frac{2p1}{qp}}`$ | $`|u_x|^{\frac{q1}{qp}},p<q,pq0,p,q1`$ |
| $`A_{4.5}^2`$ | $`\stackrel{~}{F}(u_x)`$ | $`0`$ |
| $`A_{4.6}^1`$ | $`\frac{\mathrm{exp}\left[\left(q2p\right)\mathrm{arctan}u_x\right]}{1+u_x^2}`$ | $`\sqrt{1+u_x^2}\mathrm{exp}[(qp)\mathrm{arctan}u_x],qp,p0`$ |
| $`A_{4.7}^1`$ | $`\lambda u_x,\lambda 0`$ | $`x+u_x\mathrm{ln}|u_x|`$ |
| $`A_{4.7}^2`$ | $`\pm \mathrm{exp}(2t)`$ | $`\frac{1}{2}u_x^2`$ |
| $`A_{4.8}^1(q\frac{1}{2})`$ | $`\pm |u_x|^{2q1}`$ | $`x`$ |
| $`A_{4.8}^1(q0,1)`$ | $`\lambda |u_x|^{2q1},\lambda 0`$ | $`x+|u_x|^q`$ |
| $`A_{4.8}^2`$ | $`\lambda |u_x|^1,\lambda 0`$ | $`xk\mathrm{ln}|u_x|,k0`$ |
| $`A_{4.8}^3`$ | $`\pm \mathrm{exp}(2ku_x)`$ | $`x+\mathrm{exp}(ku_x),k0`$ |
| $`A_{4.8}^4(|q|1)`$ | $`|t|^{\frac{1+q}{1q}}`$ | $`\frac{1}{2}u_x^2`$ |
| $`A_{4.9}^1(q>0)`$ | $`\pm \mathrm{exp}(2q\mathrm{arctan}t)`$ | $`\frac{t\mathrm{exp}\left(2q\mathrm{arctan}t\right)}{1+t^2}\frac{1}{2}u_x^2`$ |
| $`A_{4.10}^1`$ | $`\frac{\mathrm{exp}\left(2k\mathrm{arctan}u_x\right]}{1+u_x^2}`$ | $`\beta |t|^{\frac{1}{2}}\sqrt{1+u_x^2}\mathrm{exp}(k\mathrm{arctan}u_x),k0,\beta 0`$ |
In order to complete the group classification, we have to analyze nonlinear equations of the form (0.1) which admit five-dimensional invariance algebras. In Section III.3 we have constructed two nonlinear heat conductivity equations whose invariance algebras are five-dimensional semi-direct products of semi-simple and solvable Lie algebras. According to the results of there are three more PDEs belonging to the class (3.39), admitting the five-dimensional algebras
$`A_5^2`$ $`=`$ $`A_{4.5}^2+t_tx_u;`$
$`A_5^3`$ $`=`$ $`A_{4.5}^2+nt_tu_u,(n1,n0);`$
$`A_5^4`$ $`=`$ $`A_{4.5}^2+nt_t+u_xx_u,(n0),`$
and there is one PDE admitting the realization $`A_{4.5}^1`$. It is not difficult to verify that all the algebras $`A_5^i(i=1,\mathrm{},4)`$ are solvable five-dimensional Lie algebras. Moreover, the algebra $`A_5^1`$ is decomposable $`A_5^1A_{3.7}A_{2.2}`$ and the algebras $`A_5^i(i=2,3,4)`$ are non-isomorphic non-decomposable five-dimensional solvable Lie algebras (see, e.g., )
$`A_5^2A_{5.34}(p=2),A_5^3A_{5.33}(p=2+n,q=n)`$
$`A_5^4A_{5.35}(p=2,q=n).`$
Consequently, the PDEs invariant with respect to the above algebras are inequivalent. We give in Table 7 a complete list of nonlinear PDEs of the form (0.1) whose maximal invariance algebras are five-dimensional.
## IV Some conclusions
Surprisingly, the number of inequivalent nonlinear PDEs of the general form under consideration, and which admit non-trivial symmetry groups is reasonably small. Summarizing the results of our group classification of nonlinear heat conductivity equations of the form (0.1) we conclude that
1. There are two inequivalent nonlinear PDEs (0.1), that admit a one-dimensional invariance algebra.
2. There are five inequivalent PDEs (0.1) given in Table 1, which are invariant with respect to two-dimensional Lie algebras. Note that all two-dimensional Lie algebras are solvable.
3. Nonlinear heat conductivity equations (0.1) invariant under three-dimensional Lie algebras (note that a three-dimensional Lie algebra is either semi-simple or solvable).
1. There are six PDEs (0.1) admitting three-dimensional semi-simple invariance algebras (see Theorems 3.1 and 3.2).
2. There are twenty eight equations (0.1), given in Tables 2–4, which are invariant with respect to three-dimensional solvable Lie algebras.
4. Nonlinear heat conductivity equations (0.1) invariant under four-dimensional Lie algebras (note that there are no semi-simple four-dimensional Lie algebras).
1. There are five PDEs (0.1) admitting four-dimensional invariance algebras, that are semi-direct sums of semi-simple and solvable Lie algebras (see PDEs given at the end of Section III.3).
2. There are thirty equations (0.1) given in Tables 5,6, which are invariant with respect to four-dimensional solvable Lie algebras.
5. Nonlinear heat conductivity equations (0.1) invariant under five-dimensional Lie algebras (note that there are no semi-simple five-dimensional Lie algebras).
1. There are two PDEs (0.1) admitting five-dimensional invariance algebras, that are semi-direct sums of semi-simple and solvable Lie algebras (see PDEs given at the end of Section III.3).
2. There are four equations (0.1) given in Table 7, which are invariant with respect to five-dimensional solvable Lie algebras.
Table 7. Invariance of (0.1) under four-dimensional solvable Lie algebras
| Algebra | $`F`$ | $`G`$ |
| --- | --- | --- |
| $`A_5^1`$ | $`u_x^2`$ | $`u_x^1`$ |
| $`A_5^2`$ | $`\mathrm{exp}u_x`$ | $`0`$ |
| $`A_5^3`$ | $`u_x^n,n1,n0`$ | $`0`$ |
| $`A_5^4`$ | $`\frac{\mathrm{exp}\left(n\mathrm{arctan}u_x\right)}{1+u_x^2},n0`$ | $`0`$ |
We have shown that there are no nonlinear PDEs of the form (0.1) admitting invariance algebras of the dimension higher than five. Consequently, the classification of invariant nonlinear heat conductivity equations (0.1) presented above is complete in the sense that any PDEs of the form (0.1), which possess non-trivial Lie symmetry, can be reduced to one of the canonical forms given above.
Furthermore, we have shown that the results on group classification of particular equations from the class (0.1) obtained in , can be derived from our considerations. That is to say, for each invariant PDE (0.1) obtained in the papers enumerated above we can give an invariant equation from our list which is equivalent to it. The procedure of looking for a corresponding change of variables is purely algebraic. We start with identifying the invariance algebra by determining whether it is semi-simple, or solvable, or a semi-direct sum of semi-simple and solvable algebras, and then we find the corresponding realizations of the Lie algebras of the same dimension as the algebra under study. Comparing the two realizations, it is not difficult to find the explicit form of the change of variables connecting these realizations (and, consequently transforming the corresponding invariant equations into one another).
One more point is that our classification is in full accordance with the results of Sokolov and Magadeev . These papers discuss, in particular, estimates for the dimension of the symmetry algebras of evolution PDEs with one or more spatial variables.
Another important point is the so called quasi-local or non-local symmetries of nonlinear heat conductivity equations. One can construct a number of these kind of symmetries as indicated in and combine them with non-local transformations like the Legendre, Laplace and Bäcklund transformations . As mentioned in the introduction, there is an intriguing possibility of reducing the problem of group classification of the general second order evolution equation (0.2) to that of PDE (0.1). Moreover, the ”singular points” of this reduction are the quasi-local symmetries of (0.1) which might correspond to usual Lie symmetries of (0.2). However, these very important questions go beyond the scope of the present paper and, in fact, can be a basis of a separate paper.
Let us stress again that it is our belief that the models most adequately describing real processes should possess the highest symmetry. That is why the most probable candidates for the roles of such models are PDEs admitting four and five-dimensional invariance algebras. The corresponding list of PDEs contain the well-known equations (like the Burgers equation) and several principally new equations which certainly deserve further investigation.
The questions mentioned above are under now study and will be reported on in future publications.
Acknowledgements. R. Zhdanov thanks the Swedish Natural Sciences Research Council for financial support (grant number R-RA 521-2373/1999) and the Mathematics Department, Linköping University, for its hospitality and financial support during his visit to Sweden.
## Appendix 1: solvable Lie algebras.
Three-dimensional solvable Lie algebras $`(L=e_1,e_2,e_3)`$ over $`𝐑`$
The set of three-dimensional solvable Lie algebras consists of the following two decomposable Lie algebras:
$`A_{3.1}`$ $`=`$ $`A_1A_1A_1=3A_1;`$
$`A_{3.2}`$ $`=`$ $`A_{2.2}A_1,[e_1,e_2]=e_2,`$
and the following eight classes of non-decomposable Lie algebras:
$`A_{3.3}`$ $`:`$ $`[e_2,e_3]=e_1;`$
$`A_{3.4}`$ $`:`$ $`[e_1,e_3]=e_1,[e_2,e_3]=e_1+e_2;`$
$`A_{3.5}`$ $`:`$ $`[e_1,e_3]=e_1,[e_2,e_3]=e_2;`$
$`A_{3.6}`$ $`:`$ $`[e_1,e_3]=e_1,[e_2,e_3]=e_2;`$
$`A_{3.7}`$ $`:`$ $`[e_1,e_3]=e_1,[e_2,e_3]=qe_2(0<|q|<1);`$
$`A_{3.8}`$ $`:`$ $`[e_1,e_3]=e_2,[e_2,e_3]=e_1;`$
$`A_{3.9}`$ $`:`$ $`[e_1,e_3]=qe_1e_2,[e_2,e_3]=e_1+qe_2,(q>0).`$
We note that the algebra $`A_{3.3}`$ is nilpotent. Note also that we have $`A_{3.i}(i=3,4,\mathrm{},9)`$ such that $`e_1,e_2`$ $`=A_{2.1}=2A_1.`$
Four-dimensional solvable Lie algebras $`(L=e_1,e_2,e_3,e_4)`$ over $`𝐑`$
Amongst the four-dimensional Lie algebras there are 10 decomposable algebras: $`4A_1=A_{3.1}A_1,A_{2.2}2A_1=A_{2.2}A_{2.1},A_{2.2}A_{2.2}=2A_{2.2},A_{3.i}A_1(i=3,4,\mathrm{},9)`$; and 10 non-decomposable solvable Lie algebras:
$`A_{4.1}`$ $`:`$ $`[e_2,e_4]=e_1,[e_3,e_4]=e_2;`$
$`A_{4.2}`$ $`:`$ $`[e_1,e_4]=qe_1,[e_2,e_4]=e_2,[e_3,e_4]=e_2+e_3,q0;`$
$`A_{4.3}`$ $`:`$ $`[e_1,e_4]=e_1,[e_3,e_4]=e_2;`$
$`A_{4.4}`$ $`:`$ $`[e_1,e_4]=e_1,[e_2,e_4]=e_1+e_2,[e_3,e_4]=e_2+e_3;`$
$`A_{4.5}`$ $`:`$ $`[e_1,e_4]=e_1,[e_2,e_4]=qe_2,[e_3,e_4]=pe_3,1pq1,pq0;`$
$`A_{4.6}`$ $`:`$ $`[e_1,e_4]=qe_1,[e_2,e_4]=pe_2e_3,[e_3,e_4]=e_2+pe_3,q0,p0;`$
$`A_{4.7}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_4]=2e_1,[e_2,e_4]=e_2,[e_3,e_4]=e_2+e_3;`$
$`A_{4.8}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_4]=(1+q)e_1,[e_2,e_4]=e_2,[e_3,e_4]=qe_3,|q|1;`$
$`A_{4.9}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_4]=2qe_1,[e_2,e_4]=qe_2e_3,[e_3,e_4]=e_2+qe_3,q0;`$
$`A_{4.10}`$ $`:`$ $`[e_1,e_3]=e_1,[e_2,e_3]=e_2,[e_1,e_4]=e_2,[e_2,e_4]=e_1.`$
Five-dimensional solvable Lie algebras $`(L=e_1,e_2,\mathrm{},e_5)`$ over $`𝐑`$
The set of non-isomorphic five-dimensional Lie algebras is exhausted by 27 types of decomposable algebras: $`5A_1,A_{2.2}3A_1,2A_{2.2}A_1,A_{3.i}2A_1(i=3,4,\mathrm{},8),A_{3.i}A_{2.2}(i=3,4,\mathrm{}8),A_{4.i}A_1(i=1,\mathrm{},10)`$; and 39 non-decomposable solvable algebras:
$`A_{5.1}`$ $`:`$ $`[e_3,e_5]=e_1,[e_4,e_5]=e_2;`$
$`A_{5.2}`$ $`:`$ $`[e_2,e_5]=e_1,[e_3,e_5]=e_2,[e_4,e_5]=e_3;`$
$`A_{5.3}`$ $`:`$ $`[e_2,e_4]=e_3,[e_2,e_5]=e_1,[e_4,e_5]=e_2;`$
$`A_{5.4}`$ $`:`$ $`[e_2,e_4]=e_1,[e_3,e_5]=e_1;`$
$`A_{5.5}`$ $`:`$ $`[e_3,e_4]=e_1,[e_2,e_5]=e_1,[e_3,e_5]=e_2;`$
$`A_{5.6}`$ $`:`$ $`[e_3,e_4]=e_1,[e_2,e_5]=e_1,[e_3,e_5]=e_2,[e_4,e_5]=e_3;`$
$`A_{5.7}`$ $`:`$ $`[e_1,e_5]=e_1,[e_2,e_5]=pe_2,[e_3,e_5]=qe_3,`$
$`[e_4,e_5]=re_4,1rqp1,rpq0;`$
$`A_{5.8}`$ $`:`$ $`[e_2,e_5]=e_1,[e_3,e_5]=e_3,[e_4,e_5]=pe_4,0<|p|1;`$
$`A_{5.9}`$ $`:`$ $`[e_1,e_5]=e_1,[e_2,e_5]=e_1+e_5,[e_3,e_5]=pe_3,[e_4,e_5]=qe_4,0qp;`$
$`A_{5.10}`$ $`:`$ $`[e_2,e_5]=e_1,[e_3,e_5]=e_2,[e_4,e_5]=e_4;`$
$`A_{5.11}`$ $`:`$ $`[e_1,e_5]=e_1,[e_2,e_5]=e_1+e_2,[e_3,e_5]=e_2+e_3,[e_4,e_5]=pe_4,p0;`$
$`A_{5.12}`$ $`:`$ $`[e_1,e_5]=e_1,[e_2,e_5]=e_1+e_2,[e_3,e_5]=e_2+e_3,[e_4,e_5]=e_3+e_4;`$
$`A_{5.13}`$ $`:`$ $`[e_1,e_5]=e_1,[e_2,e_5]=pe_2,[e_3,e_5]=qe_3re_4,`$
$`[e_4,e_5]=qe_4+re_3,|p|1,pr0,q0;`$
$`A_{5.14}`$ $`:`$ $`[e_2,e_5]=e_1,[e_3,e_5]=pe_3e_4,[e_4,e_5]=e_3+pe_4,p0;`$
$`A_{5.15}`$ $`:`$ $`[e_1,e_5]=e_1,[e_2,e_5]=e_1+e_2,[e_3,e_5]=pe_3,`$
$`[e_4,e_5]=e_3+pe_4,1p1;`$
$`A_{5.16}`$ $`:`$ $`[e_1,e_5]=e_1,[e_2,e_5]=e_1+e_2,[e_3,e_5]=pe_3qe_4,`$
$`[e_4,e_5]=qe_3+pe_4,p0,q0;`$
$`A_{5.17}`$ $`:`$ $`[e_1,e_5]=pe_1e_2,[e_2,e_5]=e_1+pe_2,[e_3,e_5]=qe_3re_4,`$
$`[e_4,e_5]=re_3+qe_4,r0,p,q𝐑;`$
$`A_{5.18}`$ $`:`$ $`[e_1,e_5]=pe_1e_2,[e_2,e_5]=e_1+pe_2,[e_3,e_5]=e_1+pe_3e_4,`$
$`[e_4,e_5]=e_2+e_3pe_4,p𝐑;`$
$`A_{5.19}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=(1+p)e_1,[e_2,e_5]=e_2,[e_3,e_5]=pe_3,`$
$`[e_4,e_5]=qe_4,p𝐑,q0;`$
$`A_{5.20}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=(1+p)e_2,[e_2,e_5]=e_2,[e_3,e_5]=pe_3,`$
$`[e_4,e_5]=e_1+(1+p)e_4,p,q𝐑;`$
$`A_{5.21}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=2e_1,[e_2,e_5]=e_2+e_3,[e_3,e_5]=e_3+e_4,[e_4,e_5]=e_4;`$
$`A_{5.22}`$ $`:`$ $`[e_2,e_3]=e_1,[e_2,e_5]=e_3,[e_4,e_5]=e_4;`$
$`A_{5.23}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=2e_1,[e_2,e_5]=e_2+e_3,`$
$`[e_3,e_5]=e_3,[e_4,e_5]=pe_4,p0;`$
$`A_{5.24}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=2e_1,[e_2,e_5]=e_2+e_3,`$
$`[e_3,e_5]=e_3,[e_4,e_5]=ϵe_1+2e_4,ϵ=\pm 1;`$
$`A_{5.25}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=2pe_1,[e_2,e_5]=pe_2+e_3,[e_3,e_5]=e_2+pe_3,`$
$`[e_4,e_5]=qe_4,p𝐑,q0;`$
$`A_{5.26}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=2pe_1,[e_2,e_5]=pe_2+e_3,[e_3,e_5]=e_2+pe_3,`$
$`[e_4,e_5]=ϵe_1+2pe_4,ϵ=\pm 1,p𝐑;`$
$`A_{5.27}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=e_1,[e_3,e_5]=e_3+e_4,[e_4,e_5]=e_1+e_4;`$
$`A_{5.28}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=(1+p)e_1,[e_2,e_5]=pe_2,`$
$`[e_3,e_5]=e_3+e_4,[e_4,e_5]=e_4,p𝐑;`$
$`A_{5.29}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_5]=e_1,[e_2,e_5]=e_2,[e_3,e_5]=e_4;`$
$`A_{5.30}`$ $`:`$ $`[e_2,e_4]=e_1,[e_3,e_4]=e_2,[e_1,e_5]=(2+p)e_1,[e_2,e_5]=(1+p)e_2,`$
$`[e_3,e_5]=pe_3,[e_4,e_5]=e_4,p𝐑;`$
$`A_{5.31}`$ $`:`$ $`[e_2,e_4]=e_1,[e_3,e_4]=e_2,[e_1,e_5]=3e_1,`$
$`[e_2,e_5]=2e_2,[e_3,e_5]=e_3,[e_4,e_5]=e_3+e_4;`$
$`A_{5.32}`$ $`:`$ $`[e_2,e_4]=e_1,[e_3,e_4]=e_2,[e_1,e_5]=e_1,[e_2,e_5]=e_2,`$
$`[e_3,e_5]=pe_1+e_3,p𝐑;`$
$`A_{5.33}`$ $`:`$ $`[e_1,e_4]=e_1,[e_3,e_4]=pe_3,[e_2,e_5]=e_2,`$
$`[e_3,e_5]=qe_3,p,q𝐑,p^2+q^20;`$
$`A_{5.34}`$ $`:`$ $`[e_1,e_4]=pe_1,[e_2,e_4]=e_2,[e_3,e_4]=e_3,[e_1,e_5]=e_1,[e_3,e_5]=e_2,p𝐑;`$
$`A_{5.35}`$ $`:`$ $`[e_1,e_4]=pe_1,[e_2,e_4]=e_2,[e_3,e_4]=e_3,[e_1,e_5]=qe_1,`$
$`[e_2,e_5]=e_3,[e_3,e_5]=e_2,p,q𝐑,p^2+q^20;`$
$`A_{5.36}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_4]=e_1,[e_2,e_4]=e_2,[e_2,e_5]=e_2,[e_3,e_5]=e_3;`$
$`A_{5.37}`$ $`:`$ $`[e_2,e_3]=e_1,[e_1,e_4]=2e_1,[e_2,e_4]=e_2,`$
$`[e_3,e_4]=e_3,[e_2,e_5]=e_3,[e_3,e_5]=e_2;`$
$`A_{5.38}`$ $`:`$ $`[e_1,e_4]=e_1,[e_2,e_5]=e_2,[e_4,e_5]=e_3;`$
$`A_{5.39}`$ $`:`$ $`[e_1,e_4]=e_1,[e_2,e_4]=e_2,[e_1,e_5]=e_2,[e_2,e_5]=e_1,[e_4,e_5]=e_3.`$
## Appendix 2: Lie algebras which are semi-direct sums of semi-simple and solvable algebras.
1. Lie algebras of dimensions 5 and 6.
$`sl(2,𝐑)+A_{2.1}`$ $`:`$ $`[e_1,e_4]=e_4,[e_2,e_5]=e_4,[e_3,e_4]=e_5,[e_1,e_5]=e_5;`$
$`so(3)+A_{3.1}`$ $`:`$ $`[e_1,e_5]=e_6,[e_2,e_4]=e_6,[e_3,e_4]=e_5,[e_1,e_6]=e_5,`$
$`[e_2,e_6]=e_4,[e_3,e_5]=e_4;`$
$`sl(2,𝐑)+A_{3.i}`$ $`:`$ $`[e_1,e_4]=e_4,[e_2,e_5]=e_4,`$
$`i=3,A_{3.3}=e_6,e_4,e_5`$ $`[e_3,e_4]=e_5,[e_1,e_5]=e_5,`$
$`i=5,A_{3.5}=e_4,e_5,e_6;`$
$`sl(2,𝐑)+A_{3.1}`$ $`:`$ $`[e_1,e_4]=2e_4,[e_2,e_5]=2e_4,[e_3,e_4]=e_5,`$
$`[e_1,e_6]=2e_6,[e_2,e_6]=e_5,[e_3,e_5]=2e_6.`$
2. Lie algebras of dimension 7.
$`so(3)+A_{4.5}(p=q=1)`$ $`:`$ $`[e_1,e_5]=e_6,[e_2,e_4]=e_6,[e_3,e_4]=e_5,`$
$`[e_1,e_6]=e_5,[e_2,e_6]=e_4,[e_3,e_5]=e_4;`$
$`so(3)+4A_1`$ $`:`$ $`[e_1,e_4]={\displaystyle \frac{1}{2}}e_7,[e_2,e_4]={\displaystyle \frac{1}{2}}e_5,[e_3,e_4]={\displaystyle \frac{1}{2}}e_6,`$
$`[e_1,e_5]={\displaystyle \frac{1}{2}}e_6,[e_2,e_5]={\displaystyle \frac{1}{2}}e_4,[e_3,e_5]={\displaystyle \frac{1}{2}}e_7,`$
$`[e_1,e_6]={\displaystyle \frac{1}{2}}e_5,[e_2,e_6]={\displaystyle \frac{1}{2}}e_7,[e_3,e_6]={\displaystyle \frac{1}{2}}e_4,`$
$`[e_1,e_7]={\displaystyle \frac{1}{2}}e_4,[e_2,e_7]={\displaystyle \frac{1}{2}}e_6,[e_3,e_7]={\displaystyle \frac{1}{2}}e_5;`$
$`sl(2,𝐑)+A_{4.i}`$ $`:`$ $`[e_1,e_4]=e_4,[e_2,e_5]=e_4,`$
$`i=5:A_{4.5}(q=1),`$ $`[e_3,e_4]=e_5,[e_1,e_5]=e_5,`$
$`i=8:A_{4.8}(q=1),`$
$`A_{4.8}=e_6,e_4,e_5,e_7;`$
$`sl(2,𝐑)+A_{4.5}`$ $`:`$ $`[e_1,e_4]=2e_4,[e_2,e_5]=2e_4,[e_3,e_4]=e_5,`$
$`A_{4.5}(p=q=1)`$ $`[e_1,e_6]=2e_6,[e_2,e_6]=e_5,[e_3,e_5]=2e_6;`$
$`sl(2,𝐑)+4A_1`$ $`:`$ $`[e_1,e_4]=3e_4,[e_2,e_5]=3e_4,[e_3,e_4]=e_5,`$
$`[e_1,e_5]=e_5,[e_2,e_6]=2e_5,[e_3,e_5]=2e_6,`$
$`[e_1,e_6]=e_6,[e_2,e_7]=e_6,[e_3,e_6]=3e_7,[e_1,e_7]=3e_7;`$
$`sl(2,𝐑)+4A_1`$ $`:`$ $`[e_1,e_4]=e_4,[e_2,e_5]=e_4,[e_3,e_4]=e_5,[e_1,e_5]=e_5,`$
$`[e_1,e_6]=e_6,[e_2,e_7]=e_6,[e_3,e_6]=e_7,[e_1,e_7]=e_7.`$
3. Lie algebras of dimension 8.
$`so(3)+A_{5.7}`$ $`:`$ $`[e_1,e_5]=e_6,[e_2,e_4]=e_6,[e_3,e_4]=e_5,`$
$`A_{5.7}(p=q=1)`$ $`[e_1,e_6]=e_5,[e_2,e_6]=e_4,[e_3,e_5]=e_4;`$
$`so(3)+A_{5.i}`$ $`:`$ $`[e_1,e_4]={\displaystyle \frac{1}{2}}e_7,[e_2,e_4]={\displaystyle \frac{1}{2}}e_5,`$
$`i=4:A_{5.4}=e_8,e_4,e_7,e_5,e_6,`$ $`[e_3,e_4]={\displaystyle \frac{1}{2}}e_6,[e_1,e_5]={\displaystyle \frac{1}{2}}e_6,`$
$`i=7:A_{5.7}(p=q=r=1),`$ $`[e_2,e_5]={\displaystyle \frac{1}{2}}e_4,[e_3,e_5]={\displaystyle \frac{1}{2}}e_7,`$
$`i=17:A_{5.17}(p=q,r=1),`$ $`[e_1,e_6]={\displaystyle \frac{1}{2}}e_5,[e_2,e_6]={\displaystyle \frac{1}{2}}e_7,`$
$`A_{5.17}=e_4,e_6,e_5,e_7,e_8,`$ $`[e_3,e_6]={\displaystyle \frac{1}{2}}e_4,[e_1,e_7]={\displaystyle \frac{1}{2}}e_4,`$
$`[e_2,e_7]={\displaystyle \frac{1}{2}}e_6,[e_3,e_7]={\displaystyle \frac{1}{2}}e_5;`$
$`so(3)+5A_1`$ $`:`$ $`[e_1,e_4]={\displaystyle \frac{1}{2}}e_7,[e_1,e_5]={\displaystyle \frac{1}{2}}e_6,`$
$`[e_1,e_6]=2e_5e_8,[e_1,e_7]=2e_4,`$
$`[e_1,e_8]=3e_6,[e_2,e_4]={\displaystyle \frac{1}{2}}e_6,`$
$`[e_2,e_5]={\displaystyle \frac{1}{2}}e_7,[e_2,e_6]=2e_4,`$
$`[e_2,e_7]=2e_5e_8,[e_2,e_8]=3e_7,`$
$`[e_3,e_4]=2e_5,[e_3,e_5]=2e_4,`$
$`[e_3,e_6]=e_7,[e_3,e_7]=e_6;`$
$`sl(2,𝐑)+A_{5.i}`$ $`:`$ $`[e_1,e_4]=e_4,[e_2,e_5]=e_4,`$
$`[e_3,e_4]=e_5,[e_1,e_5]=e_5;`$
$`i=4:A_{5.4}=e_8,e_4,e_6,e_5,e_7,`$
$`i=7,8:A_i(p=1),`$
$`A_{5.8}=e_6,e_7,e_4,e_5,e_8,`$
$`i=9:A_{5.9},`$
$`i=13,19,20:A_{5.i}(p=1),`$
$`A_{5.i}(i=19,20)=e_6,e_4,e_5,e_7,e_8,`$
$`sl(2,𝐑)+A_{5.7}`$ $`:`$ $`[e_1,e_4]=2e_4,[e_2,e_5]=2e_4,[e_3,e_4]=e_5,`$
$`[e_1,e_5]=2e_6;[e_2,e_6]=e_5,[e_3,e_5]=2e_6;`$
$`sl(2,𝐑)+A_{5.i}`$ $`:`$ $`[e_1,e_4]=e_4,[e_2,e_5]=e_4,`$
$`i=4:A_{5.4}^ϵ,`$ $`[e_3,e_4]=e_5,[e_1,e_5]=e_5,`$
$`i=1:A_{5.1},`$ $`[e_1,e_6]=e_6,[e_2,e_7]=e_6,`$
$`i=3:A_{5.3},`$ $`[e_3,e_6]=e_7,[e_1,e_7]=e_7,`$
$`i=15:A_{5.15}(p=1)`$
$`i=7:A_{5.7}(p=q=1,1r1)`$
$`i=17:A_{5.17}(p=q,r=1,p0)`$
$`A_{5.i}(i=7,17)=e_4,e_6,e_5,e_7,e_8;`$
$`sl(2,𝐑)+A_{5.i}`$ $`:`$ $`[e_1,e_4]=3e_4,[e_2,e_5]=3e_4,`$
$`[e_3,e_4]=e_5,[e_1,e_5]=e_5;`$
$`i=4:A_{5.4}`$ $`[e_2,e_6]=2e_5,[e_3,e_5]=2e_6,[e_1,e_6]=e_6,`$
$`i=7:A_{5.7}(p=q=r=1)`$ $`[e_2,e_7]=e_6,[e_3,e_6]=3e_7,[e_1,e_7]=3e_7;`$
$`sl(2,𝐑)+5A_1`$ $`:`$ $`[e_1,e_4]=4e_4,[e_2,e_5]=4e_4,`$
$`[e_3,e_4]=e_5,[e_1,e_5]=2e_5;`$
$`[e_2,e_6]=3e_5,[e_3,e_5]=2e_6,`$
$`[e_1,e_7]=2e_7,[e_2,e_7]=2e_6,`$
$`[e_3,e_6]=3e_7,[e_1,e_8]=4e_8,`$
$`[e_2,e_8]=e_7,[e_3,e_7]=4e_8;`$
$`sl(2,𝐑)+5A_1`$ $`:`$ $`[e_1,e_4]=2e_4,[e_2,e_5]=2e_4,`$
$`[e_3,e_4]=e_5,[e_1,e_6]=2e_6;`$
$`[e_2,e_6]=e_5,[e_3,e_5]=2e_6,`$
$`[e_1,e_7]=e_7,[e_2,e_8]=e_7,`$
$`[e_3,e_7]=e_8,[e_1,e_8]=e_8.`$
In giving the type of the radicals, we have followed the rule that the bases of the radicals consist of the operators $`e_4,\mathrm{},e_m,`$ whenever the basis is not given explicitly; where it is given explicitly, then the basis operators are ordered as in the corresponding solvable algebra. For instance, in the algebra $`sl(2,𝐑)+A_{5.17}`$ we have written $`A_{5.17}=e_4,e_6,e_5,e_7,e_8.`$ This means that the basis operators satisfy the commutation relations which define the algebra $`A_{5.17}`$ given in the list of solvable algebras. To obtain the commutation relations for the algebra $`A_{5.17},`$ we replace the operators $`e_4,e_5,e_6,e_7,e_8`$ as follows:
$$e_4e_1,e_6e_2,e_5e_3,e_7e_4,e_8e_5.$$
Furthermore, for the five-dimensional radicals $`N=e_4,e_5,e_6,e_7,e_8`$ we use the notation
$`A_{5.9}`$ $`:`$ $`[e_4,e_8]=e_4,[e_5,e_8]=e_5,[e_6,e_8]=pe_6,`$
$`[e_7,e_8]=e_6+pe_7,p0;`$
$`A_{5.4}^ϵ`$ $`:`$ $`[e_4,e_8]=e_8,[e_6,e_7]=ϵe_8,ϵ=\pm 1;`$
$`A_{5.3}`$ $`:`$ $`[e_6,e_8]=e_4,[e_7,e_8]=e_5,[e_6,e_7]=e_8;`$
$`A_{5.4}`$ $`:`$ $`[e_4,e_7]=e_8,[e_5,e_6]=3e_8.`$
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# Decomposition of analytic measures on groups and measure spaces
## 1 Introduction
This paper is essentially providing a new approach to generalizations of the F.&M. Riesz Theorems, for example, such results as that of Helson and Lowdenslager . They showed that if $`G`$ is a compact abelian group with ordered dual, and if $`\mu `$ is an analytic measure (that is, its Fourier transform is supported on the positive elements of the dual), then it follows that the singular and absolutely continuous parts (with respect to the Haar measure) are also analytic.
Another direction is that provided by Forelli (itself a generalization of the result of de Leeuw and Glicksberg ), where one has an action of the real numbers $``$ acting on a locally compact topological space $`\mathrm{\Omega }`$, and a Baire measure $`\mu `$ on $`\mathrm{\Omega }`$ that is analytic (in a sense that we make precise below) with respect to the action. Then again, the singular and absolutely continuous parts of $`\mu `$ (with respect to any so called quasi-invariant measure) are also analytic.
Indeed common generalizations of both these ideas have been provided, for example, by Yamaguchi , considering the action of any locally compact abelian group with ordered dual, on a locally compact topological space. For more generalizations we refer the reader to Hewitt, Koshi, and Takahashi .
In the paper , a new approach to proving these kinds of results was given, providing a transference principle for spaces of measures. In that paper, the action was from a locally compact abelian group into a space of isomorphisms on the space of measures of a sigma algebra. A primary requirement that the action had to satisfy was what was called sup path attaining, a property that was satisfied, for example, by the setting of Forelli (Baire measures on a locally compact topological space). Using this transference principle, the authors were able to give an extension and a new proof of Forelli’s result. This was obtained by using a Littlewood-Paley decomposition of an analytic measure.
In this paper we wish to continue this process, applying this same transference principle to provide the common generalizations of the results of Forelli and Helson and Lowdenslager. What we provide in this paper is essentially a decomposition of an analytic measure as a sum of martingale differences with respect to a filtration defined by the order. For each martingale difference, the action of the group can be described precisely by a certain action of the group of real numbers, and so we can appeal to the results of .
In this way, we can reach the following generalization (see Theorem 6.4 below): if $`𝒫`$ is any bounded operator on the space of measures that commutes with the action (as does, for example, taking the singular part), and if $`\mu `$ is an analytic measure, then $`𝒫\mu `$ is also an analytic measure.
In the remainder of the introduction, we will establish our notation, including the notion of sup path attaining, and recall the transference principle from . In Section 2, we will describe orders on locally compact abelian groups, including the extension of Hahn’s Embedding Theorem provided in . In Section 3, we define the notions of analyticity. This somewhat technical section continues into Section 4, which examines the role of homomorphism with respect to analyticity. The technical results basically provide proofs of what is believable, and so may be skipped on first reading. It will be seen that the concept of sup path attaining comes up again and again, and may be seen to be an integral part of all our proofs.
In Section 5, we are ready to present the decomposition of analytic measures. This depends heavily on transference of martingale inequalities of Burkholder and Garling, and then using the fact that weakly unconditionally summing series are unconditionally summing in norm for any series in a space of measures . In Section 6, we then give applications of this decomposition, giving the generalizations that we alluded to above.
Throughout $`G`$ will denote a locally compact abelian group with dual group $`\mathrm{\Gamma }`$. The symbols $``$, $``$ and $``$ denote the integers, the real and complex numbers, respectively. If $`A`$ is a set, we denote the indicator function of $`A`$ by $`1_A`$. For $`1p<\mathrm{}`$, the space of Haar measurable functions $`f`$ on $`G`$ with $`_G|f|^p𝑑x<\mathrm{}`$ will be denoted by $`L^p(G)`$. The space of essentially bounded functions on $`G`$ will be denoted by $`L^{\mathrm{}}(G)`$. The expressions “locally null” and “locally almost everywhere” will have the same meanings as in \[20, Definition (11.26)\].
Let $`𝒞_0(G)`$ denote the Banach space of continuous functions on $`G`$ vanishing at infinity. The space of all complex regular Borel measures on $`G`$, denoted by $`M(G)`$, consists of all complex measures arising from bounded linear functionals on $`𝒞_0(G)`$.
Let $`(\mathrm{\Omega },\mathrm{\Sigma })`$ denote a measurable space, where $`\mathrm{\Omega }`$ is a set and $`\mathrm{\Sigma }`$ is a sigma algebra of subsets of $`\mathrm{\Omega }`$. Let $`M(\mathrm{\Sigma })`$ denote the Banach space of complex measures on $`\mathrm{\Sigma }`$ with the total variation norm, and let $`^{\mathrm{}}(\mathrm{\Sigma })`$ denote the space of measurable bounded functions on $`\mathrm{\Omega }`$.
Let $`T:tT_t`$ denote a representation of $`G`$ by isomorphisms of $`M(\mathrm{\Sigma })`$. We suppose that $`T`$ is uniformly bounded, i.e., there is a positive constant $`c`$ such that for all $`tG`$, we have
(1)
$$T_tc.$$
###### Definition 1.1
A measure $`\mu M(\mathrm{\Sigma })`$ is called weakly measurable (in symbols, $`\mu _T(\mathrm{\Sigma })`$) if for every $`A\mathrm{\Sigma }`$ the mapping $`tT_t\mu (A)`$ is Borel measurable on $`G`$.
Given a measure $`\mu _T(\mathrm{\Sigma })`$ and a Borel measure $`\nu M(G)`$, we define the ‘convolution’ $`\nu _T\mu `$ on $`\mathrm{\Sigma }`$ by
(2)
$$\nu _T\mu (A)=_GT_t\mu (A)𝑑\nu (t)$$
for all $`A\mathrm{\Sigma }`$.
We will assume throughout this paper that the representation $`T`$ commutes with the convolution (2) in the following sense: for each $`tG`$,
(3)
$$T_t(\nu _T\mu )=\nu _T(T_t\mu ).$$
Condition (3) holds if, for example, for all $`tG`$, the adjoint of $`T_t`$ maps $`^{\mathrm{}}(\mathrm{\Sigma })`$ into itself. In symbols,
(4)
$$T_t^{}:^{\mathrm{}}(\mathrm{\Sigma })^{\mathrm{}}(\mathrm{\Sigma }).$$
For proofs we refer the reader to . Using (1) and (3), it can be shown that $`\nu _T\mu `$ is a measure in $`_T(\mathrm{\Sigma })`$,
(5)
$$\nu _T\mu c\nu \mu ,$$
where $`c`$ is as in (1), and
(6)
$$\sigma _T(\nu _T\mu )=(\sigma \nu )_T\mu ,$$
for all $`\sigma ,\nu M(G)`$ and $`\mu _T(\mathrm{\Sigma })`$ (see ).
###### Definition 1.2
A representation $`T=(T_t)_{tG}`$ of a locally compact abelian group $`G`$ in $`M(\mathrm{\Sigma })`$ is said to be sup path attaining if it is uniformly bounded, satisfies property (3), and if there is a constant $`C`$ such that for every weakly measurable $`\mu _T(\mathrm{\Sigma })`$ we have
(7)
$$\mu Csup\left\{\mathrm{ess}\mathrm{sup}_{tG}\right|_\mathrm{\Omega }hd(T_t\mu )|:h^{\mathrm{}}(\mathrm{\Sigma }),h_{\mathrm{}}1\}.$$
The fact that the mapping $`t_\mathrm{\Omega }hd(T_t\mu )`$ is measurable is a simple consequence of the measurability of the mapping $`tT_t\mu (A)`$ for every $`A\mathrm{\Sigma }`$.
In were provided many examples of sup path attaining representations. Rather than give this same list again, we give a couple of examples of particular interest.
###### Example 1.3
(a) (This is the setting of Forelli’s Theorem.) Let $`G`$ be a locally compact abelian group, and $`\mathrm{\Omega }`$ be a locally compact topological space. Suppose that $`\left(T_t\right)_{tG}`$ is a group of homeomorphisms of $`\mathrm{\Omega }`$ onto itself such that the mapping
$$(t,\omega )T_t\omega $$
is jointly continuous. Then the space of Baire measures on $`\mathrm{\Omega }`$, that is, the minimal sigma algebra such that compactly supported continuous functions are measurable, is sup path attaining under the action $`T_t\mu (A)=\mu (T_t(A))`$, where $`T_t(A)=\{T_t\omega :\omega A\}`$. (Note that all Baire measures are weakly measurable.)
(b) Suppose that $`G_1`$ and $`G_2`$ are locally compact abelian groups and that $`\varphi :G_2G_1`$ is a continuous homomorphism. Define an action of $`G_2`$ on $`M(G_1)`$ (the regular Borel measures on $`G_1`$) by translation by $`\varphi `$. Hence, for $`xG_2,\mu M(G_1)`$, and any Borel subset $`AG_1`$, let $`T_x\mu (A)=\mu (A+\varphi (x))`$. Then every $`\mu M(G_1)`$ is weakly measurable, and the representation is sup path attaining with constants $`c=1`$ and $`C=1`$.
###### Proposition 1.4
Suppose that $`T`$ is sup path attaining and $`\mu `$ is weakly measurable such that for every $`A\mathrm{\Sigma }`$ we have
$$T_t\mu (A)=0$$
for locally almost all $`tG`$. Then $`\mu =0`$.
The proof is immediate (see ).
We now recall some basic definitions from spectral theory.
If $`I`$ is an ideal in $`L^1(G)`$, let
$$Z(I)=\underset{fI}{}\{\chi \mathrm{\Gamma }:\widehat{f}(\chi )=0\}.$$
The set $`Z(I)`$ is called the zero set of $`I`$. For a weakly measurable $`\mu M(\mathrm{\Sigma })`$, let
$$(\mu )=\{fL^1(G):f_T\mu =0\}.$$
When we need to be specific about the representation, we will use the symbol $`_T(\mu )`$ instead of $`(\mu )`$.
Using properties of the convolution $`_T`$, it is straightforward to show that $`(\mu )`$ is a closed ideal in $`L^1(G)`$.
###### Definition 1.5
The $`T`$-spectrum of a weakly measurable $`\mu _T(\mathrm{\Sigma })`$ is defined by
(8)
$$\mathrm{spec}_T(\mu )=\underset{f(\mu )}{}\{\chi \mathrm{\Gamma }:\widehat{f}(\chi )=0\}=Z((\mu )).$$
If $`S\mathrm{\Gamma }`$, let
$$L_S^1=L_S^1(G)=\{fL^1(G):\widehat{f}=0\text{outside of}S\}.$$
In order to state the main transference result, we introduce one more definition.
###### Definition 1.6
A subset $`S\mathrm{\Gamma }`$ is a $`𝒯`$-set if, given any compact $`KS`$, each neighborhood of $`0\mathrm{\Gamma }`$ contains a nonempty open set $`W`$ such that $`W+KS`$.
###### Example 1.7
(a) If $`\mathrm{\Gamma }`$ is a locally compact abelian group, then any open subset of $`\mathrm{\Gamma }`$ is a $`𝒯`$-set. In particular, if $`\mathrm{\Gamma }`$ is discrete then every subset of $`\mathrm{\Gamma }`$ is a $`𝒯`$-set.
(b) The set $`[a,\mathrm{})`$ is a $`𝒯`$-subset of $``$, for all $`a`$.
(c) Let $`a`$ and $`\psi :\mathrm{\Gamma }`$ be a continuous homomorphism. Then $`S=\psi ^1([a,\mathrm{}))`$ is a $`𝒯`$-set.
(d) Let $`\mathrm{\Gamma }=^2`$ and $`S=\{(x,y):y^2x\}`$. Then $`S`$ is a $`𝒯`$-subset of $`^2`$ such that there is no nonempty open set $`W^2`$ such that $`W+SS`$.
The main result of is the following transference theorem.
###### Theorem 1.8
Let $`T`$ be a sup path attaining representation of a locally compact abelian group $`G`$ by isomorphisms of $`M(\mathrm{\Sigma })`$ and let $`S`$ be a $`𝒯`$-subset of $`\mathrm{\Gamma }`$. Suppose that $`\nu `$ is a measure in $`M(G)`$ such that
(9)
$$\nu f_1f_1$$
for all $`f`$ in $`L_S^1(G)`$. Then for every weakly measurable $`\mu M(\mathrm{\Sigma })`$ with $`\mathrm{spec}_T(\mu )S`$ we have
(10)
$$\nu _T\mu c^3C\mu ,$$
where $`c`$ is as in (1) and $`C`$ is as in (7).
## 2 Orders on locally compact abelian groups
An order $`P`$ on $`\mathrm{\Gamma }`$ is a subset that satisfies the three axioms: $`P+PP`$; $`P(P)=\mathrm{\Gamma }`$; and $`P(P)\{0\}`$. We recall from the following property of orders.
###### Theorem 2.1
Let $`P`$ be a measurable order on $`\mathrm{\Gamma }`$. There are a totally ordered set $`\mathrm{\Pi }`$ with largest element $`\alpha _0`$; a chain of subgroups $`\{C_\alpha \}_{\alpha \mathrm{\Pi }}`$ of $`\mathrm{\Gamma }`$; and a collection of continuous real-valued homomorphisms $`\{\psi _\alpha \}_{\alpha \mathrm{\Pi }}`$ on $`\mathrm{\Gamma }`$ such that:
(i) for each $`\alpha \mathrm{\Pi }`$, $`C_\alpha `$ is an open subgroup of $`\mathrm{\Gamma }`$;
(ii) $`C_\alpha C_\beta `$ if $`\alpha >\beta `$.
Let $`D_\alpha =\{\chi C_\alpha :\psi _\alpha (\chi )=0\}`$. Then, for every $`\alpha \mathrm{\Pi }`$,
(iii) $`\psi _\alpha (\chi )>0`$ for every $`\chi P(C_\alpha D_\alpha )`$,
(iv) $`\psi _\alpha (\chi )<0`$ for every $`\chi (P)(C_\alpha D_\alpha ).`$
(v) When $`\mathrm{\Gamma }`$ is discrete, $`C_{\alpha _0}=\{0\}`$; and when $`\mathrm{\Gamma }`$ is not discrete, $`D_{\alpha _0}`$ has empty interior and is locally null.
When $`\mathrm{\Gamma }`$ is discrete, Theorem 2.1 can be deduced from the proof of Hahn’s Embedding Theorem for orders (see \[13, Theorem 16, p.59\]). The general case treated in Theorem 2.1 accounts for the measure theoretic aspect of orders. The proof is based on the study of orders of Hewitt and Koshi .
For $`\alpha \mathrm{\Pi }`$ with $`\alpha \alpha _0`$, let
(11) $`S_\alpha P(C_\alpha D_\alpha )`$ $`=`$ $`\{\chi C_\alpha D_\alpha :\psi _\alpha (\chi )0\}`$
(12) $`=`$ $`\{\chi C_\alpha :\psi _\alpha (\chi )>0\}.`$
For $`\alpha =\alpha _0`$, set
(13)
$$S_{\alpha _0}=\{\chi C_{\alpha _0}:\psi _{\alpha _0}(\chi )0\}.$$
Note that when $`\mathrm{\Gamma }`$ is discrete, $`C_{\alpha _0}=\{0\}`$, and so $`S_{\alpha _0}=\{0\}`$ in this case.
If $`A`$ is a subset of a topological space, we will use $`\overline{A}`$ and $`A^{}`$ to denote the closure, respectively, the interior of $`A`$.
###### Remarks 2.2
(a) It is a classical fact that a group $`\mathrm{\Gamma }`$ can be ordered if and only if it is torsion-free. Also, an order on $`\mathrm{\Gamma }`$ is any maximal positively linearly independent set. Thus, orders abound in torsion-free abelian groups, as they can be constructed using Zorn’s Lemma to obtain a maximal positively linearly independent set. (See \[18, Section 2\].) However, if we ask for measurable orders, then we are restricted in many ways in the choices of $`P`$ and also the topology on $`\mathrm{\Gamma }`$. As shown in , any measurable order on $`\mathrm{\Gamma }`$ has nonempty interior. Thus, for example, while there are infinitely many orders on $``$, only two are Lebesgue measurable: $`P=[0,\mathrm{}[`$, and $`P=]\mathrm{},0]`$. It is also shown in \[18, Theorem (3.2)\] that any order on an infinite compact torsion-free abelian group is non-Haar measurable. This effectively shows that if $`\mathrm{\Gamma }`$ contains a Haar-measurable order $`P`$, and we use the structure theorem for locally compact abelian groups to write $`\mathrm{\Gamma }`$ as $`^a\times \mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ contains a compact open subgroup \[20, Theorem (24.30)\], then either $`a`$ is a positive integer, or $`\mathrm{\Gamma }`$ is discrete. (See .)
(b) The subgroups $`(C_\alpha )`$ are characterized as being the principal convex subgroups in $`\mathrm{\Gamma }`$ and for each $`\alpha \mathrm{\Pi }`$, we have
$$D_\alpha =\underset{\beta >\alpha }{}C_\beta .$$
Consequently, we have $`C_\alpha D_\beta `$ if $`\beta <\alpha `$. By construction, the sets $`C_\alpha `$ are open. For $`\alpha <\alpha _0`$, the subgroup $`D_\alpha `$ has nonempty interior, since it contains $`C_\beta `$, with $`\alpha <\beta `$. Hence for $`\alpha \alpha _0`$, $`D_\alpha `$ is open and closed. Consequently, for $`\alpha \alpha _0`$, $`C_\alpha D_\alpha `$ is open and closed.
(c) Let $`\psi :\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ be a continuous homomorphism between two ordered groups. We say that $`\psi `$ is order-preserving if $`\psi (P_1)P_2`$. Consequently, if $`\psi `$ is continuous and order preserving, then $`\psi (\overline{P_1})\overline{P_2}`$.
For each $`\alpha \mathrm{\Pi }`$, let $`\pi _\alpha `$ denote the quotient homomorphism $`\mathrm{\Gamma }\mathrm{\Gamma }/C_\alpha `$. Because $`C_\alpha `$ is a principal subgroup, we can define an order on $`\mathrm{\Gamma }/C_\alpha `$ by setting $`\psi _\alpha (\chi )0\chi 0`$. Moreover, the principal convex subgroups in $`\mathrm{\Gamma }/C_\alpha `$ are precisely the images by $`\pi _\alpha `$ of the principal convex subgroups of $`\mathrm{\Gamma }`$ containing $`C_\alpha `$. (See \[1, Section 2\].)
We end this section with a useful property of orders.
###### Proposition 2.3
Let $`P`$ be a measurable order on $`\mathrm{\Gamma }`$. Then $`\overline{P}`$ is a $`𝒯`$-set.
Proof. If $`\mathrm{\Gamma }`$ is discrete, there is nothing to prove. If $`\mathrm{\Gamma }`$ is not discrete, the subgroup $`C_{\alpha _0}`$ is open and nonempty. Hence the set $`C_{\alpha _0}\{\chi \mathrm{\Gamma }:\psi _{\alpha _0}(\chi )>0\}`$ is nonempty, with $`0`$ as a limit point. Given an open nonempty neighborhood $`U`$ of $`0`$, let
$$W=UC_{\alpha _0}\{\chi \mathrm{\Gamma }:\psi _{\alpha _0}(\chi )>0\}.$$
Then $`W`$ is a nonempty subset of $`UP`$. Moreover, it is easy to see that $`W+\overline{P}P\overline{P}`$, and hence $`\overline{P}`$ is a $`𝒯`$-set.
## 3 Analyticity
We continue with the notation of the previous section. Using the order structure on $`\mathrm{\Gamma }`$ we define some classes of analytic functions on $`G`$:
(14) $`H^1(G)`$ $`=`$ $`\{fL^1(G):\widehat{f}=0\mathrm{on}(P)\{0\}\};`$
(15) $`H_0^1(G)`$ $`=`$ $`\{fL^1(G):\widehat{f}=0\mathrm{on}P\};`$
and
(16)
$$H^{\mathrm{}}(G)=\{fL^{\mathrm{}}(G):_Gf(x)g(x)𝑑x=0\mathrm{for}\mathrm{all}gH_0^1(G)\}.$$
We clearly have
$$H^1(G)=\{fL^1(G):\widehat{f}=0\mathrm{on}\overline{(P)\{0\}}\}.$$
We can now give the definition of analytic measures in $`_T(\mathrm{\Sigma })`$.
###### Definition 3.1
Let $`T`$ be a sup path attaining representation of $`G`$ by isomorphisms of $`M(\mathrm{\Sigma })`$. A measure $`\mu _T(\mathrm{\Sigma })`$ is called weakly analytic if the mapping $`tT_t\mu (A)`$ is in $`H^{\mathrm{}}(G)`$ for every $`A\mathrm{\Sigma }`$.
###### Definition 3.2
Recall the $`T`$-spectrum of a weakly measurable $`\mu _T(\mathrm{\Sigma })`$,
(17)
$$\mathrm{spec}_T(\mu )=\underset{f(\mu )}{}\{\chi \mathrm{\Gamma }:\widehat{f}(\chi )=0\}.$$
A measure $`\mu `$ in $`_T(\mathrm{\Sigma })`$ is called $`T`$-analytic if $`\mathrm{spec}_T(\mu )\overline{P}`$.
That the two definitions of analyticity are equivalent will be shown later in this section.
Since $`(\mu )`$ is translation-invariant, it follows readily that for all $`tG`$,
$$(T_t\mu )=(\mu ),$$
and hence
(18)
$$\mathrm{spec}_T(T_t(\mu ))=\mathrm{spec}_T(\mu ).$$
We now recall several basic results from spectral theory of bounded functions that will be needed in the sequel. Our reference is \[21, Section 40\]. If $`\varphi `$ is in $`L^{\mathrm{}}(G)`$, write $`\left[\varphi \right]`$ for the smallest weak-\* closed translation-invariant subspace of $`L^{\mathrm{}}(G)`$ containing $`\varphi `$, and let $`([\varphi ])=(\varphi )`$ denote the closed translation-invariant ideal in $`L^1(G)`$:
$$(\varphi )=\{fL^1(G):f\varphi =0\}.$$
It is clear that $`(\varphi )=\{fL^1(G):fg=0,g\left[\varphi \right]\}`$. The spectrum of $`\varphi `$, denoted by $`\sigma \left[\varphi \right]`$, is the set of all continuous characters of $`G`$ that belong to $`\left[\varphi \right]`$. This closed subset of $`\mathrm{\Gamma }`$ is also given by
(19)
$$\sigma \left[\varphi \right]=Z((\varphi )).$$
(See \[21, Theorem (40.5)\].)
Recall that a closed subset $`E`$ of $`\mathrm{\Gamma }`$ is a set of spectral synthesis for $`L^1(G)`$, or an $`S`$-set, if and only if $`([E])`$ is the only ideal in $`L^1(G)`$ whose zero set is $`E`$.
There are various equivalent definitions of $`S`$-sets. Here is one that we will use at several occasions.
A set $`E\mathrm{\Gamma }`$ is an $`S`$-set if and only if every essentially bounded function $`g`$ in $`L^{\mathrm{}}(G)`$ with $`\sigma [g]E`$ is the weak-\* limit of linear combinations of characters from $`E`$.
(See \[21, (40.23) (a)\].) This has the following immediate consequence.
###### Proposition 3.3
Suppose that $`B`$ is an $`S`$-set, $`gL^{\mathrm{}}(G)`$, and $`\mathrm{spec}(g)B`$. (i) If $`f`$ is in $`L^1(G)`$ and $`\widehat{f}=0`$ on $`B`$, then $`fg(x)=0`$ for all $`x`$ in $`G`$. In particular,
$$_Gf(x)g(x)𝑑x=0.$$
(ii) If $`\mu `$ is a measure in $`M(G)`$ with $`\widehat{\mu }=0`$ on $`B`$, then $`\mu g(x)=0`$ for almost all $`x`$ in $`G`$.
Proof. Part (i) is a simple consequence of \[21, Theorems (40.8) and (40.10)\]. We give a proof for the sake of completeness. Write $`g`$ as the weak-\* limit of trigonometric polynomials, $`_{\chi E}a_\chi \chi (x)`$, with characters in $`E`$. Then
$`{\displaystyle _G}f(x)g(yx)𝑑x`$ $`=`$ $`lim{\displaystyle _G}{\displaystyle \underset{\chi E}{}}a_\chi \chi (y)f(x)\chi (x)dx`$
$`=`$ $`lim{\displaystyle \underset{\chi E}{}}a_\chi \chi (y)\widehat{f}(\chi )=0`$
since $`\widehat{f}`$ vanishes on $`E`$.
To prove (ii), assume that $`\mu g`$ is not 0 a.e.. Then, there is $`f`$ in $`L^1(G)`$ such that $`f(\mu g)`$ is not 0 a.e.. But this contradicts (i), since $`f(\mu g)=(f\mu )g`$, $`f\mu `$ is in $`L^1(G)`$, and $`\widehat{f\mu }=0`$ on $`B`$.
The following is a converse of sorts of Proposition 3.3 and follows easily from definitions.
###### Proposition 3.4
Let $`B`$ be a nonvoid closed subset of $`\mathrm{\Gamma }`$. Suppose that $`f`$ is in $`L^{\mathrm{}}(G)`$ and
(20)
$$_Gf(x)g(x)𝑑x=0$$
for all $`g`$ in $`L^1(G)`$ such that $`\widehat{g}=0`$ on $`B`$. Then $`\sigma [f]B`$.
Proof. Let $`\chi _0`$ be any element in $`\mathrm{\Gamma }B`$. We will show that $`\chi _0`$ is not in the spectrum of $`f`$ by constructing a function $`h`$ in $`L^1(G)`$ with $`\widehat{h}(\chi _0)0`$ and $`hf=0`$. Let $`U`$ be an open neighborhood of $`\chi _0`$ not intersecting $`B`$, and let $`h`$ be in $`L^1(G)`$ such that $`\widehat{h}`$ is equal to 1 at $`\chi _0`$ and to 0 outside $`U`$. Direct computations show that the Fourier transform of the function $`g:th(xt)`$, when evaluated at $`\chi \mathrm{\Gamma }`$, gives $`\overline{\chi (x)}\widehat{h}(\chi )`$, and hence it vanishes on $`B`$. It follows from (20) that $`hf=0`$, which completes the proof.
A certain class of $`S`$-sets, known as the Calderón sets, or $`C`$-sets, is particularly useful to us. These are defined as follows. A subset $`E`$ of $`\mathrm{\Gamma }`$ is called a $`C`$-set if every $`f`$ in $`L^1(G)`$ with Fourier transform vanishing on $`E`$ can be approximated in the $`L^1`$-norm by functions of the form $`hf`$ where $`hL^1(G)`$ and $`\widehat{h}`$ vanishes on an open set containing $`E`$.
$`C`$-sets enjoy the following properties (see \[21, (39.39)\] or \[22, Section 7.5\]).
* Every $`C`$-set is an $`S`$-set.
* Every closed subgroup of $`\mathrm{\Gamma }`$ is a $`C`$-set.
* The empty set is a $`C`$-set.
* If the boundary of a set $`A`$ is a $`C`$-set, then $`A`$ is a $`C`$-set.
* Finite unions of $`C`$-sets are $`C`$-sets.
Since closed subgroups are $`C`$-sets, we conclude that $`\overline{P}\overline{(P)}`$, and $`C_\alpha `$, for all $`\alpha `$, are $`C`$-sets. ¿From the definition of $`S_{\alpha _0}`$, (13), and the fact that $`C_{\alpha _0}`$ is open and closed, it follows that the boundary of $`S_{\alpha _0}`$ is the closed subgroup $`\psi _{\alpha _0}^1(0)C_{\alpha _0}`$. Hence $`S_{\alpha _0}`$ is a $`C`$-set. For $`\alpha \alpha _0`$, the set $`S_\alpha `$ is open and closed, and so it has empty boundary, and thus it is a $`C`$-set. Likewise $`C_\alpha D_\alpha `$ is a $`C`$-set for all $`\alpha \alpha _0`$. ¿From this we conclude that arbitrary unions of $`S_\alpha `$ and $`C_\alpha D_\alpha `$ are $`C`$-sets, because an arbitrary union of such sets, not including the index $`\alpha _0`$, is open and closed, and so it is a $`C`$-set.
We summarize our findings as follows.
###### Proposition 3.5
Suppose that $`P`$ is a measurable order on $`\mathrm{\Gamma }`$. We have:
(i) $`\overline{P}`$ and $`\overline{(P)}`$ are $`C`$-sets;
(ii) $`S_\alpha `$ is a $`C`$-set for all $`\alpha `$;
(iii) arbitrary unions of $`S_\alpha `$ and $`C_\alpha D_\alpha `$ are $`C`$-sets.
As an immediate application, we have the following characterizations.
###### Corollary 3.6
Suppose that $`f`$ is in $`L^{\mathrm{}}(G)`$, then
(i) $`\sigma [f]S_\alpha `$ if and only if $`_Gf(x)g(x)𝑑x=0`$ for all $`gL^1(G)`$ such that $`\widehat{g}=0`$ on $`S_\alpha `$;
(ii) $`\sigma [f]\mathrm{\Gamma }C_\alpha `$ if and only if $`\mu _\alpha f=0`$;
(iii) $`\sigma [f]\overline{P}`$ if and only if $`fH^{\mathrm{}}(G)`$.
Proof. Assertions (i) and (iii) are clear from Propositions 3.5 and 3.4. To prove (ii), use Fubini’s Theorem to first establish that for any $`gL^1(G)`$, and any $`\mu M(G)`$, we have
$$_G(\mu f)(t)g(t)𝑑t=_Gf(t)(\mu g)(t)𝑑t.$$
Now suppose that $`\sigma [f]\mathrm{\Gamma }C_\alpha `$, and let $`g`$ be any function in $`L^1(G)`$. From Propositions 3.5 and 3.4, we have that $`_Gfg𝑑t=0`$ for all $`g`$ with Fourier transform vanishing on $`\mathrm{\Gamma }C_\alpha `$, equivalently, for all $`g=\mu _\alpha g`$. Hence, $`_Gf(\mu _\alpha g)𝑑t=_G(\mu _\alpha f)g𝑑t=0`$ for all $`g`$ in $`L^1(G)`$, from which it follows that $`\mu _\alpha f=0`$. The converse is proved similarly, and we omit the details.
Aiming for a characterization of weakly analytic measures in terms of their spectra, we present one more result.
###### Proposition 3.7
Let $`\mu `$ be weakly measurable in $`M(\mathrm{\Sigma })`$.
(i) Suppose that $`B`$ is a nonvoid closed subset of $`\mathrm{\Gamma }`$ and $`\mathrm{spec}_T\mu B`$. Then $`\sigma [tT_t\mu (A)]B`$ for all $`A\mathrm{\Sigma }`$.
(ii) Conversely, suppose that $`B`$ is an $`S`$-set in $`\mathrm{\Gamma }`$ and that $`\sigma [tT_t\mu (A)]B`$ for all $`A\mathrm{\Sigma }`$, then $`\mathrm{spec}_T\mu B`$.
Proof. We clearly have $`(\mu )([tT_t\mu (A)])`$. Hence, $`\mathrm{spec}_T\mu =Z((\mu ))Z(([tT_t\mu (A)]))=\sigma [tT_t\mu (A)],`$ and (i) follows.
Now suppose that $`B`$ is an $`S`$-set and let $`gL^1(G)`$ be such that $`\widehat{g}=0`$ on $`B`$. Then, for all $`A\mathrm{\Sigma }`$, we have from Proposition 3.4 (ii):
$$_Gg(t)T_t\mu (A)𝑑t=0.$$
Equivalently, we have that
$$_Gg(t)T_t\mu (A)𝑑t=0.$$
Since the Fourier transform of the function $`tg(t)`$ vanishes on $`B`$, we see that $`(\mu )\{f:\widehat{f}=0\mathrm{on}B\}`$. Thus $`Z((\mu ))Z(\{f:\widehat{f}=0\mathrm{on}B\})=B`$, which completes the proof.
Straightforward applications of Propositions 3.5 and 3.7 yield the desired characterization of weakly analytic measures.
###### Corollary 3.8
Suppose that $`\mu _T(\mathrm{\Sigma })`$. Then,
(i) $`\mu `$ is weakly $`T`$analytic if and only if $`\mathrm{spec}_T\mu \overline{P}`$ if and only if $`\sigma [tT_t\mu (A)]\overline{P}`$, for every $`A\mathrm{\Sigma }`$;
(ii) $`\mathrm{spec}_T\mu S_\alpha `$ if and only if $`\sigma [tT_t\mu (A)]S_\alpha `$ for every $`A\mathrm{\Sigma }`$.
(iii) $`\mathrm{spec}_T\mu C_\alpha `$ if and only if $`\sigma [tT_t\mu (A)]C_\alpha `$ for every $`A\mathrm{\Sigma }`$.
(iv) $`\mathrm{spec}_T\mu \mathrm{\Gamma }C_\alpha `$ if and only if $`\sigma [tT_t\mu (A)]\mathrm{\Gamma }C_\alpha `$ for every $`A\mathrm{\Sigma }`$.
The remaining results of this section are simple properties of measures in $`_T(\mathrm{\Sigma })`$ that will be needed later. Although the statements are direct analogues of classical facts about measures on groups, these generalization require in some places the sup path attaining property of $`T`$.
###### Proposition 3.9
Suppose that $`\mu _T(\mathrm{\Sigma })`$ and $`\nu M(G)`$. Then $`\mathrm{spec}_T\nu _T\mu `$ is contained in the support of $`\widehat{\nu }`$, and $`\mathrm{spec}_T\nu _T\mu \mathrm{spec}_T\mu `$.
Proof. Given $`\chi _0`$ not in the support of $`\widehat{\nu }`$, to conclude that it is also not in the spectrum of $`\nu _T\mu `$ it is enough to find a function $`f`$ in $`L^1(G)`$ with $`\widehat{f}(\chi _0)=1`$ and $`f_T(\nu _T\mu )=0`$. Simply choose $`f`$ with Fourier transform vanishing on the support of $`\widehat{\nu }`$ and taking value 1 at $`\chi _0`$. By Fourier inversion, we have $`f\nu =0`$, and since $`f_T(\nu _T\mu )=(f\nu )_T\mu `$, the first part of the proposition follows. For the second part, we have $`(\mu )(\nu _T\mu )`$, which implies the desired inclusion.
We next prove a property of $`L^{\mathrm{}}(G)`$ functions similar to the characterization of $`L^1`$ functions which are constant on cosets of a subgroup \[21, Theorem (28.55)\].
###### Proposition 3.10
Suppose that $`f`$ is in $`L^{\mathrm{}}(G)`$ and that $`\mathrm{\Lambda }`$ is an open subgroup of $`\mathrm{\Gamma }`$. Let $`\lambda _0`$ denote the normalized Haar measure on the compact group $`A(G,\mathrm{\Lambda })`$, the annihilator in $`G`$ of $`\mathrm{\Lambda }`$ (see \[20, (23.23)\]. Then, $`\sigma [f]\mathrm{\Lambda }`$ if and only if $`f=f\lambda _0`$ a. e. This is also the case if and only if $`f`$ is constant on cosets of $`A(G,\mathrm{\Lambda })`$.
Proof. Suppose that the spectrum of $`f`$ is contained in $`\mathrm{\Lambda }`$. Since $`\mathrm{\Lambda }`$ is an $`S`$-set, it follows that $`f`$ is the weak-\* limit of trigonometric polynomials with spectra contained in $`\mathrm{\Lambda }`$. Let $`\{f_\alpha \}`$ be a net of such trigonometric polynomials converging to $`f`$ weak-\*. Note that, for any $`\alpha `$, we have $`\lambda _0f_\alpha =f_\alpha `$. For $`g`$ in $`L^1(G)`$, we have
$$\underset{\alpha }{lim}_Gf_\alpha \overline{g}𝑑x=_Gf\overline{g}𝑑x.$$
In particular, we have
$$\underset{\alpha }{lim}_Gf_\alpha (\lambda _0\overline{g})𝑑x=_Gf(\lambda _0\overline{g})𝑑x,$$
and so
$$\underset{\alpha }{lim}_G(f_\alpha \lambda _0)\overline{g}𝑑x=_G(f\lambda _0)\overline{g}𝑑x.$$
Since this holds for any $`g`$ in $`L^1(G)`$, we conclude that $`\lambda _0f_\alpha `$ converges weak-\* to $`\lambda _0f`$. But $`\lambda _0f_\alpha =f_\alpha `$, and $`f_\alpha `$ converges weak-\* to $`f`$, hence $`f\lambda _0=f`$. The remaining assertions of the lemma are easy to prove. We omit the details.
In what follows, we use the symbol $`\mu _\alpha `$ to denote the normalized Haar measure on the compact subgroup $`A(G,C_\alpha )`$, the annihilator in $`G`$ of $`C_\alpha `$. This measure is also characterized by its Fourier transform:
$$\widehat{\mu _\alpha }=1_{C_\alpha }$$
(see \[20, (23.19)\]).
###### Corollary 3.11
Suppose that $`\mu _T(\mathrm{\Sigma })`$. Then,
(i) $`\mathrm{spec}_T\mu C_\alpha `$ if and only if $`\mu =\mu _\alpha _T\mu `$;
(ii) $`\mathrm{spec}_T\mu \mathrm{\Gamma }C_\alpha `$ if and only if $`\mu _\alpha _T\mu =0`$.
Proof. (i) If $`\mu =\mu _\alpha _T\mu `$, then, by Proposition 3.10, $`\sigma [t\mu _\alpha _TT_t\mu (A)]C_\alpha `$. Hence by Corollary 3.8, $`\mathrm{spec}_T\mu C_\alpha `$. For the other direction, suppose that $`\mathrm{spec}_T\mu C_\alpha `$. Then by Corollary 3.8 we have that the spectrum of the function $`tT_t\mu (A)`$ is contained in $`C_\alpha `$ for every $`A\mathrm{\Sigma }`$. By Proposition 3.10, we have that
$$T_t\mu (A)=_{G_\alpha }T_{ty}\mu (A)𝑑\mu _\alpha =T_t(\mu _\alpha \mu )(A)$$
for almost all $`tG`$. Since this holds for all $`A\mathrm{\Sigma }`$, the desired conclusion follows from Proposition 1.4.
Part (ii) follows from Corollary 3.6 (ii), Proposition 3.7(ii), and the fact that $`\mathrm{\Gamma }C_\alpha `$ is an $`S`$-set.
###### Corollary 3.12
Suppose that $`\mu _T(\mathrm{\Sigma })`$ and $`\mathrm{spec}_T\mu C_\alpha `$, and let $`yG_\alpha =A(G,C_\alpha )`$. Then $`T_y\mu =\mu `$.
Proof. For any $`A\mathrm{\Sigma }`$, we have from Corollary 3.11
$`T_y\mu (A)`$ $`=`$ $`T_y(\mu _\alpha \mu )(A)=\mu _\alpha T_y\mu (A)`$
$`=`$ $`{\displaystyle _{G_\alpha }}T_{yx}\mu (A)𝑑\mu _\alpha (y)`$
$`=`$ $`{\displaystyle _{G_\alpha }}T_x\mu (A)𝑑\mu _\alpha (y)=\mu _\alpha \mu (A)=\mu (A).`$
## 4 Homomorphism theorems
We continue with the notation of the previous section: $`G`$ is a locally compact abelian group, $`\mathrm{\Gamma }`$ the dual group of $`G`$, $`P`$ is a measurable order on $`\mathrm{\Gamma }`$, $`T`$ is a sup path attaining representation of $`G`$ acting on $`M(\mathrm{\Sigma })`$. Associated with $`P`$ is a collection of homomorphisms $`\psi _\alpha `$, as described by Theorem 2.1. Let $`\varphi _\alpha `$ denote the adjoint of $`\psi _\alpha `$. Thus, $`\varphi _\alpha `$ is a continuous homomorphism of $``$ into $`G`$. By composing the representation $`T`$ with the $`\varphi _\alpha `$, we define a new representation $`T_{\varphi _\alpha }`$ of $``$ acting on $`M(\mathrm{\Sigma })`$ by: $`tT_{\varphi _\alpha (t)}`$. If $`\mu `$ in $`M(\mathrm{\Sigma })`$ is weakly measurable with respect to $`T`$ then $`\mu `$ is also weakly measurable with respect to $`T_{\varphi _\alpha }`$. We will further suppose that $`T_{\varphi _\alpha }`$ is sup path attaining for each $`\alpha `$. This is the case with the representations of Example 1.7.
Our goal in this section is to relate the notion of analyticity with respect to $`T`$ to the notion of analyticity with respect to $`T_{\varphi _\alpha }`$. More generally, suppose that $`G_1`$ and $`G_2`$ are two locally compact abelian groups with dual groups $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$, respectively. Let
$$\psi :\mathrm{\Gamma }_1\mathrm{\Gamma }_2$$
be a continuous homomorphism, and let $`\varphi :G_2G_1`$ denote its adjoint homomorphism. Suppose $`\nu `$ is in $`M(G_2)`$. We define a Borel measure $`\mathrm{\Phi }(\nu )`$ in $`M(G_1)`$ on the Borel subsets $`A`$ of $`G_1`$ by:
(21)
$$\mathrm{\Phi }(\nu )(A)=_{G_2}1_A\varphi (t)𝑑\nu (t)=_{G_1}1_A𝑑\mathrm{\Phi }(\nu ),$$
where $`1_A`$ is the indicator function of $`A`$. We have $`\mathrm{\Phi }(g)_{M(G_1)}=\nu _{M(G_2)}`$ and, for every Borel measurable bounded function $`f`$ on $`G_1`$, we have
(22)
$$_{G_1}f𝑑\mathrm{\Phi }(\nu )=_{G_2}f\varphi (t)𝑑\nu (t).$$
In particular, if $`f=\chi `$, a character in $`\mathrm{\Gamma }`$, then
(23)
$$\widehat{\mathrm{\Phi }(\nu )}(\chi )=_{G_1}\overline{\chi }𝑑\mathrm{\Phi }(\nu )=_{G_2}\overline{\chi }\varphi (t)𝑑\nu (t)=_{G_2}\overline{\psi (\chi )}(t)𝑑\nu (t)=\widehat{\nu }(\psi (\chi )),$$
where $`\psi `$ is the adjoint homomorphism of $`\varphi `$. So,
(24)
$$\widehat{\mathrm{\Phi }(\nu )}=\widehat{\nu }\psi .$$
Our first result is a very useful fact from spectral synthesis of bounded functions. The proof uses in a crucial way the fact that the representation is sup path attaining, or, more precisely, satisfies the property in Proposition 1.4.
###### Lemma 4.1
Suppose that $`T`$ is a sup path attaining representation of $`G_1`$ acting on $`M(\mathrm{\Sigma })`$, $`\varphi `$ is a continuous homomorphism of $`G_2`$ into $`G_1`$ such that $`T_\varphi `$ is a sup path attaining representation of $`G_2`$. Suppose that $`B`$ is a nonempty closed $`S`$-subset of $`\mathrm{\Gamma }_1`$ and that $`\mu `$ is in $`M(\mathrm{\Sigma })`$ with $`\mathrm{spec}_T\mu B`$. Suppose that $`C`$ is an $`S`$-subset of $`\mathrm{\Gamma }_2`$ and $`\psi (B)C`$. Then $`\mathrm{spec}_{T_\varphi }\mu C`$.
Proof. Since $`C`$ is an $`S`$-subset of $`\mathrm{\Gamma }_2`$, it is enough to show that for every ‘$`A\mathrm{\Sigma }`$, $`\mathrm{spec}_{T_\varphi }(xT_{\varphi (x)}\mu (A))C`$, by Proposition 3.7. For this purpose, it is enough by \[21, Theorem (40.8)\], to show that
$$gT_{\varphi ()}\mu (A)=0$$
for every $`g`$ in $`L^1(G_2)`$ with $`\widehat{g}=0`$ on $`C`$. For $`rG_2`$ and $`xG_1`$, consider the measure
$$T_x(g_{T_\varphi }T_{\varphi (r)}\mu )=g_{T_\varphi }T_{x+\varphi (r)}\mu .$$
For $`A\mathrm{\Sigma }`$, we have
$`g_{T_\varphi }T_{x+\varphi (r)}\mu (A)`$ $`=`$ $`{\displaystyle _G}T_{t+x}(T_{\varphi (r)}\mu )(A)𝑑\mathrm{\Phi }(g)(t)`$
$`=`$ $`\mathrm{\Phi }(g)[tT_t(T_{\varphi (r)}\mu )(A)](x)`$
$`=`$ $`0`$
for almost all $`xG_1`$. To justify the last equality, we appeal to Proposition 3.3 and note that $`\widehat{\mathrm{\Phi }(g)}=\widehat{g}\psi `$ and so $`\widehat{\mathrm{\Phi }(g)}=0`$ on $`B\psi ^1(C)`$. Moreover, $`\sigma [tT_t(T_{\varphi (r)}\mu )(A)]\mathrm{spec}_T(\mu )B`$. Now, using Proposition 1.4 and the fact that, for every $`A\mathrm{\Sigma }`$,
$$T_x[g_{T_\varphi }T_{\varphi (r)}]\mu (A)=g_{T_\varphi }T_{x+\varphi (r)}\mu (A)=0,$$
for almost all $`xG_1`$, we conclude that the measure $`g_{T_\varphi }T_{\varphi (r)}\mu `$ is the zero measure, which completes the proof.
Given $`𝒞`$, a collection of elements in $`L^1(G_1)`$ or $`M(G_1)`$, let
$$Z(𝒞)=\underset{\delta 𝒞}{}\{\chi :\widehat{\delta }(\chi )=0\}.$$
This is the same notation for the zero set of an ideal in $`L^1(G)`$ that we introduced in Section 1. Given a set of measures $`𝒮`$ in $`M(G_2)`$, let
$$\mathrm{\Phi }(𝒮)=\{\mathrm{\Phi }(\nu ):\nu 𝒮\}M(G_1).$$
###### Lemma 4.2
In the above notation, if $`\mu M(\mathrm{\Sigma })`$ is weakly measurable, then
$$Z\left(\mathrm{\Phi }(_{T_\varphi }\mu )\right)=\psi ^1\left(Z(_{T_\varphi }\mu )\right)=\psi ^1\left(\mathrm{spec}_{T_\varphi }\mu \right).$$
Proof. It is enough to establish the first equality; the second one follows from definitions. We have
$`Z\left(\mathrm{\Phi }(_{T_\varphi }\mu )\right)`$ $`=`$ $`{\displaystyle \underset{\delta \mathrm{\Phi }(_{T_\varphi }(\mu ))}{}}\{\chi \mathrm{\Gamma }:\widehat{\delta }(\chi )=0\}`$
$`=`$ $`{\displaystyle \underset{g_{T_\varphi }(\mu )}{}}\{\chi \mathrm{\Gamma }:\widehat{\mathrm{\Phi }(g)}(\chi )=0\}`$
$`=`$ $`{\displaystyle \underset{g_{T_\varphi }(\mu )}{}}\{\chi \mathrm{\Gamma }:\widehat{g}(\psi (\chi ))=0\}`$
$`=`$ $`{\displaystyle \underset{g_{T_\varphi }(\mu )}{}}\psi ^1\left(Z(g)\right)`$
$`=`$ $`\psi ^1\left({\displaystyle \underset{g_{T_\varphi }(\mu )}{}}\left(Z(g)\right)\right)`$
$`=`$ $`\psi ^1\left(Z\left(_{T_\varphi }(\mu )\right)\right)=\psi ^1\left(\mathrm{spec}_{T_\varphi }\mu \right).`$
###### Lemma 4.3
Suppose that $`C`$ is a nonempty closed $`S`$-subset of $`\mathrm{\Gamma }_2`$ and that $`\psi ^1(C)`$ is an $`S`$-subset of $`\mathrm{\Gamma }_1`$. Suppose that $`\mu `$ is in $`M(\mathrm{\Sigma })`$ and $`\mathrm{spec}_{T_\varphi }(\mu )C`$. Then $`\mathrm{spec}_T\mu \psi ^1(C)`$.
Proof. We will use the notation of Lemma 4.2. If $`f_{T_\varphi }(\mu )`$ and $`tG_1`$, then $`f_{T_\varphi }(T_t\mu )`$. So, for $`A\mathrm{\Sigma }`$, we have $`f_{T_\varphi }(T_t\mu )(A)=0`$. But
$`f_{T_\varphi }(T_t\mu )(A)`$ $`=`$ $`{\displaystyle _{}}T_{t\varphi (x)}\mu (A)f(x)𝑑x`$
$`=`$ $`{\displaystyle _G}T_{tx}\mu (A)𝑑\mathrm{\Phi }(f),`$
where $`\mathrm{\Phi }(f)`$ is the homomorphic image of the measure $`f(x)dx`$. Hence, $`\mathrm{\Phi }(f)_T([tT_t\mu (A)])`$, and so $`\mathrm{\Phi }(_{T_\varphi }(\mu ))_T([tT_t\mu (A)])`$, which implies that
$$Z\left(\mathrm{\Phi }(_{T_\varphi }(\mu ))\right)Z\left(_T([tT_t\mu (A)])\right)=\mathrm{spec}_T(tT_t\mu (A)).$$
By Lemma 4.2,
$$Z\left(\mathrm{\Phi }(_{T_\varphi }(\mu ))\right)=\psi ^1\left(\mathrm{spec}_{T_\varphi }\mu \right)\psi ^1(C).$$
Hence, $`\mathrm{spec}_T(tT_t\mu (A))\psi ^1(C)`$ for all $`A\mathrm{\Sigma }`$, which by Proposition 3.7 implies that $`\mathrm{spec}_T(\mu )\psi ^1(C)`$.
Taking $`G_1=G,G_2=`$ and $`\psi =\psi _\alpha `$ to be one of the homomorphisms in Theorem 2.1, and using the fact that $`[0,\mathrm{}[`$, $`S_\alpha `$, $`C_\alpha D_\alpha `$ are all $`S`$-sets, we obtain useful relationships between different types of analyticity.
###### Theorem 4.4
Let $`G`$ be a locally compact abelian group with ordered dual group $`\mathrm{\Gamma }`$, and let $`P`$ denote a measurable order on $`\mathrm{\Gamma }`$. Suppose that $`T`$ is a sup path attaining representation of $`G`$ by isomorphisms of $`M(\mathrm{\Sigma })`$, such that $`T_{\varphi _\alpha }`$ is sup path attaining, where $`\varphi _\alpha `$ is as in Theorem 2.1.
(i) If $`\mu M(\mathrm{\Sigma })`$ and $`\mathrm{spec}_T(\mu )C_\alpha D_\alpha `$. Then
$$\mathrm{spec}_T(\mu )S_\alpha \mathrm{spec}_{T_{\varphi _\alpha }}(\mu )[0,\mathrm{}[.$$
(ii) If $`\mu M(\mathrm{\Sigma })`$ and $`\mathrm{spec}_T(\mu )C_{\alpha _0}`$. Then
$$\mathrm{spec}_T(\mu )S_{\alpha _0}\mathrm{spec}_{T_{\varphi _{\alpha _0}}}(\mu )[0,\mathrm{}[.$$
We can use the representation $`T_\varphi `$ to convolve a measure $`\nu M(G_2)`$ with $`\mu M(G_1)`$:
$$\nu _{T_\varphi }\mu (A)=_{G_2}T_{\varphi (x)}\mu (A)𝑑\nu (x)=_{G_2}\mu (A\varphi (x))𝑑\nu (x),$$
for all Borel $`AG_1`$.
Alternatively, we can convolve $`\mathrm{\Phi }(\nu )`$ in the usual sense of \[20, Definition 19.8\] with $`\mu `$ to yield another measure in $`M(G_1)`$, defined on the Borel subsets of $`G_1`$ by
$$\mathrm{\Phi }(\nu )\mu (A)=_{G_1}_{G_1}1_A(x+y)𝑑\mathrm{\Phi }(\nu )(x)𝑑\mu (y).$$
Using (22), we find that
$`\mathrm{\Phi }(\nu )\mu (A)`$ $`=`$ $`{\displaystyle _{G_1}}{\displaystyle _{G_2}}1_A(\varphi (t)+y)𝑑\nu (t)𝑑\mu (y)`$
$`=`$ $`{\displaystyle _{G_2}}\mu (A\varphi (t))𝑑\nu (t)=\nu _{T_\varphi }\mu (A).`$
Thus,
(25)
$$\mathrm{\Phi }(\nu )\mu =\nu _{T_\varphi }\mu .$$
We end the section with homomorphism theorems, which complement the well-known homomorphism theorems for $`L^p`$-multipliers (see Edwards and Gaudry \[11, Appendix B\]). In these theorems, we let $`G_1`$ act on $`M(G_1)`$ by translation. That is, if $`\mu M(G_1)`$, $`xG_1`$, and $`A`$ is a Borel subset of $`G_1`$, then
$$T_x\mu (A)=\mu (A+x).$$
Let $`\varphi :G_2G_1`$ be a continuous homomorphism. By Example 1.3, $`T`$ and $`T_\varphi `$ are sup path attaining. (Recall that if $`tG_2`$, $`\mu M(G_1)`$, then $`T_{\varphi (t)}\mu (A)=\mu (A+\varphi (t))`$.) A simple exercise with definitions shows that for $`\mu M(G_1)`$
$$\mathrm{spec}_T\mu =\mathrm{supp}\widehat{\mu }.$$
###### Theorem 4.5
Suppose that $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ contain measurable orders $`P_1`$ and $`P_2`$, respectively, and $`\psi :\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ is a continuous, order-preserving homomorphism (that is, $`\psi (\overline{P_1})\overline{P_2}`$). Suppose that there is a positive constant $`N(\nu )`$ such that for all $`fH^1(G_2)`$
(26)
$$\nu f_1N(\nu )f_1.$$
Then
(27)
$$\mathrm{\Phi }(\nu )\mu N(\nu )\mu $$
for all Borel measures in $`M(G_1)`$ such that $`\widehat{\mu }`$ is supported in $`\overline{P_1}`$.
Proof. We have $`\mathrm{\Phi }(\nu )\mu =\nu _{T_\varphi }\mu `$. Also $`\overline{P_2}`$ is a $`𝒯`$-set. So (27) will follow from Theorem 1.8 once we show that $`\mathrm{spec}_{T_\varphi }\mu \overline{P_2}`$. For that purpose, we use Lemma 4.1. We have
$$\mathrm{spec}_T\mu =\mathrm{supp}\widehat{\mu }\overline{P_1},$$
and $`\psi (\overline{P_1})\overline{P_2}`$ is an $`S`$-set. Hence $`\mathrm{spec}_{T_\varphi }\mu \overline{P_2}`$ by Lemma 4.1.
The following special case of Theorem 4.5 deserves a separate statement.
###### Theorem 4.6
With the above notation, suppose that there is a positive constant $`N(\nu )`$ such that for all $`fH^1(G_2)`$
(28)
$$\nu f_1N(\nu )f_1.$$
Then for all $`fH^1(G_1)`$ we have
(29)
$$\mathrm{\Phi }(\nu )f_1N(\nu )f_1.$$
###### Theorem 4.7
Suppose that there is a positive constant $`N(\nu )`$ such that for all $`fH^1()`$
(30)
$$\nu f_1N(\nu )f_1.$$
Then for all $`\mu M(G)`$ with support of $`\widehat{\mu }`$ contained in $`C_\alpha D_\alpha `$, where $`\alpha <\alpha _0`$, we have
(31)
$$\mathrm{\Phi }_\alpha (\nu )\mu _1N(\nu )\mu .$$
Proof. The proof is very much like the proof of Theorem 4.5. We have $`\mathrm{\Phi }_\alpha (\nu )\mu =\nu _{T_\varphi }\mu `$. Apply Theorem 1.8, taking into consideration that
$$\mathrm{spec}_T\mu =\mathrm{supp}\widehat{\mu }C_\alpha D_\alpha $$
is an $`S`$-set and so
$$\mathrm{spec}_{T_{\varphi _\alpha }}\mu \psi _\alpha (C_\alpha D_\alpha )[0,\mathrm{}[.$$
## 5 Decomposition of Analytic Measures
Define measures $`\mu _{\alpha _0}`$ and $`d_\alpha `$ by their Fourier transforms: $`\widehat{\mu _{\alpha _0}}=1_{C_{\alpha _0}}`$, and $`\widehat{d_\alpha }=1_{C_\alpha D_\alpha }`$. Then we have the following decomposition theorem.
###### Theorem 5.1
Let $`G`$ be a locally compact abelian group with an ordered dual group $`\mathrm{\Gamma }`$. Suppose that $`T`$ is a sup path attaining representation of $`G`$ in $`M(\mathrm{\Sigma })`$. Then for any weakly analytic measure $`\mu M(\mathrm{\Sigma })`$ we have that the set of $`\alpha `$ for which $`d_\alpha _T\mu 0`$ is countable, and that
(32)
$$\mu =\mu _{\alpha _0}_T\mu +\underset{\alpha }{}d_\alpha _T\mu ,$$
where the right side converges unconditionally in norm in $`M(\mathrm{\Sigma })`$. Furthermore, there is a positive constant $`c`$, depending only upon $`T`$, such that for any signs $`ϵ_\alpha =\pm 1`$ we have
(33)
$$\underset{\alpha }{}ϵ_\alpha d_\alpha _T\mu c\mu .$$
One should compare this theorem to the well-known results from Littlewood-Paley theory on $`L^p(G)`$, where $`1<p<\mathrm{}`$ (see Edwards and Gaudry ). For $`L^p(G)`$ with $`1<p<\mathrm{}`$, it is well-known that the subgroups $`(C_\alpha )`$ form a Littlewood-Paley decomposition of the group $`\mathrm{\Gamma }`$, which means that the martingale difference series
$$f=\mu _{\alpha _0}f+\underset{\alpha }{}d_\alpha f$$
converges unconditionally in $`L^p(G)`$ to $`f`$. Thus, Theorem 5.1 above may be considered as an extension of Littlewood-Paley Theory to spaces of analytic measures.
The next result, crucial to our proof of Theorem 5.1, is already known in the case that $`G=𝕋^n`$ with the lexicographic order on the dual. This is due to Garling , and is a modification of the celebrated inequalities of Burkholder. Our result can be obtained directly from the result in by combining the techniques of with the homomorphism theorem 4.5. However, we shall take a different approach, in effect reproducing Garling’s proof in this more general setting.
###### Theorem 5.2
Suppose that $`G`$ is a locally compact group with ordered dual $`\mathrm{\Gamma }`$. Then for $`fH^1(G)`$, for any set $`\{\alpha _j\}_{j=1}^n`$ of indices less than $`\alpha _0`$, and for any numbers $`ϵ_j\{0,\pm 1\}`$ ($`1jn`$), there is an absolute constant $`a>0`$ such that
(34)
$$\underset{j=1}{\overset{n}{}}ϵ_jd_{\alpha _j}f_1af_1.$$
Furthermore,
(35)
$$f=\mu _{\alpha _0}f+\underset{\alpha }{}d_\alpha f,$$
where the right hand side converges unconditionally in the norm topology on $`H^1(G)`$.
Proof. The second part of Theorem 5.2 follows easily from the first part and Fourier inversion.
Now let us show that if we have the result for compact $`G`$, then we have it for locally compact $`G`$. Let $`\pi _{\alpha _0}:\mathrm{\Gamma }\mathrm{\Gamma }/C_{\alpha _0}`$ denote the quotient homomorphism of $`\mathrm{\Gamma }`$ onto the discrete group $`\mathrm{\Gamma }/C_{\alpha _0}`$ (recall that $`C_{\alpha _0}`$ is open), and define a measurable order on $`\mathrm{\Gamma }/C_{\alpha _0}`$ to be $`\pi _{\alpha _0}(P)`$. By Remarks 2.2 (c), the decomposition of the group $`\mathrm{\Gamma }/C_{\alpha _0}`$ that we get by applying Theorem 2.1 to that group, is precisely the image by $`\pi _{\alpha _0}`$ of the decomposition of the group $`\mathrm{\Gamma }`$. Let $`G_0`$ denote the compact dual group of $`\mathrm{\Gamma }/C_{\alpha _0}`$. Thus if Theorem 5.2 holds for $`H^1(G_0)`$, then applying Theorem 4.5, we see that Theorem 5.2 holds for $`G`$.
Henceforth, let us suppose that $`G`$ is compact. We will suppose that the Haar measure on $`G`$ is normalized, so that $`G`$ with Haar measure is a probability space.
Since each one of the subgroups $`C_\alpha `$, and $`D_\alpha `$ ($`\alpha <\alpha _0`$) is open, it follows that their annihilators in $`G`$, $`G_\alpha =A(G,C_\alpha )`$, and $`A(G,D_\alpha )`$, are compact. Let $`\mu _\alpha `$ and $`\nu _\alpha `$ denote the normalized Haar measures on $`A(G,C_\alpha )`$ and $`A(G,D_\alpha )`$, respectively. We have $`\widehat{\mu }_\alpha =1_{C_\alpha }`$ (for all $`\alpha `$), and $`\widehat{\nu }_\alpha =1_{D_\alpha }`$ (for all $`\alpha \alpha _0`$), so that $`d_\alpha =\mu _\alpha \nu _\alpha `$.
For each $`\alpha `$, let $`_\alpha `$ denote the $`\sigma `$-algebra of subsets of $`G`$ of the form $`A+G_\alpha `$, where $`A`$ is a Borel subset of $`G`$. We have $`_{\alpha _1}_{\alpha _2}`$, whenever $`\alpha _1>\alpha _2`$. It is a simple matter to see that for $`fL^1(G)`$, the conditional expectation of $`f`$ with respect to $`_\alpha `$ is equal to $`\mu _\alpha f`$ (see \[11, Chapter 5, Section 2\]).
We may suppose without loss of generality that $`\alpha _1>\alpha _2>\mathrm{}>\alpha _n`$. Thus the $`\sigma `$-algebras $`_{\alpha _k}`$ form a filtration, and the sequence $`(d_{\alpha _1}f,d_{\alpha _2}f,\mathrm{},d_{\alpha _n}f)`$ is a martingale difference sequence with respect to this filtration.
In that case, we have the following result of Burkholder \[7, Inequality (1.7)\], and . If $`0<p<\mathrm{}`$, then there is a positive constant $`c`$, depending only upon $`p`$, such that
(36)
$$\underset{1kn}{sup}\left(\underset{j=1}{\overset{k}{}}ϵ_jd_{\alpha _j}f\right)_pc\underset{1kn}{sup}\left(\underset{j=1}{\overset{k}{}}d_{\alpha _j}f\right)_p.$$
###### Lemma 5.3
For any index $`\alpha `$, $`0<p<\mathrm{}`$, and $`fH^1(G)L^p(G)`$, we have almost everywhere on $`G`$
(37)
$$\left|\mu _\alpha f\right|^p\mu _\alpha \left|f\right|^p,$$
where $`\mu _\alpha `$ is the normalized Haar measure on the compact subgroup $`G_\alpha =A(G,C_\alpha )`$.
Proof. The dual group of $`G_\alpha `$ is $`\mathrm{\Gamma }/C_\alpha `$ and can be ordered by the set $`\pi _\alpha (P)`$, where $`\pi _\alpha `$ is the natural homomorphism of $`\mathrm{\Gamma }`$ onto $`\mathrm{\Gamma }/C_\alpha `$.
Next, by convolving with an approximate identity for $`L^1(G)`$ consisting of trigonometric polynomials, we may assume that $`f`$ is a trigonometric polynomial. Then we see that for each $`xG`$ that the function $`yf(x+y)`$, $`yG_\alpha `$, is in $`H^1(G_\alpha )`$. To verify this, it is sufficient to consider the case when $`f`$ is a character in $`P`$. Then
$$f(x+y)=f(x)\pi _\alpha (f)(y),$$
and by definition $`\pi _\alpha (f)`$ is in $`H^1(G_\alpha )`$.
Now we have the following generalization of Jensen’s Inequality, due to Helson and Lowdenslager \[16, Theorem 2\]. An independent proof based on the ideas of this section is given in . For all $`gH^1(G)`$
(38)
$$\left|_Gg(x)𝑑x\right|\mathrm{exp}_G\mathrm{log}|g(x)|dx.$$
Apply (38) to $`yf(x+y)`$, $`yG_\alpha `$ to obtain
$$\left|_{G_\alpha }f(x+y)𝑑\mu _\alpha (y)\right|\mathrm{exp}_{G_\alpha }\mathrm{log}|f(x+y)|d\mu _\alpha (y).$$
Extending the integrals to the whole of $`G`$ (since $`\mu _\alpha `$ is supported on $`G_\alpha `$), raising both sides to the $`p`$th power, and then applying the usual Jensen’s inequality for the logarithmic function on finite measure spaces, we obtain
$`\left|{\displaystyle _G}f(x+y)𝑑\mu _\alpha (y)\right|^p`$ $``$ $`\mathrm{exp}{\displaystyle _G}\mathrm{log}|f(x+y)|^pd\mu _\alpha (y)`$
$``$ $`{\displaystyle _G}|f(x+y)|^p𝑑\mu _\alpha (y).`$
Changing $`y`$ to $`y`$, we obtain the desired inequality.
Let us continue with the proof of Theorem 5.2. We may suppose that $`f`$ is a mean zero trigonometric polynomial, and that the spectrum of $`f`$ is contained in $`_{j=1}^nC_{\alpha _j}D_{\alpha _j}`$, that is to say
$$f=\underset{j=1}{\overset{n}{}}d_{\alpha _j}f.$$
By Lemma 5.3, we have that
(39) $`\underset{1kn}{sup}\left|\mu _{\alpha _k}f\right|`$ $`=`$ $`\left(\underset{1kn}{sup}\left|\mu _{\alpha _k}f\right|^{1/2}\right)^2`$
$``$ $`\left(\underset{1kn}{sup}\mu _{\alpha _k}|f|^{1/2}\right)^2.`$
Also, we have that $`(\mu _{\alpha _j}|f|^{1/2})_{j=1}^n`$ is a martingale with respect to the filtration $`(_j)_{j=1}^n`$. Hence, by Doob’s Maximal Inequality \[10, Theorem (3.1), p. 317\] we have that
(40) $`\underset{1kn^{}}{sup}\mu _{\beta _k}|f|^{1/2}_2^2`$ $``$ $`4\mu _{\beta _n^{}}|f|^{1/2}_2^2`$
$``$ $`4|f|^{1/2}_2^2=4f_1.`$
The desired inequality follows now upon combining Burkholder’s Inequality (36) with (39), and (40).
Proof of Theorem 5.1. Transferring inequality (34) by using Theorem 1.8, we obtain that for any set $`\{\alpha _j\}_{j=1}^n`$ of indices less than $`\alpha _0`$, and for any numbers $`ϵ_j\{0,\pm 1\}`$ ($`1jn`$), there is a positive constant $`c`$, depending only upon the representation $`T`$, such that
(41)
$$\underset{j=1}{\overset{n}{}}ϵ_jd_{\alpha _j}_T\mu c\mu .$$
Now suppose that $`\{\alpha _j\}_{j=1}^{\mathrm{}}`$ is a countable collection of indices less than $`\alpha _0`$. Then by Bessaga and Pełczyński , the series $`_{j=1}^{\mathrm{}}d_{\alpha _j}_T\mu `$ is unconditionally convergent. In particular, for any $`\delta >0`$, for only finitely many $`k`$ do we have that $`d_{\alpha _k}_T\mu >\delta `$. Since this is true for all such countable sets, we deduce that the set of $`\alpha `$ for which $`d_\alpha _T\mu 0`$ is countable.
Hence we have that $`_\alpha d_\alpha _T\mu `$ is unconditionally convergent to some measure, say $`\nu `$. Clearly $`\nu `$ is weakly measurable. To prove that $`\mu =\nu `$, it is enough by Proposition 1.4 to show that for every $`A\mathrm{\Sigma }`$, we have $`T_t\mu (A)=T_t\nu (A)`$ for almost all $`tG`$.
We first note that since for every $`fL^1(G)`$ the series $`\mu _{\alpha _0}f+_\alpha d_\alpha f`$ converges to $`f`$ in $`L^1(G)`$, it follows that, for every $`gL^{\mathrm{}}(G)`$, the series $`\mu _{\alpha _0}g+_\alpha d_\alpha g`$ converges to $`g`$ in the weak-\* topology of $`L^{\mathrm{}}(G)`$. Now on the one hand, for $`tG`$ and $`A\mathrm{\Sigma }`$, we have $`\mu _{\alpha _0}_TT_t\mu (A)+_\alpha d_\alpha _TT_t\mu (A)=T_t\nu (A)`$, because of the (unconditional) convergence of the series $`\mu _{\alpha _0}_T\mu +_\alpha d_\alpha _T\mu `$ to $`\nu `$. On the other hand, by considering the $`L^{\mathrm{}}(G)`$ function $`tT_t(A)`$, we have that $`\mu _{\alpha _0}_TT_t\mu (A)+_\alpha d_\alpha _TT_t\mu (A)=\mu _{\alpha _0}T_t\mu (A)+_\alpha d_\alpha T_t\mu (A)=T_t\mu (A)`$, weak \*. Thus $`T_t\mu (A)=T_t\nu (A)`$ for almost all $`tG`$, and the proof is complete.
## 6 Generalized F. and M. Riesz Theorems
Throughout this section, we adopt the notation of Section 5, specifically, the notation and assumptions of Theorem 5.1.
Suppose that $`T`$ is a sup path attaining representation of $``$ by isomorphisms of $`M(\mathrm{\Sigma })`$. In , we proved the following result concerning bounded operators $`𝒫`$ from $`M(\mathrm{\Sigma })`$ into $`M(\mathrm{\Sigma })`$ that commute with the representation $`T`$ in the following sense:
$$𝒫T_t=T_t𝒫$$
for all $`t`$.
###### Theorem 6.1
Suppose that $`T`$ is a representation of $``$ that is sup path attaining, and that $`𝒫`$ commutes with $`T`$. Let $`\mu M(\mathrm{\Sigma })`$ be weakly analytic. Then $`𝒫\mu `$ is also weakly analytic.
To describe an interesting application of this theorem from , let us recall the following.
###### Definition 6.2
Let $`T`$ be a sup path attaining representation of $`G`$ in $`M(\mathrm{\Sigma })`$. A weakly measurable $`\sigma `$ in $`M(\mathrm{\Sigma })`$ is called quasi-invariant if $`T_t\sigma `$ and $`\sigma `$ are mutually absolutely continuous for all $`tG`$. Hence if $`\sigma `$ is quasi-invariant and $`A\mathrm{\Sigma }`$, then $`|\sigma |(A)=0`$ if and only if $`|T_t(\sigma )|(A)=0`$ for all $`tG`$.
Using Theorem 6.1 we obtained in the following extension of results of de Leeuw-Glicksberg and Forelli , concerning quasi-invariant measures.
###### Theorem 6.3
Suppose that $`T`$ is a sup path attaining representation of $``$ by isometries of $`M(\mathrm{\Sigma })`$. Suppose that $`\mu M(\mathrm{\Sigma })`$ is weakly analytic, and $`\sigma `$ is quasi-invariant. Write $`\mu =\mu _a+\mu _s`$ for the Lebesgue decomposition of $`\mu `$ with respect to $`\sigma `$. Then both $`\mu _a`$ and $`\mu _s`$ are weakly analytic. In particular, the spectra of $`\mu _a`$ and $`\mu _s`$ are contained in $`[0,\mathrm{})`$.
Our goal in this section is to extend Theorems 6.1 above to representations of a locally compact abelian group $`G`$ with ordered dual group $`\mathrm{\Gamma }`$. More specifically, we want to prove the following theorems.
###### Theorem 6.4
Suppose that $`T`$ is a sup path attaining representation of $`G`$ by isomorphisms of $`M(\mathrm{\Sigma })`$ such that $`T_{\varphi _\alpha }`$ is sup path attaining for each $`\alpha `$. Suppose that $`𝒫`$ commutes with $`T`$ in the sense that
$$𝒫T_t=T_t𝒫$$
for all $`tG`$. Let $`\mu M(\mathrm{\Sigma })`$ be weakly analytic. Then $`𝒫\mu `$ is also weakly analytic.
As shown in \[4, Theorem (4.10)\] for the case $`G=`$, an immediate corollary of Theorem 6.4 is the following result.
###### Theorem 6.5
Suppose that $`T`$ is a sup path attaining representation of $`G`$ by isometries of $`M(\mathrm{\Sigma })`$, such that $`T_{\varphi _\alpha }`$ is sup path attaining for each $`\alpha `$. Suppose that $`\mu M(\mathrm{\Sigma })`$ is weakly analytic with respect to $`T`$, and $`\sigma `$ is quasi-invariant with respect to $`T`$. Write $`\mu =\mu _a+\mu _s`$ for the Lebesgue decomposition of $`\mu `$ with respect to $`\sigma `$. Then both $`\mu _a`$ and $`\mu _s`$ are weakly analytic with respect to $`T`$. In particular, the $`T`$-spectra of $`\mu _a`$ and $`\mu _s`$ are contained in $`\overline{P}`$.
Proof of Theorem 6.4. Write
$$\mu =\mu _{\alpha _0}_T\mu +\underset{\alpha }{}d_\alpha _T\mu ,$$
as in (5.1), where the series converges unconditionally in $`M(\mathrm{\Sigma })`$. Then
(42)
$$𝒫\mu =𝒫(\mu _{\alpha _0}_T\mu )+\underset{\alpha }{}𝒫(d_\alpha _T\mu ).$$
It is enough to show that the $`T`$-spectrum of each term is contained in $`\overline{P}`$. Consider the measure $`\mu _{\alpha _0}_T\mu `$. We have $`\mathrm{spec}_T(\mu _{\alpha _0}_T\mu )S_{\alpha _0}`$. Hence by Theorem 4.4, $`\mu _{\alpha _0}_T\mu `$ is $`T_{\varphi _{\alpha _0}}`$-analytic. Applying Theorem 6.1, we see that
(43)
$$\mathrm{spec}_{T_{\varphi _{\alpha _0}}}(𝒫(\mu _{\alpha _0}_T\mu ))[0,\mathrm{}[.$$
Since $`𝒫`$ commutes with $`T`$, it is easy to see from Proposition 3.10 and Corollary 3.11 that
$$\mathrm{spec}_T(𝒫(\mu _{\alpha _0}_T\mu ))C_{\alpha _0}.$$
Hence by (43) and Theorem 4.4,
$$\mathrm{spec}_T(𝒫(\mu _{\alpha _0}_T\mu ))S_{\alpha _0},$$
which shows the desired result for the first term of the series in (42). The other terms of the series (42) are handled similarly.
Acknowledgments The second author is grateful for financial support from the National Science Foundation (U.S.A.) and the Research Board of the University of Missouri.
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# Is the Time a Dimension of an Alien Universe?
## I Introduction
The appearance in recent years many variants of quantum gravity, including theories with change of signature sign, gravitational quantum tunnel transitions in the inflation phase, different methods of quantizing and so forth (see, e.g., a review in ) do not allow to choose the criteria to distinguish ”right” theories from ”wrong” ones. The situation is similar to that existed in the beginning of our century in physics, when Lee groups theory was invented and included as a special case the Lorenz sub-group but there were no convincing reasons to single out of it the Dirac equation. Now there is no a quite clear idea on the time nature which would permit to predict the phenomena whose experimental examination in its turn would help to judge its reliability. If there is any, it would be able to significantly narrow down the choice for possible quantum gravity theories. It is one of such ”time models” that is discussed in the present paper. Treating time as a dimension of another universe with laws similar to those of our world, it gives the way to a number of experimentally phenomena, some of which are discussed in the paper.
## II Intersecting universes hypothesis
Let us consider time to be one of dimensions of a three-dimensional space of another universe which adjoins to ours and penetrate into it by this dimension (see ). To do this, assume that at the ”Big Bang” the quantum state of the early universe gave the birth to an ensemble of inflation formations with different values of Hubble constant $`H=\frac{\dot{a}}{a}`$ ($`a`$ is the scale factor of the space elements of metric). This does not contradict to the contemporary conception of the world in the inflation universe models. Then a chains of initially three-dimensional formations also could emerge, with one dimension penetrating into another one and thus making it four-dimensional. So, restrict ourselves to only such formations. Stress, that this formations are primordial three-dimensional but in addition there exist a ”penetration” of one dimension of another formation (we shall call it ”neighboring”) into it. In other words, we are interested in only the formations having total number of dimensions being equal to four: three of them are ”own” and the fourth one penetrated from another formation and is of different, generally speaking, origin. Thus, this fourth dimension belongs to both of the neighboring formations and is common to them. It is ”own” to one formation and ”foreign” to another. We shall refer to such formations as ”universes”. In the neighboring universe (whose the fourth dimension is) the situation would be the same. To its own three dimensions (one of which is common with the first universe considered) a new ”foreign” dimension from another universe is added - either from that mentioned before or not. In the first case we have an isolated pair of universes, and in the second one have the chain of intersecting universes which can be closed and finite or open and infinite. If we are to suppose non-equal inflating of the three-dimensional volumes of this universes in regard to each other (because of different value $`a`$ ), the fourth dimension penetrated into a universe from another one can play the role of natural calibrating value while comparing the volumes and be treated as ”time”. In this case the dimension of ”foreign” universe common to ours is time to us, but in the foreign universe the role of time is played by a dimension of an alien to it universe (either new one, or, probably, ours). Note, that this model does not reduce to the York theory (see ) or to the ”scale” theories of time. Here time is not simply a scale factor but a dimension of origin and properties different from that of our three space dimensions, though it plays in its own universe role exactly the same as them. All what is needed for such treatment is a difference in expanding between the three ”own” and the fourth ”foreign” dimensions. Therefore, from this point of view, time can be considered as a dimension of another formation and common to our universe but with properties different from that of our own three spatial dimensions and, hence, expanding unlike them. The concept of the expansion rate of the universe (as well as time and rate concepts themselves) are then defined by comparing the laws of changing three-dimensional volumes of neighboring universes having one common (the fourth to us) dimension. The four-dimensional interval of special relativity $`ds^2=dr^2c^2dt^2`$ in this case, at $`ds^2=0`$ will yield the rate of changing of our spatial coordinates in relation to $`t`$ coordinate which characterize the alteration of the ”spatial” dimensions of the other universe. Light velocity $`c`$ in this terms is simply the difference between changing in own three-dimensional volumes of two neighboring universes. It appeared as a result of the Big Bang and can be exceeded in our space only at local but not at large scales. In the neighboring universe, provided validity after the phase of blowing a kind of special relativity condition like that of our world, we shall have $`ds^2=dt^2c^2dr^2`$, and, $`c^{}=c^1`$ from our point of view, in the simple case. What consequences arise from this hypothesis concerning the nature of time, based on existing of at least two such four-dimensional universes with one common dimension?
## III Anomalous Redshift Behavior
As the first consequence of the model proposed we shall discuss the red shift in the spectra of far star objects. Let the physical laws of nature be equal in the neighboring universes. Then in both of them, ours ”spatial” and the other ”temporal” the following Hubble laws can be written:
$$\frac{dx}{dt}=v_x=H_xx$$
(1)
$$\frac{dt}{dx}=v_t=H_tt$$
(2)
where $`v_x`$ and $`v_t`$ are the rates of scattering of star objects in the corresponding universes, $`H_x`$ and $`H_t`$ are the Hubble constants in them and $`x`$ and $`t`$ mean the distances between the objects in the ”spatial” and ”temporal” universes. One of three dimensions making in the ”temporal” universe vector $`t`$ penetrates into our universe and plays the role of time to us.We shall denote it as $`t`$, because only one component of $`t`$ is common. Stress, that instead of the Hubble law (1) now one should take into account in addition to it the new relation (2), that is
$$v_t=H_tt$$
(3)
The physical sense of $`v_t`$ is the rate of time alteration per distance unit and can be bound up with the time passed since the star object was born in our universe. The Doppler wavelength shift $`Z=(\lambda _0\lambda _e)/\lambda _e`$ (here $`\lambda _0`$ and $`\lambda _e`$ are wavelengths in the observer’s frame of reference and in the object’s own one) then should be calculated using (3) as well as (1). Denote it as $`Z_{DT}`$ and $`Z_{DX}`$ accordingly. If $`Z_{DT}Z_{DX}`$ or $`Z_{DT}Z_{DX}`$ then the resulting shift is simply their sum:
$$Z=Z_{DX}+Z_{DT}$$
(4)
For $`Z_{DT}`$, using (3) one easily finds
$$Z_{DT}=\frac{1+\frac{v_t}{c^{}}}{\sqrt{1\frac{v_t^2}{c^2}}}1$$
(5)
where $`c^{}=c^1`$. When $`c^{}=\alpha c^1`$, $`\alpha =const`$. Putting $`v_t`$ from (3) into (5) and assuming, in virtue of similarity of the nature’s laws in both of the universes, $`H_t=H_0c^1`$ we obtain (for $`c^{}=c^1`$)
$$v_ec^1=H_ttc^1=H_0t(c^{}c)^1=\frac{t}{t_0}$$
(6)
with $`t_0=H_0^1`$ being the time since the universe formation and $`t`$ is the difference between the birth time of the star objects of whose spectra redshift is investigated (i.e, $`t=t_et_0`$ or $`t=t_0t_e`$ depending on the sign of the expression). Then
$$Z_{DT}=\frac{1+\frac{t}{t_0}}{\sqrt{1(\frac{t}{t_0}})^2}1$$
(7)
and
$$Z=Z_{DT}+Z_{DX}\frac{1+\frac{v}{c}}{\sqrt{1(\frac{v}{c})^2}}+\frac{1+\frac{t}{t_0}}{\sqrt{1(\frac{t}{t_0})^2}}2$$
(8)
If the contributions of $`Z_{DT}`$ and $`Z_{DX}`$ into $`Z`$ are comparable by value, it is not clear which way to calculate the Doppler shift in the direction to the observer is more preferable,
$`Z=Z_{DX}+Z_{DT}`$ (9)
or
$$Z=\frac{1+A}{\sqrt{1A^2}}1$$
(10)
where
$`A={\displaystyle \frac{\frac{v_x}{c}+\frac{t}{t_0}}{1+\frac{v_xt}{ct_0}}}`$ (11)
The first case corresponds to independent contribution of $`Z_{DT}`$ and $`Z_{DX}`$, while (10) represents the attempt to add velocities of different origin (in different universes) in the direction to the observer (note, that velocity in the ”temporal” universe is a vector only in that universe, but not in ours).
## IV New Energy -Momentum relation for Alien Non-Relativistic Particles?
As the second example we put the following question: whether a particle (with rest mass $`m_0`$) can penetrate from one universe into another (say, from ours to the neighboring or into our universe from elsewhere) and if yes, which energy it should posses? To answer it, let first write two Dirac equations for a free particle in our universe and in the other one:
$$_\gamma \psi (x)=0;(x=r,it)$$
(12)
$$_\gamma \psi (x_t)=0;(x_t=t,ix)$$
(13)
where $`\gamma `$ means Dirac matrices. Energy eigenvalues for the particle then
$$E^2=p^2+E_0^2;E_0^2=m_0c^2$$
(14)
$$P_t^2=E_t^2c^2+P_{0t}^2;P_0=m_0c^{}$$
(15)
Here the particle energy and momentum are written for our and the neighboring universes accordingly. Note, that momentum corresponds to the space derivative of wave function, that is $`\mathrm{\Psi }/r`$ and energy corresponds to the time derivative $`\mathrm{\Psi }/t`$. As space and time exchange their role while turning from one universe to the other, the derivatives also will exchange their role. That is, what we usually call ”momentum” will behave like energy and vice versa. We shall keep the notations $`P`$ for $`\mathrm{\Psi }/r`$ and $`E`$ for $`\mathrm{\Psi }/t`$ though their sense in different universes may change and this fact is taken into account in (15).
If we may admit that when a particle tunneling from one universe into another its momentum remains, then $`P_t=P`$ and (14) takes the form
$$E^2=(E_t^2c^2+P_{0t}^2)c^2+E_0^2$$
(16)
But in the neighboring universe $`P_{0t}=m_0c^{}`$ and in our universe momentum can be of any quantity. Therefore $`E^t`$ turns into relativistic energy and can not be less then $`E_0=m_0c^2`$. Then from (16) follows an inequality for the particle energy
$$E>2E_0$$
(17)
This condition is necessary but not sufficient, because our assumptions require that vector momentum conservation law must realize. Probably, the validity of this law can be achieved only for a selected direction in our space or only near massive bodies which get the remainder energy and momentum. Nevertheless, the expression (17), in principle, can be examined experimentally, since it indicates on the possibility of an apparent violations of the laws of conservation of energy, charge, number of particles faster than $`\sqrt{3/4c}`$, which originates from the particles going out into another universe. For massless particles (i.g., photons) this transition may reduce to a frequency renormalization (a simultaneous for all the frequencies ”reddening”). Beside that, one can take into consideration a possibility for both mass and massless particles tunneling through a potential barrier of unknown nature which divided he universes since they was born. Discuss now some aspects of tunneling of quantum particles from one universe into another. If it is possible, no matter how small the probability of such a transition is, it seems that one can point out two observable consequences. The first of them is as follows. In the neighboring universe a free non-relativistic particle obeys the equation
$$ih\frac{}{x}\mathrm{\Psi }=\frac{h^2}{2m_0c^3}\frac{^2}{t^2}\mathrm{\Psi }$$
(18)
therefore the relation between its energy and momentum will differs from that for a particle in our universe. Namely,
$$E=\sqrt{2E_0c}\sqrt{p}$$
(19)
## V Does Intensity of Relic Radiation have a Second Maximum?
The second consequence observable appears from the relic radiation of the neighboring universe. This radiation, if tunneling through the barrier between the universes, can, in principle, be detected. Its frequency and wavelength can be calculated using the assumed earlier similarity of the physical laws in both of the universes. Thus, in the neighboring one the relic radiation wavelengths $`\lambda _{tr}`$ (measured in seconds) will coincide numerically with the wavelength obtained in our universe $`\lambda _r`$. But while turning to our universe from the ”foreign” one, the wavelength must be multiplied by the light velocity $`c`$. So, $`\lambda _r`$ of the foreign relic radiation came to us is ( in case $`c^{}=c^1`$)
$$\lambda _{rt}=\lambda c$$
(20)
Though in (20) the wavelength is very large (it corresponds to the frequencies of about 0,1 $`sec^1`$), at the intensity curve of relic radiation background there will be an anomaly in this (or another) region, observed as a maximum with fast decreasing low frequency border. Whether it is possible to distinguish this maximum from the galaxies heat noise?
## VI Is the Redshift of Arp Galaxies Consequence of Differences in Time Existence of Their Parts?
Very interesting problem is interpretation of large differences in value of redshift of different parts of Arp galaxies. All difficulties may be take of if use for explaining this the new relation (equation (7)) The hypothesis of Arp about explosion of supernova stars may be includes in this interpretation . The interpretation of redshift is: the different parts of Arp galaxies were borne in different times.
## VII Final remarks
In the frames of the model considered a number of other questions arises. Is our universe non-isotropic in relation to the direction of ”foreign” dimension penetrating? Is there only a pair of universes or every dimension can penetrate into some universe or other making an infinite series of universes? If it is possible to send a particle into another universe in what time it will find itself when returned and so on. Apparently answering such questions can stimulate producing more complicated theories (including the proposed here as a special case) and further investigations the nature of time.
In conclusion we name again the main ideas and results.
1. The model of the ”Universes” discussed treats the time as one of the dimensions of a ”neighboring” Universes born together with ours and expanding by another law. The rate of expanding is calibrated by comparing the change in ”own” three-dimensional volumes of the neighboring universes.
2. If we are to assume the similarity of the physical laws in different universes, then beside the Hubble law the analogous law is exist concerning to the rate of changing the time coursing in dependence on the time passed. This law allows to avoid some contradictions related to redshift in spectra of a number of anomalous far galaxies (including Arp galaxies).
3. The model considered gives way to a number of statements which can in principle be examined by experiment. Some of the consequences are briefly discussed: existing of relic radiation come from ”foreign” universes; a possibility for non-conservation of energy, charge and number of particles with energies greater than $`2E_0`$ (because of leaving into the neighboring universes); anomalous dependence of energy upon momentum for nun-relativistic particles of ”foreign” origin.
All the results following from the model of Universe structure proposed, though it seems to be unlikely and leading to an unusual treatment of time, allow to hope for an experimental examination. Stress,the hypothesis about nature of time doe’s non contradict the hypothesis of multifractal nature of time and space presented by author in -. Some results of multifractal theory of time and space coincide with results of hypothesis considered in this paper, but here was presented hypothesis that gives new look on the possible origin of time and its nature. As the last remark we once more pay attention on possibility to receive the results concerning the redshift if simple postulate the inhomogeneity of time flows and relation $`\frac{t}{x}=H_tt`$ without any explanations.
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# Generation and evolution of vortex-antivortex pairs in Bose-Einstein condensates
## I Introduction
Quantum liquids and gases are an excellent environment for studies of solitons and topological quantum structures such as vortices. Studies of superfluidity and vorticity have been done mostly in He II . Due to the high density and strong interactions the theoretical studies give usually only a qualitative agreement with experiments. Moreover, the size of the vortex core is very small in He II, typically $`1`$ Å. The situation is drastically different for atomic Bose-Einstein condensates . These many-body systems are dilute and weakly interacting, and as a result a quantitative agreement between theory and experiments can often be found. Also, the expected size of a vortex core is fairly large, $`0.11\mu \mathrm{m}`$, making detailed experimental studies possible.
Several theoretical approaches have been suggested for the generation of vortices in condensates, see e.g. Refs. . Among them, the stirring of a condensate with a blue-detuned laser in a single component , and the coherent interconversion between the two components of a binary condensate have been successfully implemented. In addition, the creation of vortex rings after the decay of a dark soliton has been recently reported . These vortices carry one quantum of circulation each, and they form stable, geometric patterns. However, the success in controlling atomic condensates suggests that one could create also non-equilibrium situations to investigate dynamics of vortices, and in particular, vortex collisions. Alternatively, a vortex can be generated through the “phase-imprinting” method . This approach has been successfully used to generate solitons in condensates .
Here we propose how to create and control a pair of spatially separated vortices of opposite circulation via controlled decay of dynamically unstable solitons in a toroidal trapped condensate. Due to the toroidal geometry and the employed technique for the creation of the vortices, the dynamics and circulation of the generated vortices are correlated, even when they are generated at well-separated locations.
In homogeneous condensates and in absence of dissipation a vortex–antivortex pair moves as a whole through the fluid . Such behavior can be, however, substantially modified in the case of trapped condensates. In particular, for the case of stiff potentials the vortex dynamics is altered by the constraint of vanishing normal velocity field at the boundaries . We show how vortex–antivortex dynamics in toroidal geometries become different from the equivalent dynamics in a homogeneous (not trapped) condensate.
We propose a combination of condensate splitting, phase imprinting, and subsequent merging of separated parts into a toroidal condensate, which leads to a very controlled mechanism of vortex creation. We describe the details of our proposal and numerical solution of the relevant Gross-Pitaevskii equation in Sec. II. In Sec. III we apply the method of images to explain the vortex dynamics and we give some concluding remarks in Sec. IV
## II Vortex creation and dynamics numerically
We assume initially a condensate trapped in a Gaussian optical trap. By shaking the trap rapidly we split the condensate into two spatially separated parts . First we slowly increase the amplitude of the periodic shaking along the chosen $`x`$-axis . The effective potential evolves smoothly from the Gaussian well into a trap with two separated minima; this is the time-averaged potential. The condensate follows adiabatically the changes of the time-averaged potential, splitting into two parts (Fig. 1). By adding a similar process along the $`y`$-axis the time-averaged trapping potential evolves towards a torus. Thus the two separated condensate parts combine eventually into a torus. Before merging the two halves, we imprint on one of them (with a strongly detuned laser field ) a change of phase of the wavefunction, $`\mathrm{\Delta }\varphi `$.
At sufficiently low temperatures, the system is well described by the corresponding Gross-Pitaevskii equation
$$i\mathrm{}\frac{\mathrm{\Psi }}{t}=\frac{\mathrm{}^2}{2m}^2\mathrm{\Psi }+V\mathrm{\Psi }+NU_0|\mathrm{\Psi }|^2\mathrm{\Psi }.$$
(1)
Here $`\mathrm{\Psi }`$ is the condensate wavefunction, $`V(x,y,t)`$ is the trap potential (not the time-averaged one), $`N`$ is the number of atoms and $`U_0=4\pi \mathrm{}^2a/m`$, where $`a`$ is the $`s`$-wave scattering length and $`m`$ is the atomic mass. In all simulations we have assumed a sodium condensate ($`a=2.75`$ nm) with the particle number $`N=10^5`$. The width of the trap in $`x`$\- and $`y`$-directions was $`10\mu `$m and its size in the $`z`$-direction $`L_z=2\mu `$m. The trap depth was set to $`4\mu \mathrm{K}`$. In our simulations the duration of each shaking stage was typically 50 ms. The creation of the torus can be seen in the first steps of Fig. 1. With the chosen parameters the trap ground state energy in $`z`$-direction is considerably larger than the mean field energy $`nU_0`$ at the toroidal trap ($`n`$ is the atomic density). This fact allows us to restrict ourselves to a two–dimensional geometry.
The coherence time for the applied phase shift depends on the details of the shaking. If we express the wavefunction as $`\mathrm{\Psi }=R\mathrm{exp}(i\varphi )`$ and assume the Thomas-Fermi limit, i.e., ignore terms due to kinetic energy in Eq. (1), we obtain (assuming that $`R`$ changes sufficiently slowly)
$$\frac{\varphi }{t}=\frac{V(r,t)+NU_0R(r,t)^2}{\mathrm{}}.$$
(2)
Since in the Thomas-Fermi limit the right-hand side of Eq. (2) becomes independent of $`r`$, the phase evolution is identical everywhere. During the time between applying the phase shift and the merging into a torus, the phase has changed by a constant amount over the condensate halves and thus there will be a discontinuity in the phase across the merging points, corresponding to the applied phase. If, on the other hand, the condensate exhibits sloshing, the kinetic energy term makes the phase evolution $`r`$-dependent and the applied phase shift loses its meaning. Even a small asymmetry between the two parts is then enough to mask the applied phase shift on the time-scale of $`100`$ ms, which in effect leads to a total loss of the control of the vortex production. In our simulations the density distribution followed the Thomas-Fermi limit closely and, therefore, the phase evolution was almost identical everywhere.
When the phase-shifted condensate parts form a torus, the discontinuity in phase is quickly transformed into a dark soliton, since the system can adapt to the phase discontinuity only by reducing its density. In general, such a kink-like state in a trapped condensate is dynamically stable if the mean-field energy $`nU_0`$ is roughly smaller than the characteristic trap energy in the direction given by the soliton front , which in our case corresponds to the typical energy of the trap in the radial direction. With the parameters from our setup and $`n1.710^{14}`$ cm<sup>-3</sup>, this condition is violated. Under such conditions the soliton front undergoes snake instabilities which can eventually decay into vortex-antivortex pairs (or vortex rings in 3D configurations), see e.g. . This process has been experimentally observed in non-linear optics , as well as, very recently, in matter waves . Since the front perturbation with a wavevector $`k=1/(\sqrt{2}\xi )`$, where $`\xi =1/\sqrt{4\pi na}`$ is the healing length, will grow fastest , the distance between the generated vortices is expected to be about $`d\pi /k`$. If $`k>1/\xi `$ the normal modes are dynamically stable, i.e., if the torus is narrower than about $`\pi \xi `$ the soliton is dynamically stable.
The depth of the soliton depends on the size of the phase discontinuity. It must be large enough for vortex generation. Otherwise the resulting soliton can decay into other excitations, or even remain stable. For example, a $`\pi /2`$ phase shift did not lead to generation of vortices in our simulations. Due to the $`2\pi `$ degeneracy of the phase, the phase shift $`\pi `$ corresponds to the largest possible discontinuity. But one should note that the direction of the phase gradient over the discontinuity affects the subsequent evolution. If the phase difference is equal to $`\pi `$, the $`2\pi `$ degeneracy makes the direction of the gradient ambiguous. In our simulations the motion of vortices was initially indeterminate and thus very sensitive to small disturbances, when $`\mathrm{\Delta }\varphi \pi `$. For $`\mathrm{\Delta }\varphi =0.9\pi `$ (as in Fig. 1), however, we obtain a much more predictable behavior. Due to symmetry, the mirror image of this situation is obtained with $`\mathrm{\Delta }\varphi =1.1\pi `$ (the lower vortex moves clockwise in the beginning, and the upper one counterclockwise in Fig. 1).
In our simulations each soliton decays into a pair of vortices of opposite circulation. However, one of the vortices, located close to the edge of the condensate, has a very short lifetime, and decays within a few ms at the edge. The survival of only one vortex per soliton is not surprising, considering that $`dl`$, where $`l`$ is the width of the torus. For wider torii, i. e., longer solitons, our simulations show a decay of each soliton into many vortex-antivortex pairs, as expected . Similar behaviour has been reported in simulations of condensates in rectangular boxes . In our situation the condensate wave function remains symmetric with respect to reflection to the horizontal axis (see Fig. 1). This symmetry leads to correlations between vortex production at the two phase discontinuities, located on opposite sides of the torus. We are always left with two vortices (one vortex per soliton) of opposite circulations. These vortices have also opposite vertical and equal horizontal components of the velocity, and they are therefore automatically set on a collision course. This reflects the conservation of angular momentum (we have verified that the shaking process or the adiabatic merging of the separated condensate parts do not bring any extra angular momentum into the system)
Figure 1 shows the time evolution obtained from numerical simulations. The created vortices bounce from each other when they approach. This process repeats itself if we continue simulations beyond $`t=300`$ ms. In a nondissipative homogeneous system two vortices of opposite circulation separated by a distance greater than the healing length move parallel, since each vortex will move with the velocity of the other one. Clearly, the vortex dynamics in a torus trap is more complex and completely different. It can be understood with the help of the images method.
## III Images method
The images method has been recently successfully employed to describe both a single vortex, and vortex arrays dynamics in circular 2D traps . We consider in the following a stiff torus (with infinitely high and stiff walls), i.e. we neglect the influence of the inhomogeneous trapping potential on the vortex dynamics . Additionally, we assume that the distance between the vortices is greater than the size of the vortex cores. In absence of friction the vortex velocity equals the superfluid velocity at the vortex position. Such a velocity field, which is induced by the presence of the vortices, has to fulfill the constraint of vanishing normal component at the torus boundaries. To fulfill this constraint it is enough to consider for each vortex an infinite number of fictitious image vortices, each with appropriate position and circulation. Each vortex, either real or image, contributes to the superfluid velocity field $`𝐯_{SF}(𝐫)`$ by $`\left(\kappa _i/2\pi \right)\widehat{𝐳}\times (𝐫𝐫_i)/|𝐫𝐫_i|^2`$, where $`𝐫_i`$ and $`\kappa _i`$ refer to its position and circulation, respectively.
Let us consider the case of a torus with an inner radius $`R_1`$, an outer radius $`R_2`$, and one vortex with circulation $`\kappa `$ located at $`𝐫`$ (in a frame centered at the origin of the torus). In order to cancel the normal component of the superfluid velocity on the boundaries, we have to consider two different families of image vortices. For the first family, an image vortex is placed at $`𝐫_{im1}=\frac{R_1^2}{r}\frac{𝐫}{r}`$ with circulation $`\kappa `$ as well as one image vortex with circulation $`\kappa `$ at the origin . The image vortex at $`𝐫_{im1}`$ will induce a velocity component normal on the outer circle, so in addition we need an image vortex at $`𝐫_{im2}=\frac{R_2^2}{r_{im1}}\frac{𝐫}{r}=\frac{R_2^2}{R_1^2}𝐫`$ with circulation $`\kappa `$. In turn this image will induce some normal component on the inner circle, so two more image vortices are required: one at $`𝐫_{im3}=\frac{R_1^2}{r_{im2}}\frac{𝐫}{r}=\frac{R_1^4}{R_2^2r}\frac{𝐫}{r}`$ with circulation $`\kappa `$ the other one at the origin with circulation $`\kappa `$ and so on and so forth. For the second family, one considers an image vortex at $`𝐫_{}^{}{}_{im1}{}^{}=\frac{R_2^2}{r}\frac{𝐫}{r}`$ with circulation $`\kappa `$. This image vortex gives a normal component of the velocity on the inner circle so we need in addition one image vortex at $`𝐫_{}^{}{}_{im2}{}^{}=\frac{R_1^2}{r_{im1}^{}}\frac{𝐫}{r}=\frac{R_1^2}{R_2^2}𝐫`$ with circulation $`\kappa `$ as well as one image vortex at the origin with circulation $`\kappa `$. For this we require another image vortex at $`𝐫_{}^{}{}_{im3}{}^{}=\frac{R_2^2}{r_{im2}^{}}\frac{𝐫}{r}=\frac{R_2^4}{R_1^2r}\frac{𝐫}{r}`$ with circulation $`\kappa `$ and so on. For the second real vortex in the torus we have to repeat the same formalism to get an image configuration, which ensures that the boundary condition is fulfilled.
That way we can express the total superfluid velocity field in the torus (for arbitrary circulations), in particular the velocity at the locations of the real vortices as a function of their mutual positions. Since the latter equals the vortex velocity we get the following first order equations for the vortex positions $`𝐫_i`$ that can be easily computed
$`2\pi {\displaystyle \frac{d}{dt}}\stackrel{}{r}_i`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\kappa _j\left(r_i\stackrel{}{e}_{\phi _i}r_j\left(\frac{R_2}{R_1}\right)^{2n}\stackrel{}{e}_{\phi _j}\right)}{\left(\stackrel{}{r}_i\left(\frac{R_2}{R_1}\right)^{2n}\stackrel{}{r}_j\right)^2}}+R_1R_2)`$ (3)
$``$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\kappa _j\left(r_i\stackrel{}{e}_{\phi _i}\left(\frac{R_1}{R_2}\right)^{2n}\frac{R_1^2}{r_j}\stackrel{}{e}_{\phi _j}\right)}{\left(\stackrel{}{r}_i\left(\frac{R_1}{R_2}\right)^{2n}\frac{R_1^2}{r_j^2}\stackrel{}{r}_j\right)^2}}+R_1R_2)`$ (4)
$`+`$ $`{\displaystyle \frac{\kappa _{ji}\left(r_i\stackrel{}{e}_{\phi _i}r_{ji}\stackrel{}{e}_{\phi _{ji}}\right)}{\left(\stackrel{}{r}_i\stackrel{}{r}_{ji}\right)^2}}+{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \frac{\kappa _j}{r_i}}\stackrel{}{e}_{\phi _i}i,j=1,2`$ (5)
where $`R_1R_2`$ indicates a term equal to the previous one but interchanging $`R_1`$ and $`R_2`$. The results obtained with the image method are presented in Fig. 2. As initial conditions we take two vortices located on opposite parts of the torus and displaced symmetrically from the central point between the torus boundaries. For the sake of clarity our initial conditions have been chosen such that the vortex trajectories are wider, i.e., they explore more space between the torus boundaries than in the numerical simulation from Fig. 1. But with appropriate initial vortex positions one can achieve a nearly perfect agreement with simulations of Eq. (1) concerning the shape of the trajectories, and the period of their oscillations ($`300`$ ms).
## IV Conclusions
So far, we have assumed that the condensate is at zero temperature, i.e., we have used Gross-Pitaevskii equation (1) to model the condensate, neglecting any dissipation. The presence of non-condensate atoms is expected to result in dissipative Magnus forces and vortex decay. But experimental studies so far show that a realistic time scale for vortex decay is a few seconds, which is longer than the time scale of our scenario. We tested the stability of our system by adding a small, constant imaginary term into Eq. (1). It did not modify the reported behavior unless the magnitude of the term was too strong to be considered a perturbation. Nevertheless, the problem of the effects of dissipation in our model is very challenging. The presence of dissipation might cause the vortex-antivortex pair either to collide with the trapping boundaries, or to mutually annihilate.
Another complication arises from the fact that currently the spatial resolution of the measurement of the condensate density profile is limited. Usually this problem is avoided by letting the condensate expand rapidly by removing the trapping potential. In a quasi-2D traps the optical depth will be an issue. Therefore some other detection scheme, such as matter wave interference , may have to be considered.
Summarizing, we have presented a method for creating vortices in a toroidal geometry in a controllable way. This method allows for the analysis of the interaction of vortices in such a geometry, which is significantly modified by the geometry of the trapping potential. The realm of possible phenomena that can be generated and controlled by merging the split condensate parts into a torus is not by any means limited by those discussed here, e.g. if we form the torus rapidly, other excitations are also created, which interact with the vortices and affect their dynamics.
In the analytical estimation of the vortex orbits we have neglected for simplicity the effects of the inhomogeneity of the trapping potential. Such inhomogeneity will result in an additional contribution to the velocity of the vortex in the direction perpendicular to the density gradient, as discussed in Ref. (a detailed analysis will be the subject of a later publication). On the other hand, our calculations have been constrained to quasi-2D traps. Very recently the issue of lower dimensional BEC has been subject of great interest. In particular 2D (and even 1D) condensates have been experimentally observed . Therefore the situation considered in the present paper is experimentally justified and feasible. However, the analysis of 3D geometries constitutes a challenging problem. We have made some short simulations to study the dynamics of vortex-antivortex (line) interactions in a 3D pipe geometry. The results we obtain are very similar to those presented here. In a proper 3D geometry, the solitons decay into vortex rings, as observed very recently in the case of cylindrical traps . The propagation, deformation and interaction of vortex rings will be the topic for future studies.
We thank M. Lewenstein, G. Shlyapnikov and M. Baranov for discussions. We acknowledge support of the Academy of Finland (project 43336), Center for Scientific Computing (CSC), and Deutsche Forschungsgemeinschaft (SFB 407). J.-P. M. is supported by the National Graduate School on Modern Optics and Photonics.
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# 1 Introduction
## 1 Introduction
In disordered systems, the ensemble average of the logarithm of the partition function cannot be evaluated directly in most cases. Two widely used methods, the replica trick and the supersymmetric method , have been proposed to circumvent this problem. In the replica trick the ensemble average of the logarithm of the partition function is written as
$`\mathrm{log}Z=\underset{n0}{lim}{\displaystyle \frac{Z^n1}{n}},`$ (1)
i.e., $`\mathrm{log}Z`$ is calculated by the analytic continuation of the $`n`$-dependence of $`n`$ replicated partition functions to $`n0`$. If $`Z`$ is given by a fermion determinant, $`Z^n`$ can be written as an integral over $`n`$ replicated Grassmann fields, and therefore, this method is known as the fermionic replica trick. Alternatively, $`\mathrm{log}Z`$ can be written as
$`\mathrm{log}Z=\underset{n0}{lim}{\displaystyle \frac{1Z^n}{n}}.`$ (2)
In this case, a determintal $`Z`$ can be expressed as an integral over $`n`$ replicated complex fields and this limit referred to as the bosonic replica trick. Both for the fermionic and for the bosonic replica trick, the average partition function can mapped onto a non-linear $`\sigma `$-model which is amenable to a saddle-point expansion.
In the supersymmetric method the disorder average is performed for the ratio
$`_J\mathrm{log}Z\left(J\right)_{J=0}=_J{\displaystyle \frac{Z\left(J\right)}{Z\left(J=0\right)}}|_{J=0}.`$ (3)
If $`Z\left(J\right)`$ is given by a determinant, the numerator can be expressed as a fermionic integral whereas the denominator can be written as a bosonic integral. For $`J=0`$ the partition function is thus invariant with respect to superunitary transformations that mix the fermionic and bosonic fields. Based on this symmetry, the supersymmetric partition function can be mapped onto a supersymmetric non-linear $`\sigma `$model which can be used to derive exact analytical expressions for spectral correlation functions.
he replica trick has two obvious advantages. First, it can be applied to cases where $`Z`$ cannot be expressed as a determinant (as for example is the case in the theory of spin-glasses ), and, second, it is possible to calculate the logarithm of the partition function. One disadvantage of this method is that, in order to take the replica limit, the $`n`$ dependence of the replicated partition function has to be known analytically. Therefore, the application of replica trick is limited to perturbative expansions of the partition function. A much more serious problem of this method is that the continuation of the $`n`$-dependence to $`n0`$ is not unique. For example, a term of the form $`\mathrm{sin}n\pi `$ contributing to $`Z^n`$ vanishes for integer $`n`$ but gives rise to a nonzero result in the replica limit. The failure of the replica trick was first noticed in the theory of spin glasses where a replica symmetric minimum of the free energy resulted in a negative entropy . However, in that case the problems could be resolved by means of an elaborate scheme of replica symmetry breaking . More recently, the replica trick was criticized because of its failure to reproduce the oscillatory contributions to random matrix correlation functions . Notice however that the non-oscillatory contributions to the two-point function were reproduced correctly . Another example for which the replica trick may be problematic are nonhermitian Random Matrix Theories with eigenvalues scattered in the complex plane , but we will not study such theories in this article.
The advantage of the supersymmetric method is that it is possible to derive non-perturbative analytical results. This has been shown convincingly for the calculation of spectral correlation functions of random matrix-ensembles . A disadvantage is that it requires some familiarity with supermathematics, but for perturbative expansions this method is no more complicated than the replica trick. In the example that will be discussed in this article, the exact super-symmetric calculation is actually much simpler than the perturbative replica calculation. A second disadvantage is, that because the average partition function is normalized to unity, one does not have access to the average free energy.
Recently, the replica trick was revived in an article by Kamenev and Mézard . They found that, in order to reproduce the oscillatory terms in spectral two-point correlation function, saddle points with broken replica symmetry had to be taken into account . For the Gaussian Unitary Ensemble they found the exact analytical result. As explained in an article by Zirnbauer , the reason for this miracle is a consequence of the Duistermaat-Heckman theorem which is applicable to the $`\sigma `$model for the two-point function of the Gaussian Unitary Ensemble. This theorem states the conditions under which an integral is localized on its critical points so that a saddle-point approximation becomes exact. For the Gaussian Orthogonal Ensemble and the Gaussian Symplectic Ensemble the replica trick could only reproduce the asymptotic expansion of the two-point spectral correlation function for large energy differences . We remind the reader that, unless we know the analytical properties of a function in the complex plane, it cannot be reproduced from its asymptotic series. A version of the replica trick that does not rely on the non-linear $`\sigma `$-model but instead on orthogonal polynomials was shown to reproduce the exact correlation functions for all Gaussian ensembles . In this article we will not discuss this variant of the replica trick which is not an alternative to the orthogonal polynomial method.
To investigate the replica trick we have chosen the microscopic spectral density of the QCD Dirac operator which is defined as the spectral density near zero on the scale of the average level spacing. The reason is three-fold. i) The $`n`$-fold replicated partition function is the QCD partition function with $`n`$ flavors. Its low-energy limit, relevant for the microscopic spectral density, is known analytically. Because of spontaneous breaking of chiral symmetry in QCD, it is given by a partition function of weakly interacting Goldstone bosons and, on the scale of the average level spacing, it can be reduced to a unitary matrix integral which can be evaluated analytically for any number of flavors . This, so called finite volume chiral partition function has been investigated in great detail, also in the context of one-link integrals in lattice QCD . ii) Because the eigenvalues occur in pairs $`\pm \lambda `$, the level repulsion of the eigenvalues leads to a nontrivial oscillatory behavior of the microscopic spectral density. iii) The low energy partition function can be derived as a function of two integer valued parameters, the topological charge $`\nu `$ and the number of physical flavors $`N_f`$ (in this article only $`N_f=0`$ will be discussed) and can be trivially continued to non-integer $`\nu `$. For half-integer $`\nu `$, the saddle point expansion of the $`U\left(n\right)`$-integrals terminates for finite positive integer values of $`n`$. This is closely related to Duistermaat-Heckman localization, where the leading order saddle point approximation is exact such as for the $`\sigma `$model with $`n`$ fermionic replicas of the two-point function of the Gaussian Unitary Ensemble . Based on recent work , we expect that the replica trick with replica-symmetry breaking gives the exact result in this case.
The replica limit of the finite volume chiral partition function was first studied in . It was found that the asymptotic expansion of the valence quark mass dependence of the chiral condensate in the microscopic region (i.e. for large valence quark masses in units of the average eigenvalue spacing) was reproduced by the replica trick. The small mass expansion was obtained up to logarithmic singularities which are essential for the calculation of the spectral density. Recently, the partially quenched supersymmetric chiral Lagrangian was formulated in terms of the replica trick .
In this article we analyze the oscillatory contributions to the microscopic spectral density in the framework of the replica trick. In section 2 we discuss the chiral symmetries for bosonic replicas and introduce the supersymmetric quenched low energy chiral partition function. The calculation of the microscopic spectral density by means of the supersymmetric method is given in section 3. This calculation illustrates that the compact/non-compact structure of the final result for the resolvent appears naturally in the supersymmetric method. In section 4.1 we derive the large-mass asymptotic behavior of the microscopic spectral density by means of the fermionic replica trick. The small-mass behavior is discussed in section 4.2. Bosonic replicas are discussed in section 5 and concluding remarks are made in section 6. In the Appendix we the derive trace correlators necessary for the replica calculation to fourth order in the inverse mass.
## 2 Chiral Symmetry
In this section we discuss chiral symmetry for bosonic and fermionic replicas, as well as for the supersymmetric partition function. To illustrate the difference in flavor symmetries between bosonic and fermionic replicas we present a detailed discussion for the case of one flavor or replica.
The question we wish to address in this article is whether the spectrum of the QCD Dirac operator in the sector of topological charge $`\nu `$ can be obtained from the QCD partition function with $`n`$ additional replica flavors with quark mass $`z`$. This partition function is given by
$`Z_\nu ^{\left(N_f+n\right)}\left(z\right)={\displaystyle \left[dA\right]_\nu \stackrel{n}{det}\left(i\text{ /}D+z\right)\underset{f=1}{\overset{N_f}{}}det\left(i\text{ /}D+m_f\right)e^{S_{YM}\left[A\right]}},`$ (4)
where $`m_1,\mathrm{},m_{N_f}`$ are the usual quark masses. The integral is over all gauge fields in the sector of topological charge $`\nu `$ (which is chosen positive in this article) and is weighted by the Yang-Mills action. The replica limit of the resolvent or the chiral condensate is defined by
$`\mathrm{\Sigma }\left(z\right){\displaystyle \frac{1}{V_4}}\mathrm{Tr}{\displaystyle \frac{1}{z+i\text{ /}D}}=\underset{n0}{lim}{\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{V_4}}{\displaystyle \frac{}{z}}\mathrm{ln}Z_\nu ^{\left(N_f+n\right)}\left(z\right).`$ (5)
Here, $`V_4`$ is the Euclidean 4-volume. The spectral density, which in terms of the eigenvalues $`i\lambda _k`$ of the Dirac operator is defined by
$`\rho \left(\lambda \right)={\displaystyle \underset{k}{}}\delta \left(\lambda \lambda _k\right),`$ (6)
follows from the discontinuity across the imaginary axis
$`{\displaystyle \frac{\rho \left(\lambda \right)}{V_4}}={\displaystyle \frac{1}{2\pi }}\left[\mathrm{\Sigma }\left(i\lambda +ϵ\right)\mathrm{\Sigma }\left(i\lambda ϵ\right)\right].`$ (7)
Below we only consider the dimensionless ratio
$`{\displaystyle \frac{\mathrm{\Sigma }\left(z\right)}{\mathrm{\Sigma }_0}}`$ $`=`$ $`{\displaystyle \frac{1}{V_4\mathrm{\Sigma }_0}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\lambda {\displaystyle \frac{\rho \left(\lambda \right)}{z+i\lambda }},`$ (8)
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑u{\displaystyle \frac{1}{V_4\mathrm{\Sigma }_0}}\rho \left({\displaystyle \frac{u}{V_4\mathrm{\Sigma }_0}}\right){\displaystyle \frac{1}{V_4\mathrm{\Sigma }_0z+iu}},`$
as a function of the microscopic variable $`V_4\mathrm{\Sigma }_0z`$. The chiral condensate, $`\mathrm{\Sigma }_0`$, is defined as the limit of $`\mathrm{\Sigma }\left(z\right)`$ for $`z`$ close to $`z=0`$ but many level spacings away from the center of the spectrum. As one can see from the second equality, it can be expressed as an integral over the microscopic spectral density defined by
$`\rho _s\left(u\right)={\displaystyle \frac{1}{V_4\mathrm{\Sigma }_0}}\rho \left({\displaystyle \frac{u}{V_4\mathrm{\Sigma }_0}}\right).`$ (9)
For spontaneously broken chiral symmetry, the spacing of the eigenvalues is given by $`\pi /\mathrm{\Sigma }_0V_4`$ so that (9) is stable in $`V_4`$ and can be calculated in the thermodynamic limit.
Because the resolvent $`\mathrm{\Sigma }\left(z\right)`$ has a cut along the imaginary axis it is sometimes more convenient to use the relation $`G\left(i\lambda ϵ\right)=G\left(i\lambda +ϵ\right)`$ to rewrite the discontinuity as
$`{\displaystyle \frac{\rho \left(\lambda \right)}{V_4}}={\displaystyle \frac{1}{2\pi }}\left[\mathrm{\Sigma }\left(i\lambda +ϵ\right)+\mathrm{\Sigma }\left(i\lambda +ϵ\right)\right].`$ (10)
We thus only need to calculate the resolvent in the half-plane $`\mathrm{Re}\left(z\right)>0`$. If the argument $`\mathrm{\Sigma }\left(z\right)`$ represents the microscopic variable, this relation gives us the microscopic spectral density.
The reason of working with the partition function (4) is that in the phase of spontaneously broken symmetry its low energy limit in entirely determined by chiral symmetry and is a partition function of weakly interacting Goldstone modes (or pions). We are interested in the kinematical domain
$`1/m_\pi V_4^{1/4}1/\mathrm{\Lambda }_{QCD}.`$ (11)
Because $`V_4^{1/4}1/\mathrm{\Lambda }_{QCD}`$, only the Goldstone modes contribute to the mass-dependence of the partition function . For quark masses for which the Compton wavelength of the Goldstone modes is much larger than the size of the box ( $`1/m_\pi V_4^{1/4}`$), the kinetic term of the chiral Lagrangian can be ignored, and only the constant fields contribute to the mass dependence of the low-energy partition function. Therefore, as we will see next, in the domain (11) the QCD parition function can be reduced to a unitary matrix integral.
In QCD with $`n`$ fermionic flavors the chiral symmetry group is given by $`U_V\left(n\right)\times U_A\left(n\right)`$. A $`U_A\left(1\right)`$ subgroup of the axial symmetry group is broken by the anomaly. The remaining axial symmetry group is broken spontaneously by the formation of a nonzero chiral condensate and the vector symmetry group, $`U_V\left(n\right)`$, remains unbroken. The Goldstone manifold is thus given by the axial group $`SU_A\left(n\right)`$. The low energy limit of the QCD partition function is uniquely fixed from the requirement that its transformation properties under the chiral symmetry group are the same as for full QCD. Taking into account the anomaly, one finds, for $`n`$ flavors all with mass $`m`$, in the sector of topological charge $`\nu `$, the low-energy finite volume partition function
$`Z_\nu ^{\left(n\right)}\left(x\right)={\displaystyle _{UU\left(n\right)}}𝑑U\left(detU\right)^\nu e^{\frac{x}{2}\mathrm{Tr}\left(U+U^1\right)},`$ (12)
where $`xmV_4\mathrm{\Sigma }_0`$. This partition function is valid in the range (11).
In the domain (11), the low energy partition function can also be obtained from a Random Matrix Theory with the global symmetries of the QCD partition function . In this theory, the matrix elements of the Euclidean Dirac operator are replaced by independently distributed Gaussian random variables. In the sector of topological charge $`\nu `$, the Dirac matrix is thus has given by
$`D=\left(\begin{array}{cc}0& iW\\ iW^{}& 0\end{array}\right),`$ (15)
where $`W`$ is an $`n\times \left(n+\nu \right)`$ matrix, and the integration over the gauge fields is replaced by an integration over the probability distribution of the matrix elements. If $`W`$ is complex and the probability distribution is a function of traces $`\mathrm{Tr}\left(W^{}W\right)^p`$, this ensemble is known as the chiral Unitary Ensemble (chUE) or the chiral Gaussian Unitary Ensemble (chGUE) if the distribution of the matrix elements is Gaussian. As was shown in , the statistical properties of the smallest eigenvalues of $`D`$ do not depend on the details of the probability distribution of the matrix elements.
In order to analyze the chiral symmetry for bosonic replicas let us first discuss the flavor symmetries for one flavor. The fermion determinant that occurs in the QCD partition function can be written as a integral over Grassmann variables.
$`det\left(D+m\right)={\displaystyle 𝑑\overline{\chi }𝑑\chi e^{{\scriptscriptstyle d^4x\overline{\chi }\left(D+m\right)\chi }}}.`$ (16)
The inverse determinant can be written as an integral over bosonic integration variables
$`{\displaystyle \frac{1}{det\left(D+m\right)}}={\displaystyle 𝑑\varphi ^{}𝑑\varphi e^{{\scriptscriptstyle d^4x\varphi ^{}\left(D+m\right)\varphi }}}.`$ (17)
In a well-defined theory the functional integral has to be convergent. This is automatically the case for the Grassmann integration in (16), but the bosonic integrals in (17) is only convergent for positive $`m`$. The symmetries of the partition function should be compatible with these convergence requirements. In particular, $`\varphi ^{}`$ should be identified with the complex conjugate of $`\varphi `$, and not as an independent integration variable such as the fermionic variables $`\overline{\chi }`$ and $`\chi `$.
If we decompose the spinors according to the block structure of the Dirac operator the vector symmetry of the massive theory with one fermionic replica (16) is given by
$`\left(\begin{array}{c}\chi _1\\ \chi _2\end{array}\right)e^{i\theta }\left(\begin{array}{c}\chi _1\\ \chi _2\end{array}\right),\left(\begin{array}{c}\overline{\chi }_1\\ \overline{\chi }_2\end{array}\right)e^{i\theta }\left(\begin{array}{c}\overline{\chi }_1\\ \overline{\chi }_2\end{array}\right),`$ (26)
and the axial $`U\left(1\right)`$ symmetry of the massless theory can be written as
$`\left(\begin{array}{c}\chi _1\\ \chi _2\end{array}\right)\left(\begin{array}{c}e^{i\theta }\chi _1\\ e^{i\theta }\chi _2\end{array}\right),\left(\begin{array}{c}\overline{\chi }_1\\ \overline{\chi }_2\end{array}\right)\left(\begin{array}{c}e^{i\theta }\overline{\chi }_1\\ e^{i\theta }\overline{\chi }_2\end{array}\right).`$ (35)
The $`U\left(1\right)`$ vector symmetry of the massive bosonic theory (17) is the same as for the fermionic theory
$`\left(\begin{array}{c}\varphi _1\\ \varphi _2\end{array}\right)e^{i\theta }\left(\begin{array}{c}\varphi _1\\ \varphi _2\end{array}\right),\left(\begin{array}{c}\varphi _1^{}\\ \varphi _2^{}\end{array}\right)e^{i\theta }\left(\begin{array}{c}\varphi _1^{}\\ \varphi _2^{}\end{array}\right).`$ (44)
This transformation does not affect the complex conjugation properties of $`\varphi `$. However, the axial transformation (35) applied to the bosonic fields affects their complex conjugation properties. In this case the axial transformation that is compatible with the convergence of the bosonic integral is given by
$`\left(\begin{array}{c}\varphi _1\\ \varphi _2\end{array}\right)\left(\begin{array}{c}e^s\varphi _1\\ e^s\varphi _2\end{array}\right),\left(\begin{array}{c}\varphi _1^{}\\ \varphi _2^{}\end{array}\right)\left(\begin{array}{c}e^s\varphi _1^{}\\ e^s\varphi _2^{}\end{array}\right).`$ (53)
The axial symmetry group is therefore not $`U\left(1\right)`$ but instead $`Gl\left(1\right)/U\left(1\right)`$. Of course, this axial transformation is also a symmetry of the fermionic partition function.
For $`n`$ bosonic flavors the vector symmetry is $`U\left(n\right)`$ whereas the axial symmetry is given by
$`\left(\begin{array}{c}\varphi _1\\ \varphi _2\end{array}\right)\left(\begin{array}{c}e^H\varphi _1\\ e^H\varphi _2\end{array}\right),\left(\begin{array}{c}\varphi _1^{}\\ \varphi _2^{}\end{array}\right)\left(\begin{array}{c}e^H\varphi _1^{}\\ e^H\varphi _2^{}\end{array}\right),`$ (62)
with the matrix $`H`$ containing only real elements. The axial symmetry group is thus given by the coset $`Gl\left(n\right)/U\left(n\right)`$. An explicit parameterization of this coset is given by $`AA^{}`$ with $`AGl\left(n\right)`$. We expect that the axial symmetry of the bosonic partition function
$`Z_\nu ^{\left(n\right)}={\displaystyle \frac{1}{det^n\left(D+m\right)}}_\nu `$ (63)
is broken in the same way as in the fermionic case with a $`Gl\left(1\right)/U\left(1\right)`$ coset broken explicitly by the anomaly and the remaining part of the coset broken spontaneously by the chiral condensate. In absence of explicit symmetry breaking, the Goldstone manifold is thus given by $`Gl\left(n\right)/U\left(n\right)`$.
The mass term introduced according to
$`{\displaystyle \underset{k,l=1}{\overset{n}{}}}\varphi _1^kM_{kl}\varphi _1^l+\varphi _2^kM_{kl}^{}\varphi _2^l`$ (64)
is invariant under the axial transformation (62) provided that the mass matrix is transformed at the same time as
$`Me^HMe^H,M^{}e^HM^{}e^H.`$ (65)
The bosonic partition function in the sector of topological charge $`\nu `$ transforms covariantly with (62)
$`Z_\nu ^{\left(n\right)}det\left(e^{2\nu H}\right)Z_\nu ^{\left(n\right)}.`$ (66)
The low energy limit of bosonic partition function is uniquely fixed by the requirement that its transformation properties are the same as of $`Z_\nu ^{\left(n\right)}`$. In the sector of topological charge $`\nu `$ it is given by
$`Z_\nu ^{\left(n\right)}={\displaystyle _{UGl\left(n\right)/U\left(n\right)}}𝑑U\left(detU\right)^\nu e^{\frac{\mathrm{\Sigma }_0V_4}{2}\mathrm{Tr}\left(MU+M^{}U^1\right)},`$ (67)
The measure $`dU`$ is the invariant Haar measure. Below we consider only the case with a diagonal mass matrix with all nonzero matrix elements equal to $`m`$ and use the definition $`x=mV_4\mathrm{\Sigma }_0`$. This partition function is valid in the kinematical domain (11) where constant Goldstone fields are the only relevant degrees of freedom.
Since fermionic integrals are always convergent the partition function is invariant under both compact ($`U\left(n\right)`$) and non-compact ($`Gl\left(n\right)/U\left(n\right)`$) axial transformations. However, the small mass behavior of the non-compact partition function is singular because of the volume of the non-compact group diverges. This is not the case for the fermionic partition function and therefore this $`Gl\left(n\right)/U\left(n\right)`$ is not an admissible parameterization of the Goldstone manifold. The correct parameterization is given by the compact manifold $`U\left(n\right)`$. At a more technical level, this follows from the fact that the transformations that lead to the non-compact integral are only legitimate for bosonic quarks. For fermionic quarks one necessarily finds a compact effective partition function.
In the supersymmetric method the generating function of the quenched resolvent is given by
$`Z_\nu (z,J)={\displaystyle \frac{det\left(\text{ /}D+z+J\right)}{det\left(\text{ /}D+z\right)}}_\nu .`$ (68)
It can be written as a superintegral
$`Z_\nu (M,M^{})={\displaystyle 𝑑\varphi 𝑑\varphi ^{}𝑑\chi 𝑑\overline{\chi }\mathrm{exp}[\varphi ^{}\text{ /}D\varphi +\overline{\chi }\text{ /}D\chi +(\begin{array}{c}\varphi _1^{}\\ \overline{\chi }_1\end{array})M\left(\begin{array}{c}\varphi _1\\ \chi _1\end{array}\right)+(\begin{array}{c}\varphi _2^{}\\ \overline{\chi }_2\end{array})M^{}\left(\begin{array}{c}\varphi _2\\ \chi _2\end{array}\right)]}_\nu .`$ (77)
(78)
Both $`M`$ and $`M^{}`$ are given by $`\mathrm{diag}(z+J,z)`$, but in order to study the transformation properties of the partition function, we keep them as general matrices.
With spontaneously broken axial symmetry, the bosonic part of the Goldstone manifold is $`U\left(1\right)\times Gl\left(1\right)/U\left(1\right)`$. Because, the partition function is invariant under super-unitary transformations, the full symmetry group is given by the maximum Riemannian submanifold of $`Gl\left(1|1\right)`$ (we will denote this manifold by ($`\widehat{Gl}\left(1|1\right)`$). This manifold can be parameterized as
$`U=\left(\begin{array}{cc}e^{i\theta }& \alpha \\ \beta & e^s\end{array}\right).`$ (81)
For zero topological charge the generating function (68) is invariant under $`Gl_R\left(1|1\right)\times Gl_L\left(1|1\right)`$ if at the same time the mass matrix is transformed as
$`MU_R^1MU_L,M^{}U_L^1M^{}U_R.`$ (82)
For nonzero values of $`\nu `$ the generating function (68) is not invariant under (82) but transforms according to
$`Z_\nu (z,J)\mathrm{Sdet}^\nu \left(U_R^1U_L\right)Z_\nu (z,J).`$ (83)
The low-energy partition function is obtained from the requirement that it should have the same transformation properties as the QCD partition function (68). In the sector of topological charge $`\nu `$, it is given by
$`Z_\nu (z,J)={\displaystyle _{U\widehat{Gl}\left(1|1\right)}}𝑑U\mathrm{Sdet}^\nu \left(U\right)e^{\frac{\mathrm{\Sigma }_0V_4}{2}\mathrm{Str}\left(MU+M^{}U^1\right)}.`$ (84)
The integration is over the Haar measure of $`\widehat{Gl}\left(1|1\right)`$, Below, we only the consider the case of a diagonal mass matrix with both $`M`$ and $`M^{}`$ equal to $`\mathrm{diag}(z+J,z)`$.
An amusing observation is that the topological charge in QCD partition function is discrete, whereas the number of flavors with equal mass, thought of as the power of the fermion determinant is a continuous parameter. In the low energy partition function it is just the other way round.
## 3 Supersymmetric Calculation of the Resolvent
In this section we calculate the super-integrals in (94) to obtain analytical expression for the resolvent and the microscopic spectral density. For integer $`\nu `$, this calculation was also presented in . To simplify the explicit calculations we will choose the parameterization
$`U=\left(\begin{array}{cc}e^{i\theta }& 0\\ 0& e^s\end{array}\right)\mathrm{exp}\left(\begin{array}{cc}0& \alpha \\ \beta & 0\end{array}\right).`$ (89)
In this case the Haar measure is simply given by
$`dU=d\theta dsd\alpha d\beta .`$ (90)
This results in the partition function
$`Z(z,z+J)={\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑s{\displaystyle _{C_c}}𝑑\theta 𝑑\alpha 𝑑\beta e^{i\nu \theta \nu s}\mathrm{exp}\mathrm{Str}M\left(\begin{array}{cc}\left(1+\frac{\alpha \beta }{2}\right)\mathrm{cos}\theta & \alpha \left(e^se^{i\theta }\right)\\ \beta \left(e^{i\theta }e^s\right)& \left(1\frac{\alpha \beta }{2}\right)\mathrm{cosh}s\end{array}\right),`$ (93)
(94)
Where $`z`$ and $`J`$ are now microscopic variables, i.e. they are expressed in units of $`1/V_4\mathrm{\Sigma }_0`$. The integration over $`s`$ is over the complete real axis. For integer $`\nu `$ the integration over $`\theta `$ is over the interval $`[\pi ,\pi ]`$. For non-integer $`\nu `$ the translational invariance of the $`\theta `$-integral is lost. It is recovered by extending the integration contour to include the intervals $`\pi +i\mathrm{},\pi ]`$ and $`[\pi ,\pi +i\mathrm{}`$. A picture of this integration contour, denoted by $`C_c`$, is shown in Fig. 1.
After performing the Grassmann integrations the partition function reduces to
$`Z(z,J)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑s{\displaystyle _{C_c}}𝑑\theta e^{i\nu \theta \nu s}\left(\left(z+J\right)\mathrm{cos}\theta +z\mathrm{cosh}s\right)e^{\left(z+J\right)\mathrm{cos}\theta z\mathrm{cosh}s}`$ (95)
$`=`$ $`\left(z+J\right)K_\nu \left(z\right)I_{\nu +1}\left(z+J\right)+zK_{\nu +1}\left(z\right)I_\nu \left(z+J\right).`$
The normalization of the partition function according to $`Z\left(z,J=0\right)=1`$ follows from the Wronskian identity $`zK_\nu \left(z\right)I_{\nu +1}\left(z\right)+zK_{\nu +1}\left(z\right)I_\nu \left(z\right)=1`$. The resolvent obtained by diferentiation with respect to $`J`$ is, after using some identities for Bessel functions, given by
$`\mathrm{\Sigma }\left(z\right)=z\left(K_\nu \left(z\right)I_\nu \left(z\right)+K_{\nu 1}\left(z\right)I_{\nu +1}\left(z\right)\right)+{\displaystyle \frac{\nu }{z}}.`$ (96)
The microscopic spectral density then follows from the discontinuity across the imaginary axis according to (10),
$`\rho _s\left(x\right)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left[\mathrm{\Sigma }\left(ix+ϵ\right)+\mathrm{\Sigma }\left(ix+ϵ\right)\right]`$ (97)
$`=`$ $`{\displaystyle \frac{x}{2}}\left(J_\nu \left(x\right)^2J_{\nu +1}\left(x\right)J_{\nu 1}\left(x\right)\right)+\nu \delta \left(x\right).`$
The last term is the contribution from the $`\nu `$ zero modes. For integer values of $`\nu `$ the integration contour $`C_c`$ can be replaced by the segment $`[\pi ,\pi ]`$. The restriction of the integration over $`\theta `$ to this segment for non-integer values of $`\nu `$ would have resulted in the wrong answer. This can be seen from the following representation of modified Bessel functions,
$`I_\nu \left(z\right)={\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}e^{z\mathrm{cos}\theta }e^{i\nu \theta }𝑑\theta {\displaystyle \frac{\mathrm{sin}\nu \pi }{\pi }}{\displaystyle _0^{\mathrm{}}}e^{z\mathrm{cosh}s\nu s}𝑑s.`$ (98)
The microscopic spectral density and the expression for $`\mathrm{\Sigma }\left(z\right)`$ were first obtained from chiral Random Matrix Theory by means of the orthogonal polynomial method and can also be obtained by means of the supersymmetric method . I n that case the derivation, starting from the joint probability distribution of the eigenvalues, is also correct for non-integer values of $`\nu `$ and is given by the expressions (96) and (97).
The asymptotic expansion of the imaginary part of $`\mathrm{\Sigma }\left(z\right)`$ is given by
$`i\mathrm{Im}\mathrm{\Sigma }\left(z\right)`$ $`=`$ $`{\displaystyle \frac{i\left(1\right)^\nu e^{2z}}{2z}}+{\displaystyle \frac{i\left(4\nu ^21\right)\left(1\right)^\nu e^{2z}}{8z^2}}{\displaystyle \frac{i\left(1\right)^\nu e^{2z}\left(4\nu ^21\right)\left(4\nu ^29\right)}{64z^3}}`$ (99)
$`+`$ $`{\displaystyle \frac{i\left(4\nu ^21\right)\left(4\nu ^29\right)\left(194\nu ^2\right)\left(1\right)^\nu e^{2z}}{3\times 256z^4}},`$
and for the asymptotic expansion of the spectral density we find
$`\rho _{s\nu }\left(x\right)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}[1{\displaystyle \frac{\mathrm{cos}\left(2x\pi \nu \right)}{2x}}+{\displaystyle \frac{14\nu ^2}{8x^2}}(1\mathrm{sin}(2x\pi \nu ))+(4\nu ^21)(4\nu ^29){\displaystyle \frac{\mathrm{cos}\left(2x\pi \nu \right)}{64x^3}}`$ (100)
$`+`$ $`(4\nu ^21)(4\nu ^29){\displaystyle \frac{\left(6+\left(194\nu ^2\right)\mathrm{sin}\left(2x\pi \nu \right)\right)}{x^42^73!}}+\mathrm{}].`$
The asymptotic expansion of the partition function, the resolvent, and the spectral density terminates for half-integer values of $`\nu `$. For example, for $`\nu =1/2`$ only one oscillating term in (100) is nonvanishing suggesting that it can be obtained from a leading order saddle point approximation.
Since for half-integer $`\nu `$ the integral is localized on the critical points, we expect that the asymptotic expansion generated by the replica trick reproduces the exact answer. For other values of $`\nu `$ we expect that the replica trick reproduces the asymptotic series to all orders in $`1/x`$.
## 4 Spectral Density via Fermionic Replicas
In this section, we analyze the low-energy partition in the quenched case ($`N_f=0`$). For $`n`$ replica flavors in the sector of topological charge $`\nu `$ it is given by (see (12))
$`Z_\nu ^{\left(n\right)}\left(x\right)={\displaystyle _{UU\left(n\right)}}𝑑U\stackrel{\nu }{det}Ue^{\frac{x}{2}\mathrm{Tr}\left(U+U^1\right)}.`$ (101)
In order to take the replica limit, $`n0`$, the $`n`$-dependence of (101) has to be known explicitly. The closed form of $`Z_\nu ^{\left(n\right)}\left(x\right)`$, in terms of a determinant of modified Bessel functions,
$`Z_\nu ^{\left(n\right)}\left(x\right)=det\left(I_{\nu +ji}\left(x\right)\right),i,j=1,\mathrm{},n,`$ (102)
does not provide us with an explicit $`n`$-dependence. The explicit $`n`$-dependence can only be obtained for the large-mass and the small-mass expansion of $`Z_\nu ^{\left(n\right)}\left(x\right)`$. It was obtained in for an expansion about the replica symmetric saddle point using the method of Virasoro constraints. Our new result is the asymptotic expansion of the microscopic spectral density using replica symmetry breaking à la Kamenev and Mézard which will be discussed the the second half of the next subsection. The small mass expansion of (102) was also considered in , but failed at the order for which logarithmic terms enter in the expansion. In the second subsection we derive the lowest order logarithmic term for the case of zero topological charge.
### 4.1 Large Mass Expansion
The asymptotic expansion of the microscopic spectral density is obtained from the large mass expansion of the finite volume partition function. To this end we expand the partition function (101) in powers of $`1/x`$ by means of a saddle point approximation. By diagonalizing $`U`$ it can be easily seen that the saddle points are given by unitary matrices with eigenvalues $`\pm 1`$, i.e. by unitary matrices satisfying $`U^2=1`$. The solutions of this equation are highly degenerate. They can be organized in $`n+1`$ classes, $`U=I_p`$, where $`I_p`$ is a diagonal matrix with $`p`$ elements $`1`$ and $`np`$ elements $`+1`$ (with $`0pn`$). The integrand does not depend on the submanifold $`U\left(n\right)/U\left(np\right)\times U\left(p\right)`$ of $`U\left(n\right)`$ and the integration over this coset has to be performed exactly resulting in its volume $`V_{n,p}`$. In terms of the parameterization
$`U=I_pU_0VU_0^1,`$ (103)
with $`U_0U\left(n\right)/U\left(np\right)\times U\left(p\right)`$ it is clear that the integration can be restricted to $`U\left(np\right)\times U\left(p\right)`$. Summing over all saddle points the partition function is given by
$`Z_\nu ^{\left(n\right)}\left(x\right)={\displaystyle \underset{p=0}{\overset{n}{}}}V_{n,p}{\displaystyle _{VU\left(np\right)\times U\left(p\right)}}𝑑VJ\left(V\right)\stackrel{\nu }{det}Ve^{\frac{x}{2}\mathrm{Tr}\left[I_p\left(V+V^1\right)\right]}.`$ (104)
where $`V_{n,0}=1`$ and for $`p0`$
$`V_{n,p}=\left(2\pi \right)^{p\left(np\right)}\left(\begin{array}{c}n\\ p\end{array}\right){\displaystyle \frac{_{ȷ=1}^pj!_{ȷ=1}^{np}j!}{_{ȷ=1}^nj!}}\left(2\pi \right)^{p\left(np\right)}F_n^p.`$ (107)
The integration over $`V`$ should be thought of a saddle-point integral of a formal expansion of $`V`$ about the identity to all orders. Below we will make this explicit for the different types of saddlepoints. The total number of saddle-points in the class $`p`$ is $`\frac{n!}{\left(np\right)!p!}`$. This factor is included as combinatorial factor in $`V_{n,p}`$. The volume of the coset $`U\left(n\right)/U\left(np\right)\times U\left(p\right)`$ is given by the ratio $`\frac{V\left(U\left(n\right)\right)}{V\left(U\left(p\right)\right)V\left(U\left(np\right)\right)}`$. With the volume of $`U\left(k\right)`$ given by $`\left(2\pi \right)^{\frac{k\left(k+1\right)}{2}}_{j=1}^kj!`$ we obtain the volume factor $`V_{n,p}`$.
From the exact expression of the partition function (101) in terms of modified Bessel function given in (102) it is clear that the asymptotic series of $`Z_n^\nu \left(x\right)`$ terminates for half integer $`\nu `$. Let us investigate the asymptotic expansion about the saddle point $`I_p`$ in more detail. Because the total number of degrees of freedom in $`U\left(np\right)\times U\left(p\right)`$ is equal to $`\left(np\right)^2+p^2`$, the expansion of $`Z_\nu ^{\left(n\right)}\left(x\right)`$ is of the form
$`Z_\nu ^{\left(n\right)}\left(x\right)e^{\left(n2p\right)x}\left({\displaystyle \frac{1}{x}}\right)^{\left(np\right)^2/2+p^2/2}\left(1+𝒪\left({\displaystyle \frac{1}{x}}\right)\right).`$ (108)
This result is valid for arbitrary $`\nu `$. The asymptotic series of the Bessel function $`I_{k+\frac{1}{2}}\left(x\right)`$ terminates at $`1/x^{k+1/2}`$. The result for the maximum power in $`1/x`$ occuring the expansion of the determinant (102) is particularly simple for $`\nu =\frac{1}{2}`$ and is given by
$`\left({\displaystyle \frac{1}{x}}\right)^{n^2/2}.`$ (109)
Let us consider this case in more detail. Since, as we will see below, only the saddle points for $`p=0`$ or $`p=1`$ contribute in the replica limit, we only discuss these values of $`p`$. We observe that the maximum power and the minimum power in the asymptotic series are equal for $`p=0`$ saddle point. The asymptotic series thus has only one term and we find that
$`\left(Z_{\nu =\frac{1}{2}}^{\left(n\right)}\left(x\right)\right)_{p=0}`$ $``$ $`e^{nx}\left({\displaystyle \frac{1}{x}}\right)^{n^2/2}.`$ (110)
For $`p=1`$ and larger half-integer values of $`\nu `$, more terms contribute to the asymptotic series, but it still terminates. For example, for $`p=1`$ and $`\nu =\frac{1}{2}`$ one finds from (108) and (109) that the difference between the maximum and minimum power in the expansion in $`1/\sqrt{x}`$ is $`n1`$. Therefore, the asymptotic series for $`Z_n^\nu \left(x\right)`$ cannot contain more than $`n`$ terms. From numerical examples for small values of $`n`$, one indeed finds that in this case the coefficients of all $`n`$ possible terms are nonvanishing. For the same reason as in the case of the GUE two-point function , we expect that the replica trick will give the exact result both for $`\nu =\frac{1}{2}`$ as well as larger half-integer values of $`\nu `$.
Clearly, the expression (104) makes only sense for positive integer values of $`n`$ . In order to analytically continue it we follow the work of Kamenev and Mézard . They analytically continued the factorials in $`F_n^p`$ such that $`F_n^p`$ vanishes for $`pn+1`$. Then the sum in (104) can be extended up to infinity and the replica limit $`n0`$ can be taken term by term in the sum over $`p`$. One finds
$`\underset{n0}{lim}F_n^pn^p,`$ (111)
so that only the terms $`p=0`$ and $`p=1`$ of the, to infinity continued, sum in (104) survive. In these two cases we obtain $`F_{n0}^0=1,F_{n0}^1=n`$. Notice, that the continuation of the sum over $`p`$ to infinity explicitly breaks the replica symmetry $`pnp`$. Of course, the group integral in (30) must be also continued to non-integer values of $`n`$.
The previous discussion suggests the definition
$`Z_\nu ^{\left(n\right)}\left(x\right)\left(Z_\nu ^{\left(n\right)}\left(x\right)\right)_{p=0}+n\left(Z_\nu ^{\left(n\right)}\left(x\right)\right)_{p=1}.`$ (112)
We will first consider the contribution for $`p=0`$ which originates from expanding $`U`$ around the identity matrix $`I_0`$. Although not necessary, it turned out to be convenient to parameterize $`V_{p=0}U\left(n\right)`$ according to
$`V_{p=0}={\displaystyle \frac{1+iH/2}{1iH/2}},`$ (113)
where $`H`$ is an Hermitian $`n\times n`$ matrix. From a diagonal representation of $`V_{p=0}`$ one can easily show that
$`dV_{p=0}={\displaystyle \frac{1}{det^n\left(1+H^2/4\right)}}dH.`$ (114)
In the replica limit, the Jacobian can be ignored and one simply has $`dV_{p=0}=dH`$. The $`p=0`$ contribution to the partition function is thus given by
$`\left(Z_\nu ^{\left(n\right)}\left(x\right)\right)_{p=0}={\displaystyle }`$ $`dH`$ $`det\left({\displaystyle \frac{1+iH/2}{1iH/2}}\right)^\nu e^{nx\frac{x}{2}\mathrm{Tr}H^2+x\mathrm{Tr}\left(\frac{H^4/8}{1+H^2/4}\right)}.`$ (115)
In order to keep track of the powers of $`x`$ it is convenient to rescale $`HH/\sqrt{x}`$ which leads to another Jacobian $`x^{n^2/2}`$ that also vanishes in the replica limit. In terms of $`x=m\mathrm{\Sigma }_0V`$, the mass dependence of the condensate is given by
$`{\displaystyle \frac{\mathrm{\Sigma }_{p=0}\left(x\right)}{\mathrm{\Sigma }_0}}=\underset{n0}{lim}{\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{Z_\nu ^{\left(n\right)}\left(x\right)}}{\displaystyle \frac{}{x}}\left(Z_\nu ^{\left(n\right)}\left(x\right)\right)_{p=0}.`$ (116)
so that to order $`1/x^4`$ we need to collect terms to order $`1/x^3`$ in the expansion of the partition function. Using the expressions
$`e^{\frac{x}{2}Tr\left(V_{p=0}+V_{p=0}^1\right)}`$ $`=`$ $`e^{nx\frac{1}{2}\mathrm{Tr}H^2}e^{\mathrm{Tr}\left[\frac{1}{8x}H^4\frac{1}{32x^2}H^6+\frac{1}{128x^3}H^8\mathrm{}\right]}`$
$`\left(detV_{p=0}\right)^\nu `$ $`=`$ $`e^{\nu n}e^{\frac{i\nu \mathrm{Tr}}{\sqrt{x}}\left(H\frac{2H^3}{3\times 2^3x^2}+\mathrm{}\right)}.`$ (117)
we find the result (only terms up to order $`1/x^2`$ are displayed)
$`\left(Z_\nu ^{\left(n\right)}\right)_{p=0}`$ $`=`$ $`{\displaystyle \frac{e^{nx+\nu n}}{x^{\frac{n^2}{2}}}}{\displaystyle 𝑑He^{\frac{1}{2}\mathrm{Tr}H^2}}`$ (118)
$`\times `$ $`\left(1{\displaystyle \frac{\nu ^2}{2x}}\left(\mathrm{Tr}H\right)^2+{\displaystyle \frac{\nu ^2}{12x^2}}\left(\mathrm{Tr}H\right)\left(\mathrm{Tr}H\right)^2+{\displaystyle \frac{\nu ^4}{24x^2}}\left(\mathrm{Tr}H\right)^4\right)`$
$`\times `$ $`\left(1+{\displaystyle \frac{1}{8x}}\mathrm{Tr}H^4{\displaystyle \frac{1}{32x^2}}\mathrm{Tr}H^6\right)`$
$`=`$ $`\left({\displaystyle \frac{2\pi }{x}}\right)^{n^2/2}e^{nx}\left(1n{\displaystyle \frac{\nu ^2}{2x}}+\left(2n^3+n\right){\displaystyle \frac{1}{8x}}\right)+𝒪\left({\displaystyle \frac{n^2}{x^2}}\right).`$
The Gaussian integrals have been calculated using the trace correlators given in the Appendix. To this order, the Jacobian in (114) gives rise to an extra factor $`\left(1n^3x/4\right)`$. With inclusion of this terms the $`1/x`$ corrections vanish for $`\nu =\frac{1}{2}`$. All terms of order $`1/x^2`$ are at least of order $`n^2`$. The resolvent obtained from (116)
$`{\displaystyle \frac{\mathrm{\Sigma }_{p=0}\left(x\right)}{\mathrm{\Sigma }_0}}=\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{4\nu ^21}{8x^2}}{\displaystyle \frac{\left(4\nu ^21\right)\left(4\nu ^29\right)}{128x^4}}+\mathrm{}`$ (119)
agrees with the large mass expansion of the quenched condensate obtained via other methods . We observe a cancellation of the odd powers in $`1/x`$. This is in agreement with the analytical expression for $`\mathrm{\Sigma }\left(x\right)`$ given in (96) which can be rewritten as
$`{\displaystyle \frac{\mathrm{\Sigma }\left(x\right)}{\mathrm{\Sigma }_0}}={\displaystyle \frac{x}{2}}\left(\left(2{\displaystyle \frac{4\nu ^2}{x^2}}\right)I_\nu \left(x\right)K_\nu \left(x\right)+I_{\nu +1}\left(x\right)K_{\nu +1}\left(x\right)+I_{\nu 1}\left(x\right)K_{\nu 1}\left(x\right)\right).`$ (120)
From the asymptotic behavior of the Bessel functions, one easily derives that the asymptotic series of this expression is an expansion in powers of $`1/x^2`$. Also notice the cancellation of the $`\nu /x`$ term in the large$`x`$ asymptotic expansion.
Next we consider the more subtle contribution of the saddle point given by the diagonal matrix $`I_1`$ with one element equal to $`1`$ and $`n1`$ elements equal to $`1`$. A parameterization of $`U\left(n1\right)\times U\left(1\right)`$ that is convenient for the expansion about the saddle point is given by
$`V_{p=1}=\left(\begin{array}{cc}\frac{1+iH/2}{1iH/2}& \\ & \frac{1+ih/2}{1ih/2}\end{array}\right),`$ (123)
where, now, $`H`$ is a hermitian $`\left(n1\right)\times \left(n1\right)`$ matrix and $`h`$ is a real variable. Because $`U_0U\left(n\right)/U\left(n1\right)\times U\left(1\right)`$ we have that
$`U=I_1U_0V_{p=1}U_0^1=U_0I_1V_{p=1}U_0^1.`$ (124)
The measure, in terms of the coordinates $`H`$ and $`h`$, can be obtained by diagonalizing $`U`$ and $`H`$ with unitary transformations $`U_1`$ and $`W`$, respectively. If the eigenvalues are denoted by $`e^{i\theta }`$ and $`h_k`$, in this order, we find the measure (notice the plus sign in the last factor)
$`dU`$ $`=`$ $`{\displaystyle \underset{k<l}{}}\left|e^{i\theta _k}e^{i\theta _l}\right|^2{\displaystyle \underset{k}{}}d\theta _kdU_1`$
$`=`$ $`{\displaystyle \frac{_{k<l}\left|h_kh_l\right|^2}{\left(1+h^2/4\right)_k\left(1+h_k^2/4\right)^{n1}}}{\displaystyle \underset{k}{}}\left|{\displaystyle \frac{1+ih_k/2}{1ih_k/2}}+{\displaystyle \frac{1+ih/2}{1ih/2}}\right|^2dh{\displaystyle \underset{k}{}}dh_kdU_0dV.`$
The last factor can be written as
$`{\displaystyle \underset{k}{}}\left|{\displaystyle \frac{1+ih_k/2}{1ih_k/2}}+{\displaystyle \frac{1+ih/2}{1ih/2}}\right|^2=4^{n1}{\displaystyle \frac{det^2\left(1+hH\right)}{\left(1+h^2/4\right)^{n1}det\left(1+H^2/4\right)}}.`$ (126)
Using that $`_{k<l}\left|h_kh_l\right|^2_kdh_kdV=dH`$ we thus find the measure
$`dU=4^{n1}dU_0{\displaystyle \frac{dhdHdet^2\left(1+hH\right)}{\left(1+h^2/4\right)^ndet^n\left(1+H^2/4\right)}}J(H,h)dhdHdU_0.`$ (127)
The integrand does not depend on $`U_0`$ and the integration over these variables just gives the volume of the coset which, together with the combinatorial factor, combines into the factor $`V_{n,1}`$ discussed in the first part of this subsection. The $`p=1`$ contribution to the partition function is thus given by (with one overall minus sign from orienting the $`h`$-integration from $`\mathrm{}`$ to $`\mathrm{}`$)
$`\left(Z_\nu ^{\left(n\right)}\left(x\right)\right)_{p=1}`$ $`=`$ $`\left(1\right)^\nu \left(8\pi \right)^{n1}{\displaystyle 𝑑H𝑑h\frac{det^2\left(1+hH\right)}{\left(1+h^2/4\right)^ndet^n\left(1+H^2/4\right)}}`$ (128)
$`\times \left({\displaystyle \frac{1+ih/2}{1ih/2}}\right)^\nu det\left({\displaystyle \frac{1+iH/2}{1iH/2}}\right)^\nu e^{x\mathrm{Tr}\left(\frac{1H^2/4}{1+H^2/4}\right)x\left(\frac{1h^2/4}{1+h^2/4}\right)}.`$
Since the factors in the denominator of the Jacobian can be ignored in the replica limit only the following terms in the expansion of the Jacobian contribute to order $`1/x^4`$,
$`J=4^{n1}\left[1+{\displaystyle \frac{h\mathrm{Tr}H}{2}}+{\displaystyle \frac{h^2}{16}}\left[2\left(\mathrm{Tr}H\right)^2\mathrm{Tr}H^2\right]+\mathrm{}\right].`$ (129)
It is instructive to perform the calculation to leading order in $`1/x`$. In this case the following terms should be collected
$`\left(Z_\nu ^{\left(n\right)}\left(x\right)\right)_{p=1}`$ $`=`$ $`\left(1\right)^\nu \left(8\pi \right)^{n1}e^{\left(n2\right)x+\nu n}{\displaystyle 𝑑H𝑑he^{\frac{x}{2}\mathrm{Tr}H^2+\frac{x}{2}h^2}}`$
$`\times `$ $`\left[1+{\displaystyle \frac{x}{8}}\mathrm{Tr}H^4{\displaystyle \frac{x}{8}}h^4{\displaystyle \frac{\nu ^2}{2}}\left(\left(\mathrm{Tr}H\right)^2+h^2+2ih\mathrm{Tr}H\right)+{\displaystyle \frac{h\mathrm{Tr}H}{2}}\right].`$
The last two terms vanish upon integration. Thus, the Jacobian contributes only at the next to the leading order. The saddle-point integrations can be performed conveniently by rescaling $`h`$ and $`H`$ according to $`hh/\sqrt{x};HH/\sqrt{x}`$. This results in
$`\left(Z_\nu ^{\left(n\right)}\left(x\right)\right)_{p=1}`$ $`=`$ $`i\left(1\right)^\nu \left(8\pi \right)^{n1}\left({\displaystyle \frac{2\pi }{x}}\right)^{\left(\left(n1\right)^2+1\right)/2}e^{\left(n2\right)x+\nu n}`$ (131)
$`\times `$ $`\left[1+\left(2\left(n1\right)^3+\left(n1\right)\right){\displaystyle \frac{1}{8x}}{\displaystyle \frac{3}{8x}}n\left(n2\right){\displaystyle \frac{\nu ^2}{2x}}\right].`$
Now we are in a position to calculate the contribution of the $`p=1`$ saddle point to the chiral condensate,
$`{\displaystyle \frac{\mathrm{\Sigma }_{p=1}\left(x\right)}{\mathrm{\Sigma }_0}}=\underset{n0}{lim}{\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{Z_\nu ^n\left(x\right)}}{\displaystyle \frac{}{x}}\left[n\left(Z_\nu ^{\left(n\right)}\left(x\right)\right)_{p=1}\right].`$ (132)
Using trace correlators given in the Appendix, the expansion in $`1/x`$ can be easily extended to order $`1/x^4`$. For the mass dependence of the chiral condensate we obtain
$`{\displaystyle \frac{\mathrm{\Sigma }\left(x\right)}{\mathrm{\Sigma }_0}}`$ $`=`$ $`\underset{n0}{lim}{\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{Z_\nu ^n}}\left(x\right){\displaystyle \frac{}{x}}Z_\nu ^n\left(x\right)`$ (133)
$`=`$ $`1{\displaystyle \frac{i\left(1\right)^\nu e^{2x}}{2x}}+{\displaystyle \frac{\left(4\nu ^21\right)\left(1i\left(1\right)^\nu e^{2x}\right)}{8x^2}}{\displaystyle \frac{i\left(1\right)^\nu e^{2x}\left(4\nu ^21\right)\left(4\nu ^29\right)}{64x^3}}`$
$`+`$ $`{\displaystyle \frac{\left(4\nu ^21\right)\left(4\nu ^29\right)\left[i\left(1\right)^\nu e^{2x}\left(194\nu ^2\right)6\right]}{3\times 256x^4}}.`$
All terms $`ie^{2x}`$ originate from $`p=1`$ saddle point. Before calculating the spectral density by taking the discontinuity of $`\mathrm{\Sigma }\left(x\right)`$ we wish to point out that the expression (133) has been obtained under the assumption that $`\mathrm{Re}\left(x\right)>0`$, so that we cannot calculate the discontinuity from the difference of $`\mathrm{\Sigma }\left(i\lambda +ϵ\right)`$ and $`\mathrm{\Sigma }\left(i\lambda ϵ\right)`$. The reason is that for $`x=i\lambda +ϵ`$ the dominant saddle-point is given by $`I_0`$ and for $`x=i\lambda ϵ`$ it is given by $`I_0`$. The infinitesimal increment thus breaks the replica symmetry between $`I_0`$ and $`I_0`$. The sum over $`p`$ in (104) has been extended to $`\mathrm{}`$ consistent with the breaking of the replica symmetry $`pnp`$ by the saddle point $`I_0`$. Thus, the infinitesimal increment is necessary to resolve the ambiguity as was also the case in the original calculation of . The expression for $`\mathrm{\Sigma }\left(i\lambda ϵ\right)`$ can be obtained from a replica-symmetry breaking solution for which $`I_0`$ dominates. The final result for for the asymptotic expansion of $`\rho _s\left(\lambda \right)`$ coincides with (100) obtained from the expansion of the analytical result $`\left(\lambda /2\right)\left(J_\nu ^2\left(\lambda \right)J_{\nu +1}\left(\lambda \right)J_{\nu 1}\left(\lambda \right)\right)`$ up to $`1/\lambda ^4`$ for $`\lambda >0`$.
We observe that it requires a great deal of effort to derive the asymptotic expansion of the oscillating contribution to the spectral density by means of the replica trick. This is especially true due to the lack of the Virasoro constraints for $`\nu 0`$. Those constraints are a key tool in simplifying the calculations for the mass expansions of the condensate in the sector of vanishing topological charge ($`\nu =0`$). The simplifying feature in the sector with vanishing topological charge that the partition function can be shown to belong to the universality class (see ) of the generalized Kontsevich model with potential $`𝒱\left(x\right)=1/x^2`$ and satisfies the same Virsoro constraints. For $`\nu 0`$ the one-link integral depends also on $`detJ`$ and $`detJ^{}`$ which is also the case when we have integrals over $`SU\left(N_f\right)`$ instead of $`U\left(N_f\right)`$. Therefore, a possible generalization of the Virasoro constraints to include the case $`\nu 0`$ would be certainly welcome not only to reduce our efforts in perturbative calculations but also as an identification of the universality class of the QCD finite volume partition function at nonzero topological charge.
### 4.2 Small Mass Expansion
The situation for the small mass is much more complicated. Before discussing the complications we recall that the partition function $`Z_\nu ^n\left(x\right)`$ has been extensively studied in the context of lattice QCD where it is known as the one-link integral (see for example for an updated review of the subject). In that context one considers unitary matrix integrals of the form
$`Z(J,J^{})={\displaystyle _{UU\left(n\right)}}𝑑Ue^{\mathrm{Tr}\left(JU^{}+J^{}U\right)},`$ (134)
where $`J`$ is a general $`n\times n`$ matrix. The partition function with such potential is a function of the eigenvalues $`\lambda _k`$ of the matrix $`JJ^{}`$. The one-link model exhibits two phases according to
$`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{1}{2\sqrt{\lambda _k}}}`$ $`1\left(\mathrm{weak}\mathrm{coupling}\right),`$
$`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{1}{2\sqrt{\lambda _k}}}`$ $`1\left(\mathrm{strong}\mathrm{coupling}\right).`$
In our case, with partition function given by (102), the eigenvalues of $`J^{}J`$ are given by $`x^2/4`$. From the above we see that, for $`x\mathrm{}`$, we expect to take the replica limit, $`n0`$, and remain in the weak coupling regime while, for $`x0`$, it is not clear whether the replica limit can be taken without crossing a phase boundary. These problems are reflected in logarithmic singularities of the small $`x`$ expansion of the valence quark mass dependence of the partition function and its derivatives. In the replica limit of the expansion coefficients could be derived up to the order for which terms of the form $`x^p\mathrm{log}x`$ are absent. These singular terms could be related to the presence of de Wit-’t Hooft poles .
Below we consider the replica limit of $`Z_\nu ^n\left(x\right)`$ for $`\nu =0`$. In that case logarithmic terms already enter to lowest order in the expansion. We thus consider the small $`x`$ expansion of the partition function $`Z_{\nu =0}^n\left(x\right)=det\left(I_{ij}\left(x\right)\right)`$. By expanding the Bessel functions we obtain
$`Z_0^n\left(x\right)={\displaystyle \underset{k=0}{\overset{n}{}}}\left({\displaystyle \frac{x^2}{4}}\right)^k{\displaystyle \frac{1}{k!}}+{\displaystyle \underset{k=n+1}{\overset{\mathrm{}}{}}}C_{k,n}x^{2k}.`$ (136)
The first sum can be recognized as an incomplete exponential, but we were not able to determine the coefficients $`C_{k,n}`$ in general. In order to expose the $`n`$-dependence of the first term, we rewrite the incomplete exponential as an incomplete gamma function
$`{\displaystyle \underset{k=0}{\overset{n}{}}}\left({\displaystyle \frac{x^2}{4}}\right)^k{\displaystyle \frac{1}{k!}}={\displaystyle \frac{e^{\frac{x^2}{4}}\mathrm{\Gamma }(n+1,x^2/4)}{\mathrm{\Gamma }\left(n+1\right)}}.`$ (137)
From the small mass expansion of the incomplete gamma function
$`\mathrm{\Gamma }(n+1,x^2/4)={\displaystyle _{x^2/4}^{\mathrm{}}}𝑑tt^ne^t=\mathrm{\Gamma }\left(n+1\right)+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{x^2}{4}}\right)^{k+2n}{\displaystyle \frac{\left(1\right)^k}{\left(n+k\right)\left(k1\right)!}},`$ (138)
and neglecting terms which will vanish in the replica limit, i.e.,
$`\mathrm{\Gamma }\left(n+1\right)`$ $`=`$ $`1\gamma n+𝒪\left(n^2\right)`$
$`\left({\displaystyle \frac{\mu ^2}{4}}\right)^n{\displaystyle \frac{1}{n+k}}`$ $`=`$ $`{\displaystyle \frac{1}{k}}+n\left({\displaystyle \frac{2}{k}}\mathrm{log}{\displaystyle \frac{\mu }{2}}{\displaystyle \frac{1}{k^2}}\right)+𝒪\left(n^2\right),`$ (139)
we find
$`Z_0^n\left(x\right)`$ $`=`$ $`\mathrm{\hspace{0.17em}1}+n\left[\left(1e^{x^2/4}\right)\left(2\mathrm{log}{\displaystyle \frac{x}{2}}+\gamma \right)+e^{x^2/4}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{x^2}{4}}\right)^k{\displaystyle \frac{1}{k!k}}\right]+𝒪\left(n^2\right)`$
$`+`$ $`{\displaystyle \underset{k=n+1}{\overset{\mathrm{}}{}}}C_{k,n}x^{2k}.`$
Because the coefficients $`C_{k,n}`$ are unknown only the lowest order terms of the small mass expansion can be calculated,
$`Z_0^n\left(x\right)`$ $`=`$ $`\mathrm{\hspace{0.17em}1}n{\displaystyle \frac{x^2}{2}}\mathrm{log}x+𝒪\left(x^2\right)+𝒪\left(n^2\right).`$ (141)
For the replica limit of the condensate given by,
$`\mathrm{\Sigma }\left(x\right)=\underset{n0}{lim}{\displaystyle \frac{1}{nZ_0^n\left(x\right)}}{\displaystyle \frac{Z_0^n\left(x\right)}{x}},`$ (142)
we thus obtain
$`\mathrm{\Sigma }\left(x\right)=x\mathrm{log}x+𝒪\left(x^1\right).`$ (143)
This is indeed the correct leading order term of the small mass expansion of $`x\left(I_0\left(x\right)K_0\left(x\right)+I_1\left(x\right)K_1\left(x\right)\right)`$. The linear behavior of the microscopic spectral density at the origin is reproduced by taking the discontinuity across the imaginary axis,
$`\rho _s\left(\lambda \right)={\displaystyle \frac{\lambda }{2}}+𝒪\left(\lambda ^2\right).`$ (144)
We did not succeed to generalize this calculation to arbitrary $`\nu `$, but we expect that the logarithmic terms can be obtained in a similar fashion. In particular, the first $`n+\nu `$ terms of the expansion seem to follow a simpler pattern than the coefficients of the higher powers.
## 5 Bosonic Replicas
In view of the seemingly different role played by bosonic and fermionic replicas and the fact that the supersymmetric method uses both compact and noncompact variables it is natural to try to reproduce the results of last section by introducing additional $`n`$ replicas of bosonic quarks of mass $`m`$ instead of fermionic ones. In this case, the condensate for $`N_f=0`$ is given by
$`\mathrm{\Sigma }\left(m\right)=\underset{n0}{lim}{\displaystyle \frac{1}{\left(n\right)}}{\displaystyle \frac{1}{V_4}}{\displaystyle \frac{}{m}}\mathrm{ln}Z_\nu ^{\left(n\right)}\left(x\right),`$ (145)
where $`Z_\nu ^{\left(n\right)}\left(x\right)`$ is given in (4). We will show that, for large masses, bosonic replicas can be used to reproduce the asymptotic expansion of the chiral condensate but they fail to reproduce the microscopic spectral density for a subtle reason which we will explain below.
It is convenient to express the bosonic partition function $`Z_\nu ^{\left(n\right)}\left(x\right)`$ in terms of the eigenvalues of the $`Gl\left(n\right)/U\left(n\right)`$ matrices. Up to an overall constant we have,
$`Z_\nu ^{\left(n\right)}\left(x\right)={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}ds_k{\displaystyle \underset{k<j}{}}\left(e^{s_k}e^{s_j}\right)\left(e^{s_k}e^{s_j}\right)e^{x_{k=1}^n\mathrm{cosh}s_k+\nu _{k=1}^ns_k}.`$ (146)
Notice in particular that the measure (including the Vandermonde determinant) is invariant under the symmetry $`s_ks_k+t`$ which is a remnant of the $`Gl\left(n\right)/U\left(n\right)`$ invariance. The fermionic partition function $`Z_\nu ^{\left(n\right)}\left(x\right)`$ is given by a circular ensemble (for integer $`\nu `$) ,
$`Z_\nu ^{\left(n\right)}\left(x\right)={\displaystyle _\pi ^\pi }{\displaystyle \underset{k=1}{\overset{n}{}}}d\theta _k{\displaystyle \underset{k<j}{}}\left(e^{i\theta _k}e^{i\theta _j}\right)\left(e^{i\theta _k}e^{i\theta _j}\right)e^{x_{k=1}^n\mathrm{cos}\theta _k+i\nu _{k=1}^n\theta _k}.`$ (147)
In the noncompact case, the solutions of the saddle point equation $`\mathrm{sinh}s_k=0`$ are given by $`s_k=0,\pm i\pi ,\pm i2\pi \mathrm{}`$. Thus, in principle we might have a variety of saddle points which should be all taken into account in large $`x`$ expansion. However, we will argue that only $`s_k=0`$ solution contributes to the large $`x`$ behavior of $`Z_\nu ^{\left(n\right)}\left(x\right)`$. Our discussion is based on the $`n=1`$ integral where a steepest descent analysis can be easily carried out. In this case, the bosonic partition function is given by
$`Z_\nu ^{\left(1\right)}\left(x\right)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑se^{x\mathrm{cosh}\left(s\right)+\nu s}=2K_\nu \left(x\right)`$ (148)
From the asymptotic expansion of $`K_\nu \left(x\right)`$ it is clear that only the saddle-point at $`s=0`$ contributes to the integral. The spectral density can be calculated from the resolvent at $`x=\pm i\lambda +ϵ`$ so that the integral above can be identified with the modified Bessel function $`K_\nu \left(x\right)`$.
On the other hand, for one fermionic replica the partition function is given by (assuming integer $`\nu `$)
$`Z_\nu ^{\left(1\right)}\left(x\right)=2\pi I_\nu \left(x\right).`$ (149)
In this case two saddle-points contribute to the asymptotic expansion, one $`e^x`$, and the other one $`e^x`$, which are both essential to recover the oscillatory contributions to the microscopic spectral density.
Let us analyze in detail the saddle-point calculation of the bosonic partition function for one replica. The stationary phase condition that the imaginary part of the action, $`i\lambda \mathrm{cosh}s`$, is constant results in the following curve in the complex $`s`$-plane through the saddle point at $`s=0`$. Clearly, the integration contour (the real axis) can be deformed into the steepest descent curve of fig. 2 (depending on the sign of $`\lambda `$) and no other saddle points need to be considered. A similar analysis for the compact case shows that both the saddle points at $`\theta =0`$ and at $`\theta =\pi `$ have to be taken into account.
A similar situation arises in the $`1/N`$ correction to the semicircle law of the Gaussian Unitary Ensemble (GUE) of hermitian matrices. Inside the semicircle the saddle point solutions of the $`n=1`$ fermionic replica are horizontally aligned in the complex plane, while (see also a comment in ) the solutions for the one bosonic replica are vertically aligned and only the $`p=0`$ saddle point contributes. If we insist on taking into account all vertically aligned saddle points we get an incorrect result. However, we have checked that outside the semicircle the situation is reversed (see figures in for the fermionic case). In this case the bosonic replicas correctly reproduce the semicircle law and its leading exponentially decreasing correction outside the semicircle (see also ) while the fermionic replica calculation gives an incorrect result. In the case of the microscopic spectral density analyzed in this work, we are always inside the semicircle even in the limit $`x\mathrm{}`$ which explains, assuming the same pattern of as for the GUE, the failure of the bosonic replica calculation in reproducing the oscillating part of the spectral density.
If we are just interested in the mass dependence of the chiral condensate it is sufficient to only take into account the $`p=0`$ saddle point and, one can convince oneself that bosonic and fermionic replicas produce the same large $`x`$ result as follows. The integrand of the non-compact partition function is obtained by transforming the integration variables according to $`\theta _k=is_k`$ and replacing $`xx`$. For large $`x`$, the saddle-point of both partition functions is at $`\theta _k=0`$ and $`s_k=0`$, respectively, and one finds that $`Z_\nu ^{\left(n\right)}\left(x\right)`$ and $`Z_\nu ^{\left(n\right)}\left(x\right)`$ have the same asymptotic expansion for the chiral condensate (see (119)). Clearly this proof is formal because it assumes that the replica limit can be interchanged with the operation $`xx`$. However, we have explicitly checked to the order $`1/x^3`$ that the noncompact partition function $`Z_\nu ^{\left(n\right)}\left(x\right)`$ leads to the same asymptotic expansion.
## 6 Conclusions
We have investigated the replica trick for the microscopic spectral density of the QCD Dirac operator in the quenched limit. The advantage of working with fermionic replicas is that this theory corresponds to QCD with $`n`$ flavors of equal mass $`m`$. Because the low energy properties of this theory have are well understood, the starting point of this approach has a firm basis. The valence quark mass dependence of the chiral condensate and the spectral density of the QCD Dirac operator, however, are only obtained in the limit $`n0`$. The existence of this limit has been debated for many years and the investigation of its nature has been the main topic of this article.
The alternative approach to obtain the QCD Dirac spectrum is the supersymmetric method. Although in principle rigorous, one might raise the question whether this theory with bosonic ghost quarks might have unusual properties as for example the spontaneous breaking of supersymmetry. Our results show that this is not the case. The low-energy limit of this partition function is completely dictated by chiral supersymmetry. The power of the supersymmetric method is that one obtains rigorous nonperturbative results such as, for example, the spectral density in the microscopic region.
The replica trick, on the other hand, requires an explicit $`n`$ dependence, and, up to now, only perturbative results have been obtained. Exact results have only been derived in cases where the perturbative series consists of only a finite number of terms. Our results for the microscopic spectral density confirm that the asymptotic series of its non-oscillatory part can be obtained from an expansion about the replica symmetric saddle-point. Both fermionic and bosonic replicas give the same result and are in complete agreement with the asymptotic expansion of the exact result. The same is true for the asymptotic expansion of the resolvent away from the imaginary axis (where the eigenvalues are located). Things are different for the the oscillatory part of the microscopic spectral density. In this case the correct asymptotic expansion is obtained only if a saddle point that breaks the replica symmetry is taken into account. This additional saddle point only exists for fermionic replicas. For bosonic replicas, only the replica symmetric saddle point contributes in the saddle point calculation, and therefore this approach does not reproduce the asymptotic expansion of the oscillatory part of the spectral density. A similar observation has been made for the application of the replica trick to the Wigner-Dyson ensembles.
In the asymptotic expansion of the supersymmetric partition function also two saddle-points have to be taken into account. In the boson-boson component of the Goldstone manifold only one saddle point contributes, but as is the case for one fermionic replica, we have to take into account two saddle-points in the fermion-fermion component of the Goldstone manifold.
Is it possible to go beyond perturbation theory using the replica trick? One indication in favor of an affirmative answer to this question is that we have reproduced the leading order logarithmic singularity of the small mass expansion of the resolvent. Except for these logarithmic terms, the small mass expansion of the resolvent is a convergent series which may be summed to obtain the exact result. The large mass expansion, on the other hand, is an asymptotic series which cannot be summed and cannot provide us with the exact result. Of course, for half integer $`\nu `$ when the asymptotic series terminates, an exact result is obtained from the replica trick.
The exact answer for the resolvent in the microscopic region shows a compact-noncompact dichotomy. This dichotomy is natural in the supersymmetric approach where the compact part of the Goldstone manifold is associated with the fermion-fermion sector and the non-compact part is associated with the boson-boson sector. In the fermionic replica trick this dichotomy is not at all clear but might be hidden in the $`n0`$ limit which is given by an integration over a $`1\times 1`$ matrix and and an integral over an $`\left(n1\right)\times \left(n1\right)`$ matrix. This might be another hint that it is possible to go beyond perturbation theory within the replica framework.
Finally, we hope to have convinced the reader that the supersymmetric method is the only $`\sigma `$model approach that can provide us with rigorous exact results. Even the calculation of a small number of terms in the asymptotic expansion within the replica approach requires a tremendous effort in the case of broken replica symmetry.
Acknowledgments
This work was partially supported by the US DOE grant DE-FG-88ER40388. Dominique Toublan and Poul Damgaard are thanked for useful discussions. J.J.M.V. is grateful to the Institute for Nuclear Theory at the University of Washington for its hospitality and partial support during the completion of this work. The work of D.D. is supported by FAPESP (Brazilian Agency).
Appendix
We are interested in expectation values of traces of powers of $`n\times n`$ Hermitian matrices with matrix elements distributed according to the Gaussian Unitary Ensemble (GUE). Such averages are given by
$`\mathrm{\Omega }_q(p_1,p_2,\mathrm{},p_q)`$ $`=`$ $`\mathrm{Tr}H^{p_1}\mathrm{Tr}H^{p_2}\mathrm{}\mathrm{Tr}H^{p_q}={\displaystyle \frac{𝑑He^{\frac{1}{2}\mathrm{Tr}H^2}\mathrm{Tr}H^{p_1}\mathrm{Tr}H^{p_2}\mathrm{}\mathrm{Tr}H^{p_q}}{𝑑He^{\frac{1}{2}\mathrm{Tr}H^2}}}`$ (150)
$`=`$ $`{\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}_{i=1}^n\left(dx_ie^{\frac{1}{2}x_i^2}\right)_{i<j}\left(x_ix_j\right)^2_{k_1=1}^nx_{k_1}^{p_1}_{k_2=1}^nx_{k_2}^{p_2}\mathrm{}_{k_q=1}^nx_{k_q}^{p_q}}{_{\mathrm{}}^{\mathrm{}}_{i=1}^n\left(dx_ie^{\frac{1}{2}x_i^2}\right)_{i<j}\left(x_ix_j\right)^2}}`$
$`=`$ $`\omega _{p_1}\mathrm{}\omega _{p_q},`$
where $`\omega _l\mathrm{Tr}H^l`$. All correlators $`\mathrm{\Omega }_q(p_1,\mathrm{},p_q)`$ can be calculated recursively starting from $`\mathrm{\Omega }_1\left(0\right)=n`$. The recursion relations are derived from integrals of total derivatives (Schwinger-Dyson equations),
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{n}{}}}dx_i{\displaystyle \underset{k=1}{\overset{n}{}}}_k\left(x_k^{a+1}{\displaystyle \underset{j_1=1}{\overset{n}{}}}x_{j_1}^{r_1}\mathrm{}{\displaystyle \underset{j_b=1}{\overset{n}{}}}x_{j_b}^{r_b}\mathrm{\Delta }_n^2e^{\frac{1}{2}_{i=1}^nx_i^2}\right)=0.`$ (151)
It is useful to keep in mind the identity
$`{\displaystyle \underset{k=1}{\overset{n}{}}}_k\left(x_k^{a+1}\mathrm{\Delta }_n^2\right)=\mathrm{\Delta }_n^2{\displaystyle \underset{l=1}{\overset{a}{}}}\omega _l\omega _{al}.`$ (152)
As a sample calculation let us derive $`\mathrm{\Omega }_3(1,1,4)`$. Choosing $`b=2,r_1=4,r_2=1`$ and $`a=1`$ in (151) we find
$`\mathrm{\Omega }_3(1,1,4)=4\mathrm{\Omega }_2(1,3)+n\mathrm{\Omega }_1\left(4\right).`$ (153)
Choosing now $`b=1,a=1`$ and $`r_1=3`$ or $`r_1=1`$, respectively, we get
$`\mathrm{\Omega }_2(1,3)`$ $`=`$ $`3\mathrm{\Omega }_1\left(2\right),`$
$`\mathrm{\Omega }_2(1,1)`$ $`=`$ $`n.`$ (154)
Finally, from $`b=0`$ and $`a=2`$ or $`a=0`$, respectively, we deduce
$`\mathrm{\Omega }_1\left(2\right)`$ $`=`$ $`n^2,`$
$`\mathrm{\Omega }_1\left(4\right)`$ $`=`$ $`\mathrm{\hspace{0.17em}2}n\mathrm{\Omega }_1\left(2\right)+\mathrm{\Omega }_2(1,1).`$ (155)
This system of equations results in
$`\mathrm{\Omega }_3(1,1,4)=\mathrm{\hspace{0.17em}2}n^4+13n^2.`$ (156)
For some special correlators we can easily derive a general formula, e.g., choosing $`b=2k,a=1`$ and $`r_1=r_2=\mathrm{}=r_{2k}=1`$ we find
$`\mathrm{\Omega }_{2k}(1,1,\mathrm{},1)=n^k\left(2k1\right)!!.`$ (157)
Clearly all $`\mathrm{\Omega }_q(p_1,\mathrm{},p_q)`$ with $`_ip_i`$ being odd vanish identically. Besides (157) we have used the following moments in the calculation of this paper,
$`\begin{array}{cc}\mathrm{\Omega }_1\left(0\right)=n,\hfill & \\ \mathrm{\Omega }_1\left(2\right)=n^2,\hfill & \\ \mathrm{\Omega }_1\left(4\right)=\mathrm{\hspace{0.17em}2}n^3+n,\hfill & \mathrm{\Omega }_2(2,2)=n^4+2n^2,\hfill \\ \mathrm{\Omega }_2(1,3)=\mathrm{\hspace{0.17em}3}n^2,\hfill & \mathrm{\Omega }_3(1,1,2)=n^3+2n,\hfill \\ \mathrm{\Omega }_1\left(6\right)=\mathrm{\hspace{0.17em}5}n^4+10n^2,\hfill & \mathrm{\Omega }_2(2,4)=\mathrm{\hspace{0.17em}2}n^5+9n^3+4n,\hfill \\ \mathrm{\Omega }_2(3,3)=\mathrm{\hspace{0.17em}12}n^3+3n,\hfill & \mathrm{\Omega }_2(1,5)=\mathrm{\hspace{0.17em}10}n^3+5n,\hfill \\ \mathrm{\Omega }_3(1,1,4)=\mathrm{\hspace{0.17em}2}n^4+13n^2,\hfill & \mathrm{\Omega }_4(1,1,1,3)=\mathrm{\hspace{0.17em}9}n^3+6n,\hfill \\ \mathrm{\Omega }_1\left(8\right)=\mathrm{\hspace{0.17em}14}n^5+70n^3+21n,\hfill & \mathrm{\Omega }_2(2,6)=\mathrm{\hspace{0.17em}5}n^6+40n^4+60n^2,\hfill \\ \mathrm{\Omega }_2(4,4)=\mathrm{\hspace{0.17em}4}n^6+40n^4+61n^2,\hfill & \mathrm{\Omega }_3(1,1,6)=\mathrm{\hspace{0.17em}5}n^5+70n^3+30n,\hfill \\ \mathrm{\Omega }_3(1,3,4)=\mathrm{\hspace{0.17em}6}n^5+75n^3+24n,\hfill & \mathrm{\Omega }_5(1,1,1,1,4)=\mathrm{\hspace{0.17em}6}n^5+75n^3+24n,\hfill \\ \mathrm{\Omega }_3(2,4,4)=\mathrm{\hspace{0.17em}4}n^8+72n^6+381n^4+488n^2,\hfill & \\ \mathrm{\Omega }_2(4,6)=\mathrm{\hspace{0.17em}10}n^7+169n^5+610n^3+156n,\hfill & \\ \mathrm{\Omega }_4(1,1,4,4)=\mathrm{\hspace{0.17em}4}n^7+88n^5+661n^3+192n,\hfill & \\ \mathrm{\Omega }_3(4,4,4)=\mathrm{\hspace{0.17em}8}n^9+228n^7+2202n^5+6517n^3+1440n.\hfill & \end{array}`$ (172)
The sum of the coefficients of the correlators can always be checked by means of the $`n=1`$ case where all correlators reduce to one dimensional Gaussian integrals,
$`\mathrm{\Omega }_q(p_1,\mathrm{},p_q)|_{n=1}=\left({\displaystyle \underset{i}{}}p_i\mathrm{\hspace{0.17em}1}\right)!!.`$ (173)
The coefficient of the highest order power in $`n`$ of the correlators $`\mathrm{\Omega }_1\left(2k\right)`$ can be checked as follows. First, in the limit $`n\mathrm{}`$ one can solve the loop equation (or Virasoro constraints) which, for our Gaussian potential, gives the equation
$`{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{1}{px_i}}={\displaystyle \frac{p\sqrt{p^24n}}{2}},`$ (174)
where $`p`$ is a positive number assumed to be large (outside the semi-circle). After expanding both sides of this equation in powers of $`1/p`$ we find
$`\underset{n\mathrm{}}{lim}\mathrm{\Omega }_1\left(2k\right)={\displaystyle \frac{\left(2n\right)^{k+1}\left(2k1\right)!!}{2\left(k+1\right)!}}.`$ (175)
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# Collective modes in uniaxial incommensurate-commensurate systems with the real order parameter
## I Introduction
One of the most useful insights into the properties of stable and metastable ordered states in many body systems follows from the investigations of accompanying collective modes, excitations with a coherent participation of (semi)macroscopic number of particles. The attention is usually focused on the lowest branch in the spectrum. If it is of Goldstone type, i. e. gapless (e. g. acoustic) in the long wavelength limit ($`k0`$), there is a continuous degeneracy in the characterization of ordered state, associated with the breaking of symmetry of the high temperature thermodynamic phase. Without the continuity in the degeneracy one has instead a finite gap at $`k=0`$.
Obvious extrinsic causes for the gap in the Goldstone mode are impurities, defects in the crystal structure, etc. Another cause for the gap is the presence of long range interactions . We do not consider either of these mechanisms here, but remind that, as is well known in charge density wave materials , they may play a decisive role in the collective dynamics of ordered state. Instead, we concentrate on the systems in which the above distinction regarding the degeneracy of ordered state(s) has its origin in short range interactions. Those are numerous materials that show one or more types of uniaxially modulated orderings with periodicities which may be commensurate or incommensurate with respect to the underlying crystal lattice .
Let us at the beginning invoke some simple widely accepted conclusions, accumulated through intense theoretical and experimental investigations on these incommensurate-commensurate (IC) systems in last few decades. In an ideal case of sinusoidal modulation the spectrum of collective excitations contains two types of modes, phasons and amplitudons, representing linearized fluctuations of phase and amplitude of the order parameter respectively. While the amplitudon mode has a finite gap below the critical temperature, the phason mode is acoustic if the free energy of corresponding state does not depend on the relative phase of ordered modulation and crystal lattice. In other words one has the continuous degeneracy with respect to this relative phase. It is strictly fulfilled only if the modulation is incommensurate with respect to the periodicity of crystal lattice.
For commensurate modulations the free energy depends on the relative phase, as is easily seen already from the standard Landau expansions in which the lattice discreteness is taken into account by keeping a leading Umklapp contribution. Within this standard and frequently explored model , which leads to the simple variational equation of sine-Gordon type, the phason mode acquires a gap which is finite only for the strict commensurate ordering, and diminishes rapidly (exponentially) as the order of commensurability increases. For other modulations, which may have the form of dilute soliton lattices, the Goldstone mode remains gapless, although among these modulations there are solutions with commensurate periodicities close to the exempted leading commensurability. In other words, within this model one does not distinguish between ”secondary” commensurate orderings and incommensurate orderings. This is the consequence of a crude simplification made by retaining only one Umklapp term in the free energy. The recent analysis shows that already after taking into account two leading Umklapp terms the phase diagram becomes qualitatively different . It contains a finite number of commensurate states, and shows a harmless staircase, i. e. a series of first order transitions between neighboring states. The Goldstone mode is then expected to have a finite gap for each state participating in the phase diagram.
The Landau models for the orderings with spatial modulations are generally justified providing the interactions responsible for their stabilization are weak enough, so that the variations of order parameter (defined with respect to the appropriately chosen star of wave vectors) are slow at the scale of lattice constant. Two crucial simplifications are then allowed, namely the gradient expansion and the perturbative treatment of lattice discreteness through the truncation of the sum of Umklapp contributions.
In the opposite regime of strong couplings the above spatial continuation is not allowed, and the lattice discreteness leads to qualitatively different properties of phase diagrams and related spectra of excitations, established by numerous analytical, and particularly numerical, studies of spin (e.g. Ising) , displacive (e.g. Frenkel-Kontorova) , and electron-phonon (e.g. Holstein) discrete models. Characteristically for such models, either a finite, sometimes large, number of commensurate modulations in the cases of harmless staircase, or the infinity of them in the cases of complete devil’s staircases, can participate in the phase diagram. All commensurate states then have lowest branches of collective excitations with finite gaps in the limit $`k0`$. We repeat that, while none of these possibilities can be reproduced by Landau model with one Umklapp term, the former harmless staircases with a finite number of commensurate states are realized already within extended Landau models with only two Umklapp terms taken into account .
The analysis of Frenkel-Kontorova and Holstein models established also a new type of instability that involves incommensurate modulations, the so-called transition by breaking of analyticity . Namely, by increasing the coupling constant , or by decreasing temperature , the smooth envelope of an incommensurate periodic modulation becomes nonanalytic. The free energy then depends nonanalytically on the relative phase of modulation and underlying lattice. As a consequence a finite gap opens in the Goldstone branch of collective excitations even for incommensurate modulations.
Already from the beginning of investigations on discrete models it was realized that the above complex features in phase diagrams and spectra of collective excitations have their origin in the nonintegrability of these models, i. e. in the nontrivial chaotic structures of corresponding phase spaces. In this respect it is important to emphasize that, either in their basic form or after the inclusion of further terms, Landau free energy expansions are as a rule the examples of nonintegrable functionals. For example, while the sine-Gordon model with one Umklapp term, as a basic model for class I of IC systems, is integrable, already the inclusion of another Umklapp term brings in the nonintegrability .
The situation is even more intriguing for the class II, i. e. for IC materials with modulations having the period close or equal either to the original or to the dimerized unit cell of crystal lattice. There are numerical and analytical indications that already the minimal , as well as slightly extended , models for this class are not integrable. The consequences of this nonintegrability on the phase diagram are discussed in detail in Ref. . In particular, it is shown that, in addition to simple disordered, commensurate \[i. e. (anti)ferro\] and (almost) sinusoidal incommensurate states, included into previous analyses , the phase diagram contains also an enumerable family of metastable solutions with the periodic alternations of commensurate and incommensurate sinusoidal domains.
In the present work we calculate the spectrum of collective modes for stable and metastable states in systems of class II. The corresponding Landau model is particularly convenient for the discussion of questions raised in this Introduction, since it is nonintegrable, and, in addition, the accompanying phase diagram comprises both commensurate and incommensurate (meta)stable states. Our main aim is to investigate to what extent are the collective modes influenced by the nonintegrability of, here continuous, free energy functional. Furthermore, by analyzing Goldstone modes for the modulated states of the model under consideration we also resolve some controversies present in the literature on its applicability in the description of incommensurate phases in systems of class II. The equivalent analysis for the class I, i. e. for the Landau model with two Umklapp terms, will be presented elsewhere .
The plan of the paper is as follows. The free energy functional for the class II is introduced in Sec. II. In Sec. III we perform the variational procedure up to the second order, taking care about some specific questions related to the thermodynamic minimization . The linear eigenvalue problem associated to the second order variational procedure is discussed in Sec. IV. Here we encounter a generalized Hill problem, since the systems includes four coupled first order equations (in contrast to the standard cases with two equations), and furthermore, since we are looking for the collective modes of highly multiharmonic periodic states. We therefore do not follow a standard way, appropriate for simple sinusoidal incommensurate orderings, but develop for the first time a general formalism, applicable also to other types of Landau models. This formalism enables the determination of Floquet exponents, and of corresponding Bloch basis of eigenfunctions which we consider in Sec.V. The numerical results for the collective modes of all (meta)stable states appearing in the phase diagram are presented in Sec. VI. Concluding remarks along the lines specified in the previous paragraph are given in Sec. VII.
## II Model
The free energy functional for the uniaxial incommensurate systems of class II is given by
$$\stackrel{~}{f}[\stackrel{~}{u}]=\frac{1}{2\stackrel{~}{L}}_{\stackrel{~}{L}}^{\stackrel{~}{L}}\left[d\left(\frac{d^2\stackrel{~}{u}}{d\stackrel{~}{z}^2}\right)^2+c\left(\frac{d\stackrel{~}{u}}{d\stackrel{~}{z}}\right)^2+a\stackrel{~}{u}^2+{\scriptscriptstyle \frac{1}{2}}b\stackrel{~}{u}^4\right]𝑑\stackrel{~}{z}$$
(1)
where $`\stackrel{~}{u}`$ represents the real order parameter and $`\stackrel{~}{L}`$ is the length of the system. This functional is the simplest (minimal) Landau expansion for the systems with minima of free energy density in the reciprocal space close to the center of Brillouin zone, or to the part of its border perpendicular to the uniaxial direction. Then $`c<0`$, and one has to add a highly nontrivial term with the second derivative of $`\stackrel{~}{u}`$ (and presumably positive coefficient $`d`$) in order to ensure the boundness of Landau expansion in the reciprocal space. The rest of the expansion (1) is standard, with $`b>0`$, and $`a`$ becoming negative below the critical temperature of the transition from disordered to uniform (ferro) or dimerized (antiferro) phase. We limit the further analysis to the most interesting regime characterized by $`c<0`$. It includes the incommensurate ordering and the transition to the commensurate ordering (but does not include the transition from disordered to commensurate state which takes place for $`c>0`$) . In this regime the useful dimensionless quantities are
$$z=\sqrt{\frac{c}{d}}\stackrel{~}{z},L=\sqrt{\frac{c}{d}}\stackrel{~}{L},u(z)=\frac{\sqrt{bd}}{c}\stackrel{~}{u}(\stackrel{~}{z}),f[u]=\frac{bd^2}{c^4}\stackrel{~}{f}[\stackrel{~}{u}].$$
(2)
The model (1) can be now represented as the one-parameter problem,
$$f[u]=\frac{1}{2L}_L^L\left[\left(\frac{d^2u}{dz^2}\right)^2\left(\frac{du}{dz}\right)^2+\lambda u^2+{\scriptscriptstyle \frac{1}{2}}u^4\right]𝑑z,$$
(3)
with $`\lambda ad/c^2`$. The parameterization of the phase diagram in the regime $`c<0`$ is thus very simple, since all relationships between different (meta)stable states (like phase transitions, ranges of coexistence of two or more states, etc) can be presented in the one-dimensional $`\lambda `$-space. The knowledge of the actual dependence of this parameter, as well as of the scales which enter into the reduced quantities (2), on the original physical parameters, in particular on temperature, goes together with the specification of microscopic background behind the phenomenological free energy (1). This is a necessary step in any comparison of phase diagram for the model (3) with experimental data for a given material.
In the previous works on the functional (3) we have determined thermodynamically stable states, i. e. its local minima, without taking into considerations statistical fluctuations outside these minima. This mean-field type of approximation is inappropriate for (quasi) one-dimensional systems. It is however usually sufficient for three-dimensional uniaxial systems with strong enough couplings in the perpendicular directions, on which we concentrate here.
The thermodynamic extremalization of functional (3) consists of the standard variational procedure that is equivalent to the classical mechanical one and leads to the corresponding Euler-Lagrange (EL) equation
$$\frac{d^4u}{dz^4}+\frac{d^2u}{dz^2}+\lambda u+u^3=0,$$
(4)
and of the extremalization that involves boundary conditions or some equivalent set of parameters. The general procedure that carefully takes into account the latter aspect is proposed in Ref. . The most interesting result of this approach is obtained for the functional with the kernel that is not explicitly $`z`$-dependent. Then the relation
$$F=H,$$
(5)
holds for each thermodynamic extremum $`u_0(z)`$. Here $`F`$ is the corresponding averaged free energy, and $`H`$ is the integral constant of the problem (4) which corresponds to the Hamiltonian in classical mechanics.
Early considerations of the model (1) led to the suggestion that the mean-field phase diagram contains only disordered \[$`u_d(z)=0`$\], commensurate \[$`u_c(z)=\pm \sqrt{\lambda }`$\], and (almost) sinusoidal \[$`u_s(z)2/\sqrt{3}(\sqrt{1/4\lambda })\mathrm{sin}(z/\sqrt{2})`$\] incommensurate orderings. The commensurate state is thermodynamically stable for $`\lambda <1/8`$ (and for $`a<0`$ in the range $`c>0`$). The incommensurate state is stable in the range $`2<\lambda <\lambda _{id}=1/4`$, while the first order phase transition between the commensurate and the incommensurate states occurs at $`\lambda _{ic}=1.112`$. The more precise values, obtained after taking into account corrections from higher harmonics in the sinusoidal ordering , are $`1.835<\lambda <\lambda _{id}`$ and $`\lambda _{ic}=1.177`$. Also, the wave number of this ordering, $`q`$, slightly deviates from $`1/\sqrt{2}`$ \[i. e. from $`\sqrt{c/(2d)}`$ in the original parameters of Eq. (1)\] as one approaches the left edge of instability, $`\lambda 1.835`$.
While by above solutions of EL equation (4) one exhausts all absolute minima of the free energy (1), the more involved numerical analysis showed the existence of an enumerable series of periodic solutions which are metastable in finite ranges of the parameter $`\lambda `$. The corresponding phase diagram is shown in Fig. 1 in which we ascribe to various solutions symbolic words introduced in Ref. . By their physical content the metastable solutions from Fig. 1 represent periodic trains of successive sinusoidal and uniform segments (see Fig. 1 in ), and, as domain patterns, complete in a natural way, as an inherent outcome of nonintegrable model (1), the phase diagram in the range of coexistence of two corresponding basic types of orderings.
## III Second order variational procedure
The question on which we concentrate now is the thermodynamic stability of a given state $`u(z)`$ which is a solution of EL equation (4) and fulfills additional conditions of Ref. . To this end we have to go beyond the linear terms in the extremalization procedure. Let us therefore at first extend the standard variational procedure to the second order. Later on we shall shortly consider the conditions which follow from the minimization of boundary conditions.
Let $`\eta (z)`$ be the infinitesimal variation with respect to $`u(z)`$, obeying usual conditions at the boundaries $`z=0`$ and $`z=L`$,
$$\eta (z=0)=\eta (z=L)=\eta ^{}(z=0)=\eta ^{}(z=L)=0.$$
(6)
After performing standard partial integrations,
$`{\displaystyle \frac{1}{L}}{\displaystyle _0^L}(\eta ^{})^2𝑑z`$ $`=`$ $`{\displaystyle \frac{1}{L}}\eta ^{}\eta |_0^L{\displaystyle \frac{1}{L}}{\displaystyle _0^L}\eta ^{\prime \prime }\eta 𝑑z,`$ (7)
$`{\displaystyle \frac{1}{L}}{\displaystyle _0^L}(\eta ^{\prime \prime })^2𝑑z`$ $`=`$ $`{\displaystyle \frac{1}{L}}\eta ^{\prime \prime }\eta ^{}|_0^L{\displaystyle \frac{1}{L}}\eta ^{\prime \prime \prime }\eta |_0^L+{\displaystyle \frac{1}{L}}{\displaystyle _0^L}\eta ^{IV}\eta 𝑑z,`$ (8)
the quadratic contribution to the corresponding variation of free energy functional (3) can be expressed in the form
$$\delta ^2ff[u+\eta ]f[u]=\frac{1}{L}_0^L𝑑z\eta 𝒟\eta ,$$
(9)
with
$$𝒟\frac{d^4}{dz^4}+\frac{d^2}{dz^2}+\lambda +3u^2.$$
(10)
The linear differential operator (10) defines the eigenvalue problem
$$𝒟\eta _\mathrm{\Lambda }\eta _\mathrm{\Lambda }^{\prime \prime \prime \prime }(z)+\eta _\mathrm{\Lambda }^{\prime \prime }(z)+\left[\lambda +3u^2(z)\right]\eta _\mathrm{\Lambda }(z)=\mathrm{\Lambda }\eta _\mathrm{\Lambda }(z),$$
(11)
with the boundary conditions for $`\eta (z)`$ specified by Eq. (6). The necessary condition for the thermodynamic stability of solution $`u(z)`$ is given by the requirement that the spectrum $`\mathrm{\Lambda }`$ is non-negative for all normalizable solutions $`\eta _\mathrm{\Lambda }(z)`$ of the problem (10,6).
Since the above procedure strictly respects the boundary conditions (6), it is entirely equivalent to that usually used in classical mechanics. As a consequence the obtained condition for the stability of a given solution $`u(z)`$ holds for any value of the sample length $`L`$. However, neither the extremal solution $`u(z)`$ of EL equation (4), nor the corresponding conditions of thermodynamic stability, should be sensitive to the conditions imposed on the sample surfaces in the physically relevant thermodynamic limit $`L\mathrm{}`$. Therefore the stability condition can be generalized in this limit. In particular, we may ignore the boundary conditions (6), and perform the variational procedure for any infinitesimal variation $`\eta (z)`$, noting, for later purposes, that the requirement of infinitesimality excludes variations $`\eta (z)`$ which would scale as $`|z|^\beta `$ with $`\beta >0`$.
Performing the same steps as before, but now with neglected surface terms in Eqs. (8) (which scale as $`1/L`$), we come again to the linear eigenvalue problem (11), but without a specification on boundary conditions. This means that any complete set of eigenfunctions $`\eta _\mathrm{\Lambda }(z)`$ with the eigenvalues $`\mathrm{\Lambda }`$ \[which are themselves characterized solely by the linear differential equation (11)\] can be used in the representation of a given variation $`\eta (z)`$, and in the corresponding diagonal representation of the free energy (9). The condition $`\mathrm{\Lambda }0`$ thus guarantees the thermodynamic stability of the solution $`u(z)`$ with respect to any infinitesimal variation, provided the system has the well-defined thermodynamic limit as specified above. We note that by this relaxation of boundary conditions (6) we extend the standard ”classical mechanical” second order variational procedure by including a part of, but still not all, ”thermodynamic” variations. The discussion of this question in Appendix A suggests that the criterion of thermodynamical stability is probably entirely covered by the eigenvalue problem (11).
The concise definition of the thermodynamical stability, i. e. of the stability of any (absolutely stable or metastable) local minimum of thermodynamic functional with respect to small fluctuations, is thus:
* A given solution of EL equation (4) is thermodynamically stable if and only if all solutions of Eq.(11) for any $`\mathrm{\Lambda }<0`$ are non-normalizable.
For later purpose it is appropriate to introduce here also the concept of orbital stability, relevant for the behavior of particular solutions in the phase space:
* A given solution $`u(z)`$ of EL equation (4) is orbitally stable if and only if all solutions of Eq.(11) for $`\mathrm{\Lambda }=0`$ are normalizable.
In the next section the above definitions will be used in the study of stability of homogeneous and periodic configurations $`u(z)`$. The crucial assumption in this respect is that the solutions of Eq.(11) smoothly depend on both parameters $`\lambda `$ and $`\mathrm{\Lambda }`$.
Before embarking into the calculation of the spectrum of eigenvalue problem (11), we invoke its general property which follows from the fact that the density of free energy functional (3) does not depend explicitly on the spatial coordinate $`z`$. Then there exists a normalizable solution of equation (11) with $`\mathrm{\Lambda }=0`$, namely $`\eta _0(z)u^{}(z)`$. This is the Goldstone mode that follows from the translational invariance of free energy functional (3), by which $`u(z+z_0)`$ with arbitrary $`z_0`$ is the solution of EL equation (4) if $`u(z)`$ is its solution. We note that Eq.(11) then has, together with the above Goldstone mode, another solution of the form $`\eta _1(z)=w(z)+zu^{}(z)`$, where $`w(z)`$ is some periodic function of the same period as that of the Goldstone mode. Although $`\eta _1(z)`$ is non-normalizable i.e. its norm grows as a power of $`L`$, we consider this non-normalizability as marginal. The power law growth of a solution of Eq.(11) is much easier to control than possible exponential growth of the remaining solutions, if there are any. For special values of the parameter $`\lambda `$ figuring in Eq.(11) with $`\mathrm{\Lambda }=0`$ the only normalizable solution is the Goldstone mode $`u^{}(z)`$, while other solutions have a power law growth in $`z`$, $`z^nu^{}(z)`$ with $`n3`$. These special values of $`\lambda `$ denote the edges of thermodynamical metastability of the corresponding configuration $`u(z)`$.
## IV Floquet theory
The analysis of the eigenvalue problem (11) with periodic functions $`u(z)`$ is based on general Floquet and Bloch theorems for linear differential equations with periodic coefficients. It will be performed in two stages, covered by this and the next section. The aim of the first one, based on the Floquet’s approach, is to answer the question: Whether there exists a normalizable solution $`\eta _\mathrm{\Lambda }(z)`$ for a given value of $`\mathrm{\Lambda }`$? In the second stage we calculate the set of values $`\mathrm{\Lambda }`$ for which normalizable solutions exist, i.e. the spectrum of collective modes, by using the Bloch’s wave number representation.
We begin by showing that the set of values of $`\mathrm{\Lambda }`$ for which the corresponding normalizable solutions $`\eta _\mathrm{\Lambda }(z)`$ may exist is bounded from below. To this end let us rewrite Eq. (11) in the form
$$\stackrel{~}{𝒟}^2\eta _\mathrm{\Lambda }(z)+3u(z)^2\eta _\mathrm{\Lambda }(z)=\left(\mathrm{\Lambda }+\frac{1}{4}\lambda \right)\eta _\mathrm{\Lambda }(z),\stackrel{~}{𝒟}\frac{d^2}{dz^2}+\frac{1}{2},$$
(12)
and introduce the norm of the function $`\eta _\mathrm{\Lambda }(z)`$,
$$\eta _\mathrm{\Lambda }^2\eta _\mathrm{\Lambda }^{}\eta _\mathrm{\Lambda }=\frac{1}{L}_0^L\eta _\mathrm{\Lambda }(z)^{}\eta _\mathrm{\Lambda }(z)𝑑z.$$
(13)
After multiplying Eq. (12) by $`\eta _\mathrm{\Lambda }^{}(z)`$ and integrating with respect to $`z`$ we get
$$\stackrel{~}{𝒟}\eta _\mathrm{\Lambda }(z)^2+3u(z)\eta _\mathrm{\Lambda }(z)^2=\left(\mathrm{\Lambda }+\frac{1}{4}\lambda \right)\eta _\mathrm{\Lambda }(z)^2.$$
(14)
Here it is taken into account that the operator $`\stackrel{~}{𝒟}`$ is hermitean and the function $`u(z)`$ is real. Since the left-hand side of Eq. (14) is strictly positive, we conclude that
$$\mathrm{\Lambda }\mathrm{\Lambda }_{min}=\lambda \frac{1}{4}$$
(15)
for each $`\mathrm{\Lambda }`$ for which the the norm (13) of the function $`\eta _\mathrm{\Lambda }(z)`$ exists. In particular, this means that it is sufficient to reduce a (numerical) analysis of the thermodynamic stability of a given configuration $`u(z)`$ to the search for the normalizable eigenfunctions $`\eta _\mathrm{\Lambda }(z)`$ in the finite interval of $`\mathrm{\Lambda }`$, $`\mathrm{\Lambda }_{min}\mathrm{\Lambda }<0`$.
Before considering Eq. (11) with the general periodic function $`u(z)`$, let us establish the criterion for the thermodynamic stability of the particular homogeneous (ferro or antiferro) solution $`u_c(z)=\pm \sqrt{\lambda }`$ of EL equation (4). Then Eq. (11) reduces to the linear differential equation with constant coefficients, so that the normalizable eigenfunctions must have the form $`\eta (z)e^{ikz}`$ with real values of the wave number $`k`$. The corresponding eigenvalues $`\mathrm{\Lambda }`$ are given by
$$\mathrm{\Lambda }=k^4k^22\lambda ,\lambda <0.$$
(16)
It follows that the homogeneous configuration $`u_c(z)=\pm \sqrt{\lambda }`$ is stable, i.e. that $`\mathrm{\Lambda }>0`$ for any $`k`$, provided that $`\lambda <1/8`$. Note that the latter inequality is just the condition of orbital instability of the homogeneous solution. Namely, the linearization of the EL equation with respect to this solution leads to the linear equation
$$\theta ^{\prime \prime \prime \prime }+\theta ^{\prime \prime }2\lambda \theta =0,$$
(17)
which has normalizable solutions $`\theta (z)`$ only for $`\lambda >1/8`$. Thus, we see that in this simple case the thermodynamic stability excludes the orbital stability, and vice versa, and that two stabilities ”meet” each other in one point, $`\lambda =1/8`$.
### A General Floquet’s procedure
In order to apply the well-known Floquet’s procedure to Eq. (11) with a general periodic function $`u(z)`$, we rewrite this equation in the matrix form
$$\frac{d𝚯(z)}{dz}=𝐀(z;\lambda ,\mathrm{\Lambda })𝚯(z),$$
(18)
where $`𝚯(𝐳)[\eta (z),\eta ^{}(z),\eta ^{\prime \prime }(z),\eta ^{\prime \prime \prime }(z)]^T`$, and the matrix $`\text{A}(z;\lambda ,\mathrm{\Lambda })`$ is given by
$$\text{A}(z)=\left(\begin{array}{cccc}0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\\ \mathrm{\Lambda }\lambda 3u(z)^2& 0& 1& 0\end{array}\right).$$
(19)
The system of linear equations (18) has four linearly independent solutions, $`𝚯_i(z),i=1,\mathrm{},4`$. They form the fundamental matrix
$$\text{F}(z)=[𝚯_1(z),𝚯_2(z),𝚯_3(z),𝚯_4(z)],$$
(20)
which is obviously the solution of equation
$$\frac{d\text{F}(z)}{dz}=\text{A}(z;\lambda ,\mathrm{\Lambda })\text{F}(z).$$
(21)
Without reducing generality we can always choose such initial conditions at $`z=0`$ that $`\text{F}(0)=\text{I}`$, where I is the identity matrix.
The Floquet’s theorem states that whenever the matrix $`\text{A}(z;\lambda ,\mathrm{\Lambda })`$ is a periodic function of variable $`z`$ with a period $`P`$, the fundamental matrix has the form
$$\text{F}(z)=\text{G}(z)\text{e}^{𝚺z},$$
(22)
where $`\text{G}(z)`$ is a matrix which varies periodically with $`z`$, $`\text{G}(z+P)=\text{G}(z)`$, and $`𝚺`$ is a constant matrix. Due to the periodicity of matrix $`\text{G}(x)`$ and the initial condition $`\text{G}(0)=\text{I}`$, the matrix $`𝚺`$ can be expressed in the form
$$𝚺=\frac{1}{P}\mathrm{ln}\text{F}(P).$$
(23)
The matrix $`𝐅(P)`$ is called the monodromy matrix. The eigenvalues of $`𝚺`$, $`\sigma _i`$ are Floquet’s exponents and the eigenvalues of monodromy matrix $`\text{F}(P)`$, $`\rho _i`$, are Floquet’s multipliers. Put in other words, Floquet’s theorem states that for each Floquet multiplier $`\rho _i`$ there exists a solution $`𝚯_i(z)`$ of Eq. (18) with the property
$$𝚯_i(z+P)=\rho 𝚯_i(z).$$
(24)
Floquet multipliers $`\rho _i`$ are in general complex numbers. For a normalizable solution $`𝚯_i(z)`$ the Floquet multiplier $`\rho _i`$ lies on the unit circle in the complex $`\rho `$-plane, and the corresponding Floquet exponent $`\sigma _i`$ is imaginary.
### B Poincaré-Lyapunov theorem
The problem (18) has an additional important property. After the linear transformation $`𝐙(z)[Z_1(z),Z_2(z),Z_3(z),Z_4(z)]^T=𝐓𝚯(z)`$, defined by
$$Z_1=\eta (z),Z_2=\eta ^{\prime \prime }(z),Z_3=2(\eta ^{}(z)+\eta ^{\prime \prime \prime }(z)),Z_4=2\eta ^{}(z),$$
(25)
i. e. by
$$\text{T}(z)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 2& 0& 2\\ 0& 2& 0& 0\end{array}\right),$$
(26)
Eq. (18) is transformed into the equation
$$\frac{d\text{Z}(z)}{dz}=\text{J H}(z;\lambda ,\mathrm{\Lambda })\text{Z}(z),$$
(27)
with
$$\text{H}(z)=\left(\begin{array}{cccc}2(\lambda \mathrm{\Lambda }+3v_0(z)^2)& 0& 0& 0\\ 0& 2& 0& 0\\ 0& 0& 0& \frac{1}{2}\\ 0& 0& \frac{1}{2}& \frac{1}{2}\end{array}\right),\text{J}=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right).$$
(28)
The problem (27) has the Hamiltonian form, characterized by the hermitean matrix H and the simplectic matrix J (i.e. J is antisymmetric and has the property $`\text{J}^2=\text{I}`$). The Poincaré-Lyapunov (PL) theorem for such problems states that the corresponding fundamental matrix, $`𝚽(z)`$, satisfies the relation
$$𝚽^T(z)𝐉𝚽(z)=𝚽^T(0)𝐉𝚽(0).$$
(29)
In other words, $`𝚽^T(z)𝐉𝚽(z)`$ is the ”integral of motion” for the Hamiltonian problem (27). This theorem can be easily checked by differentiating the relation (29) with respect to $`z`$, and taking into account Eq. (27). Since $`𝚽(z)=\text{T}\text{F}(z)`$, and the matrix $`\text{T}^T\text{J}\text{T}\text{J}_1`$ is also simplectic, it follows that
$$𝐅^T(z)𝐉_1𝐅(z)=𝐅^T(0)𝐉_1𝐅(0)=𝐉_1,$$
(30)
i.e. the PL theorem holds for our original fundamental matrix $`\text{F}(z)`$ as well.
From the relation (30) it follows that the matrices $`\text{F}^T(z)`$ and $`\text{F}^1(z)`$ are similar \[$`\text{F}^T(z)=\text{J}_1\text{F}^1(z)\text{J}_1^1`$\]. Thus, if $`\rho _1\rho `$ is the Floquet multiplier of $`\text{F}(P)`$, then $`\rho ^1`$ is also its Floquet multiplier. Furthermore, since in our example the matrix $`\text{F}(P)`$ is real, it follows that $`\rho ^{}`$ and $`\rho ^1`$ are Floquet multipliers as well. These simple relations link four Floquet multipliers of the problem (18) for any periodic solution of EL equation (4) and for any value of parameter $`\lambda `$. The corresponding three possible types of distributions of Floquet multipliers in the complex $`\rho `$-plane are shown in Fig. 2. Floquet multipliers are either complex (a) or real. In the latter case two pairs generally have different values and may be of the same (b) or opposite (c) signs. Fig. 2 does not include the situations with existing collective modes, i. e. when one or two pairs of solutions are normalizable, and the corresponding Floquet multipliers are on the unit circle.
Four Floquet multipliers of the problem (18) can be represented as roots of a polynomial function of fourth order. Since, due to the PL theorem, $`\rho ,\rho ^1,\rho ^{}`$ and $`\rho ^1`$ are all roots of such function, its general form is
$$P_4(\rho )=\rho ^4+a(\lambda ,\mathrm{\Lambda })\left(\rho ^3+\rho \right)+b(\lambda ,\mathrm{\Lambda })\rho ^2+1,$$
(31)
where $`a(\lambda ,\mathrm{\Lambda })`$ and $`b(\lambda ,\mathrm{\Lambda })`$ are, for a given periodic function $`u(z)`$, some smooth real functions of parameters $`\lambda `$ and $`\mathrm{\Lambda }`$ from the matrix (19).
### C Scenarios of thermodynamic (in)stabilities
Having recapitulated the Floquet theory for the Hamiltonian linear problem (18), we address the problem of thermodynamical stability for a given periodic configuration $`u(z)`$. We start by noting that Eqs. (18), (19) and (31) enable some general conclusions about the dependence of the positions of Floquet multipliers in the complex plane on the parameters $`\lambda `$ and $`\mathrm{\Lambda }`$. At first, since $`\rho =0`$ cannot be the root of $`P_4(\rho )`$, it follows that by changing continuously $`\lambda `$ and $`\mathrm{\Lambda }`$ one can come from the distributions (a) or (b) to the distribution (c) in Fig. 2 only by passing through the unit circle. Next, it is easy to determine the positions of Floquet multipliers in the limits $`\mathrm{\Lambda }\mathrm{}`$ and $`\mathrm{\Lambda }\mathrm{}`$, since then we may neglect $`\lambda `$ and $`u^2(z)`$ with respect to $`\mathrm{\Lambda }`$ in the polynomial matrix element of the matrix $`𝐀`$ (19) \[still keeping in mind that $`u(z)`$ defines the period $`P`$ which enters into the definition (23)\].
In the former limit $`\mathrm{\Lambda }\mathrm{}`$ the Floquet multipliers are given by
$$\rho _n=\text{e}^{k_nP},k_n=\mathrm{\Lambda }^{\frac{1}{4}}\text{e}^{i(2n+1)\frac{\pi }{4}},n=0,1,2,3$$
(32)
i.e. the distribution from Fig. 2(a) is realized. Note that for $`\mathrm{\Lambda }<\mathrm{\Lambda }_{min}`$ this distribution cannot pass to that from Fig. 2(c), since in this range of values of $`\mathrm{\Lambda }`$ the unit circle cannot be crossed because the problem (18) does not have normalizable solutions.
In the limit $`\mathrm{\Lambda }\mathrm{}`$ the Floquet multipliers are given by
$$\rho _1=\rho _2^1=\text{e}^{\mathrm{\Lambda }^{\frac{1}{4}}P},\rho _3=\rho _4^1=\text{e}^{i\mathrm{\Lambda }^{\frac{1}{4}}P}.$$
(33)
As is seen in Fig. 3, one then has one pair of normalizable and one pair of non-normalizable solutions.
From the other side, at $`\mathrm{\Lambda }=0`$ one particular Floquet multiplier has the value $`\rho =1`$, and corresponds to the already mentioned Goldstone mode. It has to be at least doubly degenerate, since otherwise the remaining three multipliers could not have symmetric positions required by the PL theorem. Possible distributions of Floquet multipliers for $`\mathrm{\Lambda }=0`$ are shown in Fig. 4. Aside from the possibility that the degeneracy of the Goldstone mode is complete and all four multipliers are at $`\rho =1`$ (a), one may have the remaining two multipliers either on the real axis (b,c), or on the unit circle (d).
Taking into account the above conclusions, we are now able to list possible scenarios of thermodynamic (in)stabilities for the periodic solutions of Eqs. (4,5).
1. The solution $`u(z)`$ is unstable for a given value of $`\lambda `$ if by increasing $`\mathrm{\Lambda }`$ from $`\mathrm{\Lambda }_{min}=\lambda 1/4`$ the Floquet multipliers from the distribution (a) or (b) of Fig. 2 move in such a way to come to the unit circle for some value of $`\mathrm{\Lambda }`$ in the interval $`(\lambda 1/4,0)`$.
2. If the solution $`u(z)`$ is thermodynamically stable, the Floquet multipliers for $`\lambda 1/4<\mathrm{\Lambda }<0`$ are defined either by one complex number not lying on the unit circle (Fig. 2a), or by two real numbers of the same sign ($`r_1,r_2`$) and their reciprocals (Fig. 2b). The latter case has to be realized as $`\mathrm{\Lambda }0`$ from below, since only distributions (b) and (c) from Fig. 4 represent the Goldstone mode for a thermodynamically stable configuration. Thus, at some negative value of $`\mathrm{\Lambda }`$ the distribution from Fig. 2a has to reduce to the double degenerate Floquet multiplier at the real axis ($`r_1=r_2`$), which then evolves into the distribution from Fig. 2b.
3. For special value(s) of control parameter ($`\lambda =\lambda _c`$) the thermodynamic instability of $`u(z)`$ proceeds in a particular way, realized when all four complex Floquet multipliers approach together the point $`\rho _0=1`$ as $`\mathrm{\Lambda }`$ tends to zero from below. The Goldstone mode is then completely degenerate (Fig. 4a). Putting in another way, such instability occurs when the points $`r`$ and $`r^1`$ in Fig.4b tend towards $`\rho _0=1`$ as $`\lambda \lambda _c`$. Note that the distribution of Floquet multipliers from Fig. 4c means that the instability, i.e. the crossing of Floquet multipliers with the unit circle, takes place at some negative value of $`\mathrm{\Lambda }`$. Also, the distribution from Fig. 4d signifies that the remaining non-Goldstone mode is unstable in a finite interval of values of $`\mathrm{\Lambda }`$, starting at some negative value of $`\mathrm{\Lambda }`$.
## V Bloch theory
In order to determine normalizable solutions of Eq. (18) with $`\mathrm{\Lambda }0`$, i.e. the collective modes for given periodic configuration $`u(z)`$ with the period $`P=2\pi /Q`$, we profit from the freedom in choosing boundary conditions for the solutions $`\eta (z)`$, and specify periodic (Born - von Karman) ones. By this we chose the Bloch representation,
$$\eta _k(z)=e^{ikz}\mathrm{\Psi }_k(z),\mathrm{\Psi }_k\left(z+\frac{2\pi }{Q}\right)=\mathrm{\Psi }_k(z),$$
(34)
where $`k`$ is the Bloch wave number limited to the I Brillouin zone ($`Q/2kQ/2`$). The differential equation for the periodic function $`\mathrm{\Psi }_k(z)`$ reads
$`{\displaystyle \frac{d^4\mathrm{\Psi }_k(z)}{dz^4}}+4ik{\displaystyle \frac{d^3\mathrm{\Psi }_k(z)}{dz^3}}+(16k^2){\displaystyle \frac{d^2\mathrm{\Psi }_k(z)}{dz^2}}+2ik(12k^2){\displaystyle \frac{d\mathrm{\Psi }_k(z)}{dz}}+`$ (35)
$`+\left[k^4k^2+\lambda +3u(z)^2\right]\mathrm{\Psi }_k(z)=\mathrm{\Lambda }(k)\mathrm{\Psi }_k(z),`$ (36)
and the normalizability condition is
$$\frac{Q}{2\pi }_0^{\frac{2\pi }{Q}}\mathrm{\Psi }_k^{}(z)\mathrm{\Psi }_k(z)𝑑z=1.$$
(37)
The dependence $`\mathrm{\Lambda }(k)`$, i. e. the spectrum of eigenvalue problem (11), follows from Eqs. (36, 37). Since $`k`$ is quasi-continuous in the limit $`L\mathrm{}`$, this spectrum is for a stable configuration $`u(z)`$ composed of non-negative bands. The corresponding Bloch functions $`\eta _{n,k}(z)`$, where $`n`$ enumerates bands, represent a complete orthonormal set of functions for the problem (11).
From expressions (24) and (34) it follows that the Floquet multiplier for the Bloch function $`\eta _k(z)`$ is given by $`\rho =\text{e}^{ikP}`$. The polynomial function (31) then has the form
$$P_4(\rho )=\left(\rho \text{e}^{ikP}\right)\left(\rho \text{e}^{ikP}\right)\left[\rho ^2+c_2(\lambda ,\mathrm{\Lambda })\rho +1\right],$$
(38)
where $`c_2(\lambda ,\mathrm{\Lambda })`$ is some coefficient. Comparing two representations for $`P_4(\rho )`$ we conclude that the coefficients from the expressions (31) and (38) are linked by relations
$$a(\lambda ,\mathrm{\Lambda })=c_2(\lambda ,\mathrm{\Lambda })2\mathrm{cos}(kP),b(\lambda ,\mathrm{\Lambda })=22c_2(\lambda ,\mathrm{\Lambda })\mathrm{cos}(kP),$$
(39)
i. e. that the eigenvalue $`\mathrm{\Lambda }`$ depends on the wave number $`k`$ only through the function $`\mathrm{cos}(kP)`$. This in particular means that for each band $`\mathrm{\Lambda }(k)`$ we have $`\mathrm{\Lambda }(k)=\mathrm{\Lambda }(k)`$. This is consistent with the symmetry of Eq. (36). Furthermore $`\mathrm{\Lambda }(k+\frac{2\pi }{P})=\mathrm{\Lambda }(k)`$, in accordance with the reduction of wave numbers in (34) to the I Brillouin zone.
The representation (34) is particularly convenient for the analytical discussion of long wavelength limit $`k0`$ for the Goldstone mode for which $`\mathrm{\Lambda }(k=0)=0`$ and $`\mathrm{\Psi }_{k=0}(z)=u^{}(z)`$. To this end we insert the Taylor expansions for small $`k`$,
$$\mathrm{\Psi }_k(z)=u^{}(z)+k\mathrm{\Psi }_1(z)+k^2\mathrm{\Psi }_2(z)+\mathrm{},\mathrm{\Lambda }(k)=k^2\mathrm{\Lambda }_2+k^4\mathrm{\Lambda }_4+\mathrm{},$$
(40)
into Eq.( 36). The requirement that the coefficients in front of leading powers, $`k`$ and $`k^2`$, vanish then leads to the equations
$$\left(\stackrel{~}{𝒟}^2+\lambda \frac{1}{4}+3u^2\right)\mathrm{\Psi }_1=2i\left(2u^{\prime \prime \prime \prime }+u^{\prime \prime }\right),$$
(41)
and
$$\left(\stackrel{~}{𝒟}^2+\lambda \frac{1}{4}+3u^2\right)\mathrm{\Psi }_2=(\mathrm{\Lambda }_2+1)u^{}+6u^{\prime \prime \prime }2i(2\mathrm{\Psi }_1^{\prime \prime \prime }+\mathrm{\Psi }_1^{}),$$
(42)
with the operator $`\stackrel{~}{𝒟}`$ given by Eq. (12). After multiplying equation (42) by $`v^{}`$, integrating with respect to $`z`$, using the fact that $`u^{}`$ is the Goldstone mode, and inserting $`2u^{\prime \prime \prime \prime }+u^{\prime \prime }`$ from Eq. (41), we get the expression for the coefficient $`\mathrm{\Lambda }_2`$,
$$\mathrm{\Lambda }_2=1+6\frac{u^{\prime \prime 2}}{u^2}\frac{\mathrm{\Psi }_1^{}\left(\stackrel{~}{𝒟}^2+\lambda \frac{1}{4}+3u^2\right)\mathrm{\Psi }_1}{u^2},$$
(43)
where $`\mathrm{}`$ stands for the spatial integration, like in Eq. (13). However, the general thermodynamic condition (5) for the functional (1) reads
$$\frac{u^{\prime \prime 2}}{u^2}=\frac{1}{2},$$
(44)
so that Eq. (43) can be written in a more transparent way,
$$\mathrm{\Lambda }_2=2F_2\left[u(z)\right].$$
(45)
Here we introduce the functional
$$F_2\left[u(z)\right]=\frac{\mathrm{\Psi }_1^{}\left(\stackrel{~}{𝒟}^2+\lambda \frac{1}{4}+3u^2\right)\mathrm{\Psi }_1}{u^2}.$$
(46)
Since the operator figuring in this equation just defines the eigenvalue problem (11,12), it is clear that the functional $`F_2\left[u(z)\right]`$ is positive definite for any thermodynamically stable configuration $`u(z)`$. This has two consequences.
Firstly, the common upper limit of the velocity of Goldstone mode, $`v_G=\sqrt{\mathrm{\Lambda }_2}`$, for all thermodynamically stable periodic states is $`v_{G,M}=\sqrt{2}`$. The velocity $`v_G`$ for a given (meta) stable state has the maximum value $`v_{G,max}v_{G,M}`$ when the functional (46) attains its minimum. Like the functional (3), the functional $`F_2[u(z)]`$ depends only on the parameter $`\lambda `$. Thus, taken a given solution $`u(z)`$, we can find the function $`\mathrm{\Psi }_1`$ by solving the inhomogeneous linear differential equation (41), and then determine, by calculating $`F_2[u(z)]`$, the velocity $`v_G`$ as a function of $`\lambda `$. In other words, we have a direct method for the calculation of the velocity of Goldstone mode, not related to the above Floquet-Bloch procedure (but derived from it). It can be used as an independent check of numerical results for the spectrum $`\mathrm{\Lambda }(k)`$ which follow from Eq. (36).
Secondly, the functional (46) attains its minimal value ($`F_2=0`$) if and only if the function $`\mathrm{\Psi }_1(z)`$ vanishes. As is seen from Eq. (41), this is possible only when the solution $`u(z)`$ satisfies the equation
$$2u^{\prime \prime \prime \prime }+u^{\prime \prime }=0.$$
(47)
i.e. when $`u(z)\mathrm{sin}(z/\sqrt{2})`$. The only solution from the phase diagram in Fig. 1 with this property is the almost sinusoidal incommensurate state, denoted by $`s^2`$. Since in the limit $`\lambda \lambda _{id}=1/4`$ it reduces strictly to the above simple sinusoidal dependence on $`z`$ (with the amplitude tending to zero), we conclude that just at the second order phase transition from the incommensurate to the disordered state the velocity of Goldstone mode of incommensurate state attains the maximum value $`v_{G,M}=\sqrt{2}`$. We note that other periodic (and metastable) states $`u(z)`$ from Fig. 1 cannot even approximately satisfy Eq. (47). From the other side, due to the deviations from the sinusoidal form of a given solution, the functional (46) can attain the value $`F_2=2`$, in which case the velocity of Goldstone mode vanishes. As will be seen from the numerical results in the next section, this is indeed the case for all periodic solutions (including the almost sinusoidal configuration $`s^2`$) at the edges of their local thermodynamic stabilities.
Let us conclude this general discussion with the remark on the class of periodic solutions $`u(z)`$ which in addition have the property $`u(z+P/2)=u(z)`$. Since only $`u^2(z)`$, which then has the period $`P/2`$ (and not $`P`$), enters into the problem (18), the corresponding spectrum $`\mathrm{\Lambda }(k)`$ and the Floquet multipliers can be calculated with respect to the former period. The I Brillouin zone is then doubled ($`QkQ`$), and the number of branches of collective modes is halved. In particular, within this choice the value of Floquet multiplier for the Goldstone mode $`u^{}(z)`$, defined by Eq. (24), is $`1`$ and not $`1`$. The approaching of $`\mathrm{\Lambda }=0`$ from below for the stable solution $`u(z)`$ then proceeds like in the point (ii) of Subsection IV.A, but with one pair of Floquet multipliers tending towards the point $`\rho _0=1`$, and the other pair placed at the negative real semiaxis (Fig. 4c). Correspondingly, the Bloch representation of Goldstone solution of Eq. (18) is $`\eta =\text{e}^{\pm iQz}\mathrm{\Psi }(z)`$ with $`\mathrm{\Psi }(z)=\text{e}^{iQz}u^{}(z)`$, i.e. the Goldstone mode is placed at the border of the doubled Brillouin zone. However, the propagation of collective modes takes place in the periodic structure determined by the configuration $`u(z)`$ (and the period $`P`$). Thus the above doubled Brillouin zone has to be folded once to get the physical one, $`Q/2kQ/2`$. In other words, the wave numbers $`k=Q`$ and $`k=0`$ coincide, so that the Goldstone mode is realized as the long-wavelength one for such solutions as well. Furthermore, we note that after this folding the above Taylor expansion (40) and subsequent conclusions on the velocity of Goldstone mode \[Eqs. (45-47)\] follow in the same way for states with the property $`u(z+P/2)=u(z)`$ as well.
## VI Collective modes for systems of class II
In order to derive collective modes for configurations participating in the phase diagram from Fig. 1, we extend the numerical method developed in Refs. to the calculation of eigenvalues and Bloch solutions (34) of the linear problem (11, 18). In the further discussion we shall mostly use the notation $`\mathrm{\Omega }(k)\sqrt{\mathrm{\Lambda }(k)}`$, where $`\mathrm{\Omega }(k)`$ has the meaning of the frequency of collective mode. Note that the energy scale for $`\mathrm{\Omega }(k)`$ \[as well as that for averaged free energies in Fig. 1\] is defined by the last expression in Eq. (2).
The periodic solutions from the phase diagram were determined by solving a system of algebraic equations for coefficients of their Fourier expansions. These Fourier sums were truncated at finite degrees, high enough to ensure a sufficient precision for $`u(z)`$, as well as for corresponding wave number $`Q`$, averaged free energy, etc. The limitation of this method comes from the increase of number of non-negligible Fourier coefficients as the period $`2\pi /Q`$ increases, and the corresponding solutions $`u(z)`$ contain more and more elementary sinusoidal and uniform segments.
Representing the function $`u^2(z)`$ in Eq. (36) by corresponding truncated Fourier series, and writing the function $`\mathrm{\Psi }_k(z)`$ in the same manner,
$$\mathrm{\Psi }_k(z)=a_0+\sqrt{2}\underset{n=1}{\overset{N}{}}\left[a_n\mathrm{cos}(nQz)+b_n\mathrm{sin}(nQz)\right],$$
(48)
we come to the homogeneous linear algebraic system for the coefficients $`a_0,a_1,..,a_N`$ and $`b_1,b_2,\mathrm{},b_N`$. In order to calculate collective modes $`\mathrm{\Lambda }(k)`$, it remains to diagonalize the corresponding ($`2N+1`$) - dimensional matrix. This matrix is generally complex and hermitean. Again, one has to keep a sufficient number of Fourier components in the expansion (48) to get a reliable result for at least two lowest branches in the spectrum $`\mathrm{\Lambda }(k)`$. In actual calculations the truncation at a given number of coefficients $`N`$ is taken as acceptable if for a given branch $`\mathrm{\Lambda }(k)`$ one fulfills to a certain degree of approximation the equality $`\mathrm{\Lambda }(k=0)=\mathrm{\Lambda }(k=Q)`$ \[i.e. the equality $`\mathrm{\Lambda }(k=0)=\mathrm{\Lambda }(k=2Q)`$ for the solutions with the property $`u(z+P/2)=u(z)`$\].
### A Collective modes of states $`u_s(z)`$, $`u_c(z)`$ and $`u_d(z)`$
In Fig. 5 we present the spectrum of collective modes for the incommensurate almost sinusoidal state $`u_s(z)`$, denoted by $`s^2`$ in Fig. 1, choosing few characteristic values of the parameter $`\lambda `$. As announced above, we use here the reduced Brillouin zone, $`Q/2<k<Q/2`$, for all values of $`\lambda `$, except for those for which the periodic modulation is absent ($`\lambda =0.3>\lambda _{id}`$ in Fig. 5c). At the very second order transition from the incommensurate state to the disordered state $`u_d(z)=0`$ ($`\lambda =\lambda _{id}=1/4`$) we present the spectrum in both, reduced and extended, zone schemes (Fig. 5c). Note that due to the additional symmetry of $`s^2`$ state, $`u_s(z+\pi /Q)=u_s(z)`$, the subsequent branches in Figs. 5a-b are not separated by gaps at the zone edges $`k=\pm Q/2`$.
At first, we see that the lowest branch has the property of Goldstone mode \[$`\mathrm{\Omega }(k)k`$ for $`k0`$\] in the whole range of stability of the configuration $`s^2`$. For $`\lambda `$ well below the critical value $`\lambda _{id}`$ ($`\lambda =0.1`$ in Fig. 5a) the subsequent pairs of branches defined in such way are separated by gaps at $`k=0`$. In other words, the general property obtained before in the limit $`\mathrm{\Lambda }=\mathrm{\Omega }^2\mathrm{}`$ by which only one pair of Floquet multipliers is on the unit circle (Fig. 3), is here realized for all values of $`\mathrm{\Lambda }`$. However, as $`\lambda `$ increases the gap between two lowest pairs of branches decreases, and finally disappears for $`\lambda 0.05`$, as is seen in Fig. 5a. Then one has an overlap of branches in a finite range of values of $`\mathrm{\Omega }`$, i.e. all four Floquet multipliers are on the unit circle. This overlap increases, and the minimum of higher branch tends towards $`0`$, as $`\lambda `$ approaches the critical value $`\lambda _{id}`$ ($`\lambda =0.249`$ in Fig. 5b). At $`\lambda =\lambda _{id}`$ this minimum has the value $`\mathrm{\Omega }=0`$, while the slope $`d\mathrm{\Omega }(k)/dk`$ has a finite value which coincides with that of already existing Goldstone branch (Fig. 5c). In other words, just at the second order phase transitions one has two acoustic modes, which, although with same phase velocities $`vlim_{k0}\frac{d\mathrm{\Omega }(k)}{dk}`$, have different dispersions at finite values of $`k`$. For $`\lambda >\lambda _{id}`$ these two branches combine into a single mode which has minima at $`k=\pm Q`$ with a finite value $`\mathrm{\Omega }(Q)`$, and a maximum at $`k=0`$ (Fig. 5c), as it follows directly from the quadratic part of Landau expansion (3).
The dependence of the phase velocity of Goldstone mode, $`v_G`$, on the parameter $`\lambda `$ is shown in Fig. 6. It is finite at $`\lambda =\lambda _{id}`$, decreases as the amplitude of the incommensurate state increases, and vanishes at the metastability edge for the $`s^2`$ state, $`\lambda =1.835`$. This dependence is in accordance with the analytic results (43-47) on the asymptotic behavior of the Goldstone mode.
The spectrum of collective modes for the commensurate state $`u_c(z)=\pm \sqrt{\lambda }`$ \[i.e. $`\stackrel{~}{u}_c(\stackrel{~}{z})=\pm \sqrt{a/b}`$ in the original notation of Eq. (1)\] follows from Eq. (16). This state is thermodynamically stable in the range $`a<0`$ for $`c>0`$ and $`a<c^2/(8d)`$ for $`c<0`$, which comprises positive values of $`c`$, excluded from the analysis after the transformation (2). In order to cover the whole range of stability of $`\stackrel{~}{u}_c(\stackrel{~}{z})`$, we rewrite Eq. (16) in the original notation,
$$\stackrel{~}{\mathrm{\Omega }}^2=d\stackrel{~}{k}^4+c\stackrel{~}{k}^22a=d\left(\stackrel{~}{k}^2+\frac{c}{2d}\right)^2\frac{c^2}{4d}2a$$
(49)
\[with $`\stackrel{~}{k}\sqrt{c/d}k,\stackrel{~}{\mathrm{\Omega }}\sqrt{c^2/d}\mathrm{\Omega }`$ in the range $`c<0`$\]. The second equality in the expression (49) shows that for $`c<0`$ the dispersion curve has minima at $`\stackrel{~}{k}=\pm \sqrt{c/(2d)}`$ (i.e. at $`k=\pm 1/\sqrt{2}`$), with the gap $`\stackrel{~}{\mathrm{\Omega }}(\stackrel{~}{k})`$ equal to $`\sqrt{2ac^2/(4d)}`$ \[i.e. $`\sqrt{2(\lambda 1/8)}`$ in the reduced scale $`\mathrm{\Omega }`$\]. As for the range $`c>0`$, it follows from the first equality in Eq. (16) that the collective mode has the minimum at $`\stackrel{~}{k}=0`$, with the gap $`\stackrel{~}{\mathrm{\Omega }}(0)=\sqrt{2a}`$. The gap vanishes at $`a=0`$, i.e. at the line of second order transition from the commensurate state $`\stackrel{~}{u}_c(\stackrel{~}{z})`$ to the disordered state $`\stackrel{~}{u}_d(\stackrel{~}{z})=0`$.
The commensurate solutions $`\stackrel{~}{u}_c(\stackrel{~}{z})=\pm \sqrt{a/b}`$, which here represent the uniform or dimerized ordering for the Landau expansions (1) around the center or the border of the original Brillouin zone respectively, have the same symmetry properties as the solution for the disordered state, $`\stackrel{~}{u}_d(\stackrel{~}{z})=0`$. The only collective excitations with finite activated frequencies are fluctuations of the amplitude $`\stackrel{~}{u}(\stackrel{~}{z})`$ with the above dispersion relation (49). Note that the mode of Goldstone (acoustic) type is absent. Since the solutions $`\stackrel{~}{u}_c(\stackrel{~}{z})`$ possess, as constants, a trivial translational degeneracy, we prefer to associate this absence of acoustic branch with its reduction to the trivial dependence $`\stackrel{~}{\mathrm{\Omega }}(\stackrel{~}{k})=0`$. The purpose of this interpretation will become clear in the next subsection.
Finally, as it follows directly from the expression (1), the disordered state $`\stackrel{~}{u}_d(\stackrel{~}{z})=0`$ which is stable in the range $`a>0,c>\sqrt{4ad}`$, has a branch of collective excitations with the minimum at $`\stackrel{~}{k}=0`$ for $`c>0`$, and with two minima at $`\stackrel{~}{k}=\pm \sqrt{c/2d}`$ for $`c<0`$. The respective gaps at these minima are equal to $`\sqrt{a}`$ (for $`c>0`$), and to $`\sqrt{ac^2/4d}`$ (for $`c<0`$).
### B Collective modes of periodic metastable states
An illustration of spectra of collective modes for metastable states is shown in Fig. 7. We take the state $`sd`$, chose the value of control parameter somewhere in the middle of corresponding region of stability from Fig. 1 ($`\lambda =1`$), and plot four lowest branches of collective modes. Spectra for all other metastable states from Fig. 1 have the same qualitative properties, and therefore are not plotted. More specifically, for all states, and for all values of $`\lambda `$ within the respective ranges of stability, the subsequent branches are separated by finite gaps, i. e. there is no branch overlap, like that obtained for the state $`s^2`$ \[Figs. 5a,b\].
Furthermore, the lowest branch for all states is the Goldstone mode with the dispersion $`\mathrm{\Omega }(k)v_Gk`$ for $`k0`$, and with the corresponding phase velocity $`v_G`$ vanishing for the values of parameter $`\lambda `$ at the edges of stability. The dependence of $`v_G`$ on $`\lambda `$ for all metastable states from Fig. 1 is shown in Fig. 8. The characteristic scales for these velocities, given by maxima $`v_{G,max}`$ of curves $`v_G(\lambda )`$ for each metastable state, are situated in the range of values ($`0.40.5`$ in dimensionless units of Fig. 8). This is to be compared with the maximum value of about $`1.4`$ of $`v_G`$ for the configuration $`u_s(z)`$ (Fig. 6). In this respect one may recognize a rough tendency by which $`v_{G,max}`$ decreases as the proportion of incommensurate ($`s`$) domains decreases. This decrease is particularly evident as one compares $`s^2`$ with $`s^9d^3`$, and with other configurations from Fig. 1 of Ref. , in which the proportion of commensurate domains $`d`$ is larger and both types of domains become rather short. The decrease of $`v_{G,max}`$ is then saturated, i. e. values of $`v_{G,max}`$ are roughly concentrated in the narrow range $`0.470.48`$.
The above tendency can be plausibly interpreted along the lines from the preceding subsection. The metastable states are in fact domain trains, built as successions of segments with local sinusoidal ($`u_s`$) and commensurate ($`u_c`$) orderings. As was already stated, it is plausible to associate to a commensurate segment a Goldstone mode with the vanishing frequency (and the vanishing velocity as well). The total Goldstone mode, which is some hybrid of these vanishing contributions and the contributions from the local sinusoidal ordering, tends to be softer and softer as the train has more and more commensurate domains. As a consequence, the velocity $`v_{G,max}`$ gradually decreases as the proportion of commensurate segments in metastable states increases.
## VII Conclusions
The results presented in Sec.5 show that the spectra of collective excitations for all periodic states, stable and metastable, from the phase diagram of the model (1,3) \[Fig. 1\] have Goldstone branches with a linear dispersion $`\mathrm{\Omega }=v_Gk`$ in the long wavelength limit. Thus, although these spectra belong to the nonintegrable model, they have standard characteristics that essentially follow from the absence of an explicit $`x`$-dependence of free energy density in Eq. (1). The latter property of the free energy in turn ensures the translational degeneracy of all solutions of EL equation (4), including those participating in the phase diagram. In this respect the present spectrum does not differ qualitatively from those of integrable models with the same property.
The fact that the chaotic content of the phase space for nonintegrable models, like for that defined by Eq. (1 or for other examples , does not have as substantial impact on the spectrum of collective excitations as it has on the thermodynamic phase diagram, can be interpreted in the following way. The states from the phase diagram belong to the subset of solutions of EL equation defined by conditions like Eq. (5). They are localized in the orbitally unstable chaotic layers which cover the phase space, have the measure zero in this space, and are mutually separated by topological barriers with characteristic heights given by averaged free energies of these layers . These barriers do not allow for smooth changes from one state to others, and as such represent an intrinsic mechanism for frequently observed phenomena like memory effects and thermal hysteresis, as discussed in detail in Ref.. In general, nonintegrable free energy functionals have more complex phase diagrams than integrable ones.
On the other hand, collective modes belong to another space of states, denser than the phase space, i. e. to that defined by the second order variational procedure and the corresponding eigenvalue problem (11). All states in this space are realizable as thermodynamic fluctuations. They have usual properties of double periodic linear systems, although the corresponding Bloch functions $`\mathrm{\Psi }_k(z)`$ in Eq. (34) may be far from a simple sinusoidal form. These properties are not essentially dependent on the level of integrability of free energy functional.
In order to resolve the eigenvalue problem (11) for the model (1, 3), we formulate here a method based on the general Floquet-Bloch formalism, applicable to any IC system showing stable multiharmonic (i.e. non-sinusoidal) periodic ordering(s). Beside being a basis for the numerical calculations of eigenvalues and eigenfunctions (Sec. 5), this approach clearly indicates that for more complex models and orderings the traditional notions of phasons and amplitudons are not appropriate. In particular, it was often claimed that, being an expansion in terms of real order parameter, the functional (1) itself is insufficient for the stabilization of modulated states in the systems of II class, since incommensurate states, in particular those with soliton lattice like modulations, should have to be described by at least a two-dimensional order parameter . Also, the absence of phase variable in Eq.(1) caused a belief that the states which emerge from this functional do not have an acoustic (phason-like) collective mode. However, while the previous study led to the conclusion that almost sinusoidal and highly non-sinusoidal configurations are among (meta)stable states of the model (1) (as is seen in Fig.1), the present analysis shows that Goldstone modes are well defined for all these configurations. From the other side, all dispersive modes for the homogeneous ($`u=`$const.) states are massive, i.e. have finite gaps.
The gap of the lowest such mode tends to zero at continuous (2nd order) phase transitions from one homogeneous state to another, or to some periodic ordering. The examples are the lines ($`c>0,a=0`$) and ($`c<0,\lambda ad/c^2=1/4`$), representing the transitions from the disordered state to the commensurate and incommensurate states respectively. As is shown in Fig. 5, the situation is qualitatively different at the transition from the disordered state to the incommensurate, almost sinusoidal, one. The reason is the specific behavior of Goldstone mode in the incommensurate state. By approaching the transition from the incommensurate side the phase velocity of this mode, $`v_G`$, remains finite, while, as is shown elsewhere , its oscillatory strength tends to zero. In fact, the above behavior of Goldstone mode for the state $`s^2`$ at the second order transition to the disordered state is exceptional. Namely, the Goldstone modes in the (meta)stable states behave critically at the edges of stabilities for these states, including the lower edge of state $`s^2`$ at $`\lambda =1.835`$. At these edges the phase velocities $`v_G`$ vanish. All these specific properties of collective modes, particularly of the most interesting Goldstone modes, are expected to be directly experimentally observable in X-ray and neutron scatterings, as well as in optical and similar measurements. The particular discussion on the role of these collective modes in the dielectric response, and the comparison with measurements on some materials of the II class, is given in Ref..
Finally, we comment on the general property of Goldstone modes for metastable periodic states by which they become softer and softer as the period of these states increases. This tendency, shown in Fig. 8, has its origin in the elastic nature of Goldstone modes in the long wavelength limit. More specifically, as the segments of local sinusoidal order become more and more dilute in the underlying commensurate background, the slight variations in their mutual distances cost less and less energy, i. e. the corresponding effective elastic constant decreases. In this interpretation, which holds for dilute soliton lattices as well, the commensurate ordering is by assumption perfectly elastic, i. e. the notion of relative distance has no sense since the lattice discreteness is neglected. The only possible deformations are those invoking the variations of amplitude, and resulting in the massive collective modes. The lattice discreteness introduces, through an ”external” potential of Peierls-Nabarro type, the finite stiffness of the local commensurate ordering, or even opens the gap in the Goldstone mode for dilute incommensurate states at the transition by breaking of analyticity .
ACKNOWLEDGMENTS
The work was supported by the Ministry of Science and Technology of the Republic of Croatia through project No. 119201.
## A
The procedure from Sec. 3 takes into account, after relaxing boundary conditions (6), all infinitesimal variations of the order parameter $`u(z)`$. This generalization includes some thermodynamic variations, like e. g. those, specified by the scaling $`u(z)su(z)`$ with $`s1`$, responsible for the condition obeyed by $`u(z)`$ at boundaries $`z=0`$ and $`z=L`$ (condition B in Ref. ). However, by this procedure the analysis of thermodynamic stability is still not completed, since there remain variations which invoke infinitesimal relative changes in the configuration $`u(z)`$, but are not infinitesimal at the absolute scale. An example is the scaling
$$u(z)u[(1+ϵ)z],ϵ0,$$
(A1)
which leads to the condition (5). The variation that corresponds to this scaling is not infinitesimal. Indeed, after the transformation $`(1+ϵ)zz`$ in the integral (3), it follows that this variation behaves as $`z`$ and therefore does not fulfill the criterion of infinitesimality specified in Sec. 3. Thus the above procedure has to be enlarged by including the expansion of the free energy with respect to $`ϵ`$ up to the quadratic terms. While the requirement that the linear term vanishes gives the condition (5), the second order variation reads
$`\delta ^2f`$ $``$ $`f\left[u((1+ϵ)z)+\eta \right]f[u(z)]=`$ (A2)
$`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle _0^L}𝑑z\left[\eta (z)𝒟\eta (z)+2(u^{}(z))^2ϵ^2+4(2u^{\prime \prime \prime \prime }(z)+u^{\prime \prime }(z))\eta (z)ϵ\right],`$ (A3)
i.e. the expression (9) is extended by the term quadratic in $`ϵ`$, and the term representing the bilinear coupling between $`\eta (z)`$ and $`ϵ`$.
The previous analysis of the model (1,3) led to the conclusion that all solutions $`u(z)`$ of the EL equation (4) that participate in the thermodynamic phase diagram as stable or metastable configurations, are simple periodic. The analysis in Sec.4 shows that the corresponding eigenfunctions $`\eta _\mathrm{\Lambda }(z)`$ of the problem (11) are then double periodic. This means that for periodic extrema $`u(z)`$ the bilinear coupling in the expression (A3) vanishes, i. e. the fluctuations in $`ϵ`$ are decoupled from $`\eta (z)`$ fluctuations. The remaining $`ϵ^2`$ term is positively definite, i. e. all periodic configurations satisfying the EL equation (4) and the condition (5) are also stable with respect to the variation defined by the scaling (A1), irrespectively to the value of the control parameter $`\lambda `$ figuring in the functional (3).
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# References
Transition Strength Sums and Quantum Chaos in Shell Model States
V.K.B. Kota<sup>a,</sup> <sup>1</sup><sup>1</sup>1Invited talk in the National Seminars on Nuclear Physics held at Institute of Physics, Bhubaneswar during July 26-29, 1999., R. Sahu<sup>b</sup>, K. Kar<sup>c</sup>, J.M.G. Gómez<sup>d</sup> and J. Retamosa<sup>d</sup>
<sup>a</sup>Physical Research Laboratory, Ahmedabad 380 009, India
<sup>b</sup>Physics Department, Berhampur University, Berhampur 760 007, India
<sup>c</sup>Saha Institute of Nuclear Physics, 1/AF Bidhannagar,
Calcutta 700 064, India
<sup>d</sup>Departamento de F$`\stackrel{´}{i}`$sica At$`\stackrel{´}{o}`$mica, Molecular y Nuclear, Universidad Complutense de Madrid, E-28040 Madrid, Spain
ABSTRACT: For the embedded Gaussian orthogonal ensemble (EGOE) of random matrices, the strength sums generated by a transition operator acting on an eigenstate vary with the excitation energy as the ratio of two Gaussians. This general result is compared to exact shell model calculations, with realistic interactions, of spherical orbit occupancies and Gamow-Teller strength sums in some $`\left(ds\right)`$ and $`\left(fp\right)`$ shell examples. In order to confirm that EGOE operates in the chaotic domain of the shell model spectrum, calculations are carried out using two different interpolating hamiltonians generating order-chaos transitions. Good agreement is obtained in the chaotic domain of the spectrum, and strong deviations are observed as nuclear motion approaches a regular regime (transition strength sums appear to follow the Dyson’s $`\mathrm{\Delta }_3`$ statistic). More importantly, they shed new light on the newly emerging understanding that in the chaotic domain of isolated finite interacting many particle systems smoothed densities (they include strength functions) define the statistical description of these systems and these densities follow from embedded random matrix ensembles; some EGOE calculations to this end are presented.
1. Introduction
In the last fifteen years there has been an explosive growth in the use of random matrix theories for quantum systems particularly in the context of quntum chaos . Recently, working with the aim of developing a statistical theory for finite interacting many particle systems, such as atoms, molecules, nuclei, atomic clusters, metallic quantum dots etc., by incorporating the ideas of random matrices and chaos, several research groups recognized the importance of investigating the embedded random matrix ensembles in detail, i.e. the embedded Gaussian orthogonal ensemble of random matrices of $`k`$-body interactions (EGOE($`k`$)) and their various deformations \[2-7\]. The EGOE is introduced in the context of nuclear shell model studies ; The EGOE($`k`$) is defined in $`m`$-particle spaces (i.e. in the $`\left(\begin{array}{c}𝒩\\ m\end{array}\right)`$ dimensional space generated by distributing $`m`$ fermions over $`𝒩`$ single particle states) with a GOE representation in $`k`$-particle space for $`k`$-body operators (usually $`k<<m`$). This ensemble and its deformed versions are studied using Monte-Carlo methods with EGOE constructed in occupation number representation and alternatively assuming that realistic nuclear (similarly atomic) shell model hamiltonians are typical members of EGOE(2) and its deformations, generic features of these ensembles are inferred from exact shell model calculations. For finite interacting particle systems in the chaotic domain there are now several studies of the statistical properties (both smoothed forms and fluctuations) of energy levels, wavefunction amplitudes or equivalently transition strengths generated by action of a transition operator on an eigenstate and strength functions \[1-11\]. However, only in the last 2-3 years similar studies (i.e. in the context of quantum chaos) on expectation values of operators, which measure transition strength sums, as function of excitation energy have began \[3,8,12-14\]. Given an operator $`K=𝒪^{}𝒪`$, the expectation values $`K^E`$ are the diagonal elements of $`K`$ in the hamiltonian ($`H`$) diagonal basis (a more precise definition is given ahead in Eq. (3)); they give strength sums generated by the transition operator $`𝒪`$ acting on the eigenstate with energy $`E`$. Two major examples are single particle transfer strength sums which are expectation values of number operators that give occupancies of single particle states (they determine thermodynamic behaviour) and Gamow-Teller (GT) strength sums as function of excitation energy which are relevant in astrophysics (presupernova evolution and stellar collapse). It is expected that the smoothed $`K^E`$ vs $`E`$ will give information about order-chaos transitions just as energies, wavefunction amplitudes and transition strengths.
Two important results given by EGOE are that in strongly interacting shell model spaces (essentially in $`0\mathrm{}\omega `$ spaces) (i) the state densities take Gaussian form and (ii) the bivariate strength densities take bivariate Gaussian form . These results have their basis in the EGOE representation of the hamiltonian $`H`$ (which is in general one plus two-body in nuclear case) and transition operators $`𝒪`$. The eigenvalue density $`I(E)`$ or its normalized version $`\rho (E)`$ is defined by
$$\begin{array}{c}I\left(E\right)=\delta \left(HE\right)=d\delta \left(HE\right)=d\rho \left(E\right);\\ \rho \left(E\right)\stackrel{\text{EGOE}}{}\overline{\rho \left(E\right)}=\rho _𝒢\left(E\right)=\frac{1}{\sqrt{2\pi }\sigma }exp\frac{1}{2}\left(\frac{Eϵ}{\sigma }\right)^2\end{array}$$
In (1) $`\mathrm{}`$ denotes trace (similary $`\mathrm{}`$ denotes average), the $`ϵ`$, $`\sigma `$ and $`d`$ are centroid, width ($`\sigma ^2`$ is variance) and dimensionality respectively. Note that $`ϵ=H`$, $`\sigma ^2=(Hϵ)^2`$, ‘$`𝒢`$’ stands for Gaussian and the bar over $`\rho (E)`$ indicates ensemble average (smoothed) with respect to EGOE. The strength $`R(E,E^{})`$ generated by a transition operator $`𝒪`$ in the $`H`$-diagonal basis is $`R(E,E^{})`$ = $`E^{}𝒪E^2`$. Correspondingly the bivariate strength density $`I_{biv;𝒪}(E,E^{})`$ or $`\rho _{biv;𝒪}(E,E^{})`$ which is positive definite and normalized to unity is
$$\begin{array}{c}\begin{array}{ccc}\hfill I_{biv;𝒪}(E,E^{})& =& 𝒪^{}\delta \left(HE^{}\right)𝒪\delta \left(HE\right)\hfill \\ & =& I^{}\left(E^{}\right)E^{}𝒪E^2I\left(E\right)\hfill \\ & =& 𝒪^{}𝒪\rho _{biv;𝒪}(E,E^{});\hfill \end{array}\\ \begin{array}{c}\rho _{biv;𝒪}(E,E^{})\stackrel{\text{EGOE}}{}\overline{\rho _{biv;𝒪}(E,E^{})}=\rho _{biv𝒢;𝒪}(E,E^{})\end{array}\end{array}$$
The bivariate Gaussian $`\rho _{biv𝒢;𝒪}(E,E^{})`$ in (2) is defined by the centroids and variances of its marginal densities and a bivariate correlation coefficient. Though the EGOE forms in (1,2) are derived by evaluating the averages over fixed-$`m`$ spaces, in large number of numerical shell model examples it is verified that they apply equally well in fixed- $`m`$, $`mT`$ and $`mJT`$ spaces. In practice the so-called Edgeworth corrections are added to the Gaussian forms in (1,2). Before proceeding further two remarks are in order:
1. Level and strength fluctuations for EGOE follow the GOE fluctuations, i.e. nearest spacing distribution is Wigner distribution, Dyson-Mehta $`\mathrm{\Delta }_3\left(L\right)`$ statistic follow the GOE $`ln\left(L\right)`$ behaviour for large $`L`$ and strength fluctuations are of Porter-Thomas (P-T) type. The $`\mathrm{\Delta }_3`$ form is tested ahead in Fig. 3 and recent tests of P-T form for EGOE are given in Fig. 1.
2. EGOE smoothed forms (1,2) (and (3) ahead) gave birth to the so called statistical nuclear spectroscopy (see \[9,10,15-17\] and references therein) and there are recent studies of this in atoms , molecules and solids and mesoscopic systems \[3-5\].
The pupose of this paper is to first point out that EGOE, via (1,2) gives rise to a statistical theory for the smoothed forms for transition strength sums and the theory operates in the chaotic domain of the spectrum. Shell model tests in $`(ds)^6`$ space for occupancies and GT strength sums in $`(fp)^6`$ space are carried out. The EGOE theory and the shell model studies are described in Sect. 2. In order to confirm that the agreement between shell model and EGOE theory is a consequence of chaoticity of the shell model spectrum, GT strength sums and occupancies are calculated using two different interpolating hamiltonian that generate order-chaos transitions. Results of these calculations form Sect. 3. Results in Sects. 2, 3 give nuclear physics examples for the newly emerging understanding that in the chaotic domain of isolated finite interacting many particle systems smoothed densities, defined by EGOE, give rise to the statistical description of these systems. Further remarks on this important generic result for quantum chaos in isolated finite interacting particle systems are given (together with the results from a EGOE calculation for occupancies) in the concluding Section 4.
2. EGOE theory for transition strength sums and shell model tests
One important by product of (2) is that the transition strength sum density $`𝒪^{}𝒪\delta (HE)`$, which is a marginal density of the bivariate strength density, takes a Gaussian form as the marginal of a bivariate Gaussian is a Gaussian. Therefore, using (1), it is immediately seen that transition strength sums vary with excitation energy as ratio of Gaussians. Given $`K=𝒪^{}𝒪`$, the transition strength sum density is the expectation value density denoted by $`\rho _K(E)`$ and then ,
$$\begin{array}{c}K^E=\left[d\left(m\right)\rho \left(E\right)\right]^1\left[\underset{\alpha \epsilon E}{}E\alpha KE\alpha \right]=I_K\left(E\right)/I\left(E\right)=\rho _K\left(E\right)/\rho \left(E\right)\\ \stackrel{\text{EGOE}}{}\overline{\rho _K\left(E\right)}/\overline{\rho \left(E\right)}=\rho _{K:𝒢}\left(E\right)/\rho _𝒢\left(E\right),\\ \rho _K\left(E\right)=K\delta \left(HE\right)=d^1I_K\left(E\right)=d^1K\delta \left(HE\right);K=𝒪^{}𝒪.\end{array}$$
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Fig. 1. (a) Distribution of renormalized shell model transition strengths compared with the Porter-Thomas (P-T) form. The locally averaged strengths $`\overline{R(Ei,Ef)}`$ are calculated using the EGOE bivariate Gaussian form (2). Shell model (SM) calculations are in 223 dimensional $`\left(ds\right)^{m=5,J=5/2,T=1/2}`$ space with hamiltonian defined by Wildenthal interaction and the transition operator is the two-body part of the hamiltonian obtained after substracting the configuration isospin centroid part. The SM+EGOE result is in good agreement with P-T. (b) Number of principal components (NPC) and information entropy ($`S^{info}`$) versus energy ($`E`$) for a strength distribution in 307 dimensional $`\left(ds\right)^{m=6,J=2,T=0}`$ space. The hamiltonian is Kuo interaction with <sup>17</sup>O single particle energies and the transition operator is same as in (a) but for the Kuo interaction. The exact shell model results are compared with the GOE and EGOE predictions; the EGOE formulas (they use P-T) are given in . (c) Same as (b) but for wavefunctions.
In deriving (3) it is assumed that the smoothed form of $`\rho _K\left(E\right)/\rho \left(E\right)`$ reduces to ratio of smoothed forms of $`\rho _K\left(E\right)`$ and $`\rho \left(E\right)`$. This result ignores the fluctuations in both $`\rho _K\left(E\right)`$ and $`\rho \left(E\right)`$ and the r.m.s error due to neglect of the fluctuations is given in terms of the number of principal components (NPC) or the inverse participation ratio for the transition operator $`𝒪`$ . Note that (3) takes into account $`(K,H)`$ and $`(K,H^2)`$ correlations which define the centroid $`ϵ_K`$ and width $`\sigma _K`$ of $`\rho _K\left(E\right)`$; $`ϵ_K=<KH>/<K>`$ and $`\sigma _K^2=<KH^2>/<K>ϵ_K^2`$. First discussions of the EGOE result in (3) are in . Statistical models that are inappropriate, are applied recently in the study of GT strength sums as function of excitation enegy in $`\left(ds\right)`$-shell although there are several studies of GT strengths and strength sums in statistical nuclear spectroscopy .
In order to study the domain of validity of (3), for the occupancies $`n_\alpha ^E`$ shell model calculations in 307 dimensional $`(2s1d)^{m=6,J=2,T=0}`$ space are carried out using the Rochester - Oak Ridge shell model code and for the GT strength sums in 814 dimensional $`(1f2p)^{m=6,J=0,T=0}`$ space using the NATHAN code of the Strasbourg-Madrid group. In the $`(ds)`$-shell studies the hamiltonian employed is $`H=h(1)+V(2)`$ defined by Kuo’s two-body matrix elements ($`V(2)`$) and <sup>17</sup>O single particle energies ($`h(1)ϵ_{d_{5/2}}=4.15`$ MeV, $`ϵ_{d_{3/2}}=0.93`$ MeV, $`ϵ_{s_{1/2}}=3.28`$ MeV). Similarly in the $`(fp)`$ shell studies the so called KB3 interaction is employed with $`h(1)ϵ_{f_{7/2}}=0.0`$ MeV, $`ϵ_{f_{5/2}}=6.5`$ MeV, $`ϵ_{p_{3/2}}=2`$ MeV, $`ϵ_{p_{1/2}}=4`$ MeV). The expectation value density $`\rho _{n_\alpha :𝒢}`$ for the number operators $`n_\alpha `$ in $`(ds)`$ shell and $`\rho _{K(GT):𝒢}`$ for the $`K(GT)`$ operator that generates GT strength sums in $`(fp)`$ shell are constructed in terms of their centriods and widths and similarly the state density Gaussian. Then using (3) the smoothed form of GT strength sum as function of excitation energy is constructed and compared with exact shell model results. From Fig. 2 it is seen that the EGOE result (3) describes very well the shell model results except at the edges of the spectra. The reason for the deviation at the edges is well known - here the states are not chaotic (sufficiently complex). In the $`(ds)`$ shell example the $`K`$-density centroid $`ϵ_K`$, width $`\sigma _K`$, skewness $`\gamma _{1:K}`$ and excess $`\gamma _{2:K}`$ (the $`\gamma _1`$ and $`\gamma _2`$ are shape parameters that measure deviations from the Gaussian form) are $`ϵ_K=35.08`$MeV, $`\sigma _K=9.63`$MeV, $`\gamma _{1:K}=0.08`$ and $`\gamma _{2:K}=0.09`$ for $`d_{5/2}`$ density and ($`29.50`$ MeV, 10.24 MeV, $`0.08`$, $`0.17`$) respectively for the $`d_{3/2}`$ density. Similarly for the state density the parameters are $`ϵ=32.78`$MeV, $`\sigma =10.24`$MeV, $`\gamma _1=0.05`$ and $`\gamma _2=0.18`$. In the $`(fp)`$ shell example, $`ϵ_K=11.12`$MeV, $`\sigma _K=8.65`$MeV, $`\gamma _{1:K}=0.09`$, $`\gamma _{2:K}=0.18`$, $`ϵ=9.51`$MeV, $`\sigma =8.62`$MeV, $`\gamma _1=0.10`$ and $`\gamma _2=0.19`$. Firstly the $`|\gamma _1|`$ and $`|\gamma _2|`$ values (being much less than $`0.3`$) clearly show that all the densities are close to Gaussian. Moreover $`\widehat{\sigma }=\sigma _K/\sigma 1`$. The scaled centroid shifts $`\widehat{\mathrm{\Delta }}_K=(ϵ_Kϵ)/\sigma `$ are $`0.225`$ and 0.32 for the $`d_{5/2}`$ and $`d_{3/2}`$ densities while it is 0.187 for the $`(fp)`$ shell GT example. With $`K^E`$ given as ratio of Gausians and that $`\widehat{\sigma }1`$ imply that in the middle of the spectrum $`K^EK\left\{1+\widehat{\mathrm{\Delta }}_K\widehat{E}\right\}`$; $`\widehat{E}=(Eϵ)/\sigma `$, i.e $`K^E`$ is linear in energy. This linear form and its polynomial extensions are used in the past . The value of $`\widehat{E}`$ for the ground state are $`2.6`$ and $`3.1`$ for $`(ds)`$ and $`(fp)`$ shell examples respectively. The range of validity (from the center of the spectrum) of the linear form for $`K^E`$ is much smaller for $`(ds)`$-shell than in $`(fp)`$-shell as $`\widehat{\mathrm{\Delta }}_K`$ for the former is larger than for the later; see . Fig. 2 showing good agreement between shell model results for $`n_\alpha ^E`$ and $`K(\text{GT})^E`$ with realistic interactions and the EGOE prediction (3), gives a justification to our assertion that EGOE starts operating from the region of the onset of chaos for the strength sums as well. Next we study this question via $`n_\alpha ^E`$ and $`K(\text{GT})^E`$ when the hamiltonian changes thorough a parameter from a symmetry preserving, i.e. regular hamiltonian to a chaotic one.
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Fig. 2. (a) Occupation numbers for $`d_{5/2}`$ and $`d_{3/2}`$ orbits vs excitation energy ($`E`$) in the same example as in Fig. 1b. Shown are also the shell model state and occupation number densities (histograms) compared to the EGOE Gaussian forms (continuous curves) given by (1,3). (b) GT strength sum versus excitation energy ($`E`$) for the 814 dimensional six particle $`\left(fp\right)`$-shell space with $`J=0`$, $`T=0`$. The exact shell-model results for the realistic KB3 interaction are compared with the EGOE predictions given by (3).
3. Results with MF and SU(4) interpolating hamiltonians
For further confirming the conclusions from Fig. 2, shell model calculations are performed in the 325 dimensional $`(ds)^{m=8,J=T=0}`$ space for occupancies and GT strength sums with two different interpolating hamiltonians . First set of calculations use the spherical shell model mean-field (MF) hamiltonian $`h(1)`$ as the unperturbed hamiltonian $`H_0`$ and in this case the occupation number operators commute with $`H_0`$,
$$H_\lambda \left(MF\right)=h\left(1\right)+\lambda V\left(2\right)=H_0+\lambda \left(H_{SM}H_0\right)$$
(1)
Note that $`H_{SM}=h(1)+V(2)`$; the $`h(1)`$ is defined by <sup>17</sup>O single particle energies and $`V(2)`$ by Kuo’s two-body matrix elements as in the Sect. 2 $`(ds)`$-shell example. In the figures in Figs. 3a-d the calculations with (4) are denoted by MF. It is easily seen that spherical configurations (generated by distributing the nucleons in the three $`(ds)`$ shell orbits) are eigenstates for $`H_0`$. Therefore for $`\lambda =0`$ in (4), the spectrum will have degeneracies. In the second set of calculations the $`SU(4)ST`$ scalar part $`H_{SU(4)ST}`$ of $`H_{SM}`$ is used as the unperturbed hamiltonian,
$$H_\lambda \left(SU\left(4\right)\right)=H_{SU\left(4\right)ST:scalar}+\lambda \left(H_{SM}H_{SU\left(4\right)ST:scalar}\right)=H_0+\lambda \left(H_{SM}H_0\right)$$
(2)
In the figures in Figs. 3a-d calculations with (5) are denoted by SU(4). Note that the GT operator commutes with the SU(4) hamiltonian $`H_{SU(4)ST}`$. For $`H=H_{SU(4)ST}`$, the eigenvalues and eigenvectors are given easily by $`SU(4)ST`$ algebra. The eigenstates are labelled, for a given number $`m`$ of valence nucleons, by $`L`$, $`S`$, $`J`$, $`T`$ and the $`U(4)`$ irreducible representations (irreps) $`\{f\}=\{f_1f_2f_3f_4\}`$ or the $`SU(4)`$ irreps $`\{F\}=\{F_1F_2F_3\}`$ where $`f_1+f_2+f_3+f_4=m`$, $`f_1f_2f_3f_40`$, $`F_1=(f_1+f_2f_3f_4)/2`$, $`F_2=(f_1f_2+f_3f_4)/2`$, $`F_3=(f_1f_2f_3+f_4)/2`$ with additional restrictions on $`f_i`$’s coming from the spatial part. In addition, as GT operator is a generator of the $`SU(4)`$ group, the $`SU(4)`$ algebra gives directly $`K(\text{GT})^E`$ in terms of $`SU(4)ST`$ quantum numbers. Methods for constructing $`H_{SU(4)ST}`$ part of a given $`H`$ are given in ,
$$\begin{array}{ccc}\hfill H_{SU\left(4\right)ST}& =& \frac{1}{2}\left(n1\right)\left(n2\right)E(0,\left\{0\right\}00)n\left(n2\right)E(1,\left\{1\right\}\frac{1}{2}\frac{1}{2})\hfill \\ & & +\frac{1}{8}\left[n\left(2n9\right)+G_2+2\left(S^2+T^2\right)\right]E(2,\left\{2\right\}11)\hfill \\ & & \frac{1}{8}\left[nG_2+2\left(S^2+T^2\right)\right]E(2,\left\{2\right\}00)\hfill \\ & & +\frac{1}{8}\left[n\left(n+3\right)G_2+2\left(S^2T^2\right)\right]E(2,\left\{11\right\}10)\hfill \\ & & +\frac{1}{8}\left[n\left(n+3\right)G_22\left(S^2T^2\right)\right]E(2,\left\{11\right\}01)\hfill \end{array}$$
In (6) $`n`$ is number operator, $`G_2`$ is $`U(4)`$ quadratic Casimir operator with eigenvalues $`G_2^{\{f\}}=f_1(f_1+3)+f_2(f_2+1)+f_3(f_31)+f_4(f_43)`$ and $`S^2`$ and $`T^2`$ are operators with eigenvalues $`S(S+1)`$ and $`T(T+1)`$. Construction of $`H_{SU(4)ST}`$ requires the values for the centriod energies (determined by the given $`H`$) $`E(m,\{f\}ST)=H^{m,\{f\}ST}`$. For the $`(ds)`$-shell $`H_{SM}`$ hamiltonian (defined after (4)), they are $`E(0,\{0\}00)=0`$, $`E(1,\{1\}\frac{1}{2}\frac{1}{2})=2.302`$MeV, $`E(2,\{2\}11)=4.176`$MeV, $`E(2,\{2\}00)=2.975`$MeV, $`E(2,\{11\}10)=8.360`$MeV and $`E(2,\{11\}01)=7.048`$MeV. Using these, $`H_\lambda `$ is constructed and then occupancies and $`K(\text{GT})^E`$ are calculated for various $`\lambda `$ values. The $`SU(4)ST`$ algebra gives for $`\lambda =0`$, $`K(\text{GT})^E=\frac{2}{3}(C_2(SU(4))^{\{F\}}S(S+1))`$; $`C_2(SU(4))^{\{F\}}`$ = $`F_1(F_1+4)+F_2(F_2+2)+F_3^2`$. In addition, the allowed $`\{f\}S`$ values for the $`(ds)`$-shell example are: $`\{2222\}0`$; $`\{3221\}1`$; $`\{3311\}2,0`$; $`\{332\}1`$; $`\{4211\}1`$; $`\{422\}2,0`$; $`\{431\}1,2,3`$; $`\{44\}0,2,4`$; $`\{5111\}0`$; $`\{521\}1,2`$; $`\{53\}1,3`$; $`\{611\}1`$; $`\{62\}2`$. With these results, the GT curve for $`\lambda =0`$ case is constructed as a check of the shell model calculations. Each of the $`\{f\}S`$ states will have several $`L=S`$ states and therefore the states here have degeneracies.
It is clearly seen from the results in Figs. 3a,b that the EGOE smooth form is not a good approximation to the exact results in the case of regular motion. For $`\lambda 0`$ there are several (approximately) good quantum numbers with nearby levels carrying different sets of quantum numbers and therefore expectation values show large fluctuations as a function of excitation energy. The order-chaos transition as $`\lambda `$ increases is clearly illustrated by the spectral regidity $`\mathrm{\Delta }_3`$ in Figs. 3c,d (also by the distribution of nearest neighbour level spacings as shown in ). It is seen that occupancies and GT strength sums behave rather like the $`\mathrm{\Delta }_3`$ statistic, approching more slowly the EGOE and GOE limits respectively. This similarity is probably due to the fact that both $`\mathrm{\Delta }_3`$ and strength sums are related to long-range correlations between the energy levels or wavefunctions. Thus, in the quantum chaotic domain transition strength sums, independent of the hamiltonian, follow EGOE forms and we have statistical spectroscopy in the chaotic domain.
There is another important observation that follows from Figs. 3a-d, i.e. as the interacting particle system becomes chaotic, expectation values take smoothed forms (within GOE fluctuations ) and hence described by the smoothed densities. For $`\lambda >>\lambda _c`$ ($`\lambda _c`$ corresponds to order-chaos border and we determine $`\lambda _c`$ by using $`\chi _{(\lambda )}^2`$ which is the mean square deviation of the exact shell model strength sum from the EGOE smoothed form ; see Fig. 4) they take the EGOE form given by (3). This generic result is of central interest in quantum chaos studies of finite interacting particle systems as discussed in the next section.
4. Further results and concluding remarks
Besides the nuclear shell model results for GT strength sums and occupation numbers presented in Figs. 2-4 (GT results for the SU(4) case are reported in recently), the behaviour of occupancies as a many particle hamiltonian makes order-chaos transitions is studied recently by several groups: (i) using a 20 member EGOE(1+2) in 330 dimensional $`𝒩=11,m=4`$ space with the MF hamiltonian (4); $`h(1)`$ is defined by the single particle enegies $`ϵ_i=i+(1/i);i=1,2,\mathrm{},11`$ and $`V(2)`$ is EGOE(2) ; (ii) using the four interacting electrons Ce atom ; (iii) using a symmetrized coupled two-rotor model ; (iv) using EGOE(1+2) as in (i) but in the 924 dimensional $`𝒩=12,m=6`$ space with 25 members . Most significant conclusion of all these studies is that transition strength sums show quite different behaviour in regular and chaotic domains of the spectrum. In order to make this argument clear results of (iv) are shown in Fig. 5; here occupancies for the single particle states are calculated for various values of the interpolating parameter $`\lambda `$. It is clearly seen, for example from Fig. 5, that below the region of onset of chaos
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Fig. 3. (a) Occupation numbers for $`d_{5/2}`$ and $`d_{3/2}`$ orbits vs excitation energy ($`E`$) for the MF and SU(4) interpolating hamiltonians given by (4) and (5). Shell model results are compared with the EGOE predictions (continuous curves) given by (3). (b) same as (a) but for GT strength sums. (c) Averaged spectral rigidity $`\mathrm{\Delta }_3\left(L\right)`$ for the eigenvalues of the MF hamiltonian (4). Error bars give the standard deviation of the $`\mathrm{\Delta }_3`$ average over overlaping intervals of length $`L`$. The dashed curves are for Poisson and the continuous curves are for GOE. (d) same as (c) but for the SU(4) hamiltonian (5). For the $`\lambda =1`$ case the parameters $`(ϵ,\sigma ,\gamma _1,\gamma _2)`$ for the state, $`d_{5/2}`$, $`d_{3/2}`$ and GT densities are ($`52.59`$ MeV, 13.15 MeV, 0.10, 0.03), ($`55.30`$ MeV, 12.31 MeV, 0.01, $`0.06`$), ($`48.70`$ MeV, 13.42 MeV, 0.07, 0.05) and ($`48.37`$ MeV, 12.64 MeV, 0.06, $`0.1`$) respectively.
Fig. 4. $`R\left(\lambda \right)=\chi _{\left(\lambda \right)}^2/\chi _{\left(1\right)}^2`$ vs $`\lambda `$. Results for MF and SU(4) calculations are shown for $`d_{5/2}`$ occupancies and GT strength sums. From the observations in the studies in , a plausiable definition that the $`H_\lambda `$ system is chaotic is given by the condition $`R\left(\lambda \right)1.5`$. Then, $`\lambda _c`$ is defined by $`R\left(\lambda _c\right)=1.5`$. The $`\lambda _c`$ values for the four cases are shown in the figure. The difference in the two $`\lambda _c`$ values in the case of occupancies is easily understood in terms of the norms of $`H_0`$ and $`H_1=H_{SM}H_0`$; however this is not the case with GT strength sums. Thus it appears that a complete theory for $`\lambda _c`$, in case of strength sums, may come from the study of $`\mathrm{\Delta }_3`$, strength fluctuations, NPC and information entropy for the hamiltonian and the transition operator involved.
Fig. 5. Occupation numbers for a 25 member EE(1+2) ensemble, defined by the hamiltonian $`h\left(1\right)+\lambda \left\{V\left(2\right)\right\}`$, in the 924 dimensional $`𝒩=12`$, $`m=6`$ space (see text); in the figure $`𝒩`$ is denoted as $`N`$. Details of matrix construction etc. are given in . Results are shown for the lowest 5 single particle states and for six values of $`\lambda `$. In the calculations occupation numbers are averaged over a bin size of 0.1 in $`\left(Eϵ\right)/\sigma `$; $`ϵ`$ is centroid and $`\sigma `$ is width. The spectra of all the ensemble members are first zero centered and scaled to unit width and then the ensemble average is carried out. The estimate of gives $`\lambda _c0.05`$ for order-chaos border in the present EE(1+2) example. It is clearly seen that once chaos sets in, the occupation numbers take stable smoothed forms. See for further details.
occupation numbers show strong fluctuations (in the regular ground state domain perturbation theory applies). In this region there is no equilibrium distribution for $`I_K\left(E\right)`$, $`I\left(E\right)`$ and other densities. However in the chaotic domain (see and also Fig. 4 for methods of determining $`\lambda _c`$) the densities can be replaced by their smoothed forms as in (1)-(3). Therefore there is a statistical mechanics, defined by various smoothed densities, operating in the quantum chaotic domain of finite interacting particle systems (this is also the essence of statistical nuclear spectroscopy \[9,10,15-17\]; see for arguments for statistical spectroscopy in atoms). In fact in favourable situations, it is possible to introduce thermodynamic concepts (effective temperatures and chemical potencials etc.) in the chaotic domain. In order to firmly establish the universality of these results, it is essential to carry out numericl studies for a wide variety of interacting particle systems and investigate various deformed EGOE’s in detail.
Acknowledgements This work is partially supported by DGES (Spain) Project No. PB96-0604 and DST(India).
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# High energy QCD and hard diffraction at HERA versus Tevatron
## 1 Forward jets at HERA and Mueller-Navelet jets at Tevatron
### 1.1 Forward jets at HERA
The study of forward jets at $`ep`$ colliders is considered as the milestone of QCD studies at high energies, since it provides a direct way of testing the perturbative resummations of soft gluon radiation. It is similar to the previous proposal of studying two jets separated by a large rapidity interval in hadronic colliders , for which preliminary results are available .
The cross-section for forward jet production at HERA in the dipole model reads :
$`{\displaystyle \frac{d^{(4)}\sigma }{dxdQ^2dx_Jdk_T^2d\mathrm{\Phi }}}={\displaystyle \frac{\pi N_C\alpha ^2\alpha _S(k_T^2)}{Q^4k_T^2}}f_{eff}(x,\mu _f^2)\mathrm{\Sigma }e_Q^2{\displaystyle _{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}}{\displaystyle \frac{d\gamma }{2i\pi }}\left({\displaystyle \frac{Q^2}{k_T^2}}\right)^\gamma \times `$ (1)
$`\times \mathrm{exp}\{ϵ(\gamma ,0)Y\}\left[{\displaystyle \frac{h_T(\gamma )+h_L(\gamma )}{\gamma }}(1y)+{\displaystyle \frac{h_T(\gamma )}{\gamma }}{\displaystyle \frac{y^2}{2}}\right]`$
where
$`Y`$ $`=`$ $`\mathrm{ln}{\displaystyle \frac{x_J}{x}}`$ (2)
$`ϵ(\gamma ,p)`$ $`=`$ $`\overline{\alpha }\left[2\psi (1)\psi (p+1\gamma )\psi (p+\gamma )\right]`$ (3)
$`f_{eff}(x,\mu _f^2)`$ $`=`$ $`G(x,\mu _f^2)+{\displaystyle \frac{4}{9}}\mathrm{\Sigma }(Q_f+\overline{Q_f})`$ (4)
$`\mu _f^2`$ $``$ $`k_T^2,`$ (5)
are, respectively, $`Y`$ the rapidity interval between the photon probe and the jet, $`ϵ(\gamma ,p)`$ the BFKL kernel eigenvalues, $`f_{eff}`$ the effective structure function combination, and $`\mu _f`$ the corresponding factorization scale. The main BFKL parameter is $`\overline{\alpha },`$ which is the (fixed) value of the effective strong coupling constant in LO-BFKL formulae.
The so-called “impact factors” $`h_T`$ and $`h_L`$ are obtained from the $`k_T`$ factorization properties of the coupling of the BFKL amplitudes to external hard probes. The same factors can be related to the photon wave functions within the equivalent context of the QCD dipole model .
The main problem to solve is to investigate the effect of the experimental cuts on the determination of the integration variables leading to a prediction for $`d\sigma /dx`$ from the given theoretical formula for $`d^{(4)}\sigma `$ as given in formula (1). The effect appears as bin-per-bin correction factors to be multiplied to the theoretical cross-sections for average values of the kinematic variables for a given $`x`$-bin before comparing to data .
The experimental correction factors have been determined using a toy Monte-Carlo designed as follows. We generate flat distributions in the variables $`k_T^2/Q^2`$, $`1/Q^2`$, $`x_J,`$ using reference intervals which include the whole of the experimental phase-space (we chose the variables which minimize the variation of the cross-section over the measured kinematical range. The correction factors are given in Reference . We perform a fit to the H1 and ZEUS data with only two free parameters. these are the effective strong coupling constant in LO BFKL formulae $`\overline{\alpha }`$ corresponding to the effective Lipatov intercept $`\alpha _P=1+4\mathrm{log}2\overline{\alpha }N_C/\pi `$, and the cross-section normalisation. The obtained values of the parameters and the $`\chi ^2`$ of the fit are given in Table I for a fit to the H1 and ZEUS data separately, and then to the H1 + ZEUS data together.
| fit | $`\overline{\alpha }`$ | Norm. | $`\chi ^2(/dof)`$ |
| --- | --- | --- | --- |
| H1 | 0.17 $`\pm `$ 0.02 $`\pm `$ 0.01 | 29.4 $`\pm `$ 4.8 $`\pm `$ 5.2 | 5.7 (/9) |
| ZEUS | 0.20 $`\pm `$ 0.02 $`\pm `$ 0.01 | 26.4 $`\pm `$ 3.9 $`\pm `$ 4.7 | 2.0 (/2) |
| H1+ZEUS | 0.16 $`\pm `$ 0.01 $`\pm `$ 0.01 | 30.7 $`\pm `$ 2.9 $`\pm `$ 3.5 | 12.0 (/13) |
Table I- Fit results
### 1.2 Mueller Navelet Jets at Tevatron (D0)
It is fruitful to compare our results with the effective intercept we obtain from recent preliminary dijet data obtained by the D0 Collaboration at Tevatron . The measurement consists in the ratio $`R=\sigma _{1800}/\sigma _{630}`$ where $`\sigma `$ is the dijet cross-section at large rapidity interval $`Y\mathrm{\Delta }\eta `$ for two center-of-mass energies (630 and 1800 GeV), $`\mathrm{\Delta }\eta _{1800}=4.6`$, $`\mathrm{\Delta }\eta _{630}=2.4.`$ The experimental measurement is $`R=2.9\pm 0.3`$ (stat.) $`\pm 0.3`$ (syst.). Using the Mueller-Navelet formula , this measurement allows us to get a value of the effective intercept for this process. We get $`\alpha _P`$=1.65 $`\pm `$ 0.05 (stat.) $`\pm `$ 0.05 (syst.), in agreement with the value obtained by D0 using a saddle-point approximation . This intercept is higher than the one obtained in the forward jet study. The question arises to interpret the different values of the effective intercept. It could reasonably come from the differences in higher order QCD corrections for the BFKL kernel and/or in the impact factors depending on the initial probes .
## 2 Diffraction at HERA versus Tevatron
### 2.1 Diffraction at HERA
The data taken in 1992 by the H1 and ZEUS experiments showed some new interesting events with an interval of rapidity around the incident proton direction devoid of any hadronic activity. This means that a colourless exchange (”pomeron”) must have occured since there is no colour connection between the remnant proton and the struck quark. The selection used to tag these events is mainly by asking a gap in rapidity in the forward proton direction <sup>1</sup><sup>1</sup>1There are other selection used by H1 and ZEUS collaborations, namely tagging a proton in the final state or using the $`M_X`$ subtraction method .. These events represent about 10% of the total deep inelastic scattering events.
The statistics obtained in 1994 allowed a measurement of the proton diffractive structure function defined in analogy to the standard proton structure function in a wide kinematical domain . The data accumulated between 1995 and 1997 allowed to extend the measurement to lower and higher $`Q^2`$. In addition to the usual deep inelastic variable $`x`$ (the momentum fraction of the interacting quark), and $`Q^2`$ (the transfered energy squared between the electron and the interacting quark), two other kinematical variables are used, namely $`\beta `$, the momentum fraction of the colourless exchanged object, and $`x_{\mathrm{l}\mathrm{P}}=x/\beta `$, the momentum fraction of the parton inside this object if we assume it has a partonic structure.
The diffractive structure function $`F_2^{D(3)}`$ can be investigated in the framework of Regge phenomenology when both a leading ($`\mathrm{l}\mathrm{P}`$) and a sub-leading ($`\mathrm{l}\mathrm{R}`$) trajectory are considered, such that
$`F_2^{D(3)}(Q^2,\beta ,x_{\mathrm{l}\mathrm{P}})=f_{\mathrm{l}\mathrm{P}/\mathrm{p}}(x_{\mathrm{l}\mathrm{P}})F_2^{\mathrm{l}\mathrm{P}}(Q^2,\beta )+f_{\mathrm{l}\mathrm{R}/\mathrm{p}}(x_{\mathrm{l}\mathrm{P}})F_2^{\mathrm{l}\mathrm{R}}(Q^2,\beta ).`$ (6)
In this parameterisation, $`F_2^{\mathrm{l}\mathrm{P}}`$ can be interpreted as the structure function of the pomeron. The values of $`F_2^{\mathrm{l}\mathrm{R}}(Q^2,\beta )`$ are taken from a parameterisation of the pion structure function , with a single free normalisation. The pomeron flux is assumed to follow a Regge behaviour with a linear trajectory $`\alpha _{\mathrm{l}\mathrm{P}}(t)=\alpha _{\mathrm{l}\mathrm{P}}(0)+\alpha _{\mathrm{l}\mathrm{P}}^{^{}}t`$.
Using H1 1994 data , the resulting value of $`\alpha _{\mathrm{l}\mathrm{P}}(0)`$ is $`\alpha _{\mathrm{l}\mathrm{P}}(0)=1.20\pm 0.09`$ and is significantly larger than values extracted from soft hadronic data ($`\alpha _{\mathrm{l}\mathrm{P}}1.08`$). Also, we find $`\alpha _{\mathrm{l}\mathrm{R}}(0)=0.62\pm 0.03`$. Using ZEUS data , we find a pomeron intercept $`\alpha _{\mathrm{l}\mathrm{P}}=1.127\pm 0.040`$, lower than the H1 value <sup>2</sup><sup>2</sup>2The selection of the diffractive sample in ZEUS is different from H1 and uses the so called $`M_X`$ method ..
The $`Q^2`$ evolution of the pomeron structure function may be understood in terms of parton dynamics and therefore perturbative QCD where parton densities are evolved according to DGLAP equations. We assign parton distribution functions to the pomeron and to the reggeon. A simple prescription is adopted in which the parton distributions of both the pomeron and the reggeon are parameterised in terms of non-perturbative input distributions at some low scale $`Q_0^2=3`$ GeV<sup>2</sup> .
The resulting parton densities of the pomeron are presented in figure 2 as a function of $`z`$, the fractional momentum of the pomeron carried by the struck parton. We find one possible fit quoted here as fit 1. Fit 1 shows a large gluonic content. The quark contribution is much smaller compared to the gluon one. We also note that we find an other much less favoured fit quoted here as fit 2 with a peaked gluon at high $`\beta `$ \[ourapf2d\].
We have redone this QCD analysis with ZEUS 1994 diffractive cross-section measurements applying the same cuts. The gluon density is found to be lower by about a factor 2, but the error bar on the ZEUS gluon density is large (about 50%) as only 30 data points are included in the fit. The quark component is found to be similar. This difference in the gluon density can however lead to differences in the charm structure function as it is very much sensible to the gluon structure function .
### 2.2 Diffraction at Tevatron
Diffractive events selection at Tevatron is basically the same as at HERA. A rapidity gap is asked, either between the proton (antiproton) direction and the jets inside the detector, or a rapidity central gap is asked between two jets inside the main part of the detector.
Hard single diffraction
D0 and CDF collaborations have studied the fraction of central and forward jet events with rapidity gap at two different center-of-mass energies of 630 and 1800 GeV. The D0 results <sup>3</sup><sup>3</sup>3 CDF results are compatible with the D0 ones. are given in Table II and are compared to MC simulations with different pomeron structure functions (hard gluon, $`f(\beta )=\beta (1\beta )`$, soft gluon, $`f(\beta )=(1\beta )^5`$, and quarks). It can be first noticed that the gap fractions at 630 GeV are larger than the ones at 1800 GeV, and are much smaller than at HERA (about 1% to be compared with about 10% at HERA), which can be explained because there is no factorisation between both experiments . The flat and hard gluon predictions are too high compared to the data, whereas the quark scenario could work (it has however been shown that this would lead to predict an excesive rate of diffractive W production). A combination of soft and hard gluon could be possible to describe this measurement, whereas the HERA data have the tendency to favour the hard gluon scenario. The HERA and Tevatron kinematical domains in $`\beta `$ and $`W`$ are however different, and a more precise study including the QCD evolution of parton distributions measured at HERA in the Tevatron domain would be of great interest.
| Sample | Data | Hard Gluon | Flat Gluon | Quark |
| --- | --- | --- | --- | --- |
| $`1800`$ GeV $`|\eta |>1.6`$ | $`(0.65\pm 0.04)\%`$ | $`(2.2\pm 0.3)\%`$ | $`(2.2\pm 0.3)\%`$ | $`(0.79\pm 0.12)\%`$ |
| $`1800`$ GeV $`|\eta |<1.0`$ | $`(0.22\pm 0.05\%`$ | $`(2.5\pm 0.4)\%`$ | $`(3.5\pm 0.5)\%`$ | $`(0.49\pm 0.06)\%`$ |
| $`630`$ GeV $`|\eta |>1.6`$ | $`(1.19\pm 0.08)\%`$ | $`(3.9\pm 0.9)\%`$ | $`(3.1\pm 0.8)\%`$ | $`(2.2\pm 0.5)\%`$ |
| $`630`$ GeV $`|\eta |<1.0`$ | $`(0.90\pm 0.06)\%`$ | $`(5.2\pm 0.7)\%`$ | $`(6.3\pm 0.9)\%`$ | $`(1.6\pm 0.2)\%`$ |
Table II- Measured and predicted gap fractions and their ratios (D0).
Hard color singlet exchange
At Tevatron (like at HERA), it is also possible to study events with a central rapidity gap between jets. This class of events could come from pomeron exchange at large momentum transfer $`t>>100`$ GeV<sup>2</sup>, compared to the previous measurement which was at low $`t`$ ($`t`$0). The fraction of dijet events with a central rapidity gap is again about 1% , to be compared with 10% at HERA. An increase of gap fraction is observed with jet transverse energy and rapidity.
Diffractive dijets with a leading antiproton
The CDF collaboration installed roman-pot detectors allowing to tag $`\overline{p}`$ in the final state. They select a diffractive dijet subsample of it requiring two jets of an energy greater than 7 GeV .
Figure 3 shows the ratio of single diffractive events over the non diffractive ones, where the two data samples are normalised to the same luminosity, in six $`x_{\mathrm{l}\mathrm{P}}`$ bins (at Tevatron, $`x_{\mathrm{l}\mathrm{P}}`$ is called $`\xi `$). This ratio does not show any $`x_{\mathrm{l}\mathrm{P}}`$ dependence and the $`x_{\mathrm{l}\mathrm{P}}`$ dependence can be factorised out from the $`x`$ one, which is similar to what has been found at HERA except at high $`x_{\mathrm{l}\mathrm{P}}`$ and low $`\beta `$ where secondary reggeons are necessary to fit H1 data.
Knowing the ratio of diffarctive to non-diffractive events, it is quite easy to extract the pomeron structure function at Tevatron, by multiplying this ratio by the proton structure function. For this sake, the CDF collaboration chose to use the GRV parametrisation . The result together with extrapolations of H1 fits is given in Figure 4. The results disagree both in normalisation and shapes. This disagreement represents a breaking of factorisation as expected . It is quite challenging to analyze this breaking of factorisation which might be $`\beta `$ dependent using run I and run II data. In run II, the D0 collaboration will have roman pots in both sides of the detector allowing to tag both proton and antiproton in the final state, it will be quite interesting to be able to test factorisation directly.
## 3 Conclusion
We have presented and discussed two related topics at HERA and Tevatron, namely forward jets and diffraction.
Using a new method to disantangle the effects of the kinematic cuts from the genuine dynamical values of the forward jet cross-sections at HERA, we find that the effective pomeron intercept is $`\alpha _P=1.43\pm 0.025`$ (stat.) $`\pm 0.025`$ (syst.). It is much higher than the soft pomeron intercept, and, among those determined in hard processes, it is intermediate between $`\gamma ^{}\gamma ^{}`$ interactions at LEP and dijet productions with large rapidity intervals at Tevatron.
Diffraction at HERA and Tevatron give quite different results. The D0 and CDF collaborations measure about 1% diffractive events whereas at HERA, it is close to 10%. QCD analysis of HERA data lead to a pomeron made of gluons, and a hard structure function is favored whereas the D0 collaboration showed that a combination of hard and soft structure functions is needed. The CDF collaboration was able to make for the first time a direct measurement of the antiproton diffractive structure function, allowing a direct comparison with HERA, by using their roman pot data. The data show large discrepancies both in shape and in normalisation. A QCD analysis of run I Tevatron data and run II data where D0 will have roman pot detectors on each side are quite challenging.
## Acknowledgments
I thank J.Bartels, G.Contreras, H.Jung, R.Peschanski and L.Schoeffel for collaboration.
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# Nonlocal Boundary Dynamics of Traveling Spots in a Reaction-Diffusion System
## Acknowledgement.
This work has been supported by the German–Israeli Science Foundation.
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# Invariant mass dependence of Λ polarization in 𝑝𝑝→𝑝Λ𝐾⁺𝜋⁺𝜋⁻𝜋⁺𝜋⁻
We show that there is a correlation between the invariant mass of the produced $`\mathrm{\Lambda }K^+`$, $`\mathrm{\Lambda }K^+\pi ^+\pi ^{}`$ or $`\mathrm{\Lambda }K^+\pi ^+\pi ^{}\pi ^+\pi ^{}`$ system in the exclusive reaction $`ppp\mathrm{\Lambda }K^+\pi ^+\pi ^{}\pi ^+\pi ^{}`$ and the longitudinal or transverse momentum of $`\mathrm{\Lambda }`$. Together with the longitudinal and transverse momentum dependence of $`\mathrm{\Lambda }`$ polarization observed in inclusive reactions, such a correlation implies a dependence of $`\mathrm{\Lambda }`$ polarization on these invariant masses. The qualitative features of this dependence are consistent with the recent observation by E766 collaboration at BNL. A quantitative estimation has been made using an event generator for $`pp`$ collisions. A detailed comparison with the data is made.
Since the discovery of hyperon polarization ($`P_H`$) in inclusive production processes at high energies, there has been constant interest in studying the origin of this effect, both experimentally and theoretically \[References-References\]. It is now an established experimental fact that, in high energy hadron-hadron or hadron-nucleus collisions, the produced hyperons are polarized transversely to the production plane, although neither the projectiles nor the targets are polarized before the collisions. A large number of experiments on the inclusive production processes show that the polarization is significantly different from zero in the fragmentation regions for moderately large transverse momenta. The magnitude of the polarization increase with increasing $`x_F`$ and $`p_{}`$, where $`x_F2p_{}/\sqrt{s}`$ is the Feynman-$`x`$, $`s`$ is the total center of mass energy squared, $`p_{}`$ and $`p_{}`$ are respectively the longitudinal and transverse component of the momentum of the produced hyperon.
Recently, further progresses have been achieved in experimental studies. An interesting program has been established by E766 Collaboration at Brookheaven National Laboratory (BNL) to study $`\mathrm{\Lambda }`$ polarization $`P_\mathrm{\Lambda }`$ in different specific reaction channels. The first experiment has already been carried out on $`P_\mathrm{\Lambda }`$ in $`ppp\mathrm{\Lambda }K^+\pi ^+\pi ^{}\pi ^+\pi ^{}`$. They found in particular the following interesting phenomenon: $`P_\mathrm{\Lambda }`$ in this channel depends not only on $`x_F`$ and $`p_{}`$ of the produced $`\mathrm{\Lambda }`$ (as observed earlier in inclusive experiments), but also on the invariant masses $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$, and $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$, of the particle systems $`\mathrm{\Lambda }K^+`$, $`\mathrm{\Lambda }K^+\pi ^+\pi ^{}`$ and $`\mathrm{\Lambda }K^+\pi ^+\pi ^{}\pi ^+\pi ^{}`$ respectively. The magnitude of $`P_\mathrm{\Lambda }`$ increases with the increasing of these invariant masses $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$, and $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$.
We recall that, compared with those on inclusive processes, the experiments on exclusive processes have the advantage to study the dependence on kinematic variables other than $`x_F`$ or $`p_{}`$. It offers the opportunity to investigate different correlations which usually provide us with information that can not be obtained in inclusive experiments. Such information gives often deeper insight into the physics behind the data. Hence, it is conceivable that the above mentioned new experimental finding provide further important tests of the different models and give us some new clue in the searching of the origin of $`\mathrm{\Lambda }`$ polarization in high energy hadron-hadron collisions. However, kinematic analysis readily shows that, due to energy-momentum conservation, there has to be a correlation between the above-mentioned invariant mass $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$, or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$ and $`x_F`$ or $`p_{}`$ of $`\mathrm{\Lambda }`$. On the average, both $`x_F`$ and $`p_{}`$ increase with the increasing of these invariant masses $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$, or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$. Together with the $`x_F`$ and $`p_{}`$ dependences of $`P_\mathrm{\Lambda }`$ observed earlier in the inclusive experiments, such correlations lead already to some increase of $`|P_\mathrm{\Lambda }|`$ with the increasing of $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$, or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$. Hence, we are led naturally to the following questions: How large are these correlations? Can they already account for the observed invariant mass dependence of $`P_\mathrm{\Lambda }`$? Do the observed increase of $`P_\mathrm{\Lambda }`$ with increasing $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$, or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$ give us some further tests of the different models or they are just different manifestations of the $`x_F`$ and $`p_{}`$ dependences observed earlier in inclusive experiments? These are questions which we would like to investigate in this note.
We recall that the invariant mass of $`\mathrm{\Lambda }K^+`$-system is defined as the total energy in their center of mass (c.m.) system, i.e.,
$$M_{\mathrm{\Lambda }K}=\sqrt{m_\mathrm{\Lambda }^2+p^2}+\sqrt{m_K^2+p^2},$$
(1)
where $`m_K`$ and $`m_\mathrm{\Lambda }`$ are their masses, $`p^{}`$ is their momentum in the c.m. frame of this particle system. It is obvious that, $`M_{\mathrm{\Lambda }K}`$ increases with increasing $`p^{}`$. Large value of $`p^{}`$ implies that the difference between the momentum of $`\mathrm{\Lambda }`$ and that of $`K^+`$ is large, also in the c.m. frame of the colliding $`pp`$-system. Hence, we expect large difference between $`x_F`$ (or $`p_{}`$) of $`\mathrm{\Lambda }`$ and that of $`K^+`$ for large $`M_{\mathrm{\Lambda }K}`$. On the other hand, if $`x_F`$ is very large (say, larger than $`0.5`$), the longitudinal momentum of $`\mathrm{\Lambda }`$ is very large. In this case, $`\mathrm{\Lambda }`$ carries already a very large part of the momentum of the whole system. According to energy-momentum conservation, the momentum of $`K^+`$ cannot be the same as that for $`\mathrm{\Lambda }`$ since the sum of them cannot exceed $`\sqrt{s}/2`$. The longitudinal momentum for $`K^+`$ has to be much smaller. This implies a large difference between them thus a large $`M_{\mathrm{\Lambda }K}`$. Hence, we expect that $`M_{\mathrm{\Lambda }K}`$ increases with increasing $`x_F`$ for large $`x_F`$. Similarly, if $`p_{}`$ of $`\mathrm{\Lambda }`$ is large, there should be a large probability that $`p_{}`$ of $`K^+`$ is large and in the opposite direction to guarantee the transverse momentum conservation. This leads to a large difference in transverse momenta for $`\mathrm{\Lambda }`$ and $`K^+`$ thus also a large $`M_{\mathrm{\Lambda }K}`$. Hence, we expect $`M_{\mathrm{\Lambda }K}`$ increases also with increasing $`p_{}`$. We are naturally led to the following questions: How fast does $`M_{\mathrm{\Lambda }K}`$ increase with increasing $`x_F`$ or $`p_{}`$? Are there similar correlations between $`M_{\mathrm{\Lambda }K\pi \pi }`$ (or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$) and $`x_F`$ or $`p_{}`$ of $`\mathrm{\Lambda }`$? Apparently, the answers to these questions are determined by the momentum distributions of the particles produced in the collision processes and these distributions can be obtained in unpolarized reactions.
In order to study these questions quantitatively, we used a Monte-Carlo event generator PYTHIA . We recall that PYTHIA is an event generator for unpolarized high energy hadronic reactions based on Lund fragmentation model . It describes most (if not all) of the different features of the data for particle production in unpolarized reactions. We therefore expect that we should be able to obtain a reasonable description of the correlations between $`x_F`$ or $`p_{}`$ and the invariant masses mentioned above. We generated 100,000 $`ppp\mathrm{\Lambda }K^+\pi ^+\pi ^{}\pi ^+\pi ^{}`$ events using PYTHIA. From these events, we calculated the average values of $`x_F`$ and those of $`p_{}`$ for different $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$ or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$ bins. The obtained results are shown in Figs.1 and 2. From Figs.1(a) and 2(a), we see clearly that both $`x_F`$ and $`p_{}`$ increase with increasing $`M_{\mathrm{\Lambda }K}`$. This is consistent with our qualitative expectations. From Figs.1(b),1(c),2(b) and 2(c), we see also similar correlations between $`M_{\mathrm{\Lambda }K\pi \pi }`$ or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$ and $`x_F`$ or $`p_{}`$. The average values $`x_F`$ and $`p_{}`$ increase also with increasing $`M_{\mathrm{\Lambda }K\pi \pi }`$ or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$. But, we also see that the magnitudes of these correlations are smaller than that between $`x_F`$ or $`p_{}`$ and $`M_{\mathrm{\Lambda }K}`$.
Since the magnitude of $`\mathrm{\Lambda }`$ polarization $`P_\mathrm{\Lambda }`$ increases with increasing $`x_F`$ or increasing $`p_{}`$, we expect that $`|P_\mathrm{\Lambda }|`$ has to increase also with $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$, or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$ because of the above mentioned correlations. This qualitative feature is in agreement with the data. To see whether this effect explains also quantitatively the observed dependences of $`P_\mathrm{\Lambda }`$ on $`M_{\mathrm{\Lambda }K\pi \pi }`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$ or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$, we did the following calculations. We assume that $`P_\mathrm{\Lambda }`$ is determined completely by the $`x_F`$ and $`p_{}`$ of $`\mathrm{\Lambda }`$ and use the $`x_F`$ and $`p_{}`$ dependences of $`P_\mathrm{\Lambda }`$ obtained in experiments on inclusive reactions as input to calculate $`P_\mathrm{\Lambda }`$ as a function of $`M_{\mathrm{\Lambda }K}`$, that of $`M_{\mathrm{\Lambda }K\pi \pi }`$ and that of $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$ respectively. The results for $`P_\mathrm{\Lambda }`$ as a function of $`x_F`$ and $`p_{}`$ obtained earlier in inclusive experiments can be parametrized as
$$P_\mathrm{\Lambda }(x_F,p_{})=1.5(c_1x_F+c_2x_F^3)(1e^{c_3p_{}^2})$$
(2)
where $`c_1=0.268\pm 0.003`$, $`c_2=0.338\pm 0.015`$, and $`c_3=4.5\pm 0.6`$, are constants determined by the data. Using this parameterization for $`P_\mathrm{\Lambda }(x_F,p_{})`$, we obtain $`P_\mathrm{\Lambda }`$ as functions of $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$ and $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$ shown in Fig.3.
From Fig.3, we see that, $`|P_\mathrm{\Lambda }|`$ increases indeed with increasing $`M_{\mathrm{\Lambda }K}`$, $`M_{\mathrm{\Lambda }K\pi \pi }`$ or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$. This qualitative feature agrees with the data. We see also that, the $`M_{\mathrm{\Lambda }K}`$ dependence of $`P_\mathrm{\Lambda }`$ agrees with the data even quantitatively. This shows that the above mentioned kinematic effect together with the early observed $`x_F`$ and $`p_{}`$ dependences of $`P_\mathrm{\Lambda }`$ can already explain this $`M_{\mathrm{\Lambda }K}`$ dependence. The qualitative features of $`P_\mathrm{\Lambda }`$ as a function of $`M_{\mathrm{\Lambda }K\pi \pi }`$ and that of $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$ agree also with the data. But, quantitatively, the increase of $`|P_\mathrm{\Lambda }|`$ is too slow and is not enough to account for the observed $`M_{\mathrm{\Lambda }K\pi \pi }`$ or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$ dependence. In particular, $`|P_\mathrm{\Lambda }|`$ seems too large near the threshold of $`M_{\mathrm{\Lambda }K\pi \pi }`$ or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$. This means that other dynamic effects have to be introduced to account for this effect. This also implies that these dependences should be other independent tests for the different theoretical models \[References-References\].
In summary, we made a kinematic analysis of the dependence of $`P_\mathrm{\Lambda }`$ in $`ppp\mathrm{\Lambda }K^+\pi ^+\pi ^{}\pi ^+\pi ^{}`$ on the invariant masses of the produced $`\mathrm{\Lambda }K^+`$, $`\mathrm{\Lambda }K^+\pi ^+\pi ^{}`$ and $`\mathrm{\Lambda }K^+\pi ^+\pi ^{}\pi ^+\pi ^{}`$ systems. We showed that there is a correlation between these invariant masses and $`x_F`$ or $`p_{}`$ of $`\mathrm{\Lambda }`$. This correlation gives already a reasonable explanation of the increase of $`|P_\mathrm{\Lambda }|`$ with increasing $`M_{\mathrm{\Lambda }K}`$ but not that with increasing $`M_{\mathrm{\Lambda }K\pi \pi }`$ or $`M_{\mathrm{\Lambda }K\pi \pi \pi \pi }`$.
We thank Xie Qu-bing, Wang Qun, Si Zong-guo for helpful discussions. This work is supported in part by the National Natural Science Foundation (NSFC) and the Education Ministry of China.
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# Interaction in Quantum Communication Complexity
## 1 Introduction
A recurring theme in quantum information processing has been the idea of exploiting the exponential resources afforded by quantum states to encode information in very non-obvious ways. Perhaps the most representative result of this kind is due to Ambainis, Schulman, Ta-Shma, Vazirani and Wigderson , which shows that it is possible to deal a random set of $`\sqrt{N}`$ cards each from a set of $`N`$ by the exchange of $`O(\mathrm{log}N)`$ quantum bits between two players. Raz gives a communication problem where the information storage capacity of quantum states is exploited more explicitly. Both are examples of problems for which exponentially fewer quantum bits are required to accomplish a communication task, as compared to classical bits.
The protocols presented by also share the feature that they require minimal interaction between the communicating players. For example, in the first protocol, one player prepares a set of qubits in a certain state and sends half of them across as the message, after which both players measure their qubits to obtain the result. On the other hand, efficient quantum protocols for problems such as checking set disjointness (DISJ) seem to require much more interaction: Buhrman, Cleve and Wigderson give an $`O(\sqrt{N}\mathrm{log}N)`$ qubit protocol for DISJ that takes $`O(\sqrt{N})`$ message exchanges. This represents quadratic savings in communication cost, but also an unbounded increase in the number of messages exchanged (from one message to $`\sqrt{N}`$), as compared to classical protocols. Can we exploit the features of quantum communication and always reduce interaction while maintaining the same communication cost? In other words, do all efficient quantum protocols have the simple structure shared by those of ?
In this paper, we study the effect of interaction on the quantum communication complexity of problems. We show that for any constant $`k`$, allowing even one more message may lead to an exponential decrease in the communication complexity of a problem, thus answering the above question in the negative. More formally,
###### Theorem 1.1
For any constant $`k`$, there is a problem such that any quantum protocol with only $`k`$ messages and constant probability of error requires $`\mathrm{\Omega }(N^{1/(k+1)})`$ communication qubits, whereas it can be solved with $`k+1`$ messages by a deterministic protocol with $`O(\mathrm{log}N)`$ bits.
Klauck states a relationship between the bounded message complexity of Pointer Jumping and DISJ. Together with our result, this implies an $`\mathrm{\Omega }(N^{1/k(k+1)})`$ lower bound for $`k`$ message protocols for DISJ, for any constant $`k`$.
The role of interaction in classical communication is well-studied, especially in the context of the pointer jumping function . In fact, the problem we study in this paper is the subproblem of Pointer Jumping singled out in . Our analysis has the same gross structure as that in (also explained in ), but relies on entirely new ideas from quantum information theory.
In the context of quantum communication, it was observed by Buhrman and de Wolf (based on a lower bound of Nayak ) that any one message quantum protocol for DISJ has linear communication complexity. Thus, allowing more interaction leads to a quadratic improvement in communication cost. The lower bound of immediately implies a much stronger separation: it shows that the two message complexity of a problem may be exponentially smaller than its one message complexity (see also ). Our result subsumes all these.
Our interest in the role of interaction in quantum communication also springs from the need to better understand the ways in which we can access and manipulate information encoded in quantum states. We develop information-theoretic techniques that expose some of the limitations of quantum communication. More specifically, we present a new primitive in quantum encoding, as suggested by the following theorem.
###### Theorem 1.2 (Average encoding theorem)
Let $`x\sigma _x`$ be a quantum encoding mapping $`m`$ bit strings $`x\{0,1\}^m`$ into mixed states $`\sigma _x`$. Let $`X`$ be distributed uniformly over $`\{0,1\}^m`$, let $`Q`$ be the encoding of $`X`$ according to this map, and let $`\sigma =\frac{1}{2^m}_x\sigma _x`$. Then,
$`{\displaystyle \frac{1}{2^m}}{\displaystyle \underset{x}{}}\sigma \sigma _x_\mathrm{t}`$ $``$ $`2\sqrt{I(Q:X)}.`$
In other words, if an encoding $`Q`$ is only weakly correlated to a random variable $`X`$, then the “average encoding” $`\sigma `$ (corresponding to a random string) is on average a good approximation of any encoded state. Thus, in certain situations, we may dispense with the encoding altogether, and use the single state $`\sigma `$ instead.
We also use another primitive derived from the work of Lo and Chau and Mayers which combines results of Jozsa , and Fuchs and van de Graaf . Consider two bi-partite pure states such that one party sharing the states cannot locally distinguish well between the two states with good probability. Then the other party can locally transform any of the states close to the other.
###### Theorem 1.3 (Local transition theorem)
(based on ) Let $`\rho _1,\rho _2`$ be two mixed states with support in a Hilbert space $``$, $`𝒦`$ any Hilbert space of dimension at least $`dim()`$, and $`|\varphi _i`$ any purifications of $`\rho _i`$ in $`𝒦`$. Then, there is a local unitary transformation $`U`$ on $`𝒦`$ that maps $`|\varphi _2`$ to $`|\varphi _2^{}=IU|\varphi _2`$ such that
$$|\varphi _1\varphi _1||\varphi _2^{}\varphi _2^{}|_\mathrm{t}\mathrm{\hspace{0.33em}\hspace{0.33em}2}\rho _1\rho _2_\mathrm{t}^{\frac{1}{2}}.$$
This may be of significance in cryptographic applications as well.
## 2 Preliminaries
### 2.1 The communication complexity model
In the quantum communication complexity model , Alice and Bob hold qubits. When the game starts Alice holds a superposition $`|x`$ and Bob holds $`|y`$ (representing the input to the two players), and so the initial joint state is simply $`|x|y`$. The two parties then play in turns. Suppose it is Alice’s turn to play. Alice can do an arbitrary unitary transformation on her qubits and then send one or more qubits to Bob. Sending qubits does not change the overall superposition, but rather changes the ownership of the qubits, allowing Bob to apply his next unitary transformation on the newly received qubits. At the end of the protocol, one player makes a measurement and declares that as the result of the protocol.
In general, each player may also (partially) measure her qubits during her turn. However, we assume (by invoking the principle of safe storage ) that all such measurements are postponed to the end. We also assume that the two players do not modify the qubits holding the input superposition during the protocol. Neither of these affects the aspect of communication we focus on in this paper.
The complexity of a quantum (or classical) protocol is the number of qubits (respectively, bits) exchanged between the two players. We say a protocol computes a function $`f:𝒳\times 𝒴\{0,1\}`$ with $`ϵ0`$ error, if for any input $`x𝒳,y𝒴`$ the probability that the two players compute $`f(x,y)`$ is at least $`1ϵ`$. $`Q_ϵ(f)`$ denotes the complexity of the best quantum protocol that computes $`f`$ with at most $`ϵ`$ error.
For a player $`P\{\mathrm{Alice},\mathrm{Bob}\}`$, $`Q_ϵ^{c,P}(f)`$ denotes the complexity of the best quantum protocol that computes $`f`$ with at most $`ϵ`$ error with only $`c`$ messages, where the first message is sent by $`P`$. If the name of the player is omitted from the superscript, either player is allowed to start the protocol.
We say a protocol $`𝒫`$ computes $`f`$ with $`ϵ`$ error with respect to a distribution $`\mu `$ on $`𝒳\times 𝒴`$, if
$$\mathrm{Prob}_{(x,y)\mu ,𝒫}(𝒫(x,y)=f(x,y))\mathrm{\hspace{0.33em}\hspace{0.33em}1}ϵ.$$
$`Q_{\mu ,ϵ}^{c,P}(f)`$ is the complexity of computing $`f`$ with at most $`ϵ`$ error with respect to $`\mu `$, with only $`c`$ messages where the first message is sent by player $`P`$. The following is immediate.
###### Fact 2.1
For any distribution $`\mu `$, number of messages $`c`$ and player $`P`$$`Q_{\mu ,ϵ}^{c,P}(f)Q_ϵ^{c,P}(f)`$.
### 2.2 Classical entropy and mutual information
The Shannon entropy $`S(X)`$ of a classical random variable $`X`$ quantifies the amount of randomness in it. If $`X`$ takes values $`x`$ in some finite set with probability $`p_x`$, its Shannon entropy is defined as $`S(X)=_xp_x\mathrm{log}p_x`$. The mutual information $`I(X:Y)`$ of a pair of random variables $`X,Y`$ is defined by $`I(X:Y)=S(X)+S(Y)S(XY)`$. It is a measure of how correlated the two random variables are.
The following are some basic facts about the mutual information function that we use in the paper. For any random variables $`X,Y,Z`$,
$`I(X:YZ)`$ $`=`$ $`I(X:Y)+I(XY:Z)I(Y:Z)`$ (1)
$`I(X:YZ)`$ $``$ $`I(X:Y).`$ (2)
Fano’s inequality states that if $`Y`$ can predict another random variable $`X`$ with an advantage, then $`X`$ and $`Y`$ have large mutual information. We use it only in the following simple form.
###### Fact 2.2 (Fano’s inequality)
Let $`X`$ be a uniformly distributed boolean random variable, and let $`Y`$ be a boolean random variable such that $`\mathrm{Prob}(X=Y)\frac{1}{2}+\delta `$, where $`\delta 0`$. Then $`I(X:Y)1H(\frac{1}{2}+\delta )`$.
For other equivalent definitions and properties of these concepts, we refer the reader to a standard text (such as ) on information theory. Finally, we give a simple bound on the deviation of the binary entropy function $`H(p)`$ from $`1`$ as $`p`$ deviates from $`1/2`$.
###### Fact 2.3
For $`\delta [\frac{1}{2},\frac{1}{2}]`$, we have $`H(\frac{1}{2}+\delta )\mathrm{\hspace{0.33em}\hspace{0.33em}1}\delta ^2`$.
Proof: From the definition of the binary entropy function, we have
$`H({\displaystyle \frac{1}{2}}+\delta )`$ $`=`$ $`1{\displaystyle \frac{1}{2}}[(1+2\delta )\mathrm{log}(1+2\delta )+(12\delta )\mathrm{log}(12\delta )].`$
Using the expansion $`\mathrm{ln}(1+x)=x\frac{x^2}{2}+\frac{x^3}{3}\frac{x^4}{4}+\mathrm{}`$ for $`\left|x\right|<1`$, and simplifying, we get
$`H({\displaystyle \frac{1}{2}}+\delta )`$ $`=`$ $`1(\mathrm{log}\mathrm{e})\left[\left(2{\displaystyle \frac{2^2}{22}}\right)\delta ^2+\left({\displaystyle \frac{2^3}{3}}{\displaystyle \frac{2^4}{24}}\right)\delta ^4+\left({\displaystyle \frac{2^5}{5}}{\displaystyle \frac{2^6}{26}}\right)\delta ^6+\mathrm{}\right]`$
$``$ $`1\delta ^2,`$
which is the claimed bound.
### 2.3 The density matrix and the trace norm
The quantum mechanical analogue of a random variable is a probability distribution over superpositions, also called a mixed state. Consider the mixed state $`X=\{p_i,|\varphi _i\}`$, where the superposition $`|\varphi _i`$ is drawn with probability $`p_i`$. The density matrix of the mixed state $`X`$ is $`\rho _X=_ip_i|\varphi _i\varphi _i|`$. The following properties of density matrices are immediate from the definition: every density matrix $`\rho `$ is Hermitian, i.e., $`\rho =\rho ^{}`$, has unit trace, i.e., $`\mathrm{Tr}(\rho )=_i\rho (i,i)=1`$, and is positive semi-definite, i.e., $`\psi \left|\rho \right|\psi 0`$ for all $`|\psi `$. Thus, every density matrix is unitarily diagonalizable and has non-negative real eigenvalues that sum up to $`1`$.
Given a quantum system in a mixed state with density matrix $`\rho `$ and a (general) measurement $`𝒪`$ on it, let $`\rho ^𝒪`$ denote the classical distribution on the possible results that we get by measuring $`\rho `$ according to $`𝒪`$. Suppose that it is some classical distribution $`p_1,\mathrm{},p_k`$ where we get result $`i`$ with probability $`p_i`$. Given two different mixed states, we can ask how well one can distinguish between the two mixtures, or equivalently, how different the distributions resulting from a measurement may be. To quantify this, we consider the $`\mathrm{}_1`$ metric: if $`p=(p_1,\mathrm{},p_k)`$ and $`q=(q_1,\mathrm{},q_k)`$ are two probability distributions over $`\{1,\mathrm{},k\}`$, then the $`\mathrm{}_1`$ distance between them is $`pq_1=_i\left|p_iq_i\right|`$. A fundamental theorem about distinguishing density matrices (see ) tells us:
###### Theorem 2.4
Let $`\rho _1,\rho _2`$ be two density matrices on the same space $``$. Then for any (general) measurement $`𝒪`$
$$\rho _1^𝒪\rho _2^𝒪_1\mathrm{Tr}\sqrt{A^{}A},$$
where $`A=\rho _1\rho _2`$. Furthermore, the bound is tight, and the orthogonal measurement $`𝒪`$ that projects a state on the eigenvectors of $`\rho _1\rho _2`$ achieves this bound.
Theorem 2.4 shows that the density matrix captures all the accessible information that a quantum state contains. If two different mixtures have the same density matrix (which is indeed possible) then even though they are two distinct distributions, they are physically, and thus from a computational point of view, indistinguishable. As the behavior of a mixed state is completely characterized by its density matrix we often identify a mixed state with its density matrix.
The quantity $`\mathrm{Tr}\sqrt{A^{}A}`$ is of independent interest. (Note that this is compact notation for the sum of the (magnitudes of the) singular values of $`A`$.) If we define $`A_\mathrm{t}=\mathrm{Tr}\sqrt{A^{}A}`$ then $`_\mathrm{t}`$ defines a norm (the trace norm), and has some additional properties such as $`AB_\mathrm{t}=A_\mathrm{t}B_\mathrm{t}`$, $`A_\mathrm{t}=1`$ for any density matrix $`A`$ and $`AB_\mathrm{t},BA_\mathrm{t}A_\mathrm{t}B_\mathrm{t}`$. (See for more details.) We single out the following fact for later use.
###### Fact 2.5
If $`|\varphi _1,|\varphi _2`$ are two pure states, and $`\rho _i`$ is the density matrix of $`|\varphi _i`$, then
$$\rho _1\rho _2_\mathrm{t}=\mathrm{\hspace{0.33em}\hspace{0.33em}2}\sqrt{1|\varphi _1|\varphi _2|^2}.$$
### 2.4 The fidelity measure
A useful alternative to the trace metric as a measure of closeness of density matrices is fidelity, which is defined in terms of the pure states that can give rise to those density matrices. A purification of a mixed state $`\rho `$ with support in a Hilbert space $``$ is any pure state $`|\varphi `$ in an extended Hilbert space $`𝒦`$ such that $`\mathrm{Tr}_𝒦|\varphi \varphi |=\rho `$. Given two density matrices $`\rho _1,\rho _2`$ on the same Hilbert space $``$, their fidelity is defined as
$$F(\rho _1,\rho _2)=sup\left|\varphi _1\varphi _2\right|^2,$$
where the supremum is taken over all purifications $`|\varphi _i`$ of $`\rho _i`$ in the same Hilbert space . We state a few properties of this measure: $`0F(\rho _1,\rho _2)1`$, $`F(\rho _1,\rho _2)=1\rho _1=\rho _2`$ and if $`\rho _1=|\varphi _1\varphi _1|`$, then we have $`F(\rho _1,\rho _2)=\varphi _1\left|\rho _2\right|\varphi _1`$. Jozsa proved that the optimum is always achieved when finite dimensional density matrices are considered.
###### Theorem 2.6 (Jozsa)
Let $`\rho _1,\rho _2`$ be any two mixed states with support in a finite dimensional Hilbert space $``$, $`𝒦`$ a Hilbert space of dimension at least $`dim()`$, and $`|\varphi _1`$ any purification of $`\rho _1`$ in $`𝒦`$. Then there exists a purification $`|\varphi _2𝒦`$ of $`\rho _2`$ such that $`\left|\varphi _1\varphi _2\right|^2=F(\rho _1,\rho _2)`$.
Jozsa also gave a simple proof (again for the finite dimensional case) of the following remarkable equivalence first established by Uhlmann .
$$F(\rho _1,\rho _2)=\left[\mathrm{Tr}\left(\sqrt{\rho _1}\rho _2\sqrt{\rho _1}\right)^{\frac{1}{2}}\right]^2=\sqrt{\rho _1}\sqrt{\rho _2}_\mathrm{t}^2.$$
Using this equivalence, Fuchs and van de Graaf show that the fidelity and the trace measures of distance between density matrices are closely related. They prove:
###### Theorem 2.7 (Fuchs, van de Graaf)
For any two mixed states $`\rho _1,\rho _2`$,
$$1\sqrt{F(\rho _1,\rho _2)}\frac{1}{2}\rho _1\rho _2_\mathrm{t}\sqrt{1F(\rho _1,\rho _2)}.$$
### 2.5 Von Neumann entropy and quantum mutual information
As mentioned earlier, the eigenvalues of a density matrix are all real, non-negative and sum up to one. Thus, they induce a probability distribution on the corresponding eigenvectors. Since the eigenvectors are all orthogonal, this is essentially a classical distribution. Every mixed state with the same density matrix is physically equivalent to such a canonical classical distribution. It is thus natural to define the entropy of a mixed state as the Shannon entropy of this distribution. Formally, the von Neumann entropy $`S(\rho )`$ of a density matrix $`\rho `$ is defined as $`S(\rho )=_i\lambda _i\mathrm{log}\lambda _i`$, where $`\{\lambda _i\}`$ is the multi-set of all the eigenvalues of $`\rho `$. More compactly, $`S(\rho )=\mathrm{Tr}\rho \mathrm{log}\rho `$.
Not all properties of classical Shannon entropy carry over to the quantum case. For example it is quite possible that $`S(XY)<S(X)`$ as can be seen by considering the pure state $`\frac{1}{\sqrt{2}}(|0_X|0_Y+|1_X|1_Y)`$. Nonetheless, some of the classical properties do carry over, e.g., $`S(X)0`$, $`S(X)`$ is concave and $`S(XY)S(X)+S(Y)`$. A property of interest to us is the following, which also generalizes a classical assertion.
###### Fact 2.8
Suppose a quantum system $`A`$ is in mixed state $`\{p_i,|i\}`$, where $`\left\{|i\right\}`$ are orthogonal, and $`\sigma _i`$ are density matrices for another system $`B`$, then $`S(_ip_i|ii|\sigma _i)=H(A)+_ip_iS(\sigma _i)`$.
The density matrix corresponding to a mixed state with superpositions drawn from a Hilbert space $``$ is said to have support in $``$. A density matrix with support in a Hilbert space of dimension $`d`$, has $`d`$ eigenvalues, hence the entropy of any such distribution is at most $`\mathrm{log}d`$. A pure-state has zero entropy. Measuring a pure-state may result in a non-trivial mixture and positive entropy. In general, orthogonal measurements increase the entropy. For a comprehensive introduction to this concept and its properties see, for instance, .
We define the “mutual information” $`I(X:Y)`$ of two disjoint systems $`X,Y`$ in analogy with classical mutual information: $`I(X:Y)=S(X)+S(Y)S(XY)`$, where $`XY`$ is density matrix of the system that includes the qubits of both systems. Again, not all properties of classical mutual information carry over to the quantum case. For example, it is not true in general that $`I(X:Y)S(X)`$. Nonetheless, some of the intuition we have about mutual information still applies. Equation (1) still holds, as follows immediately from the definition. Equation (2) also continues to be true, but its proof is much more involved. It is in fact equivalent to the strong sub-additivity property of von Neumann entropy. An important consequence of this property is that local measurements can only decrease the amount of mutual information. A special case of this is the classic Holevo theorem from quantum information theory, which bounds the amount of information we can extract from a quantum encoding of classical bits.
###### Theorem 2.9 (Holevo)
Let $`x\sigma _x`$ be any quantum encoding of bit strings into density matrices. let $`X`$ be a random variable with a distribution given by $`\mathrm{Prob}(X=x)=p_x`$, let $`Q`$ be the quantum encoding of $`X`$ according to this map, and let $`\sigma =_xp_x\sigma _x`$. If $`Y`$ is any random variable obtained by performing a measurement on the encoding, then
$$I(X:Y)I(X:Q)=S(\sigma )\underset{x}{}p_xS(\sigma _x).$$
In analogy with classical conditional entropy, we define $`S(Y|X)=_xp_xS(\sigma _x)`$, where $`X`$ is a classical random variable and $`Y`$ is a quantum encoding of it given by $`x\sigma _x`$. We similarly define conditional von Neumann entropy and mutual information with respect to a classical event. Thus, for example, $`I(X:Y)=S(Y)S(Y|X)`$.
## 3 The average encoding theorem
The average encoding theorem asserts that if a quantum encoding has little correlation with the encoded classical information then the encoded states are essentially indistinguishable. In particular, they are all “close” to the average encoding. This theorem formalizes a very intuitive idea and might seem to be immediate from Holevo’s theorem. However, there is a subtle difference: in Holevo’s theorem one is interested in a single measurement that simultaneously distinguishes all the states, whereas in our case we are interested in the pairwise distinguishability of the encoded states. We first prove:
###### Theorem 3.1
Let $`x\sigma _x`$ be a quantum encoding mapping $`m`$ bit strings $`x\{0,1\}^m`$ into mixed states $`\sigma _x`$. Let $`X`$ be distributed uniformly over $`\{0,1\}^m`$ and let $`Q`$ be the encoding of $`X`$ according to this map. Denote $`\mathrm{\Delta }=\frac{1}{2^{2m}}_{x_1,x_2\{0,1\}^m}\sigma _{x_1}\sigma _{x_2}_\mathrm{t}`$. Then $`I(X:Q)1H(\frac{1+\mathrm{\Delta }}{2})`$.
Proof: We start with the special case of $`m=1`$. By Theorem 2.4, there is a measurement $`𝒪`$ on $`Q`$ that realizes the trace norm distance $`t=\sigma _0\sigma _1_\mathrm{t}`$ between $`\sigma _0`$ and $`\sigma _1`$. Using Bayes’ strategy (see, for example, ), the resulting distributions can be distinguished with probability $`\frac{1}{2}+\frac{t}{4}`$. Let $`Y`$ denote the classical random variable holding the result of this entire procedure. We have $`\mathrm{Prob}(Y=X)=\frac{1}{2}+\frac{t}{4}`$. Thus, by Fano’s Inequality,
$$I(X:Y)\mathrm{\hspace{0.33em}\hspace{0.33em}1}H(\frac{1}{2}+\frac{t}{4})$$
We complete the proof for $`m=1`$ by noticing that measurements can only reduce the entropy, hence $`I(X:Q)I(X:Y)`$, and that $`\mathrm{\Delta }=\frac{t}{2}`$.
To prove the theorem for general $`m`$ we reduce it to the $`m=1`$ case. We do this by partitioning the set of strings into pairs with “easily” distinguishable encoding.
###### Lemma 3.2
There is a set of $`2^m/2`$ disjoint pairs $`(x_{2i1},x_{2i})`$ which together cover $`\{0,1\}^m`$ such that
$$\frac{2}{2^m}\underset{i}{}\sigma _{x_{2i1}}\sigma _{x_{2i}}_\mathrm{t}\mathrm{\Delta }.$$
Proof: The expectation of the LHS over a random pairing is $`\frac{2^m}{2^m1}\mathrm{\Delta }`$; hence there is a pairing that achieves this $`\mathrm{\Delta }`$.
We now fix this pairing. Let $`Z_i`$ denote the set of elements in the $`i`$’th pair, i.e., $`Z_i=\{x_{2i1},x_{2i}\}`$ and $`\mathrm{\Delta }_i=\sigma _{x_{2i1}}\sigma _{x_{2i}}_\mathrm{t}`$. We know that $`\frac{2}{2^m}\mathrm{\Delta }_i\mathrm{\Delta }`$. Let us also denote $`f(\delta )=1H(\frac{1+\delta }{2})`$. From the base case $`m=1`$, we know that for any $`i=1,\mathrm{},2^m/2`$, $`I(X:Q|XZ_i)f(\mathrm{\Delta }_i)`$. Thus we get:
$$S(Q|XZ_i)\frac{1}{2}[S(\sigma _{x_{2i}})S(\sigma _{x_{2i+1}})]f(\mathrm{\Delta }_i).$$
Averaging all the $`2^m/2`$ equations yields:
$`{\displaystyle \frac{2}{2^m}}{\displaystyle \underset{i}{}}S(Q|XZ_i){\displaystyle \frac{1}{2^m}}{\displaystyle \underset{x}{}}S(\sigma _x)`$ $``$ $`{\displaystyle \frac{2}{2^m}}{\displaystyle \underset{i}{}}f(\mathrm{\Delta }_i)`$
By the concavity of the entropy function, $`S(Q)\frac{2}{2^m}_iS(Q|XZ_i)`$, and by definition $`\frac{1}{2^m}_xS(\sigma _x)=S(Q|X)`$. Therefore,
$$I(X:Q)=S(Q)S(Q|X)\frac{2}{2^m}\underset{i}{}f(\mathrm{\Delta }_i).$$
Since $`f`$ is convex, $`\frac{2}{2^m}_if(\mathrm{\Delta }_i)f(\frac{2}{2^m}_i\mathrm{\Delta }_i)`$. Also, $`f(\delta )`$ is monotone increasing for $`0\delta \frac{1}{2}`$, so $`f(\frac{2}{2^m}_i\mathrm{\Delta }_i)f(\mathrm{\Delta })`$. Together this yields $`I(X:Q)f(\mathrm{\Delta })`$, as required.
Now, we can easily deduce Theorem 1.2.
Proof of Theorem 1.2: Let $`\mathrm{\Delta }^{}=\frac{1}{2^m}_{x_1}\sigma _{x_1}\sigma _\mathrm{t}`$. We have:
$`\mathrm{\Delta }^{}={\displaystyle \frac{1}{2^m}}{\displaystyle \underset{x_1}{}}\sigma _{x_1}\sigma _\mathrm{t}`$ $`=`$ $`{\displaystyle \frac{1}{2^m}}{\displaystyle \underset{x_1}{}}{\displaystyle \frac{1}{2^m}}{\displaystyle \underset{x_2}{}}(\sigma _{x_1}\sigma _{x_2})_\mathrm{t}{\displaystyle \frac{1}{2^{2m}}}{\displaystyle \underset{x_1,x_2}{}}\sigma _{x_1}\sigma _{x_2}_\mathrm{t}\mathrm{\Delta }`$
By Theorem 3.1, $`I(X:Q)1H(\frac{1+\mathrm{\Delta }}{2})`$, and by Fact 2.3 we have $`1H(\frac{1+\mathrm{\Delta }}{2})1(1(\frac{\mathrm{\Delta }}{2})^2)=\frac{\mathrm{\Delta }^2}{4}`$. Thus, $`\mathrm{\Delta }^{}\mathrm{\Delta }2\sqrt{I(X:Q)}`$.
## 4 Local transition between bipartite states
Lo and Chau and Mayers proved:
###### Theorem 4.1 (Lo and Chau; Mayers)
Suppose $`|\varphi _1`$ and $`|\varphi _2`$ are two pure states in the Hilbert space $`𝒦`$, such that $`\mathrm{Tr}_𝒦|\varphi _2\varphi _2|=\mathrm{Tr}_𝒦|\varphi _1\varphi _1|`$, i.e., the reduced density matrix of $`|\varphi _2`$ to $``$ is the same as the reduced density matrix of $`|\varphi _1`$ to $``$. Then, there is a local unitary transformation $`U`$ on $`𝒦`$ such that $`IU|\varphi _2=|\varphi _1`$.
The theorem follows by examining the Schmidt decomposition of the two states.
A natural generalization of this is to the case where the reduced density matrices are close to each other but not quite the same, which is what appears in Theorem 1.3. Lo and Chau and Mayers considered this case as well. Theorem 1.3 follows from their work by using the newer results of stated in Theorem 2.7.
Proof of Theorem 1.3: By Theorem 2.6, there exists a purification $`|\varphi _2^{}𝒦`$ of $`\rho _2`$ such that $`\left|\varphi _1\varphi _2^{}\right|^2=F(\rho _1,\rho _2)`$. Since $`|\varphi _2`$ and $`|\varphi _2^{}`$ have the same reduced density matrix in $``$, by Theorem 4.1, there is a (local) unitary transformation $`U`$ on $`𝒦`$ such that $`IU|\varphi _2=|\varphi _2^{}`$. Moreover, by Fact 2.5 we have
$$|\varphi _1\varphi _1||\varphi _2^{}\varphi _2^{}|_\mathrm{t}=\mathrm{\hspace{0.33em}\hspace{0.33em}2}\sqrt{1\left|\varphi _1\varphi _2^{}\right|^2}=\mathrm{\hspace{0.33em}\hspace{0.33em}2}\sqrt{1F(\rho _1,\rho _2)}.$$
By Theorem 2.7$`\sqrt{F(\rho _1,\rho _2)}\mathrm{\hspace{0.33em}1}\frac{1}{2}\rho _1\rho _2_\mathrm{t}`$, so
$$1F(\rho _1,\rho _2)\mathrm{\hspace{0.33em}\hspace{0.33em}1}\left(1\frac{1}{2}\rho _1\rho _2_\mathrm{t}\right)^2\rho _1\rho _2_\mathrm{t}.$$
This, when combined with the earlier bound on the trace distance between $`|\varphi _1,|\varphi _2^{}`$ gives us the required result.
## 5 The role of interaction in quantum communication
In this section, we prove that allowing more interaction between two players in a quantum communication game can substantially reduce the amount of communication required. We first define a communication problem and state our results formally (giving an overview of the proof), and then give the details of the proofs.
### 5.1 The communication problem and its complexity
In this section, we give the main components of the proof of Theorem 1.1. We define a sequence of problems $`S_0,S_1,\mathrm{},S_k,\mathrm{}`$ by induction. The problem $`S_1`$ is the index function, i.e., Alice has a $`n`$-bit string $`x𝒳_1=\{0,1\}^n`$, Bob has an index $`i𝒴_1=[n]`$ and the desired output is $`S_1(x,i)=x_i`$. Suppose we have already defined the function $`S_{k1}:𝒳_{k1}\times 𝒴_{k1}\{0,1\}`$. In the problem $`S_k`$, Alice has as input her part of $`n`$ independent instances of $`S_{k1}`$, i.e., $`x𝒳_{k1}^n`$, Bob has his share of $`n`$ independent instances of $`S_{k1}`$, i.e., $`y𝒴_{k1}^n`$, and in addition, there is an extra input $`a[n]`$ which is given to Alice if $`k`$ is even and to Bob if $`k`$ is odd. The output we seek is the solution to the $`a`$th instance of $`S_{k1}`$. In other words, $`S_k(x_1,\mathrm{},x_n,a,y_1,\mathrm{},y_n)=S_{k1}(x_a,y_a)`$.
Note that the input size to the problem $`S_k`$ is $`N=\mathrm{\Theta }(n^k)`$. If we allow $`k`$ message exchanges for solving the problem, it can be solved by exchanging $`\mathrm{\Theta }(\mathrm{log}N)=\mathrm{\Theta }(k\mathrm{log}n)`$ bits: for $`k=1`$, Bob sends Alice the index $`i`$ and Alice then knows the answer; for $`k>1`$, the player with the index $`a`$ sends it to the other player and then they recursively solve for $`S_{k1}(x_a,y_a)`$. However, we show that if we allow one less message, then no quantum protocol can compute $`S_k`$ as efficiently. In fact, no quantum protocol can compute the function as efficiently even if we require small probability of error only on average.
###### Theorem 5.1
For all constant $`k1`$, $`0ϵ<\frac{1}{2}`$, $`Q_{U,ϵ}^k(S_{k+1})\mathrm{\Omega }\left(N^{1/(k+1)}\right).`$
In fact, we prove a stronger intermediate claim. Let $`P_1`$ be Alice, and for $`k2`$, let $`P_k`$ denote the player that holds the index $`a`$ in an instance of $`S_k`$ ($`a`$ indicates which of the $`n`$ instances of $`S_{k1}`$ to solve). Let $`\overline{P}_k`$ denote the other player. We refer to $`\overline{P}_k`$ as the “wrong” player to start a protocol for $`S_k`$. The stronger claim is that any $`k`$ message protocol for $`S_k`$ in which the wrong player starts is exponentially inefficient as compared to the $`\mathrm{log}N`$ protocol described above.
###### Theorem 5.2
For all constant $`k1`$, $`0ϵ<\frac{1}{2}`$, $`Q_{U,ϵ}^{k,\overline{P}_k}(S_k)\mathrm{\Omega }\left(N^{1/k}\right).`$
In fact, there is a classical $`k`$-message protocol in which the wrong player starts with complexity $`O(n)`$, so our lower bound is optimal.
Theorem 5.1 now follows directly.
Proof of Theorem 5.1: It is enough to show the lower bound for the two cases when the protocol starts either with $`P_{k+1}`$ or with the other player.
Let $`P_{k+1}`$ be the player to start. Note that if we set $`a`$ to a fixed value, say $`1`$, then we get an instance of $`S_k`$. So $`Q_{U,ϵ}^{k,P_{k+1}}(S_k)Q_{U,ϵ}^{k,P_{k+1}}(S_{k+1})`$. But $`P_{k+1}=\overline{P}_k`$, so the bound of Theorem 5.2 applies.
Let player $`\overline{P}_{k+1}`$ be the one to start. Then, observe that if we allow one more message (i.e., $`k+1`$ messages in all), the complexity of the problem only decreases: $`Q_{U,ϵ}^{k+1,\overline{P}_{k+1}}(S_{k+1})Q_{U,ϵ}^{k,\overline{P}_{k+1}}(S_{k+1})`$. So we again get the same bound from Theorem 5.2.
We prove Theorem 5.2 by induction. First, we show that the index function is hard to solve with one message if the wrong player starts. This essentially follows from the lower bound for random access codes in . The only difference is that we seek a lower bound for a protocol that has low error probability on average rather than in the worst case, so we need a refinement of the original argument. We give this in the next section.
###### Lemma 5.3
For any $`0ϵ1`$, $`Q_{U,ϵ}^{1,A}(S_1)(1H(ϵ))n`$.
Next, we show that if we can solve $`S_k`$ with $`k`$ messages with the wrong player starting, then we can also solve $`S_{k1}`$ with only $`k1`$ messages of almost the same total length, again with the wrong player starting, at the cost of a slight increase in the average probability of error.
###### Lemma 5.4
For all $`k2`$, $`0ϵ<\frac{1}{2}`$, $`Q_{U,ϵ^{}}^{k1,\overline{P}_{k1}}(S_{k1})\mathrm{}+\mathrm{log}n,`$ where $`\mathrm{}=Q_{U,ϵ}^{k,\overline{P}_k}(S_k)`$, and $`ϵ^{}=ϵ+4(\mathrm{}/n)^{1/4}`$.
We defer the proof of this lemma to a later section, but show how it implies Theorem 5.2 above.
Proof of Theorem 5.2: We prove the theorem by induction on $`k`$. The case $`k=1`$ is handled by Lemma 5.3. Suppose the theorem holds for $`k1`$. We prove by contradiction that it holds for $`k`$ as well.
If $`Q_{U,ϵ}^{k,\overline{P}_k}(S_k)=o(n)`$, then by Lemma 5.4 there is a $`k1`$ message protocol for $`S_{k1}`$ with the wrong player starting, with error $`ϵ^{}=ϵ+o(1)<\frac{1}{2}`$, and with the same communication complexity $`o(n)`$. This contradicts the induction hypothesis.
### 5.2 Hardness of the index function
We now prove the average case hardness of the index function.
Proof of Lemma 5.3: Let $`Q`$ denote the message sent by Alice. For a prefix $`y\{0,1\}^i`$ of length $`i0`$, let $`Q_y`$ be the encoding which is prepared by first fixing $`x_1=y_1,\mathrm{},x_i=y_i`$ and then choosing $`x_{i+1},\mathrm{},x_m`$ at random and sending the state $`\sigma _x`$. Its density matrix is given by
$$\sigma _y=\frac{1}{2^{mi}}\underset{z\{0,1\}^{mi}}{}\sigma _{yz}.$$
On the one hand, $`I(Q:X)\mathrm{}`$, the number of qubits in $`Q`$. On the other hand, for $`y\{0,1\}^j`$, let $`ϵ_y`$ be the error probability when $`x_l=y_l`$, $`lj`$, and the index $`i=j+1`$. Note that $`ϵ=\frac{1}{n}_{j=0}^{n1}\frac{1}{2^j}_{y\{0,1\}^j}ϵ_y`$. Moreover, we have $`I(Q_y:X_{j+1})1H(ϵ_y)`$, since Bob has a measurement that predicts $`X_{j+1}`$ with probability $`1ϵ_y`$ given $`Q_y`$. We now claim that
###### Lemma 5.5
$`\frac{1}{2^m}_x_{i=0}^{m1}I(Q_{x_1\mathrm{}x_i}:X_{i+1})I(Q:X)`$.
By this lemma,
$$I(Q:X)\underset{j=0}{\overset{n1}{}}\frac{1}{2^j}\underset{y\{0,1\}^j}{}I(Q_y:X_{j+1})(1H(ϵ))n,$$
using the concavity of the entropy function.
Proof of Lemma 5.5: By the definition of mutual information, and using Fact 2.8,
$`I(QX_1\mathrm{}X_i:X_{i+1})`$ $`=`$ $`S(QX_1\mathrm{}X_i)+S(X_{i+1})S(QX_1\mathrm{}X_{i+1})`$
$`=`$ $`[i+{\displaystyle \frac{1}{2^i}}{\displaystyle \underset{y\{0,1\}^i}{}}S(\sigma _y)]+[1][(i+1)+{\displaystyle \frac{1}{2^{i+1}}}{\displaystyle \underset{y\{0,1\}^{i+1}}{}}S(\sigma _y)]`$
$`=`$ $`{\displaystyle \frac{1}{2^i}}{\displaystyle \underset{y\{0,1\}^i}{}}[S(\sigma _y){\displaystyle \frac{1}{2}}(S(\sigma _{y0})+S(\sigma _{y1}))]`$
$`=`$ $`{\displaystyle \frac{1}{2^i}}{\displaystyle \underset{y\{0,1\}^i}{}}I(Q_y:X_{i+1}).`$
Moreover, from Properties (1) and (2),
$`I(Q:X)`$ $``$ $`{\displaystyle \underset{i=0}{\overset{m1}{}}}I(Q,X_1,\mathrm{},X_i:X_{i+1})`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{m1}{}}}{\displaystyle \frac{1}{2^i}}{\displaystyle \underset{y\{0,1\}^i}{}}I(Q_y:X_{i+1})`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{m1}{}}}{\displaystyle \frac{1}{2^m}}{\displaystyle \underset{y\{0,1\}^m}{}}I(Q_{y_1\mathrm{}y_i}:X_{i+1}),`$
which proves the claim.
### 5.3 The reduction step
In this section, we show how an efficient protocol for $`S_k`$ gives rise to an efficient protocol for $`S_{k1}`$. The gross structure of the argument is the same as in . However, we use entirely new techniques from quantum information theory, as developed in Section 3 and 4 and also get better bounds in the process.
Proof of Lemma 5.4: For concreteness, we assume that $`k`$ is even, so that $`\overline{P}_k`$ is Bob. Let $`𝒫`$ be a protocol that solves $`S_k`$ with respect to $`U`$ with $`\mathrm{}`$ message qubits, error $`ϵ`$, and $`k`$ messages starting with Bob. We would like to concentrate on inputs where $`a`$ is fixed to a particular value in $`[n]`$. This would give rise to an instance of $`S_{k1}`$ that is also solved by $`𝒫`$, but with $`k`$ messages. An easy argument shows the first message carries almost no information about $`y_a`$, and we would like to argue that it is not relevant for solving $`S_{k1}`$. However, the correctness of the protocol relies on the message, so we try to reconstruct the message with Alice starting the protocol instead. We give the details below.
We first derive a protocol $`𝒫^{}`$ which has low error on an input for $`S_k`$ generated as below (we call the resulting distribution $`U_{a=j}`$): $`x_1,\mathrm{},x_n`$ are chosen uniformly at random from $`𝒳_{k1}`$, $`a`$ is set to $`j`$, $`y_j`$ is chosen uniformly at random from $`𝒴_{k1}`$, and for all $`ij`$, register $`Y_i`$ is initialized to the state $`_{z𝒴_{k1}}|z`$ (normalized).
Let $`ϵ_j`$ denote the error of $`𝒫`$ with respect to the distribution $`U_{a=j}`$. Note that $`\frac{1}{n}_iϵ_iϵ`$, since having the $`Y_i`$ in a uniform superposition over all possible inputs has the same effect on the result of the protocol as having it randomly distributed over the inputs (recall that we require that the input registers are not changed during a quantum protocol). Let $`\mu _j`$ be the mutual information $`I(M:Y_j)`$ in the protocol $`𝒫`$ when run on the mixed state $`U_{a=j}`$ with $`y_j`$ being chosen randomly.
###### Lemma 5.6
There is a protocol $`𝒫^{}`$ which solves $`S_k`$ with respect to the distribution $`U_{a=j}`$ with error $`\delta _j=ϵ_j+4\mu _j^{1/4}`$ error, $`\mathrm{}`$ message qubits and $`k`$ rounds starting with Bob, such that $`I(M:Y_j)=0`$.
The protocol $`𝒫^{}`$ is obtained by slightly modifying the first message in protocol $`𝒫`$ so that it is completely independent of $`Y_j`$. This only affects the average probability of error. Moreover, in $`𝒫^{}`$ the first message does not carry any information about $`y_j`$ and is therefore completely independent of it. Intuitively this means that Alice does not need to get that message at all, or equivalently that she can recreate it herself. This gives a protocol for solving $`S_{k1}(x_j,y_j)`$ with $`k1`$ messages and with Alice starting.
###### Lemma 5.7
There is a protocol $`𝒫^{\prime \prime }`$ that solves $`S_{k1}`$ with respect to $`U`$ with $`ϵ^{}`$ error, $`\mathrm{}+\mathrm{log}n`$ message qubits and $`k1`$ messages starting with Alice.
Together we get $`Q_{U,ϵ^{}}^{k1,A}(S_{k1})\mathrm{}+\mathrm{log}n`$ as claimed.
### 5.4 Proof of Lemmas 5.6 and 5.7
Proof of Lemma 5.6: First consider the case when $`Y_j`$ is fixed to some $`z`$, but the rest of the inputs are as in $`U_{a=j}`$. In protocol $`𝒫`$ Bob applies a unitary transformation $`V`$ on his qubits and computes $`|\varphi (z)=V|\overline{0},Y_1,\mathrm{},Y_n`$ in register $`M`$ (for the message) and $`B`$ (for Bob’s ancilla and input). In $`𝒫^{}`$ the message computation is slightly different. Instead of computing $`|\varphi (z)`$, Bob computes $`|\varphi ^{}=V|\overline{0},Y_1,\mathrm{},Y_{j1}|\psi |Y_{j+1},\mathrm{},Y_n`$, where $`|\psi `$ is the uniform superposition over $`𝒴_{k1}`$. Clearly, in $`𝒫^{}`$ the state $`|\varphi ^{}`$ and hence the message $`M`$ does not depend on $`y_j=z`$, hence $`I(M:Y_j)=0`$ when $`y_j`$ is chosen randomly.
Let us denote by $`\rho _M(z)`$ the reduced density matrix of the message register $`M`$ in $`𝒫`$ when the input is drawn according to $`U_{a=j}`$ and $`y_j=z`$, and let the corresponding density matrix for $`𝒫^{}`$ be $`\rho _M`$. Clearly, $`\rho _M=\frac{1}{\left|𝒴_{k1}\right|}_{z𝒴_{k1}}\rho _M(z)`$. Let $`t_z=\rho _M\rho _M(z)_\mathrm{t}`$. By Theorem 1.2 we know that $`E_zt_z2\sqrt{\mu _j}`$.
Protocol $`𝒫^{}`$ generates the pure state $`|\varphi ^{}`$, while the desired pure state is $`|\varphi (z)`$. Bob, who knows $`y_j=z`$ knows both $`|\varphi (z)`$ and $`|\varphi ^{}`$. By Theorem 1.3 there is a local unitary transformation $`T_z`$ acting on register $`B`$ alone, such that
$`|T_z\varphi ^{}T_z\varphi ^{}||\varphi (z)\varphi (z)|_\mathrm{t}`$ $``$ $`2\sqrt{t_z}.`$
The next step in protocol $`𝒫^{}`$ is that Bob applies the transformation $`T_z`$ to his register $`B`$. After that, protocol $`𝒫^{}`$ proceeds exactly as in $`𝒫`$. Therefore, for a given $`z`$, the probability that $`𝒫`$ and $`𝒫^{}`$ disagree on the result is at most $`2\sqrt{t_z}`$, and the error probability of $`𝒫^{}`$ on $`U_{a=j}`$ is at most
$$\delta _j=ϵ_j+2E_z\sqrt{t_z}ϵ_j+2\sqrt{E_zt_z}ϵ_j+4\mu _j^{1/4},$$
where the second step follows from Jensen’s inequality.
Proof of Lemma 5.7: Protocol $`𝒫^{\prime \prime }`$ solves an instance of $`S_{k1}`$. Alice is given an input $`\widehat{x}_R𝒳_{k1}`$ and Bob is given an input $`\widehat{y}_R𝒴_{k1}`$. The protocol proceeds as follows. Alice and Bob first reduce the problem to an $`S_k`$ instance taken from the distribution $`U_{a=j}`$ for a random $`j`$. To do that, Alice picks $`j[n]`$ at random, sets $`a=j`$ and sends it to Bob; Alice sets $`x_j=\widehat{x}`$ and Bob sets $`y_j=\widehat{y}`$; Alice picks $`x_i_R𝒳_{k1}`$ for $`ij`$; and Bob initializes each register $`Y_i`$ for $`ij`$ with $`_{z𝒴_{k1}}|z`$ (normalized).
Notice that if Alice and Bob run the protocol $`𝒫^{}`$ over this input, then they get the answer $`S_{k1}(x,y)`$ with probability of error at most
$$ϵ^{}=\frac{1}{n}\underset{i=1}{\overset{n}{}}\delta _i\frac{1}{n}\underset{i=1}{\overset{n}{}}ϵ_i+4\frac{1}{n}\underset{i=1}{\overset{n}{}}\mu _i^{1/4}ϵ+4\left[\frac{1}{n}\underset{i=1}{\overset{n}{}}\mu _i\right]^{\frac{1}{4}}.$$
We claim that
###### Claim 5.8
$`_i\mu _i\mathrm{}_1`$, where $`\mathrm{}_1`$ is the length of the message $`M`$.
Hence $`ϵ^{}ϵ+4(\mathrm{}/n)^{1/4}`$.
Alice and Bob do not run the protocol $`𝒫^{}`$ itself, but a modification of it in which Alice sends the first message instead of Bob, thus reducing the number of rounds to $`k1`$.
Let $`\rho _M`$ be the reduced density matrix of register $`M`$ holding the first message that Bob sends to Alice in $`𝒫^{}`$, for the input given above. By Lemma 5.6, we know that $`\rho _M`$ does not depend on $`y_j=\widehat{y}`$. So $`\rho _M`$ is known in advance to Alice. Alice starts the protocol $`P^{\prime \prime }`$ by purifying $`\rho _M`$. More specifically, let $`\left\{|e_i\right\}`$ be an eigenvector basis for $`\rho _M`$ with real and positive eigenvalues $`\lambda _i`$. Alice constructs the superposition $`_i\sqrt{\lambda _i}|e_i,i_{MB}`$ over two registers $`M`$ (containing the eigenvectors) and $`B`$ (containing the index $`i`$), and sends register $`B`$ to Bob. The state of the system after this message in $`𝒫^{\prime \prime }`$ is
$`|\xi `$ $`=`$ $`|x_1,\mathrm{},x_n_A{\displaystyle \underset{i}{}}\sqrt{\lambda _i}|e_i_M|i_B`$
whereas in $`𝒫^{}`$ it is
$`|\chi (y)`$ $`=`$ $`|x_1,\mathrm{},x_n_A|T_y\varphi ^{}_{MB}.`$
The reduced density matrix of $`|\xi `$ to registers $`AM`$ is the same as the reduced density matrix of $`|\chi (y)`$ to registers $`AM`$. By Theorem 4.1, Bob has a local unitary transformation $`V_y`$ (operating on his register $`B`$) that transforms $`|\xi `$ to $`|\chi (y)`$. Bob applies $`V_y`$, and Alice and Bob then simulate the rest of the protocol $`𝒫^{}`$. From this stage on, the runs of the protocols $`𝒫^{}`$ and $`𝒫^{\prime \prime }`$ are identical have the same communication complexity and success probability.
Proof of Claim 5.8: Note that $`\mu _j`$ is the same as the mutual information $`I(M:Y_j)`$ when $`𝒫`$ is run on the uniform distribution on $`𝒳_{k1}^n\times 𝒴_{k1}^n`$. So we prove the claim for the latter.
For any $`i`$, $`I(Y_i:Y_1\mathrm{}Y_{i1}Y_{i+1}\mathrm{}Y_n)=0`$. Therefore by Properties (1) and (2) (cf. Section 2) we have
$$I(M:Y_1\mathrm{}Y_n)\underset{i=1}{\overset{n}{}}I(MY_1\mathrm{}Y_{i1}:Y_i)\underset{i=1}{\overset{n}{}}I(M:Y_i)=\underset{i}{}\mu _i$$
As the first message $`M`$ contains only $`\mathrm{}_1`$ qubits, we have $`_i\mu _iI(M:Y_1\mathrm{}Y_n)\mathrm{}_1`$.
### Acknowledgements
We thank Jaikumar Radhakrishnan and Venkatesh Srinivasan for their input on the classical communication complexity of the pointer jumping and the subproblem we study in this paper, and Dorit Aharonov for helpful comments.
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# Optically Driven Qubits in Artificial Molecules
## I Introduction
Since useful algorithms such as exhaustive search and factorization have been proposed, there has been considerable interest in quantum gates which are the basic units for quantum computation and information processing. Quantum gates manipulate quantum states through the unitary transformation with an externally driven signal and it was shown that two kinds of a quantum gate are enough for quantum computational procedure; a single-bit and two-bit gates. A single-bit gate controls a single-particle state, called a qubit which is the superposition of two orthogonal states namely $`0`$ and $`1`$; mathematically a qubit is written as $`\psi =\mathrm{cos}\theta e^{i\phi }0+\mathrm{sin}\theta 1`$. Thus, with controllable parameters $`\theta `$ and $`\phi `$, $`\psi `$ spans the entire Hilbert space of $`0`$ and $`1`$ by a unitary transformation, including a reversible NOT operation; $`\psi \psi _r=\mathrm{cos}\theta e^{i\phi }1+\mathrm{sin}\theta 0`$ on the single-bit gate. A two-bit gate (controlled NOT gate) functions on two single-qubits called a target bit $`\psi `$ and a control bit $`\chi =\mathrm{cos}\theta _ce^{i\phi _c}0_c+\mathrm{sin}\theta _c1_c`$ and evolves the target bit conditionally, i.e., depending on the state of the control bit. In detail, if the control bit is $`1_c`$, the gate performs the reversible NOT operation on the target bit. Otherwise, there is no change. Hence, for a initial state of $`\psi \chi `$, we obtain, by the controlled NOT operation,
$`\psi \chi \mathrm{cos}\theta _ce^{i\phi _c}\psi 0_c+\mathrm{sin}\theta _c\psi _r1_c`$ (1)
where the parameters $`\theta _c`$ and $`\phi _c`$ characterize the state of the control bit.
Various quantum gates have been explored for realizing a quantum computer, including trapped ions, spins on a nuclear, and cooper-pairing states in Joshepson box. Especially, quantum bits based on the semiconductor quantum dots draw attention relative to others because advance in nanosemiconductor technology makes it possible to tailor the quantum dots. Both spin and spatial parts of wavefunctions for an electron confined in quantum dot can be exploited as a qubit. As shown by several authors, up-and down-spin states of an electron were found to be a good basis for a qubit which is controlled by a magnetic field. It was also shown that the operation of the controlled NOT gate is easily modeled using the Heisenberg exchange interaction between two-electron spins.
The spatial part of an electronic wavefunction is also very interesting for a qubit because its energy can be tailored easily by gate electrodes. Moreover, a recent experiment succeeded in controlling the spatial part of wavefunctions by a light pulse. A typical model along this scheme uses the superposition of the ground and the first excited states of a single quantum dot as a qubit. However, since electrons at the excited state are relaxed rapidly to the ground state by the phonon process, the coherent time of the case is estimated to be shorter than that of the spin state. More advanced model is proposed by Openov recently where the superposition of each ground states of two separated quantum dots is viewed as a qubit. Since two quantum dots are assumed to interact each other only by their excited states, one expects that an electron at the ground states has longer coherence time than that of the spin case once it is defined. It is also shown that by optical illumination, an electron can transfer between dots coherently and one can complete the reversible NOT operation with appropriate frequency and strength of the light pulse. However, even though such a qubit is very feasible, models of single-bit and two-bit gates are not suggested and detailed discussion for the coherent time is still lacking.
In this work, we suggest novel models of quantum gates based on coupled quantum dots or artificial molecules. By solving the time-dependent Schrődinger equation, we show that simple artificial molecules associated with light pulses serve as the quantum gates. In addition, by calculating the relaxation rate by phonon process, we also discuss the decoherence of a quantum state during the unitary operation. For the quantum dots, we consider two-dimensional or disk-like shape with the lateral size much larger than the extent in the growth direction ($`z`$-direction) by patterning isolated metallic gates or etching vertically quantum wells. Since the size of the quantum dot is comparable to the effective Bohr radius of a host semiconductor, discrete energy levels are formed in the quantum dot where the number of electrons and confinement potential are controlled artificially. Thus, since the number of quantum states in the quantum dot depends on the confinement potential or the radius of the quantum dot, we assume that one or two levels are bound at each isolated quantum dot; the ground and first excited states. Furthermore, in our case, since the quantum dot is circular symmetric about $`z`$-axis, the ground state has an angular momentum $`l=0`$($`s`$ state) and the first excited state with $`l=\pm 1`$($`p`$ states) is degenerate.
## II Single-bit gate
In our model, the single-bit gate consists of two larger quantum dots and one smaller dot which are embedded in a barrier material with a potential energy $`V_b`$ as shown in Fig. 1. Depending on its radius, each dot is assumed to have the different number of atomic states with itself if isolated from others. Two larger quantum dots named $`A`$ and $`B`$ have both the ground ($`l=0`$) and the first excited states ($`l=\pm 1`$) with energies $`ϵ_s`$ and $`ϵ_p`$ respectively while a smaller quantum named $`C`$ has only the ground state $`l=0`$ with a energy $`ϵ_s^{}`$. Further, the energies $`ϵ_p`$ and $`ϵ_s^{}`$ of the atomic states at the dot $`A`$, $`B`$, and $`C`$ are assumed to lie close to the barrier potential energy $`V_b`$. As a result, wavefunctions of these states have a large extension spatially and the quantum dots can interact each other via these states to be an artificial molecule. However, the atomic ground states of the dot $`A`$ and $`B`$ are supposed to have deep energies and be well-isolated from states at the other quantum dots. This means that an excess electron occupied at the ground states has a large coherence time as discussed by Openov. For this reason, we exploit the atomic ground states of the dots $`A`$ and $`B`$ as the basis of a qubit, i.e. with the ground states for the excess electron in the system of Fig. 1, we write a qubit as,
$`\psi =\mathrm{cos}\theta e^{i\phi }0+\mathrm{sin}\theta 1=\mathrm{cos}\theta e^{i\phi }0;A+\mathrm{sin}\theta 0;B`$ (2)
where $`00;A`$, $`10;B`$, and $`l;\alpha `$ is a state located at the quantum dot $`\alpha `$ with an angular momentum $`l`$.
Then, single-particle states of the artificial molecule in Fig. 1 are determined by the Hamiltonian $`H_0`$,
$`H_0`$ $`=`$ $`ϵ_s(d_{0A}^{}d_{0A}+d_{0B}^{}d_{0B})+ϵ_p(d_{\sigma A}^{}d_{\sigma A}+d_{\sigma B}^{}d_{\sigma B}+d_{\pi A}^{}d_{\pi A}+d_{\pi B}^{}d_{\pi B}+d_{0C}^{}d_{0C})`$ (4)
$`+\{Vd_{0C}^{}d_{\sigma A}+Vd_{0C}^{}d_{\sigma B}+\mathrm{H}.\mathrm{c}.\}`$
where $`ϵ_s^{}=ϵ_p`$ for simplicity and $`d_{l\alpha }`$($`d_{l\alpha }^{}`$) is the electron annihilation(creation) operators for an electron in states $`l;\alpha `$ ($`\alpha =A,B,C;l=0,\pm 1`$). Here, operators $`d_{\alpha \sigma }`$ and $`d_{\alpha \pi }`$ describe two orthogonal states localized at the dot $`\alpha `$ and are defined as
$`d_{\sigma \alpha }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(e^{i\varphi _\alpha }d_{1\alpha }+e^{i\varphi _\alpha }d_{1\alpha })`$ (5)
$`d_{\pi \alpha }`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{2}}}(e^{i\varphi _\alpha }d_{1\alpha }e^{i\varphi _\alpha }d_{1\alpha })`$ (6)
with $`\varphi _\alpha `$ an angle between the $`x`$-axis and the dot $`\alpha `$. These states have $`p`$-like wavefunctions and, as shown in Fig. 1-(b), the $`\sigma ;\alpha `$ states have their globes of wavefunctions directed to the dot C from the dot $`\alpha `$ whereas those of the $`\pi ;\alpha `$ states are normal to the line connecting the dot C with the dot $`\alpha `$. The last term of Eq. (4) represents the interaction between excited states of the quantum dots. Due to a geometircal distance, atomic states of the dots $`A`$ and $`B`$ are supposed not to be coupled to each other, but to be coupled to the ground state of the dot $`C`$ with a strength $`V`$. It is noted that $`\pi ;\alpha `$ is not coupled to $`0;C`$ because of a geometrical symmetry. The Hamiltonian $`H_0`$ is easily diagonalized and its molecular orbitals $`k`$ are given as,
$`0`$ $`=`$ $`0;A,1=0;B,`$ (7)
$`2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\sqrt{2}0;C\sigma ;B\sigma ;A),`$ (8)
$`3`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\sigma ;B\sigma ;A),4=\pi ;A,5=\pi ;B,`$ (9)
$`6`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\sqrt{2}0;C+\sigma ;B+\sigma ;A)`$ (10)
with their eigenenergies $`ϵ_0=ϵ_1=ϵ_s`$, $`ϵ_2=ϵ_p\sqrt{2}V`$, $`ϵ_3=ϵ_4=ϵ_5=ϵ_p`$, and $`ϵ_6=ϵ_p+\sqrt{2}V`$. As noted previously, both $`0;A`$ and $`0;B`$ are still eigenstates of the artificial molecule in Fig. 1 and the qubit of Eq. (2) is not evolved without a perturbation.
In order to evolve the qubit, we use optical transitions between molecular states $`k`$ of Eq. (10). Under the classical field approximation, an excess electron interacts with a light by a Hamiltonian $`H_{light}(t)`$,
$`H_{light}(t)={\displaystyle \frac{e}{2m^{}c}}(\stackrel{}{A}\stackrel{}{p}+\stackrel{}{p}\stackrel{}{A})`$ (11)
where $`\stackrel{}{A}`$ is the vector potential of the light, $`\stackrel{}{p}`$ is a momentum operator of the electron, and the $`\stackrel{}{A}^2`$ term is omitted for simplicity. If the light propagates along the $`z`$-axis and its electric field is polarized by an angle $`\varphi `$ from the $`x`$-axis, the vector potential $`\stackrel{}{A}`$ is given by,
$`\stackrel{}{A}={\displaystyle \frac{cE_0}{\omega }}(\widehat{x}\mathrm{cos}\varphi +\widehat{y}\mathrm{sin}\varphi )\mathrm{cos}(\omega z/c\omega t)`$ (12)
with an electric field strength $`E_0`$ and a frequency $`\omega `$. Hence, at the $`z`$-plane on which the quantum dots reside, the optical transitions of the ground states to or from excited states of the molecule are governed by,
$`H_{light}(t)`$ $`=`$ $`{\displaystyle \frac{i(ϵ_pϵ_s)E_0\xi }{\sqrt{2}\mathrm{}\omega }}\mathrm{cos}(\omega t)[\mathrm{cos}(\varphi _B+\varphi )d_0^{}d_2\mathrm{cos}(\varphi _B\varphi )d_1^{}d_2`$ (15)
$`+\sqrt{2}\mathrm{cos}(\varphi _B+\varphi )d_0^{}d_3+\sqrt{2}\mathrm{cos}(\varphi _B\varphi )d_1^{}d_32\mathrm{sin}(\varphi _B+\varphi )d_0^{}d_4`$
$`2\mathrm{sin}(\varphi _B\varphi )d_1^{}d_5\mathrm{cos}(\varphi _B+\varphi )d_0^{}d_6+\mathrm{cos}(\varphi _B\varphi )d_1^{}d_6]+\mathrm{H}.\mathrm{c}.`$
where $`\xi 0;\alpha ex1;\alpha `$ is the electric dipole moment and the optical transitions between excited states are omitted.
To study the evolution of the wavefunction caused by the illumination of the light, we write the wavefunction of the excess electron as,
$`\mathrm{\Psi }(t)={\displaystyle \underset{k=0}{\overset{6}{}}}S_k(t)\mathrm{exp}\{iϵ_kt/\mathrm{}\}k.`$ (16)
Here, the expansion coefficient $`S_k`$ denotes the probability amplitude of the state $`k`$ and is determined by the time-dependent Schődinger equation,
$`i\mathrm{}{\displaystyle \frac{\mathrm{\Psi }(t)}{t}}=H(t)\mathrm{\Psi }(t)`$ (17)
with $`H=H_0+H_{light}(t)`$, and the values of $`S_0(0)`$ and $`S_1(0)`$ define an initial qubit with $`S_k(0)=0(k=2,3,4,5,6)`$.
For an effective control of the excess electron, we exploit the optical transition through the first excited state $`2`$ by tuning the frequency $`\omega `$ much more closer to $`(ϵ_2ϵ_0)/\mathrm{}`$ than $`(ϵ_nϵ_0)/\mathrm{}`$ ($`n=3,4,5,6`$). Then, the ground and the first excited states are nearly resonant with the light while the higher excited states are out of resonance. In this work, to a good approximation, we neglect the non-resonant optical transitions for the range of the frequency,
$`\mathrm{}\delta ϵ_3ϵ_2=\sqrt{2}V`$ (18)
with $`\delta (ϵ_2ϵ_0)/\mathrm{}\omega `$. Then, the Hamiltonian $`H`$ governing the wavefunction of Eq. (16) is replaced with $`H_r`$,
$`H_r={\displaystyle \underset{k=0}{\overset{6}{}}}ϵ_kd_k^{}d_k+\{{\displaystyle \frac{i\mathrm{}\mathrm{\Omega }}{4}}e^{i\omega t}[ud_0^{}d_2vd_1^{}d_2]+\mathrm{H}.\mathrm{c}.\}`$ (19)
with $`u=\mathrm{cos}(\varphi _B+\varphi )`$, $`v=\mathrm{cos}(\varphi _B\varphi )`$, and $`\mathrm{}\mathrm{\Omega }\sqrt{2}E_0\xi (ϵ_pϵ_s)/(ϵ_2ϵ_0)`$.
The simple form of the Hamiltonian $`H_r`$ in Eq. (19) makes it possible to obtain an analytic solution of the wavefunction Eq. (16). To do this, first we make use of the unitary transformation,
$`U(t)=\mathrm{exp}\{i\omega t/2(d_0^{}d_0+d_1^{}d_1d_2^{}d_2)\}.`$ (20)
Substituting $`\mathrm{\Psi }(t)=U(t)\stackrel{~}{\mathrm{\Psi }}(t)`$ into the Schrődinger equation Eq. (17), we obtain another equation for $`\stackrel{~}{\mathrm{\Psi }}(t)`$:
$`i\mathrm{}{\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}(t)}{t}}=\stackrel{~}{H}\stackrel{~}{\mathrm{\Psi }}(t)`$ (21)
with the rotated Hamiltonian $`\stackrel{~}{H}`$
$`\stackrel{~}{H}`$ $`=`$ $`U^{}(t)H_rU(t)i\mathrm{}U^{}(t){\displaystyle \frac{U(t)}{t}}`$ (22)
$`=`$ $`(ϵ_0+\mathrm{}\omega /2)(d_0^{}d_0+d_1^{}d_1)+(ϵ_2\mathrm{}\omega /2)d_2^{}d_2+{\displaystyle \underset{k=3}{\overset{6}{}}}ϵ_kd_k^{}d_k`$ (24)
$`+\{{\displaystyle \frac{i\mathrm{}\mathrm{\Omega }}{4}}[ud_0^{}d_2vd_1^{}d_2]+\mathrm{H}.\mathrm{c}.\}.`$
Since the Hamiltonian $`\stackrel{~}{H}`$ is independent on time, the solution for $`\stackrel{~}{\mathrm{\Psi }}(t)`$ is given by;
$`\stackrel{~}{\mathrm{\Psi }}(t)=\mathrm{exp}[i\stackrel{~}{H}t/\mathrm{}]\stackrel{~}{\mathrm{\Psi }}(0).`$ (25)
Then, as soon as the light pulse is illuminated at the instance $`t=0`$, the probability amplitude $`S_k(t)`$ in Eq. (16) evolves as,
$`S_k(t)={\displaystyle \underset{m,n=0}{\overset{2}{}}}e^{iϵ_kt/\mathrm{}}kU(t)mm\mathrm{exp}[i\stackrel{~}{H}t/\mathrm{}]nS_n(0),k=0,1,2`$ (26)
and $`S_k(t)=0(k=3,4,5,6)`$. Since the evolution operator $`\mathrm{exp}[i\stackrel{~}{H}t/\mathrm{}]`$ is diagonal for eigenvectors of the Hamiltonian $`\stackrel{~}{H}`$, it is calculated, on the basis $`k(k=0,1,2)`$, as
$`k^{}`$ $``$ $`\mathrm{exp}[{\displaystyle \frac{(2\stackrel{~}{H}ϵ_0ϵ_2)t}{2i\mathrm{}}}]k={\displaystyle \frac{e^{i\delta t/2}}{u^2+v^2}}\left(\begin{array}{ccc}v^2& uv& 0\\ uv& u^2& 0\\ 0& 0& 0\end{array}\right)+{\displaystyle \frac{\mathrm{cos}(\mathrm{\Omega }_tt/2)}{u^2+v^2}}\left(\begin{array}{ccc}u^2& uv& 0\\ uv& v^2& 0\\ 0& 0& u^2+v^2\end{array}\right)`$ (33)
$`+`$ $`{\displaystyle \frac{\mathrm{sin}(\mathrm{\Omega }_tt/2)}{2\mathrm{\Omega }_t(u^2+v^2)}}\left(\begin{array}{ccc}2i\delta u^2& 2i\delta uv& \mathrm{\Omega }u(u^2+v^2)\\ 2i\delta uv& 2i\delta v^2& \mathrm{\Omega }v(u^2+v^2)\\ \mathrm{\Omega }u(u^2+v^2)& \mathrm{\Omega }v(u^2+v^2)& 2i\delta (u^2+v^2)\end{array}\right)`$ (37)
with $`\mathrm{\Omega }_t(\delta )=\sqrt{\delta ^2+\mathrm{\Omega }^2(u^2+v^2)/4}`$. Substituting Eq. (37) into Eq. (25), we can see that the off-diagonal components of Eq. (37) describes the transfer of the electron between quantum states because the operator $`U(t)`$ is diagonal on the basis $`k`$. Especially, the transfer of the electron from or to the state $`2`$ is determined by the last matrix of Eq. (37). Hence the occupancy of the state $`2`$ is oscillating proportional to $`\mathrm{sin}(\mathrm{\Omega }_tt/2)`$ during the illumination of the light. For this reason, we choose the duration $`\tau `$ of the light pulse to be $`\mathrm{sin}(\mathrm{\Omega }_t\tau /2)=0`$ or
$`\tau =2N\pi /\mathrm{\Omega }_t(\delta )`$ (38)
with an integer $`N`$. Then, after a pulse is completed, the excess electron still occupies the ground states of the dots $`A`$ and $`B`$, i.e., for $`t\tau `$, the wavefunction is given by
$`\mathrm{\Psi }(t)`$ $`=`$ $`e^{iϵ_0t}[S_0(\tau )0+S_1(\tau )1)]`$ (39)
where, from Eq. (26),
$`\left(\begin{array}{c}S_0(\tau )\\ S_1(\tau )\end{array}\right)`$ $`=`$ $`(\varphi ,\delta )\left(\begin{array}{c}S_0(0)\\ S_1(0)\end{array}\right)`$ (44)
$`(\varphi ,\delta )`$ $`=`$ $`e^{i\delta \tau /4}\left(\begin{array}{cc}\mathrm{cos}(2\mathrm{\Delta }\theta )\mathrm{cos}(\frac{\delta \tau }{4})+i\mathrm{sin}(\frac{\delta \tau }{4})& \mathrm{sin}(2\mathrm{\Delta }\theta )\mathrm{cos}(\frac{\delta \tau }{4})\\ \mathrm{sin}(2\mathrm{\Delta }\theta )\mathrm{cos}(\frac{\delta \tau }{4})& \mathrm{cos}(2\mathrm{\Delta }\theta )\mathrm{cos}(\frac{\delta \tau }{4})+i\mathrm{sin}(\frac{\delta \tau }{4})\end{array}\right)`$ (47)
for an odd integer $`N`$ with $`\mathrm{cos}(\mathrm{\Delta }\theta )u/\sqrt{u^2+v^2}`$ and $`\mathrm{sin}(\mathrm{\Delta }\theta )v/\sqrt{u^2+v^2}`$.
Thus, varying the polarization or (and) the frequency of the light pulse we can control the single qubit. For instance, the relative phase between $`0`$ and $`1`$ in the qubit can be manipulated by the frequency of the light pulse with the polarization angle $`\varphi =\varphi _B+\pi /2`$, i.e., $`v=0`$. From Eq. (47), we find that the operator $`(\varphi _B+\pi /2,\delta )`$ changes the phase of the state $`0`$ by $`\pi \delta \tau /2`$ relative to the state $`1`$,
$`S_0(0)0+S_1(0)1\stackrel{}{}^{𝒾(\pi \delta \tau /\mathcal{2})}𝒮_\mathcal{0}(\mathcal{0})\mathcal{0}+𝒮_\mathcal{1}(\mathcal{0})\mathcal{1}.`$ (48)
So, by varying the tuning frequency $`\delta `$ in the range of,
$`{\displaystyle \frac{\mathrm{\Omega }\mathrm{sin}(2\varphi _B)}{\sqrt{N^21}}}\delta {\displaystyle \frac{\mathrm{\Omega }\mathrm{sin}(2\varphi _B)}{\sqrt{N^21}}},`$ (49)
we get the relative phase changing from a zero to $`2\pi `$ for $`N3`$. This range of the tuning frequency $`\delta `$ is also small enough to satisfy the resonant approximation of Eq. (18) because of usually $`\mathrm{}\mathrm{\Omega }V`$. It is noted in Eq. (49) that $`\varphi _B0`$ for the control of the relative phase or, in other words, the artificial molecule should be bent on the $`z`$-plane like a water molecule.
On the other hand, the light pulse described by $`(\varphi ,0)`$ \[or $`(0,\delta )`$\] is very useful to control an amount of the excess electron at the quantum dots because its rotation matrix is given by,
$`(\varphi ,0)=\left(\begin{array}{cc}\mathrm{cos}(2\mathrm{\Delta }\theta )& \mathrm{sin}(2\mathrm{\Delta }\theta )\\ \mathrm{sin}(2\mathrm{\Delta }\theta )& \mathrm{cos}(2\mathrm{\Delta }\theta )\end{array}\right).`$ (52)
If this light pulse is illuminated on the qubit of $`\psi _0=\mathrm{cos}\theta _00+\mathrm{sin}\theta _01`$ with a zero relative phase, the operation $`(\varphi ,0)`$ transforms $`\theta _0`$ to $`\theta _0+2\mathrm{\Delta }\theta +\pi `$. Using this fact, we can prepare an arbitrary qubit of Eq. (2) by applying two successive light pulses of $`(\varphi _B+\pi /2,\delta )(\varphi ,0)`$ on the excess electron located initially at the state $`0`$. Then, we obtain the final qubit $`\psi `$ with $`\theta =2\mathrm{\Delta }\theta +\pi `$ and $`\phi =\delta \tau /2`$. As a special case of the operation $`(\varphi ,0)`$, the light pulse described by $`(0,0)`$ serves as the reversible NOT operation, as also found by Openov;
$`\psi =S_0(0)0+S_1(0)1\stackrel{}{}\psi _𝓇=𝒮_\mathcal{1}(\mathcal{0})\mathcal{0}+𝒮_\mathcal{0}(\mathcal{0})\mathcal{1}.`$ (53)
## III Controlled NOT gate
For the controlled NOT operation, we consider the arrangement of the quantum dots as shown in Fig. 2 where a target(control) bit is regarded as the superposition of atomic ground states at two larger quantum dots located at the upper(lower) side. In detail, the first electron ”1” expresses the target bit by occupying the atomic ground states of the dot $`F`$ and $`G`$ as,
$`\psi =\mathrm{cos}\theta e^{i\phi }0_t+\mathrm{sin}\theta 1_t=\mathrm{cos}\theta e^{i\phi }0;F_1+\mathrm{sin}\theta 0;G_1`$ (54)
and the occupation of the second electron ”2” on the dots $`J`$ and $`K`$ defines the control bit as,
$`\chi =\mathrm{cos}\theta _ce^{i\phi _c}0_c+\mathrm{sin}\theta _c1_c=\mathrm{cos}\theta _ce^{i\phi _c}0;J_2+\mathrm{sin}\theta _c0;K_2.`$ (55)
As the case of a single-bit gate, we assume that the interaction among the larger quantum dots is mediated via a single ground state of a smaller quantum dot $`I`$. Furthermore, due to a special location of the smaller quantum dot $`I`$ in Fig. 2, atomic excited states($`l=\pm 1`$) of the three larger quantum dots $`F,G,J`$ are coupled to the ground state of the dot $`I`$ while the dot $`K`$ is supposed to be isolated from others. For this system, the single-particle Hamiltonian is given as
$`H_0^C`$ $`=`$ $`ϵ_s(d_{0F}^{}d_{0F}+d_{0G}^{}d_{0G}+d_{0J}^{}d_{0J}+d_{0K}^{}d_{0K})`$ (59)
$`+ϵ_p(d_{\sigma F}^{}d_{\sigma F}+d_{\sigma G}^{}d_{\sigma G}+d_{\sigma J}^{}d_{\sigma J}+d_{\sigma K}^{}d_{\sigma K})`$
$`+ϵ_p(d_{\pi F}^{}d_{\pi F}+d_{\pi G}^{}d_{\pi G}+d_{\pi J}^{}d_{\pi J}+d_{\pi K}^{}d_{\pi K})+ϵ_pd_{0I}^{}d_{0I}`$
$`+\{Vd_{0I}^{}d_{\sigma F}+Vd_{0I}^{}d_{\sigma G}+V^{}d_{0I}^{}d_{\sigma J}+\mathrm{H}.\mathrm{c}.\}`$
where we denote the coupling strength between $`\sigma ;J`$ and $`0;I`$ as $`V^{}`$ which can be different from the value of $`V`$ between $`0;I`$ and $`\sigma ;F,G`$ depending on the dot-dot distance.
In fact, since the system contains two excess electrons representing the control and the target bits, respectively, there is in general the electron-electron interaction. Then, this interaction may lead to the hybridization of the ground states at each larger quantum dots with other states, so that the basis of a quantum bit may be no longer defined as the superposition of the atomic ground states. To resolve this problem, we incorporate a metal layer into the substrate, i.e. below the quantum dots. Then, since the excess electrons in the quantum dots are screened by charges in the metal layer, the electron-electron interaction could be short-ranged. If the distance of the quantum dots to the metal layer is sufficiently small, it is reasonable to assume that the Coulomb interaction is neglected for two electrons located at different quantum dots, respectively. Under this condition, we have still four single-particle ground states whose wavefunctions are localized at each larger quantum dots and the basis of a qubit used in the single-bit gate is well-defined.
Including the fact that the electron-electron interaction is short-ranged, the Hamiltonian of the two excess electrons in the controlled NOT gate of Fig. 2 can be written as,
$`H^C=H_0^C(1)+H_0^C(2)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha ,i,j,k,l}{}}V_{ijkl}d_{i\alpha }^{}(1)d_{j\alpha }^{}(2)d_{k\alpha }(2)d_{l\alpha }(1)`$ (60)
where $`\alpha `$ runs over all the quantum dots and $`i,j,k,`$ and $`l`$ designate an angular momentum of a state at the quantum dot $`\alpha `$. Here, $`V_{ijkl}`$ is the matrix element of the Coulomb potential which has a non-zero value only when two electrons occupy a quantum dot simultaneously. As one expects, the ground states of the Hamiltonian Eq. (60) are just the direct product of two single-particle ground states located at different larger quantum dots as,
$`0;\alpha _10;\beta _2(\alpha \beta \alpha ,\beta =\mathrm{larger}\mathrm{dots}).`$ (61)
because there is no the Coulomb interaction between them. Excited states of two excess electrons in the artificial molecule of Fig. 2 are in general expressed in the linear combination of all possible states $`i;\alpha _1j;\beta _2`$ except for the states in Eq. (61) and can be obtained in a numerical way. However, for low-lying excited states, an analytic form of a wavefunction can be derived. Since two electrons occupying the same site of the quantum dots require a large charging energy, to a good approximation, we expect that states such as $`i;\alpha _1j;\alpha _2`$ are not hybridized to the low-lying excited states. Then, through a simple calculation, we find that several low-lying excited states of two electrons in the controlled NOT gate have a form of
$`0;\alpha _1\overline{\alpha }_2\mathrm{or}\overline{\alpha }_10;\alpha _2\alpha =F,G,J,K.`$ (62)
Here, $`\overline{\alpha }`$ is the lowest excited state of a single particle in Eq. (59) if the quantum dot $`\alpha `$ is not exist in Fig. 2 and is given as, for each dot $`\alpha `$,
$`\overline{K}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\eta ^2+4}}}\left\{\sigma ;F+\sigma ;G+\eta \sigma ;J\sqrt{2+\eta ^2}0;I\right\},`$ (63)
$`\overline{J}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\sigma ;F+\sigma ;G\sqrt{2}0;I),`$ (64)
$`\overline{G}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\eta ^2+2}}}\left\{\sigma ;F+\eta \sigma ;J\sqrt{1+\eta ^2}0;I\right\},`$ (65)
$`\overline{F}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\eta ^2+2}}}\left\{\sigma ;G+\eta \sigma ;J\sqrt{1+\eta ^2}0;I\right\}`$ (66)
where their eigenenergies are $`ϵ_p\sqrt{2+\eta ^2}V`$, $`ϵ_p\sqrt{2}V`$, $`ϵ_p\sqrt{1+\eta ^2}V`$, and $`ϵ_p\sqrt{1+\eta ^2}V`$, respectively, with $`\eta V^{}/V`$. In Fig. 3, we show the energy spectrum of the two electrons at the controlled NOT gate resulted from Eqs. (61) and (62).
To complete the controlled NOT operation, we drive a two-particle state by illuminating a monochromatic light pulse polarized along the $`x`$-direction. Using Eqs. (11) and (12) we derive the interaction of an electron in the system of Fig. 2 with the light as,
$`H_{light}^C`$ $`=`$ $`{\displaystyle \frac{i\sqrt{2}(ϵ_pϵ_s)\xi E_0}{\mathrm{}\omega }}\mathrm{cos}(\omega t)[\mathrm{cos}(\varphi _G)d_{0F}^{}d_{\sigma F}+\mathrm{sin}(\varphi _G)d_{0F}^{}d_{\pi F}\mathrm{cos}(\varphi _G)d_{0G}^{}d_{\sigma G}`$ (68)
$`+\mathrm{sin}(\varphi _G)d_{0G}^{}d_{\pi G}d_{0J}^{}d_{\pi J}\mathrm{cos}(\varphi _K)d_{0K}^{}d_{\sigma K}+\mathrm{sin}(\varphi _K)d_{0K}^{}d_{\pi K}]+\mathrm{H}.\mathrm{c}.`$
where $`\varphi _G`$($`\varphi _K`$) is an angle between the $`x`$-axis and the dot $`G`$($`K`$). It is noted that in the dot $`J`$ the state $`\sigma ;J`$ is inactive optically because the polarization of the light is normal to a globe of its wavefunction. Then, the total Hamiltonian governing the motion of two particles in the controlled NOT gate becomes
$`H_{total}^C=H^C+H_{light}^C(1)+H_{light}^C(2)`$ (69)
from Eq. (60).
In order to demonstrate the controlled NOT operation of the Hamiltonian $`H_{total}^C`$, we expand the two-particle wavefunction in terms of states of Eqs. (61) and (62) as,
$`\mathrm{\Psi }_2(t)`$ $`=`$ $`\mathrm{exp}\{iH^Ct/\mathrm{}\}{\displaystyle \underset{\alpha }{}}[{\displaystyle \underset{\beta \alpha }{}}S_{\alpha \beta }(t)0;\alpha _10;\beta _2`$ (71)
$`+Y_\alpha (t)\overline{\alpha }_10;\alpha _2+Z_\alpha (t)0;\alpha _1\overline{\alpha }_2]`$
and then, solve the time-dependent Schődinger equation of the Hamiltonian $`H_{total}^C`$ for an initial two-qubit, $`\mathrm{\Psi }_2(0)=\psi \chi `$ defined in Eqs. (54) and (55). In this case, the initial values of $`S_{\alpha \beta }(t)`$ are given by,
$`S_{FJ}(0)`$ $`=`$ $`\mathrm{cos}\theta e^{i\phi }\mathrm{cos}\theta _ce^{i\phi _c},S_{GJ}(0)=\mathrm{sin}\theta \mathrm{cos}\theta _ce^{i\phi _c},`$ (72)
$`S_{FK}(0)`$ $`=`$ $`\mathrm{cos}\theta e^{i\phi }\mathrm{sin}\theta _c,S_{GK}(0)=\mathrm{sin}\theta \mathrm{sin}\theta _c`$ (73)
and all the others are zero. Considering the initial condition that the electron ”2” is located at the ground states of the dots $`J`$ and $`K`$ for $`t0`$, we can simplify the problem further. That is, the electron ”2” can not be evolved to the state $`\overline{\alpha }`$ because both the states $`0;J`$ and $`0;K`$ are not coupled to any excited states of Eq. (62) according to Eq. (68). So, $`Z_\alpha (t)=0`$ and Eq. (71) is reduced to
$`\mathrm{\Psi }_2(t)`$ $`=`$ $`\mathrm{exp}\{2iϵ_st/\mathrm{}\}\left[S_{FK}(t)0;F_1+S_{GK}(t)0;G_1+Y_K(t)e^{i\omega _Kt}\overline{K}_1\right]0;K_2+`$ (75)
$`\mathrm{exp}\{2iϵ_st/\mathrm{}\}\left[S_{FJ}(t)0;F_1+S_{GJ}(t)0;G_1+Y_J(t)e^{i\omega _Jt}\overline{J}_1\right]0;J_2`$
where $`\mathrm{}\omega _J=ϵ_pϵ_s\sqrt{2}V`$ and $`\mathrm{}\omega _K=ϵ_pϵ_s\sqrt{2+\eta ^2}V`$. Thus, we can see that, depending on the state of the electron ”2”, the optical transition of the electron ”1” occurs via $`\overline{K}_1`$ or $`\overline{J}_1`$.
Now, for the controlled NOT gate, we tune the frequency of the light equal to $`\omega _K`$ corresponding to the energy difference between the ground state of Eq. (61) and the first excited state $`\overline{K}_10;K_2`$. Then, the state $`\overline{J}_10;J_2`$ is out of resonance and the optical transition to it is suppressed because its energy is larger than that of the first excited state by $`(\sqrt{2+\eta ^2}\sqrt{2})V`$. Under the resonant approximation, we neglect the optical transition between the ground state and the state $`\overline{J}_10;J_2`$. As a result, only in the case of the electron ”2” at the state $`0;K_2`$, the electron ”1” evolves and its motion is determined by the Hamiltonian $`H_r^C`$ projected from $`H_{total}^C`$,
$`H_r^C`$ $`=`$ $`ϵ_s(d_{0F}^{}d_{0F}+d_{0G}^{}d_{0G})+(ϵ_s+\mathrm{}\omega _K)d_{\overline{K}}^{}d_{\overline{K}}`$ (76)
$`+`$ $`{\displaystyle \frac{\mathrm{}\mathrm{\Omega }_C}{\sqrt{8}}}\{ie^{i\omega _Kt}[d_{0F}^{}d_{\overline{K}}d_{0G}^{}d_{\overline{K}}]+\mathrm{H}.\mathrm{c}.\}`$ (77)
with
$`\mathrm{\Omega }_C={\displaystyle \frac{\sqrt{2}(ϵ_pϵ_s)\xi E_0\mathrm{cos}(\varphi _G)}{\mathrm{}^2\omega _K\sqrt{2+\eta ^2}}}.`$ (78)
Since the above Hamiltonian $`H_r^C`$ is similar to Eq. (19), through the same procedure done in the case of the single-bit gate, we obtain the expansion coefficients of Eq. (71) as a function of a time,
$`\left(\begin{array}{c}S_{FK}(t)\\ S_{GK}(t)\\ Y_K(t)\end{array}\right)=\left(\begin{array}{ccc}\mathrm{cos}^2(\mathrm{\Omega }_Ct/4)& \mathrm{sin}^2(\mathrm{\Omega }_Ct/4)& \frac{1}{\sqrt{2}}\mathrm{sin}(\mathrm{\Omega }_Ct/2)\\ \mathrm{sin}^2(\mathrm{\Omega }_Ct/4)& \mathrm{cos}^2(\mathrm{\Omega }_Ct/4)& \frac{1}{\sqrt{2}}\mathrm{sin}(\mathrm{\Omega }_Ct/2)\\ \frac{1}{\sqrt{2}}\mathrm{sin}(\mathrm{\Omega }_Ct/2)& \frac{1}{\sqrt{2}}\mathrm{sin}(\mathrm{\Omega }_Ct/2)& \mathrm{cos}(\mathrm{\Omega }_Ct/2)\end{array}\right)\left(\begin{array}{c}S_{FK}(0)\\ S_{GK}(0)\\ Y_K(0)\end{array}\right).`$ (88)
Hence, if we choose the duration of the light pulse as,
$`\tau _C=2N\pi /\mathrm{\Omega }_C,`$ (89)
we obtain the wavefunction of Eq. (71) for $`t\tau _C`$ as
$`\mathrm{\Psi }_2(t)`$ $`=`$ $`\left[S_{GK}(0)0;F_1+S_{FK}(0)0;G_1\right]0;K_2`$ (90)
$`+`$ $`\left[S_{FJ}(0)0;F_1+S_{GJ}(0)0;G_1\right]0;J_2`$ (91)
$`=`$ $`\left[S_{GK}(0)0_t+S_{FK}(0)1_t\right]1_c+\left[S_{FJ}(0)0_t+S_{GJ}(0)1_t\right]0_c`$ (92)
in which we use the fact that $`S_{FJ}(t)`$, $`S_{GJ}(t)`$, and $`Y_J(t)`$ are independent of time. Substituting the initial condition of Eq. (73), we can see that the above equation represents exactly the controlled NOT operation of Eq. (1).
## IV Decoherence of quantum states
So far, we assume that the line-width broadening of a state in a quantum dot is zero, i.e. an electronic life time at the state is infinite because the line width is proportional to the scattering rate. In reality, the level of the state is broadened due to various mechanisms such as impurities, structural imperfection, and phonons. Especially, the broadening from phonons is important for the practical quantum gates because it is inherent to the solid-state device while others can be controlled by improved technology. The scattering of an electron from both longitudinal-acoustic(LA) and longitudinal-optic(LO) phonons has been studied extensively for low-dimensional systems such as quantum wells, wire, and dots. In quantum wells, the dominant scattering of the electron results from the LO phonons via the Frőhlich interaction. In a quantum dot, however, this process is forbidden due to the discrete nature of the levels, unless the level separation equals to the LO phonon energy $`\mathrm{}\omega _{LO}`$.
Now, we examine the electron-phonon scattering mainly contributed by LA phonons. The scattering rate is calculated in first-order perturbation theory using the Fermi golden rule,
$`\mathrm{\Gamma }={\displaystyle \frac{2\pi }{\mathrm{}}}{\displaystyle \underset{f,\stackrel{}{q}}{}}M(q)^2\psi _fe^{i\stackrel{}{q}\stackrel{}{r}}\psi _i^2\delta (E_fE_i\pm \mathrm{}\omega _q)[n_B+{\displaystyle \frac{1}{2}}\pm {\displaystyle \frac{1}{2}}]`$ (93)
where the upper(lower) signs account for emission(absorption) of phonons by the transition of an electron from an initial state $`\psi _i`$ to a final state $`\psi _f`$. The sum extends over all possible final quantum states $`\psi _f`$ and phonon wave vector $`\stackrel{}{q}`$. $`n_B`$ stands for the Bose distribution function $`n_B=[e^{\mathrm{}\omega _q/kT}1]^1`$ with $`\omega _q=c_sq`$ and a longitudinal velocity of sound $`c_s`$. For a given deformation potential $`\mathrm{\Xi }`$, the coupling strength of the electron to LA phonons is given by
$`M(q)^2={\displaystyle \frac{\mathrm{\Xi }^2}{2\rho c_s\mathrm{\Omega }_v}}\mathrm{}q`$ (94)
with a mass density $`\rho `$ and a system volume $`\mathrm{\Omega }_v`$. To calculate the quantity $`\psi _fe^{i\stackrel{}{q}\stackrel{}{r}}\psi _i`$, we use wavefunctions as,
$`\stackrel{}{r}0;\alpha `$ $`=`$ $`\left({\displaystyle \frac{1}{\pi ^3\lambda _p^4\lambda _z^2}}\right)^{1/4}e^{(x^2+y^2)/2\lambda _p^2z^2/2\lambda _z^2}`$ (95)
$`\stackrel{}{r}\pm 1;\alpha `$ $`=`$ $`{\displaystyle \frac{x\pm iy}{\lambda _p}}\stackrel{}{r}0;\alpha `$ (96)
which are eigenstates of the quantum dot with a parabolic confinement potential proportional to $`(x^2/\lambda _p^2+y^2/\lambda _p^2+z^2\lambda _p^2/\lambda _z^4)`$. Here, we choose $`\lambda _p\lambda _z`$ to model a disk-like quantum dot.
Since the electron occupies the ground states or the first excited state in the quantum gate during the quantum operation, it is important to examine the scattering from those states. First, in the case of the electron initially at the ground state $`0;\alpha `$, the scattering occurs to the excited state $`n`$ of Eq. (10) by absorbing phonons and its rate is given as
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \underset{n}{}}P_n\mathrm{\Gamma }_0(q_n)`$ (97)
$`\mathrm{\Gamma }_0(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_v}{4\pi \mathrm{}^2c_s}}q^2M(q)^2n_B(q){\displaystyle _1^1}𝑑u\mathrm{exp}\{q^2\lambda _p^2(1u^2)q^2\lambda _z^2u^2\}(1u^2)`$ (98)
where $`q_n=(ϵ_nϵ_0)/\mathrm{}c_s`$ and $`P_n`$ is the probability of finding $`\sigma ;\alpha `$ or $`\pi ;\alpha `$ in the excited state $`n`$. For a GaAs quantum dot with parameters $`\lambda _p=100\AA `$ and $`\lambda _z=20\AA `$, we show $`\mathrm{\Gamma }_0(q)`$ in Fig. 4 with a solid line as a function of an energy difference $`E=ϵ_nϵ_0`$. We find that the scattering rate strongly depends on the energy difference $`E`$. Especially, for a high energy $`E2meV`$, the scattering is estimated to be very rare. This means that, if the energy difference between the first excited and the ground states in the quantum gates is sufficiently large and not close to $`\mathrm{}\omega _{LO}`$, the electron at the ground state is rarely scattered to the excited states and has a long coherent time in the quantum gates.
Through a similar argument, we can see that the electron initially at the first excited states is not relaxed to the ground state by emitting phonons, but frequently scattered to adjacent excited states by absorbing phonons because their energy difference is relatively small. For example, the electron at the first excited state in the single-bit gate is scattered to the second one with the rate of,
$`\mathrm{\Gamma }_2={\displaystyle \frac{\mathrm{\Omega }_vq^6\lambda _p^4}{8\pi \mathrm{}}}M(q)^2n_B(q){\displaystyle _1^1}𝑑u\mathrm{exp}\{q^2\lambda _p^2(1u^2)q^2\lambda _z^2u^2\}(1u^2)^2`$ (99)
where $`q=(ϵ_3ϵ_2)/\mathrm{}c_s`$. Plotting the result as a function of the energy difference $`E=ϵ_3ϵ_2`$ as shown in Fig. 4, we can see that the electron at the first excited state is scattered frequently over a more wide range of $`E`$ than that at the ground states. If the coupling strength $`V`$ is $`1meV`$ or $`E=\sqrt{2}meV`$, the electron of the first excited state is scattered with the rate of $`10^{11}/sec`$. Therefore, for the single-bit gate to work well, the duration $`\tau `$ of a light pulse should be smaller than the inverse of the scattering rate $`\mathrm{\Gamma }_2`$, i.e.,
$`{\displaystyle \frac{2N\pi }{\mathrm{\Omega }\mathrm{cos}(\varphi _B)/\sqrt{2}}}{\displaystyle \frac{1}{\mathrm{\Gamma }}}\mathrm{or}\mathrm{\Omega }N\mathrm{\Gamma }.`$ (100)
For a GaAs dot, $`\mathrm{}\mathrm{\Omega }`$ should have a larger value than $`0.4meV`$ at 300 K for $`\mathrm{\Gamma }=10^{11}/sec`$ or the field strength of the light $`E_0400V/cm`$ for a 10nm disk-like dot from Eq. (38).
In summary, we present novel models of quantum gates based on coupled quantum dots or artificial molecules. By varying the size and the location of each quantum dot in the artificial molecule, we assume that well-localized ground states are present while excited states form molecular orbits with a particular geometrical symmetry. First, for the single-bit gate, we locate a smaller quantum dot between two larger quantum dots as shown in Fig. 1. Since two larger quantum dots in Fig. 1 are separated enough, their ground states are well-localized and we define a qubit as the superposition of them. To drive the qubit, we exploit the optical transition between the ground states and the first excited state of the artificial molecule. Since the wavefunction of the first excited state extends over all three quantum dots arranged with a ”V” shape, we show that, by solving the time-dependent Schődinger equation, the light pulse rotates the qubit coherently depending on its frequency and polarization to demonstrate the single-bit operation. Secondly, for the controlled NOT operation, we examine the artificial molecule shown in Fig. 2 where two larger quantum dots containing a control bit are added to the single-bit gate with a target bit. Under the illumination of the light pulse, we show that the electron representing the target bit evolves conditionally, i.e. depending on the state of the control bit because of the asymmetrical location of two added quantum dots. Furthermore, when the electron-electron interaction is short-ranged, we find that the light pulse with the resonant frequency between the ground and the first excited states severs as the controlled NOT operation. Finally, to examine the decoherence of quantum states, we discuss electronic relaxation contributed mainly by LA phonon processes. By calculating the scattering rate using the Fermi-golden rule, we estimate the duration and the field strength of the light pulse.
###### Acknowledgements.
This work was supported by the Korean Ministry of Science and Technology through the Creative Research Initiatives Program under Contract No. 98-CR-01-01-A-20.
Figure Captions
Fig. 1. An artificial molecule consisted of three disk-like quantum dots on the substrate is shown for the single-bit gate in (a). The ground state of each larger quantum dot is assumed to be well-localized and is used as the basis of a qubit. However, excited states such as $`\sigma ;A,B`$ \[solid line in (b)\] extend their wavefunctions over the dot $`C`$ and form molecular states.
Fig. 2. We show the arrangement of the quantum dots for the controlled NOT gate. Here, a target(control) bit is regarded as the superposition of the grounds states at two larger quantum dots located at the upper(lower) side and their interaction is mediated by the quantum dot $`I`$.
Fig. 3. We show a schematic structure of two-particle energy levels. The conjugate state $`\beta _1\alpha _2`$ to the state $`\alpha _1\beta _2`$ is also possible, however, it is not shown here for simplicity.
Fig. 4. Log plot of the phonon scattering rate is shown as a function of an energy difference $`E`$ between two relevant states. $`\mathrm{\Gamma }_0`$ ($`\mathrm{\Gamma }_2`$) with a solid (dotted) line describes the scattering rate when an electron is initially at the ground (first excited) state of the single-bit gate by absorbing phonons. Used parameters are $`\mathrm{\Xi }=6.8eV`$, $`\rho =5.36g/cm^3`$, and $`c_s=5150m/sec`$ for GaAs.
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# Low-energy Limits on the Antisymmetric Tensor Field Background on the Brane and on the Non-commutative Scale
## I Introduction
One of the initial motivations for field theories on non-commutative spaces was their intrinsically more convergent behavior in the ultraviolet regime than the one observed for ordinary field theories.
Field theories on noncommutative spaces (NCFT) can be defined as theories in their own right, independent of string theory. The coordinates in these spaces are represented by self-adjoint operators acting on some Hilbert space $``$ and satisfying the following commutation relations:
$$[\widehat{x}^\mu ,\widehat{x}^\nu ]=i\theta ^{\mu \nu }[\theta ^{\mu \nu },x^\rho ]=0$$
(1)
Consequently, fields on such a space are replaced by operators. To each such operator one can associate an ordinary field on a commutative space as follows:
$$\varphi (x)=\frac{1}{(2\pi )^{d/2}}d^dk\mathrm{e}^{ik_\mu x^\mu }\text{Tr}[\widehat{\varphi }(\widehat{x})\mathrm{e}^{ik_\mu \widehat{x}^\mu }]$$
(2)
where the trace is taken in the Hilbert space $``$. By $`\varphi `$ here we denote a generic field; we can associate to it some space-time indices as it will be the case for gauge fields or fermions.
An action defined on $``$ must be naturally writable in terms of traces; furthermore, we want that in the limit $`\theta ^{\mu \nu }0`$ the expression reduces to an ordinary action on an ordinary, commutative space. Then, the generic form of the action is
$$S=\text{Tr}[(\theta _{\mu \nu }^1[\widehat{x}^\nu ,\widehat{\varphi }(\widehat{x})])^2+P(\widehat{\varphi })]$$
(3)
where $`P`$ is some polynomial in $`\widehat{\varphi }(\widehat{x})`$. It is not difficult to see that, using (2), the commutation relations for $`\widehat{x}`$ and the Baker-Campbell-Hausdorff formula, this action reduces to:
$$S=d^Dx(_\mu \varphi (x))^2+P_{}(\varphi )$$
(4)
where by $`P_{}(\varphi )`$ we mean that in $`P(\widehat{\varphi })`$ we replace $`\widehat{\varphi }`$ by $`\varphi `$ and the product of fields is the Moyal product, given by:
$$\varphi _1\varphi _2(x)=e^{i\frac{1}{2}\theta ^{\mu \nu }\frac{}{\xi ^\mu }\frac{}{\zeta ^\nu }}\varphi _1(x+\xi )\varphi _2(x+\zeta )|_{\xi =\zeta =0}$$
(5)
Noncommutative field theories became popular among string theorists with the work of Connes, Douglas and Schwarz who argued that M-theory in constant 3-form background is equivalent to the supersymmetric Yang-Mills theory, defined on a non-commutative torus. A second wave of interest was generated by the paper of Seiberg and Witten which summarized and extended earlier ideas about the appearance of noncommutative geometry in string theory with constant NS-NS $`B`$ field background.
NCFT are constructed starting from string theory in much the same way as the usual field theories. In particular, one computes the string theory S-matrix elements and writes a low energy effective action that reproduces them at the tree level. Furthermore, the only difference from the usual computation is that the world sheet propagator is modified by the presence of $`B`$. If the world sheet has no boundaries, then a constant $`B`$ field can be gauged away. Thus, a constant B flux manifests itself only in the presence of world sheet boundaries, i.e. in the presence of D-branes. Moreover, the same argument shows that only the components of $`B`$ parallel to the D-brane can be non-zero. The world sheet propagator restricted to the boundary is modified by the addition of a term $`\frac{i}{2}\theta ^{\mu \nu }ϵ(\tau \tau ^{})`$, where $`ϵ(\tau \tau ^{})`$ is the step function and $`\theta ^{\mu \nu }=(2\pi \alpha ^{})^2\left(\frac{1}{g+2\pi \alpha ^{}B}B\frac{1}{g2\pi \alpha ^{}B}\right)^{\mu \nu }`$. Since all vertex operators contain factors of the type $`exp(ikx)`$, it is easy to see that all correlation functions will get the extra factor
$$e^{i\theta ^{\mu \nu }\underset{i<j=1}{\overset{n}{}}k_\mu ^ik_\nu ^j}$$
(6)
which is just the $``$-product defined in (5) written in momentum space. Thus, the effective action in the presence of B-field has the interpretation of a field theory on a noncommutative space with noncommutativity parameter given by $`\theta `$.
Transition from the open string theory to the non-commutative field theory becomes explicit in a zero slope limit, $`\alpha ^{}ϵ^{1/2}0`$, $`g_{\mu \nu }ϵ0`$. For example, in this limit, the dynamics of Yang-Mills fields living on the brane is governed by the following action:
$$S=\frac{1}{2g_{YM}^2}Tr_{U(N)}F_{\mu \nu }^{nc}F^{nc,\mu \nu }$$
(7)
where we have taken the open string metric to be $`\eta _{\mu \nu }`$ and $`F_{\mu \nu }^{nc}`$ is given by:
$$F_{\mu \nu }^{nc}=_\mu A_\nu ^{nc}_\nu A_\mu ^{nc}+iA_\mu ^{nc}A_\nu ^{nc}iA_\nu ^{nc}A_\mu ^{nc}.$$
(8)
Comparing this with equation (4) we see that this can be interpreted as the Yang-Mills action on a noncommutative space with noncommutativity parameter $`\theta `$. Up to now the structure of $`B_{\mu \nu }`$ in the directions parallel to the brane was not really important. It should be noted, however, that large electric-like background $`B_{0i}`$ creates various problems in the zero slope limit , . A solution to these problems was proposed in and it leads to (7).
It became a colloquial wisdom that the parameter $`\theta _{\mu \nu }`$ does not necessarily have to be of the order of the inverse Plank scale squared . For example, $`\theta _{\mu \nu }`$ could be significantly larger in the “brane-world” proposal , in which the string scale $`M_s`$ is much smaller than four-dimensional Plank scale due to the large volume of extra dimensions, $`M_{\mathrm{Pl}}^2=M_s^{n+2}V_n`$. In this case a “natural” scale for $`\theta _{\mu \nu }`$ could be $`M_s^2`$. Another example is the open string realization of the non-commutative field theories on the brane in the zero slope limit described above, which presumes the non-commutative scale to be fixed, while gravity is decoupled (i.e. $`M_{\mathrm{Pl}}\mathrm{}`$). If instead we choose to fix the gravitational scale, zero slope limit would correspond to a large non-commutative parameter in units of $`M_{\mathrm{Pl}}^2`$.
All these cases pose one interesting phenomenological question as to how large $`\theta _{\mu \nu }`$ could be without contradicting existing experimental data. To answer this question we have to adopt certain calculational framework, incorporating together Standard Model fields and $`\theta _{\mu \nu }`$ background. The simplest way of doing this is to assume that SM is realized on the brane and the external background of $`B_{\mu \nu }`$ (or $`\theta _{\mu \nu }`$) field is included via Moyal product. In this paper we take space and time independent background which could be a bad approximation at large distances. Indeed, due to the interaction on the brane, $`B_{\mu \nu }(x)`$ could become a massive field so that the minimum of energy correspond to the vanishing vev. In this case the constraints on $`\theta _{\mu \nu }`$ that we are aiming to produce will be trivially satisfied since there is no $`\theta `$ to begin with and the theory is the usual one. Despite this possibility, we believe that the question of experimental constraints on $`B_{\mu \nu }`$ background/noncommutativity parameter $`\theta _{\mu \nu }`$ deserves special investigation.
Our strategy for the rest of this paper is quite straightforward. We take the (\*)-modified Standard Model and expand it once in the external $`\theta `$ parameter,
$$fg=fg+\frac{i}{2}\theta _{\mu \nu }_\mu f_\nu g,$$
(9)
so that SM is extended by the series of dimension 6 operators, composed from three or more fields:
$$SM()=SM+\underset{i}{}\theta _{\mu \nu }O_{\mu \nu }^{(i)}.$$
(10)
Here the summation is performed over different types of operators, full list of which is outside the scope of the present paper. First order in $`\theta `$ is sufficient for our purposes. Such a procedure of constructing an effective action does not depend on taking the zero slope limit and holds for a generic situation .
It is also to our advantage that in the effective Lagrangian approach we do not have to worry about the renormalizability of this theory. Although the ultimate resolution to the question of renormalizability of non-commutative field theories is very interesting and very important , here we can simply assume that all loop divergences are regularized at momenta comparable to $`M_s`$. Moreover, since the natural scale for $`\theta `$ is $`M_s^2`$, the iteration of $`\theta _{\mu \nu }O_{\mu \nu }`$-interactions does not create any problems as long as external momenta are much smaller than $`\sqrt{1/\theta }`$.
## II Experimental limits on $`\theta _{\mu \nu }`$
High-energy limits
At first glance, it is advantageous to use high-energy processes to obtain the most stringent limits on $`\theta `$. Indeed, the higher the energy/momentum transfer is, the larger the effect of $`O_{\mu \nu }`$ will be. Let us consider such a well-studied process as the Z-boson decay. We shall profit from the fact that $`O_{\mu \nu }`$ violate Lorentz invariance and use a decay channel which is strictly forbidden at the SM level and allowed when $`\theta _{\mu \nu }0`$. A good candidate for this channel will be the decay of Z into a pair of photons, forbidden by Lorentz invariance and Bose statistics in SM. The experimental limit on the branching ratio for this decay is also very good, $`Br(Z\gamma \gamma )510^5`$ .
To calculate this decay width we need to evaluate the $`\theta _{\mu \nu }O_{\mu \nu }`$-expansion of the $`U(1)`$ gauge sector, which in the presence of external $`\theta _{\mu \nu }`$-background could be taken in the following form:
$$S_{nc}=\frac{1}{4}d^4xF_{\mu \nu }^{nc}F^{nc,\mu \nu }$$
(11)
Here $`F_{nc}`$ denotes “non-commutative” field strength given by
$$F_{\mu \nu }^{nc}=_\mu B_\nu _\nu B_\mu g^{}\theta ^{\alpha \beta }_\alpha B_\mu _\beta B_\nu .$$
(12)
Expanding Eq. (11) to first order in $`\theta _{\mu \nu }`$ we get
$$S_{nc}=S+\frac{1}{2}g^{}d^4x\theta ^{\alpha \beta }(_\mu B_\nu _\nu B_\mu )_\alpha B^\mu _\beta B^\nu .$$
(13)
Going to the physical basis, we obtain the $`Z\gamma \gamma `$ interaction term:
$$S_{int}=\frac{1}{2}g^{^{}}\mathrm{sin}\theta _W\mathrm{cos}^2\theta _Wd^4x\theta _{\rho \sigma }[_\mu Z_\nu _\rho A_\mu _\sigma A_\nu +(_\mu A_\nu _\nu A_\mu )_\rho Z_\mu _\sigma A_\nu ].$$
(14)
It should be noted here that the expansion of the $`SU(2)`$ sector has quadratic terms in $`\theta _{\mu \nu }`$, but not linear ones and thus does not contribute into (14).
Finally, we arrive at the following answer for the two-photon decay width of the Z-boson, induced by the $`\theta _{\mu \nu }`$-background:
$$\mathrm{\Gamma }_{Z\gamma \gamma }=\frac{\alpha }{144}\mathrm{cos}^4\theta _WM_Z^5\underset{i}{}\theta _{0i}^2.$$
(15)
The decay width is evaluated in the Lorentz system in which Z is produced at rest.
Comparing this result with the experimental limits on the branching ratio, we conclude that the sensitivity to $`1/\sqrt{\theta }`$ is not better than 250 GeV. This is a very modest limit, which could perhaps be improved had we considered $`\theta `$-induced corrections to other high-energy processes such as $`e^+e^{}`$ cross sections, forward-backward asymmetry and so on. At any rate, a significant improvement beyond the level of 250 GeV is not possible. Thus, we conclude that the high-energy processes cannot produce sufficiently strong bounds on $`\theta _{\mu \nu }`$, and turn to the low-energy limits on this parameter.
Low-energy limits
The most notable feature of the effective Lagrangian (10) is the explicit violation of Lorentz invariance . This violation is, of course, controlled by the size of $`\theta _{\mu \nu }`$ and could be made arbitrarily small “by hand”.
The extension of Standard Model by Lorentz-noninvariant operators of dimension $``$ 4, has been actively studied in the past . Some of the limits obtained in are extremely strong. A priori, there are several problems with such a generic description. If we believe that the Lorentz non-invariant terms originate from short distances, it is not clear why dimension 3 and 4 operators should be suppressed. Another problem is of rather technical nature, as the number of free parameters in such extensions of SM is generally very large. These problems do not exist in our approach. The breaking of Lorentz invariance is given only by $`\theta _{\mu \nu }`$, and its effect first shows up in dimension 6 operators.
The qualitative understanding of the role of $`\theta _{\mu \nu }`$ for low-energy physics comes from the considerations of the non-relativistic limit for the $`d^4xe\overline{\psi }A_\mu \gamma ^\mu \psi `$ interaction term taken in the external Coulomb field:
$$V=\frac{Z\alpha }{r}\frac{Z\alpha }{2r^3}(\theta _B𝐋)\frac{m_eZ\alpha }{r^3}(\theta _E𝐫).$$
(16)
Here we split $`\theta _{\mu \nu }`$ into two three-vectors in analogy with the electromagnetic field, $`ϵ_{ijk}(\theta _B)_k\theta _{ij}`$, $`(\theta _E)_i\theta _{0i}`$. The diagonal matrix element of the last operator in Eq. (16), taken over a wave function of the discrete spectrum, is equal to zero, unless CP is broken. The second term, however, may produce interesting phenomena, as it gives an effective coupling of the angular momentum with an external vector $`\theta _B`$,
$$V=\kappa (\theta _B𝐉).$$
(17)
The effective coupling constant $`\kappa `$ has obvious $`L,J`$ dependence and, more importantly, is determined by the third power of the characteristic atomic momentum. For an outer atomic electron this coupling is of the order of $`Z^2\alpha (m_e\alpha )^3`$. The coupling to $`\theta _B`$ creates a Zeeman-like splitting of atomic orbitals with respect to an external vector and thus can be tested in high-precision atomic experiments.
At this point it becomes clear that similar effects ($`\theta _B`$ angular momentum) will appear in the hadronic physics, when we consider the (\*)-extended interaction of quarks and gluons. This interaction will lead to the effective coupling of the nucleon spin with $`\theta `$,
$$N|_{QCD}()|N=\theta \mathrm{independent}\mathrm{terms}+\frac{d_\theta }{2}\theta _{\mu \nu }\overline{N}\sigma _{\mu \nu }N,$$
(18)
which in non-relativistic limit simply becomes $`d_\theta (\theta _B\frac{𝐒}{S})`$. The size of $`d_\theta `$ is given by a cube of a characteristic hadronic scale.
To estimate $`d_\theta `$ and to convince ourselves that such an effect exists, we perform the following exercise, a simplified version of the nucleon three-point function QCD sum rules . We calculate the operator product expansion (OPE) of two nucleon currents in the presence of the external $`\theta _{\mu \nu }`$-background and observe the non-zero result. Then we take the ”phenomenological” part of the QCD sum rule and saturate it with the nucleon double-pole contribution. Matching the two sides at 1 GeV, we obtain an estimate of $`d_\theta `$ for nucleons.
On the OPE side we can use an asymptotically free description, and thus include the $`\theta _{\mu \nu }`$-piece as the correction to a free massless quark propagator:
$$S(x,0)=\frac{ix_\mu \gamma _\mu }{2\pi ^2x^4}\frac{ix_\mu \gamma _\beta }{2\pi ^2x^4}t^aG_{\nu \beta }^a\theta _{\mu \nu }.$$
(19)
This correction originates from the (\*)-extended quark-gluon interaction, whereas pure gluonic sector does not contain terms linear in theta. Using this propagator it is straightforward to calculate the OPE of two nucleon currents, expressing the result as the combination of a pure perturbative piece and non-perturbative condensates such as $`\overline{q}\sigma _{\mu \nu }q_{\theta _{\mu \nu }}`$, $`\overline{q}\sigma _{\mu \nu }t^aG_{\mu \nu }^aq`$, and so on. Deferring further details for more extended publication, we present here the final estimate for $`d_\theta :`$
$$d_\theta 0.1\mathrm{GeV}^3$$
(20)
Further progress in refining this estimate can be achieved by including the anomalous dimensions, and fixing the size of ”susceptibility” condensate $`\overline{q}\sigma _{\mu \nu }q_{\theta _{\mu \nu }}`$ by comparing the stability of different sum rule channels.
As we remarked earlier, the coupling of the nucleon spin to an external vector will lead to the Zeeman-like splitting of the hyperfine structure in atoms. Since the precision in measuring hyperfine splittings is at the level of mHz, this should lead to a strong bound on $`\theta _B.`$ One of the most sensitive systems in this respect is the hydrogen maser, where the effects of Lorentz violating terms have been searched for in a recent experiment . In this double-resonance experiment the absence of sidereal variations of the maser frequency leads to 1 mHz limit on the product of $`d_\theta `$ and $`\theta _B`$, component of $`\theta _B`$ perpendicular to earth’s axis. This puts the limit on the magnetic component of $`\theta _{\mu \nu }`$ at the level $`1/\sqrt{\theta }>510^{12}`$ GeV.
Another even more stringent limit could be extracted from the experiment which compares magnetic field measured by Cs and Hg atoms . The non-vanishing $`\theta _{\mu \nu }`$-background affects primarily nuclear spin. Thus the magnetic field, measured by mercury atom, will be corrected due to the interaction of the nuclear spin with $`\theta _B`$ to a larger extent than the magnetic field measured by Cs. The absence of sidereal variations in the difference of magnetic field measured by Cs and Hg is verified at 100 nHz level. This translates to the following, extremely tight bound on $`\theta `$:
$$\frac{1}{\sqrt{\theta }}>510^{14}\mathrm{GeV}$$
(21)
This is the main result of the present paper.
## III Discussion
The limit on $`\theta _{\mu \nu }`$ obtained in this paper is quite strong, signaling the sensitivity of low-energy experiments to the scales comparable to standard GUT/string scales. Of course, this sensitivity is the consequence of the assumption about non-vanishing $`\theta _{\mu \nu }`$-dependent background, which breaks explicitly Lorentz invariance.
Limit (21) is derived from the atomic experiments, performed in a laboratory. It would be interesting to explore whether similar or stronger bounds could be obtained from non-observation of the anisotropy of the Universe at large scales. Indeed, the coupling of baryon spins to $`\theta _{\mu \nu }`$ at certain level should lead to the anisotropy of matter distribution and polarization of interstellar medium which could, in principle, result in strong bounds on $`\theta _{\mu \nu }`$ .
Returning to the brane-world scenario it is fair to question a “natural” value for $`\theta _{\mu \nu }`$ on the brane. For the energies much higher than the inverse radius of extra dimensions, a natural “string” value of $`\theta _{\mu \nu }`$ is given by the square of the inverse string scale. <sup>§</sup><sup>§</sup>§ Throughout our discussion we do not assume zero slope limit, so that $`M_s`$ is the only natural dimensional high-energy parameter. For $`M_s`$ 1 TeV this is in apparent contradiction with the limit (21).
Does our result (21) pose another “naturalness” problem for the low-energy string scale models? The answer to this question depends on the dynamical properties of $`\theta _{\mu \nu }(x,y)`$, where $`y`$ indicates the dependence on extra space-like coordinates. If the lowest mode of the Kaluza-Klein expansion of $`\theta _{\mu \nu }(x,y)`$ remains essentially massless, as it would happen in theories with extra dimensions of infinite volume , then the vev of $`\theta _{\mu \nu }(x,y)`$ may freeze at a large value, which would be in contradiction to the observational evidence for absence of $`\theta _{\mu \nu }`$ at the level (21). When the volume of extra dimensions is finite, it may turn out that the ultimate value of $`\theta _{\mu \nu }(x)`$ is zero. Indeed, every Kaluza-Klein copy of this tensor field may receive a mass term due to the interaction with brane fields and/or details of mechanisms responsible for compactification. The mass term induced by the interaction $`\theta _{\mu \nu }O_{\mu \nu }`$ on the brane could be estimated when the volume of extra dimensions is large. If we assume that the correlator of two $`O_{\mu \nu }`$ currents is saturated at the string scale $`M_s`$,
$$O(x),O(0)d^4x(\mathrm{two}\mathrm{loop}\mathrm{factor})\times M_s^4,$$
(22)
and the two-loop phase space factor originates from the fact that by construction every $`O_{\mu \nu }`$ has at least three fields. The kinetic term for every Kaluza-Klein mode of $`\theta _{\mu \nu }(x,y)`$ receives a volume enhancement factor, so that the effective mass for the lowest mode, canonically normalized in four dimensions is
$$m_{eff}^2(\mathrm{two}\mathrm{loop}\mathrm{factor})\times \frac{M_s^2}{M_s^nV_n}10^4\times \frac{M_s^4}{M_{\mathrm{Pl}}^2}(10^5\mathrm{eV})^2\frac{M_s^2}{1\mathrm{T}\mathrm{e}\mathrm{V}^2}$$
(23)
It is clear that for a TeV range $`M_s^2`$, $`m_{eff}`$ is very low, which could create cosmological problems, similar to those arising in the standard axion relaxation mechanism when the axion mass is small. In this case the “$`\theta _{\mu \nu }`$-problem” is replaced by a new moduli problem which indicates that the choice $`M_s1`$ TeV is unnatural. However, the ultimate answer to the question about the mass of $`\theta _{\mu \nu }(x)`$ and its vev, surviving until present cosmological times, needs other physical inputs such as information about the compactification mechanisms and supersymmetry breaking.
Acknowledgments M.P. thanks A. Lossev for very helpful discussions and interest taken in this work. I.M. and R.R. thank I. Chepelev for discussions. This work was supported in part by the Department of Energy under Grant No. DE-FG02-94ER40823 and NSF grant PHY-9722101.
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# Asteroids in the Inner Solar System I – Existence
## 1 Introduction
Lagrange’s (1772) triangular solution of the three body problem was long thought to be just an elegant mathematical curiosity. The three bodies occupy the vertices of an equilateral triangle. Any two of the bodies trace out elliptical paths with the same eccentricity about the third body as a focus (see e.g., Whittaker 1904, Pars 1965). The detection of 588 Achilles near Jupiter’s Lagrange point in 1906 by Wolf changed matters. This object librates about the Sun-Jupiter $`L_4`$ Lagrange point, which is $`60^{}`$ ahead of the mean orbital longitude of Jupiter (e.g., Érdi 1997). About 470 Jovian Trojans are now known (see “http://cfa-www.harvard.edu/iau/lists/Trojans.html”), though the total population exceeding 15 km in diameter may be as high as $`2500`$ (Shoemaker, Shoemaker & Wolfe 1989; French et al. 1989). Roughly $`80\%`$ of the known Trojans are in the $`L_4`$ swarm. The remaining $`20\%`$ librate about the $`L_5`$ Lagrange point, which trails $`60^{}`$ behind the mean orbital longitude of Jupiter. There are also Trojan configurations amongst the Saturnian moons. The Pioneer 11 and the Voyager 1 and 2 flybys of Saturn discovered five small moons on tadpole or horseshoe orbits. The small moon Helene librates about the Saturn-Dione $`L_4`$ point. The large moon Tethys has two smaller Trojan moons called Telesto and Calypso, one of which librates about the Saturn-Tethys $`L_4`$ and the other about the Saturn-Tethys $`L_5`$ points. Finally, the two small moons Janus and Epimetheus follow horseshoe orbits coorbiting with Saturn (see e.g., Smith et al. 1983; Yoder et al. 1983; Yoder, Synnott & Salo 1989). Another example of a Trojan configuration closer to home is provided by the extensive dust clouds in the neighbourhood of the $`L_5`$ point of the Earth-Moon system claimed by Winiarski (1989).
This paper is concerned with the existence of coorbiting asteroids near the triangular Lagrange points of the four terrestrial planets. For non-Jovian Trojans, the disturbing forces due to the other planets are typically larger than those caused by the primary planet itself. For this reason, it was formerly considered unlikely that long-lived Trojans of the terrestrial planets could survive. Over the past decade, two lines of evidence have suggested that this reasoning is incorrect. The first is the direct discovery of inclined asteroids librating about the Sun-Mars $`L_5`$ Lagrange point. The second is numerical integrations, which have steadily increased in duration and sophistication.
The first non-Jovian Trojan asteroid, 5261 Eureka, was discovered by Holt & Levy (1990) near the $`L_5`$ point of Mars. Surprisingly, the orbit of 5261 Eureka is inclined to the plane of the ecliptic by $`20.3^{}`$. The determination of Eureka’s orbital elements and a preliminary analysis of its orbit were published in Mikkola et al. (1994). Numerical integrations were performed for several dozen Trojan test particles with different initial inclinations for $`4`$ Myr by Mikkola & Innanen (1994), who claimed that long-term stability of Martian Trojans was possible only in well-defined inclination windows, namely $`15^{}i30^{}`$ and $`32^{}i44^{}`$ with respect to Jupiter’s orbit. The discovery of a second Mars Trojan, 1998 VF31, soon followed (see e.g., Minor Planet Circular 33763; Tabachnik & Evans 1999). Perhaps the most important point to take from the observational discoveries is that if a comparatively puny body like Mars possesses Trojans, it is quite likely indeed that the more massive planets also harbour such satellites.
The main argument for the possible existence of Trojans of the Earth, Venus and Mercury comes from numerical test particle surveys. Much of the credit for reviving modern interest in the problem belongs to Zhang & Innanen (1988a,b,c). Using the framework of the planar elliptic restricted four and five-body problem, their integrations extended from 2000 to $`10^5`$ years, though their model was not entirely self-consistent as mutual interactions between the planets were not taken into account. Further important results on the stability of the Trojans of the terrestrial planets were obtained in a series of papers by Mikkola & Innanen (1990, 1992, 1994, 1995). The orbits of the planets from Venus to Jupiter were computed at first using a Bulirsch-Stoer integration and later with a Wisdom-Holman (1991) symplectic integrator. In the case of Mercury, test particles placed at the Lagrange point exhibited a strongly unstable behaviour rather rapidly. Conversely, the mean librational motion of all the Trojan test particles near Venus and the Earth’s Lagrange points appeared extremely stable. Thus far, the longest available integrations are the 6 Myr survey of test particles near the triangular $`L_4`$ point of Venus and the Earth (Mikkola & Innanen 1995). The initial inclinations ranged from $`0^{}`$ to $`40^{}`$ with respect to the orbital plane of the primary planet. The Trojans of Venus and the Earth persisted on stable orbits for small inclinations ($`i18^{}`$ and $`i11^{}`$ respectively).
These lines of reasoning suggest that a complete survey of the Lagrange points of the terrestrial planets is warranted. It is the purpose of the present paper to map out the zones in which coorbital asteroids of the terrestrial planets can survive for timescales up to 100 million years. Section 2 provides a theoretical introduction to the dynamics of coorbital satellites within the framework of the elliptic restricted three body problem. Our fully numerical survey is discussed in Section 3 and takes into account the effects of all the planets (except Pluto), as well as the most important post-Newtonian corrections and the quadrupole moment of the Moon. The results of the survey are presented for each of the terrestrial planets in Sections 4 to 7 (Mercury, Venus, the Earth, and Mars). Finally, a companion paper in this issue of Monthly Notices discusses the observable properties of the asteroids.
## 2 Theoretical Framework
Consider the restricted three-body problem, defined by a system of two massive bodies on Keplerian orbits attracting a massless test particle that does not perturb them in return. The five Lagrange points are the stationary points of the effective potential (e.g., Danby 1988). The collinear Lagrange points $`L_1`$, $`L_2`$ and $`L_3`$ are unstable, whereas the triangular Lagrange points $`L_4`$ and $`L_5`$ form equilateral triangles with the two massive bodies, as illustrated in Fig. 1. Motion around $`L_4`$ and $`L_5`$ can be stable. The planar orbits that corotate with a planet are classified as either tadpoles or horseshoes. Tadpole orbits perform a simple libration about either the Lagrange point $`L_4`$ or $`L_5`$, whereas the horseshoes perform a double libration about both $`L_4`$ and $`L_5`$. The tadpole and horseshoe orbits are divided by a separatrix, which corresponds to an orbit of infinite period that passes through $`L_3`$. On average, both tadpoles and horseshoes move about the Sun at the same rate as the parent planet.
It is helpful to develop an approximate treatment of coorbital motion. Within the framework of the elliptic restricted problem of three bodies (Sun, planet and asteroid), the Hamiltonian of the asteroid in a non-rotating heliocentric coordinate system reads:
$$H=\frac{k^2}{2a}m_\mathrm{p}k^2R,$$
(1)
where $`m_\mathrm{p}`$ is the mass of the planet in Solar masses and $`k`$ is the Gaussian gravitational constant. The disturbing function $`R`$ is defined by (e.g., Danby 1988)
$`R`$ $`=`$ $`{\displaystyle \frac{1}{(r^2+r_{\mathrm{p}}^{}{}_{}{}^{2}2rr_\mathrm{p}\mathrm{cos}S)^{1/2}}}{\displaystyle \frac{r\mathrm{cos}S}{r_{\mathrm{p}}^{}{}_{}{}^{2}}}.`$ (2)
Here, the radius vector $`r`$ and $`r_\mathrm{p}`$ refer to the asteroid and the planet respectively, $`S`$ being the elongation of the asteroid from the parent planet. Following Brouwer & Clemence (1961), $`\mathrm{cos}S`$ can be expressed in function of the mutual inclination of the two orbits $`i`$, the true anomalies $`\nu `$ and $`\nu _\mathrm{p}`$, the argument of perihelion $`\omega `$ and the longitude of the ascending node $`\mathrm{\Omega }`$ of the asteroid:
$`\mathrm{cos}S`$ $`=`$ $`\mathrm{cos}^2\frac{i}{2}\mathrm{cos}(\nu \nu _\mathrm{p}+\omega +\mathrm{\Omega })`$ (3)
$`+`$ $`\mathrm{sin}^2\frac{i}{2}\mathrm{cos}(\nu +\nu _\mathrm{p}+\omega \mathrm{\Omega }).`$
The disturbing function is then expressed in terms of the mean synodic longitudes $`\lambda _\mathrm{p}=M_\mathrm{p}`$ and $`\lambda =M+\omega +\mathrm{\Omega }`$. Here, $`M_\mathrm{p}`$ and $`M`$ are the mean anomalies of the planet and asteroid. In other words, we are using a coordinate system for which $`(x,y)`$ lies in the orbital plane of Mars and the $`x`$-axis points towards Mars’ perihelion. The disturbing function is expanded to second order in the eccentricities and to fourth order in the inclination. Finally, we make a change of variables $`\varphi =\lambda _\mathrm{p}\lambda `$ and average $`R`$ over the mean anomaly of the planet $`M_\mathrm{p}`$:
$$R=\frac{1}{2\pi }_0^{2\pi }R𝑑M_\mathrm{p}=U_0+U_2+O(e^3,e_\mathrm{p}^3,i^5).$$
(4)
The zeroth order term in the averaged disturbing function is
$$U_0=\frac{1}{(a^2+a_\mathrm{p}^22aa_\mathrm{p}\mathrm{cos}^2\frac{i}{2}\mathrm{cos}\varphi )^{1/2}}\frac{a}{a_\mathrm{p}^2}\mathrm{cos}^2\frac{i}{2}\mathrm{cos}\varphi .$$
(5)
The term that is second order in the eccentricities and inclinations has three parts
$`U_2=U_{2,0}+{\displaystyle \frac{U_{2,3}}{(a^2+a_\mathrm{p}^22aa_\mathrm{p}\mathrm{cos}^2\frac{i}{2}\mathrm{cos}\varphi )^{3/2}}}`$
$`+{\displaystyle \frac{U_{2,5}}{(a^2+a_\mathrm{p}^22aa_\mathrm{p}\mathrm{cos}^2\frac{i}{2}\mathrm{cos}\varphi )^{5/2}}},`$ (6)
where $`U_{2,0},U_{2,3}`$ and $`U_{2,5}`$ are given in Appendix A. Still lengthier expressions accurate to the fourth order in the eccentricities and inclinations are given in Tabachnik (1999). It is interesting to note that the second order expressions depend only on the longitude of perihelion, while the fourth order contains terms depending on both the longitude of perihelion and the longitude of the ascending node. Some of these averaged disturbing functions are shown in Fig. 2 and compared to the numerically computed and averaged disturbing function $`R_{\mathrm{Num}}`$.
For small and moderate inclinations, $`R`$ has a double-welled structure. The libration about the Lagrange points can be viewed as the trapping of the test particle in the well, with the tadpole orbits located in the deeper part of the well and the horseshoe orbits in the shallower part. The two orbital families are divided by a separatrix. The slight eccentricities and inclinations of the planets cause small deviations of the Lagrange points from their classical values (c.f. Namouni & Murray 1999). These deviations can be calculated by finding the local minima of the secular potential (4). The results for each terrestrial planet are recorded in Table 1.
The position of the separatrix is found by locating the local maximum of the secular potential (4). The separatrix terminates at a differential longitude from the planet $`\varphi _{}`$. This angle is the closest any tadpole orbit can come to the planet. In the restricted circular problem, this angle is $`\varphi _{}=2\mathrm{a}\mathrm{s}\mathrm{i}\mathrm{n}(\frac{1}{\sqrt{2}}\frac{1}{2})=23.906^{}`$ (see e.g., Brown & Shook 1933). The deviation from this classical value caused by the eccentricity and inclination of the planets is also recorded in Table 1.
The secular potential (4) can also be used to work out the approximate period $`P`$ of small angle librations about the Lagrange points. For simplicity, we take only the zeroth order term (5) with $`a=a_\mathrm{p}`$ and obtain the minimum of the secular potential at
$$\varphi =\mathrm{acos}\left(\frac{1}{2\mathrm{cos}^2\frac{i}{2}}\right).$$
(7)
The period of small librations about this minimum is
$$P=\frac{2}{3}\frac{P_\mathrm{p}}{\sqrt{m_\mathrm{p}(4\mathrm{cos}^4(\frac{i}{2})1)}}$$
(8)
where $`P_\mathrm{p}`$ is the orbital period of the planet. The period of libration increases with increasing inclination.
Expansion of the disturbing function is a good guide to the dynamics provided the inclinations and eccentricities are low or moderate. At high inclinations and eccentricities, the classification of the corotating orbits has only recently been undertaken by Namouni (1999).
## 3 Numerical Method
For each of the terrestrial planets, we carry out numerical surveys of in-plane and inclined corotating orbits. Section 3.1 introduces the integration method, while Section 3.2 discusses the numerical procedure.
### 3.1 Mixed Variable Symplectic Integrators
Our model includes all the planets (except Pluto whose contribution is negligible in the inner Solar System). The asteroids are represented as test particles with infinitesimal mass. They are perturbed by the planets but they do not perturb them in return.
The orbits of the planets are integrated using a mixed variable symplectic integrator scheme (Wisdom & Holman 1991; Kinoshita, Yoshida & Nakai 1991) with individual time steps (Saha & Tremaine 1994) which takes into account post-Newtonian corrections and the quadrupole moment of the Sun’s attraction on the barycentre of the Earth-Moon system (Quinn, Tremaine & Duncan 1991). The two latter contributions must be included for calculations in the inner Solar System, whereas the cumulative effect of the asteroids, satellites, galactic tidal acceleration, passing stars, solar mass loss and oblateness is believed to be smaller than $`10^{10}`$ and is neglected (Quinn et al. 1991). Mixed variable symplectic integrators exploit the fact that the Hamiltonian written in Jacobi coordinates (Plummer 1960; Wisdom & Holman 1991) is dominated by a nearly Keplerian term. The mixed variable symplectic integrators are so-called because they evaluate the planetary disturbing forces in Cartesian coordinates while using the elements to advance the orbits. Fast algorithms, like Gauss’ $`f`$ and $`g`$ functions, exist to perform the latter task (e.g., Danby 1988; Wisdom & Holman 1991).
The orbit of any of the test particles is derived from the Hamiltonian
$$H_{\mathrm{tp}}=H_{\mathrm{kep}}+H_{\mathrm{int}}.$$
(9)
Here, the test particle is given the first Jacobi index and the planets carry the higher Jacobi indices. As Wisdom & Holman (1991) point out, this is an advantageous choice as it gives the simplest interaction Hamiltonian. Denoting the Jacobi position of the test particle as $`𝐱_1`$ and its velocity as $`𝐯_1`$, then
$$H_{\mathrm{kep}}=\frac{v_1^2}{2}\frac{k^2}{r_1},H_{\mathrm{int}}=k^2\underset{i=2}{\overset{N+1}{}}m_i\left(\frac{𝐱_1𝐱_i}{r_i^3}\frac{1}{r_{1i}}\right),$$
(10)
where $`k`$ is the Gaussian gravitational constant and $`N`$ is the number of planets included in the model ($`N=8`$ for our calculations). We have used the notation $`r_{1i}=|𝐱_1𝐱_i|`$ as the distance from the test particle to the $`i`$th body. As mentioned by Wisdom and Holman (1991), the intuitive interpretation of $`H_{\mathrm{int}}`$ is to note that the attraction of the Sun on the test particle equals the difference between the direct acceleration of the massive planets on the test particle and the gravitational pull of the planets on the Sun.
The most important general relativistic effects can be included by modifying the test particle Hamiltonian to
$$H_{\mathrm{tp}}=H_{\mathrm{kep}}+H_{\mathrm{int}}+H_{\mathrm{PN}}.$$
(11)
A clever device for incorporating the most important post-Newtonian effects into mixed variable symplectic integrators is given by Saha & Tremaine (1994). The post-Newtonian Hamiltonian is recast as
$$H_{\mathrm{PN}}=\frac{1}{c^2}\left(\frac{3}{2}H_{\mathrm{kep}}^2\frac{k^4}{r_1^2}\frac{v_1^4}{2}\right).$$
(12)
The last expression contains three terms which are each integrable individually. The Keplerian part of eq. (12) can be incorporated into the usual Keplerian orbital advance using
$$\mathrm{exp}\left(\tau \{,H_{\mathrm{kep}}+\frac{3}{2c^2}H_{\mathrm{kep}}^2\}\right).$$
(13)
The Keplerian Hamiltonian is conserved with time and equals $`\frac{1}{2}k^2/a_{\mathrm{tp}}`$, $`a_{\mathrm{tp}}`$ being the semi-major axis of the test particle. After some straightforward algebra, eq. (13) becomes
$$\mathrm{exp}\left(\tau ^{}\{,H_{\mathrm{kep}}\}\right),\tau ^{}=(1\frac{3k^2}{2c^2a_{\mathrm{tp}}})\tau .$$
(14)
The second and third terms within the bracket of eq. (12) imply modifications of the test particle position and velocity before and after advancing the Keplerian orbit:
$$\frac{d𝐯_1}{dt}=\frac{2k^4}{c^2}\frac{𝐫_1}{r_1^4},\frac{d𝐱_1}{dt}=\left(\frac{2}{c^2}v_1^2\right)𝐯_1$$
(15)
Further details of this neat formulation are given in Saha & Tremaine (1994).
### 3.2 The Numerical Procedure
For each of the terrestrial planets, we carry out two surveys. The first is restricted to test particles in the orbital plane of the planet. The test particles are given the same eccentricity $`e`$, inclination $`i`$, longitude of the ascending node $`\mathrm{\Omega }`$ and mean anomaly $`M`$ as the planet. The argument of pericentre $`\omega `$ is varied from $`0^{}`$ to $`360^{}`$ in steps of $`5^{}`$. The initial semimajor axis is equal to the semimajor axis of the planet multiplied by a semimajor axis factor (c.f., Innanen & Mikkola’s (1989) investigation of Saturnian Trojans). The second survey is restricted to test particles with the same semimajor axis as the planet. The initial inclinations of the test particles (with respect to the plane of the planet’s orbit) are spaced every $`2^{}`$ and the initial arguments of pericentre are spaced every $`15^{}`$. The eccentricities of the asteroids are inherited from the parent planet.
The initial conditions come from the JPL Planetary and Lunar Ephemerides, DE405 which is available at “http://ssd.jpl.nasa.gov/” (Chamberlain et al. 1997). The starting epoch of the integration is JED2440400.5 (28 June 1969). The standard units used for the integration are the astronomical unit, the day and the Gaussian gravitational constant $`k^2=GM_{}`$. The Earth to Moon mass ratio is $`M_{}/M_\mathrm{L}=81.3`$. For most of the computations described below, the timestep for Mercury is $`14.27`$ days. The timesteps of the planets are in the ratio $`1:2:2:4:8:8:64:64`$ for Mercury moving outwards, so that Neptune has a timestep of $`2.5`$ years. The test particles all have the same timestep as Mercury. These values were chosen after some experimentation to ensure the relative energy error has a peak amplitude of $`10^6`$ over the tens of million year integration timespans (c.f., Holman & Wisdom 1993; Saha & Tremaine 1994). Individual planetary stepsizes do introduce additional stepsize resonances and it is important to check that the resonances do not degrade the accuracy of the numerical results. Referring to Figure 3 of Wisdom & Holman (1992), the stepsize resonances begin to overlap when the ratio of semimajor axes is $`>0.8`$ and stepsize is $`>0.2`$ times the orbital period. Bearing this in mind, it seems that the stepsize resonances are not a serious concern even for the test particles near Mercury. Nonetheless, Mercury does provide severe challenges for long-term integrations, and there is some evidence that round-off error may be affecting our results for the highly inclined test particles around Mercury.
As the test particles’ orbits are integrated, they are examined at each time step. If their trajectories become parabolic or hyperbolic orbits, they are removed from the survey. In addition, test particles which experience close encounters with a massive planet or the Sun are also terminated. The sphere of influence is defined as the surface around a planet at which the perturbation of the planet on the two-body heliocentric orbit is equal to that of the Sun on the two-body planetocentric orbit (e.g., Roy 1988). If the planet’s mass is much less than that of the Sun, this surface is roughly spherical with radius
$$r_\mathrm{s}=a_\mathrm{p}m_\mathrm{p}^{2/5},$$
(16)
where $`a_\mathrm{p}`$ is the semi-major axis of the planet and $`m_\mathrm{p}`$ is the planet’s mass in solar mass units. Since the algorithm we use does not allow the variable step size necessary for the treatment of close encounters, the exact size of the sphere of influence is not of great importance. Furthermore, test particles which enter the sphere of influence are typically ejected from the solar system in another 1-10 Myr (e.g., Holman 1997). In the case of the Sun, a close encounter is defined to be passage within 10 solar radii ($`0.005`$ AU). This general procedure is inherited from a number of recent studies on the stability of test particles in the Solar System (see e.g., Gladman & Duncan 1990; Holman & Wisdom 1993; Holman 1997; Evans & Tabachnik 1999)
We discuss the detailed results for each of the terrestrial planets in turn. A broad overview of the results is given in three tables. The fates of the test particles in the in-plane and inclined surveys are summarised in Tables 2 and 3 respectively. These provide the number of survivors on tadpole and horseshoe orbits, as well as the numbers suffering close encounters with each planet. Table 4 provides the average eccentric-ties and inclinations of the survivors at the end of the simulations.
## 4 Mercurian Surveys
For the Mercurian in-plane survey, the semimajor axis factor is chosen between $`0.998`$ and $`1.002`$ in steps of $`0.0004`$. If the semimajor axis factor is exactly unity and the argument of pericentre is displaced by $`60^{}`$ or $`300^{}`$, then the test particle is at the classical Lagrange point and can remain there on a stable orbit if the perturbations from the rest of the Solar System are neglected. Fig. 3 shows the time evolution of this array of test particles. The regions occupied by the test particles remaining after 1, 5, 20 and 100 million years are shown in red, green, blue and yellow respectively. After 100 million years, only 53 out of the original 792 test particles remain. The locations of the survivors are shown in close-up in Fig. 4. The most striking point to notice is that the stable zones do not include the classical Lagrange points themselves. In fact, all the survivors follow horseshoe orbits and there are no surviving tadpole orbits. There are no long-lived Mercurian Trojans. There are two possible reasons for this. First, Mercury is the least massive of the terrestrial planets and therefore the potential wells in which any long-lived Trojans inhabit are less deep than for the other terrestrial planets. Second, Mercury is the most eccentric of the terrestrial planets and this also encourages the erosion of the test particles. During the course of the 100 million years simulation, Mercury’s eccentricity fluctuates between $`0.1`$ and $`0.3`$.
The top three panels of Fig. 5 show the evolution of the semimajor axis, inclination and eccentricity for Mercury, together with a stable and an unstable test particle. The distributions of the orbital elements of the survivors suggest that they belong to a low inclination, low eccentricity family of horseshoe orbits. At the end of the simulation, the average eccentricity of the survivors is $`0.186`$ and their average inclination is $`6.968^{}`$. The lowest panel shows the evolution of the amplitude of libration $`D`$ (in degrees) about the Lagrange point for the two test particles. The stable particle starts off at a semimajor axis factor of 1.001 and an initial longitude of $`15^{}`$. The unstable particle starts off at exactly the $`L_5`$ Lagrange point. Notice that the eccentricity and inclination variations of the stable particle closely follow those of Mercury. The unstable test particle maintains its tadpole character only for some $`1.2\times 10^7`$ years, before passing through a brief horseshoe phase for $`0.5\times 10^7`$ years. After this, its semimajor axis increases to $`0.42`$ AU and finally decreases to $`0.35`$ AU before entering the sphere of influence of Mercury. The stable test particle follows a horseshoe orbit. This is obvious on examining the lowest panel, which shows that its angular amplitude of libration about the Lagrange point is $`320^{}`$.
For the Mercurian inclined survey, the orbits of 1104 test particles around Mercury are followed for 100 million years. Only thirteen test particles survived until the end of the 100 million year integration. Seven of these were low inclination ($`i<6^{}`$) test particles started off at arguments of pericentre ($`\omega =15^{}`$ or $`345^{}`$) very close to that of Mercury. The remaining six survivors were all high inclination objects. On repeating these calculations on similar machines with identical roundoff, but an updated ephemerides, all six of these highly inclined orbits were terminated before the end of the 100 million years (mostly because they entered the sphere of influence of Mercury), although the low inclination results were reproduced successfully. Clearly, it is not possible to draw definitive conclusions about the longevity of highly inclined asteroids near Mercury, although it seems likely that any stable zones must be small and depend sensitively on initial conditions.
## 5 Venusian Surveys
Perhaps one of the likeliest planets in the inner Solar System to harbour undiscovered Trojans is Venus. Fig. 6 shows the results of the in-plane survey of Venusian test particles. The starting semimajor axes are scaled by a fraction of the planet’s semimajor axis in steps of $`0.0012`$, while the argument of pericentre is offset from that of the planet by $`5^{}`$ steps. The test particles are integrated for 25 million years and the survivors recorded in the Figure. A filled circle represents a tadpole orbit, an open circle represents a horseshoe orbit. Notice that there are some long-lived survivors on horseshoe orbits even around the conjunction point. The tadpole orbits survive around $`L_4`$ for starting longitudes between $`15^{}`$ and $`160^{}`$ and around $`L_5`$ for starting longitudes between $`195^{}`$ and $`345^{}`$. The offset in the semimajor axes of the survivors $`\mathrm{\Delta }a`$ compared to the parent planet $`a_\mathrm{p}`$ satisfy $`\mathrm{\Delta }a/a_\mathrm{p}<0.72\%`$. There are 792 test particles at the beginning of the simulation, but only 407 persist till the end, of which 168 are on tadpole orbits. Horseshoes and tadpoles are of course divided by a separatrix in phase space (see Section 2). The break-up of the separatrix is associated with a chaotic layer, and this is responsible for erosion between the filled and open circles in the stable zones. At the edge of the figure, we show the stability boundary inferred from the elliptic restricted three body problem. The diamonds represent unstable test particles in the elliptic restricted three body problem (comprising the Sun, Venus and the massless asteroid). Everything within this outer boundary of diamonds is stable at the level of the restricted three body problem. Fig. 7 shows an at first sight surprising regularity of the orbits of the survivors. The extrema of the semimajor axis $`a_{\mathrm{ext}}`$ of the Trojan test particles are plotted against the initial longitude and fall on a one-parameter family of curves according to the semimajor axis factor. The heliocentric Hamiltonian for a test particle in the frame rotating with the mean motion of the planet may be written
$$H=\frac{k^2}{2a}m_\mathrm{p}k^2Rn_\mathrm{p}\sqrt{(1+m_\mathrm{p})a}.$$
(17)
Here, $`n_\mathrm{p}=\sqrt{(1+m_\mathrm{p})/a_\mathrm{p}^3}`$ is the mean motion of the planet using Kepler’s third law, and $`R`$ is approximated by the zeroth order term in the disturbing function
$$R=\frac{1}{a_\mathrm{p}}\left[\frac{1}{2|\mathrm{sin}\frac{1}{2}\varphi |}\mathrm{cos}\varphi \right],$$
(18)
where $`\varphi `$ is the difference in longitude between the planet and the asteroid. This follows on setting $`a=a_\mathrm{p}`$ in eq. (5). This Hamiltonian depends on time only through the slow variation of the planet’s orbital elements. These take place on a timescale much longer than the libration of the Trojan, and so the Hamiltonian is effectively constant. Expanding in the difference between the semimajor axis of the planet and the Trojan, we readily deduce that a test particle with a semimajor axis factor $`f`$ and which starts at a differential longitude $`\theta `$ has extrema $`a_{\mathrm{ext}}`$ satisfying
$$\frac{a_{\mathrm{ext}}}{a_\mathrm{p}}=1+\left[(f1)^2+\frac{8\mu }{3}\left(\frac{1}{2|\mathrm{sin}\frac{1}{2}\theta |}\mathrm{cos}\theta \frac{1}{2}\right)\right]^{{\scriptscriptstyle \frac{1}{2}}},$$
(19)
where $`\mu =m_\mathrm{p}/(1+m_\mathrm{p})`$ is the reduced mass. The extrema of the semimajor axis of the Trojans fall on this one-parameter family of curves, as depicted in Fig. 7.
Fig. 8 shows distributions of the eccentricities and inclinations of the survivors plotted against their initial semimajor axis and longitude from Venus. In all the panels, the maximum values attained during the course of the 25 million year orbit integrations are marked by crosses, the instantaneous values are marked by circles. There are two striking features of these diagrams. First, the eccentricities and the inclinations of the survivors remain very low indeed. The mean eccentricity of the sample after 25 million years is $`0.027`$, while the mean inclination is $`0.706^{}`$. This is a very stable family. Second, most of the survivors seem to occupy very nearly the same regions of the plots. The striking horizontal line of crosses in the rightmost panels of Fig. 8 suggest that these test particles do lie on similar orbits belonging to similar families exploring similar regions of phase space.
Fig. 9 shows the results of the Venusian inclined survey. Here, the orbits of 1104 test particles around Venus are integrated for 75 million years. The initial inclinations of the test particles (with respect to the plane of Venus’ orbit) are spaced every $`2^{}`$ and the initial longitudes (again with respect to Venus) are spaced every $`15^{}`$. The test particles are colour-coded according to whether they survive till the end of the 5, 20, 50 and 75 million year integration timespans. The 137 survivors after 75 million years all have smallish inclinations ($`i<16^{}`$). Notice, too, that the stable zones around the Lagrange points are still connected by some surviving test particles at conjunction. Bodies trapped around the Lagrange points should be sought at all longitudes close to the orbital plane. Fig. 10 shows a close-up of the stable zones, with tadpole orbits represented by filled circles and horseshoe orbits by open circles. Most of the objects that do survive are true Trojans, in that there orbits are recognisably of a tadpole character. There are just 8 surviving horseshoe orbits. The ensemble has a mean eccentricity of $`0.041`$ and a mean inclination of $`6.941^{}`$.
Assuming that they are primordial, we can estimate the number of coorbiting Venusian satellites by extrapolating from the number of Main Belt asteroids (c.f., Holman 1997, Evans & Tabachnik 1999). The number of Main Belt asteroids $`N_{\mathrm{MB}}`$ is $`N_{\mathrm{MB}}<\mathrm{\Sigma }_{\mathrm{MB}}A_{\mathrm{MB}}f`$, where $`A_{\mathrm{MB}}`$ is the area of the Main Belt, $`\mathrm{\Sigma }_{\mathrm{MB}}`$ is the surface density of the proto-planetary disk and $`f`$ is the fraction of primordial objects that survive ejection (which we assume to be a universal constant). Let us take the Main Belt to be centred on $`2.75`$ AU with a width of $`1.5`$ AU. Fig. 6 suggests that the belt of Venusian Trojans is centred on $`0.723`$ AU and has a width of $`<0.008`$ AU. If the primordial surface density falls off inversely proportional to distance, then the number of coorbiting Venusian asteroids $`N_\mathrm{V}`$ is
$$N_\mathrm{V}<\left(\frac{2.75}{0.723}\right)\left(\frac{0.723\times 0.008}{2.75\times 1.5}\right)N_{\mathrm{MB}}0.0053N_{\mathrm{MB}}.$$
(20)
The number of Main Belt asteroids with diameters $`>1`$ km is $`40000`$, which suggests that the number of Venusian Trojans is $`100`$ with perhaps a further $`100`$ coorbiting companions on horseshoe orbits.
## 6 Terrestrial Surveys
The Earth is slightly more massive than Venus, and this augurs well for the existence of coorbiting satellite companions. The Earth has no known Trojans (Whiteley & Tholen 1998), but it does possess the asteroidal companion 3753 Cruithne which moves on a temporary horseshoe orbit (Wiegert, Innanen & Mikkola 1997; Namouni, Christou & Murray 1999). The asteroid will persist on this horseshoe orbit for a few thousand years.
Fig. 11 shows the surviving in-plane test particles near the Earth after their orbits have been integrated for 50 million years. Again, filled circles represent the tadpole orbits, open circles the horseshoe orbits. On comparison with Fig. 6, we see that the stable zones of the Earth are more extensive and the number of survivors is greater than for Venus. Tadpole orbits survive for $`\mathrm{\Delta }a/a_\mathrm{p}<0.48\%`$ and horseshoes for $`\mathrm{\Delta }a/a_\mathrm{p}<1.20\%`$. Although the number of survivors is greater, the number of true Trojans on tadpole orbits is much less than for Venus. Specifically, of the 792 original test particles, only 509 persist till the end of the simulation and of these just 95 are on tadpole orbits. In the case of the Earth, just $`19\%`$ of the survivors are tadpoles, as opposed to $`41\%`$ for Venus. The visual consequence of this is that the holes in the stable regions for the Earth are more pronounced than for Venus (compare also the equivalent diagram for Saturn provided by Holman & Wisdom 1993). The survivors have a mean eccentricity of $`0.038`$ and a mean inclination of $`1.349^{}`$, consistent with a long-lived population.
Fig. 12 shows the erosion of an ensemble of 1104 inclined test particles positioned at the same semimajor axis as the Earth but varying in longitude. Again, the the initial inclinations of the test particles (with respect to the plane of the Earth’s orbit) are spaced every $`2^{}`$ and the initial longitudes are spaced every $`15^{}`$. The test particles are colour-coded according to their survival times – those surviving after 1, 5, 10 and 25 million years are shown in green, blue, yellow and red respectively. The ensemble after 25 million years is shown in Fig. 13 with tadpoles represented as filled circles and horseshoes as open circles. Again, the first thing to note is the number of survivors – there are 200 in total, almost all of which are moving on tadpole orbits. There seem to be two bands of stability, one at low starting inclinations ($`i<16^{}`$) and one at moderate starting inclinations ($`24^{}<i<34^{}`$). On careful inspection of our earlier Fig. 9, it is possible to discern in blue and green a similar band of stable trajectories at moderate inclination for Venus. These though are swept out much more quickly than for the Earth; they are all gone after just 5 million years. The survivors in Fig. 13 have a low mean eccentricity of $`0.064`$. The distribution of inclinations is strikingly bimodal as is evident from Fig. 14. The mean inclination at the end of the simulation is $`15.778^{}`$.
Let us remark that the Earth has more surviving test particles in both of our surveys than Venus, as well as more extensive stable zones. This suggests that the asteroid 3753 Cruithne may not be unique, but the first member of a larger class of coorbiting terrestrial companions (Namouni, Christou & Murray 1999). Making the same assumptions as in (20), our estimate for the number of terrestrial companions is
$$N_\mathrm{E}<\left(\frac{2.75}{1.00}\right)\left(\frac{1.00\times 0.01}{2.75\times 1.5}\right)N_{\mathrm{MB}}260.$$
(21)
Here, we have used Fig. 11 to set the width of the stable zone around the Earth as $`0.005`$ AU. The total number of coorbiting companions will be much higher, if one includes transient objects like 3753 Cruithne.
## 7 Martian Surveys
Mars is the only terrestrial planet already known to possess Trojan asteroids. These are 5261 Eureka and 1998 VF31 (see, for example, Mikkola et al. 1994, Tabachnik & Evans 1999 and the references therein). Both have moderate inclinations to the ecliptic, namely $`20.3^{}`$ and $`31.3^{}`$ respectively.
The in-plane Martian survey is presented in Fig. 15. The test particles are colour-coded according to their survival times. Most of the particles are already swept out after 5 million years. After 50 million years, there is only one test particle remaining, and it too is removed shortly afterwards. There are no survivors after 60 million years. This confirms the earlier suspicions of Mikkola & Innanen (1994) that Trojans in the orbital plane of Mars are not long-lived. As is evident from Table 2, the most common fate of the test particles is to enter the sphere of influence of Mars. The inclined survey is shown in Fig. 16, together with the positions of the two Trojans, 5261 Eureka (marked by a square) and 1998 VF31 (asterisk). Two further asteroids – 1998 QH56 (triangle) and 1998 SD4 (diamond) – have been suggested as Trojan candidates, although improved orbital elements together with detailed numerical simulations (Tabachnik & Evans 1999) now make this seem rather unlikely. The result of the inclined survey is to show stable zones for inclinations between $`14^{}`$ and $`40^{}`$ for timespans of 25 million years. The stable zones are strongly eroded at $`29^{}`$.
Following this, we conducted another experiment in which inclined Martian Trojans are simulated for 100 Myrs. The initial conditions are inherited from Mars, except for the argument of pericentre which is offset by $`60^{}`$ ($`L_4`$) and $`300^{}`$ ($`L_5`$), and the inclinations which are selected in the range from $`0^{}40^{}`$ from Mars’ orbital plane. To examine in more detail the effects of the other planet’s perturbations, a 1-degree step in inclination is chosen. Fig. 17 shows the results of this exercise. The upper panel expresses the termination time versus the initial inclinations of the test particles at both Lagrange points. Not surprisingly, low inclination Trojans ($`i<5^{}`$) enter Mars’ sphere of influence on a 10 million year timescale. The stable inclination windows are also recovered with a strong disturbing mechanism at $`29^{}`$ for $`L_5`$ and $`30^{}`$ for $`L_4`$. The four lower panels give the eccentricities and inclinations of the remaining test particles at the end of the integration. Crosses identify the maximum quantities over the entire timespan, while open circles show the instantaneous values at 100 Myr. The general trend is to have stable orbits ($`e_{\mathrm{max}}<0.2`$) in the range $`15^{}34^{}`$ in the case of $`L_4`$ and $`9^{}36^{}`$ in the case of $`L_5`$. Interestingly, the two securely known Trojans, namely 5261 Eureka and 1998 VF31, occupy positions at $`L_5`$ corresponding to the two local minima of the maximum eccentricity curve.
The results for the inclined Martian Trojans are unusual, and it is natural to seek an explanation in terms of secular resonances. Large disturbances can occur when there is a secular resonance, that is, when the averaged precession frequency of the asteroid’s longitude of pericentre $`\dot{\varpi }`$ or longitude of node $`\dot{\mathrm{\Omega }}`$ becomes nearly equal to an eigenfrequency of the planetary system (e.g., Brouwer & Clemence 1961; Williams & Faulkner 1981; Scholl et al. 1989). The secular precession frequencies in linear theory are usually labelled $`g_j`$ and $`s_j`$ ($`j=1,\mathrm{}8`$ for Mercury to Neptune) for the longitude of pericentre and longitude of node respectively. Their mean values computed over 200 million years are listed in Laskar (1990).
Fig. 18 can be used to infer the positions of some of the principal secular resonances as a function of inclination in the vicinity of each Lagrange point. The vertical axis is the rate of variation of various angles, the horizontal axis is the inclination. A resonance occurs whenever the rate of variation vanishes. The blue curve shows the frequency $`\dot{\varpi }g_5`$, which vanishes at inclinations $`2830^{}`$ for the $`L_4`$ point and $`29^{}`$ at the $`L_5`$ point. Notice that the resonance is much broader at the $`L_4`$ point. The yellow curves show $`3\dot{\mathrm{\Omega }}+2\dot{\omega }g_5`$. This frequency vanishes at a range of inclinations between $`611^{}`$, again with slight differences noticeable at the two Lagrange points. As these are resonances with Jupiter, they are expected to be the most substantial. At this semimajor axis, there is just one resonance with Saturn. The black curves show the frequency $`2\dot{\mathrm{\Omega }}+3\dot{\omega }g_6`$. This resonance occurs at inclinations of $`15^{}`$. Lastly, there are two weaker resonances with the Earth. These may be tracked down using the red curve, which shows $`\dot{\varpi }g_3`$, and the green curve, which shows $`6\dot{\mathrm{\Omega }}+\dot{\omega }s_3`$. This completes the list of the main resonances in the vicinity of Mars. The importance of the Jovian resonances in particular has been pointed out before by Mikkola & Innanen (1994).
Fig. 19 illustrates the evolution of a few arbitrarily selected orbits. The left panels refer to the $`L_4`$ Lagrange point, the right panels to the $`L_5`$ point. The panels are labelled according to the initial inclination with respect to Mars’ orbit. They plot the evolution of the eccentricity $`e`$ and the inclination with respect to Jupiter $`i_{\mathrm{jup}}`$ for a typical test particle (full curve) and Mars (broken curve). The inclinations over which the two Jovian resonances operate are shown as shaded bands. The uppermost two panels refer to orbits that start out in the orbital plane of Mars. Their inclination is initially increased, and this takes the orbits into the régime in which one of the Jovian resonances is dominant. The eccentricity of the orbit is pumped whenever it lingers in the shaded inclination band. This makes the orbit Mars-crossing and the test particle is terminated. The middle two panels show the fates of orbits at intermediate inclinations of $`12^{}`$ at the $`L_4`$ point and $`5^{}`$ at the $`L_5`$ point. Although the behaviour is quite complex, the final increase in eccentricity in both cases coincides with prolonged stays in the resonant region. We conclude that the low inclinations test particles are destabilised by this secular resonance with Jupiter. The bottom two panels show fates of test particles starting off at $`30^{}`$ at the $`L_4`$ point and $`29^{}`$ at the $`L_5`$ point. In both cases, there is a rapid and pronounced increase in the eccentricity, which takes it onto a Mars-crossing path. This destabilisation occurs only for a very narrow range of inclinations. This is manifest in the erosion in Fig. 16, especially at the $`L_4`$ point near inclinations of $`30^{}`$.
Mikkola & Innanen (1994) suggested that the instability of low inclination Martian Trojans was due to a secular resonance with Mars driving the inclination upwards. In their picture, this continues until a critical inclination of $`12^{}`$ is reached when $`3\mathrm{\Omega }+2\omega `$ resonates with Jupiter. The difficulty with this is that it is not clear whether the claimed Martian secular resonance – which is really just equivalent to the statement that the test particle is coorbiting – is responsible for the inclination increase. Our Fig. 18 seems to show that $`3\dot{\mathrm{\Omega }}+2\dot{\omega }g_5`$ nearly vanishes over a broad range of inclinations and we suspect it may be able to cause the damage on its own.
## 8 Conclusions
The possible existence of long-lived coorbiting satellites of the terrestrial planets has been examined using numerical simulations of the Solar System. Of course, integrations in the inner Solar System are laborious as much smaller timesteps are required to follow the orbits of the satellites of Mercury as opposed to the giant planets like Jupiter. Our numerical surveys have been pursued for timescales up to 100 million years – typically an order of magnitude greater than previous computations in the inner Solar System. The numerical algorithm is a symplectic integrator with individual timesteps that incorporates the most important post-Newtonian corrections (Wisdom & Holman 1991; Saha & Tremaine 1994).
One worry concerning our integrations is that they extend at most to 100 million years, which is a small fraction of the age of the Solar System ($`5`$ Gyrs). To gain an idea of the possible effects of longer integration times, we can use the approximate device of fitting our existing data and extrapolating. Graphs of the number of surviving particles against time are shown in Figs. 2021 for the in-plane and inclined surveys respectively. In the former case, the data are generally well-fit by a logarithmic decay law of form
$$N(t)=a+b\mathrm{log}_{10}\left(t[\mathrm{yrs}]\right)$$
(22)
Table 5 shows the best fitting values of $`a`$ and $`b`$. It also gives the extrapolated number of test particles after 1 and 5 Gyrs. For both Venus and the Earth, this suggests that several hundred test particles remain, even if the simulations are run for the age of the Solar System. In the inclined case, the data is not well-fitted by logarithmic decay laws. Instead, Table 6 gives the results of fitting the data between 10 and 100 Myr to a power-law decay of form
$$N(t)=\frac{10^c}{\left(t[\mathrm{yrs}]\right)^d}$$
(23)
Extrapolation suggests that some inclined test particles remain for Venus, the Earth and Mars even after 5 Gyr. However, these numbers are highly speculative for the inclined survey, as only the tail of the distribution is fitted.
The results of our surveys for Venus and the Earth are somewhat similar. Long-lived coorbiting satellites can persist in the orbital planes of both planets. The stable zones of the Earth are larger than those Venus, although the Earth retains fewer true Trojans or tadpoles. The semimajor axes of the stable test particles $`a`$, as compared to the parent planet $`a_\mathrm{p}`$, satisfy $`\mathrm{\Delta }a/a_\mathrm{p}<0.72\%`$ for Venus and $`\mathrm{\Delta }a/a_\mathrm{p}<1.2\%`$ for the Earth. Both Venus and the Earth have low inclination régimes in which long-lived test particles survive for timescales of tens of millions of years, despite the disturbing perturbations from the remainder of the Solar System. The stable zones satisfy $`i<16^{}`$ for Venus. For the Earth, there are two bands of stability, one at low inclinations ($`i<16^{}`$) and one at moderate inclinations ($`24^{}<i<34^{}`$). However, there is a hint that the higher inclination band may be further eroded with still longer timespan integrations. The inclined test particles that survive primarily move on tadpole orbits. It is possible to make very crude estimates of numbers by extrapolation from the Main Belt. These suggest that there may be some hundreds of asteroids in the coorbital regions of Venus and the Earth.
In the case of Mercury, very few of the test particles – whether starting in the orbital plane or at higher inclinations – survive till the end of the integrations. This seems reasonable, as both the low mass and the high eccentricity of Mercury militate against stable zones. A population of numerous, long-lived Mercurian Trojans seems rather unlikely. Recently, Namouni et al. (1999) have speculated that highly inclined coorbiting Mercurian satellites (“Vulcanoids”) may exist. Our survey does hint at the survival of a handful of test particles at high inclinations. On re-simulating these results over the same timespan using an updated ephemerides, these test particles did not survive but were ejected. So, if there are stable zones of inclined coorbiting satellites, then they must be narrow and rather sensitive to the detailed initial conditions.
From the viewpoint of dynamics, Mars offers perhaps the greatest interest of all. Here, the test particles within the orbital plane are all ejected on 60 million year timescales. The inclined survey shows stable zones for inclinations between $`14^{}`$ and $`40^{}`$. This is certainly consistent with the two known Martian Trojans (5261 Eureka and 1998 VF31) which have orbits moderately inclined to the ecliptic ($`20.3^{}`$ and $`31.3^{}`$ respectively) about the $`L_5`$ point. The survey shows strong erosion at a narrow band of inclinations concentrated around $`2830^{}`$ at $`L_4`$ and $`29^{}`$ at $`L_5`$. This may be traced to a strong, but narrow, Jovian resonance. The averaged precession frequency of the test particle’s longitude of pericentre $`\dot{\varpi }`$ is equal to the secular precession frequency of the longitude of pericentre of Jupiter $`g_5`$. We believe that the destabilisation of the low inclination and in-plane orbits is also due to Jupiter. Here, the combination of frequencies $`3\dot{\mathrm{\Omega }}+2\dot{\omega }g_5`$ vanishes over a broad range of inclinations between $`611^{}`$. As they evolve, low inclination test particles enter this band of inclinations, their eccentricity is increased and they become Mars-crossing.
## Acknowledgments
NWE is supported by the Royal Society, while ST acknowledges financial help from the European Community. We wish to thank John Chambers, Jane Luu, Seppo Mikkola, Fathi Namouni, Prasenjit Saha and Scott Tremaine for helpful comments and suggestions. Jane Luu and Prasenjit Saha provided critical readings of the manuscript.
## Appendix A Auxiliary Formulae
This appendix lists auxiliary functions for the averaged disturbing function
$`U_{2,0}`$ $`=`$ $`{\displaystyle \frac{a}{a_\mathrm{p}^2}}\left[\frac{1}{2}(e^2+e_\mathrm{p}^2)\mathrm{cos}\varphi ee_\mathrm{p}\mathrm{cos}(\omega +\mathrm{\Omega }+2\varphi )\right]\mathrm{cos}^2\frac{i}{2},`$
$`U_{2,3}`$ $`=`$ $`[\frac{1}{4}ee_\mathrm{p}aa_\mathrm{p}\mathrm{cos}(\omega +\mathrm{\Omega }+2\varphi )\frac{1}{2}(e^2+e_\mathrm{p}^2)aa_\mathrm{p}\mathrm{cos}\varphi `$
$`+\frac{9}{4}ee_\mathrm{p}aa_\mathrm{p}\mathrm{cos}(\omega +\mathrm{\Omega })]\mathrm{cos}^2\frac{i}{2}\frac{3}{4}(a^2e^2+a_\mathrm{p}^2e_\mathrm{p}^2),`$
$`U_{2,5}`$ $`=`$ $`[\frac{15}{8}(e^2+e_\mathrm{p}^2)a^2a_\mathrm{p}^2+\frac{27}{8}ee_\mathrm{p}a^2a_\mathrm{p}^2\mathrm{cos}(\omega +\mathrm{\Omega }\varphi )`$
$`\frac{9}{4}ee_\mathrm{p}a^2a_\mathrm{p}^2\mathrm{cos}(\omega +\mathrm{\Omega }+\varphi )\frac{9}{8}(e^2+e_\mathrm{p}^2)a^2a_\mathrm{p}^2\mathrm{cos}2\varphi `$
$`+\frac{3}{8}ee_\mathrm{p}a^2a_\mathrm{p}^2\mathrm{cos}(\omega +\mathrm{\Omega }+3\varphi )]\mathrm{cos}^4\frac{i}{2}+`$
$`[\frac{3}{4}ee_\mathrm{p}aa_\mathrm{p}(a^2+a_\mathrm{p}^2)\mathrm{cos}(\omega +\mathrm{\Omega }+2\varphi )\frac{3}{2}aa_\mathrm{p}(a^2e^2+`$
$`a_\mathrm{p}^2e_\mathrm{p}^2)\mathrm{cos}\varphi \frac{9}{4}ee_\mathrm{p}aa_\mathrm{p}(a^2+a_\mathrm{p}^2)\mathrm{cos}(\omega +\mathrm{\Omega })]\mathrm{cos}^2\frac{i}{2}`$
$`+\frac{3}{4}(a^4e^2+a_\mathrm{p}^4e_\mathrm{p}^2)+\frac{3}{2}ee_\mathrm{p}a^2a_\mathrm{p}^2\mathrm{cos}(\omega +\mathrm{\Omega }+\varphi )`$
$`+\frac{3}{4}a^2a_\mathrm{p}^2\mathrm{sin}^4\frac{i}{2}`$
We are indebted to Seppo Mikkola, who kindly confirmed the correctness of these expressions for us. The formula given in Mikkola et al. (1994) is erroneous. Expansions to the fourth order in both eccentricity and inclination are presented in Tabachnik (1999).
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# 1 Introduction
## 1 Introduction
Any Noether current associated to a gauge symmetry can be rewritten on-shell as the divergence of a superpotential. This is a very general property and does not depend on the gauge invariant theory we are considering . Let us be more precise. Suppose that a gauge symmetry variation is given locally by
$$\delta _\xi 𝝋^i=𝒅\xi ^\alpha \mathbf{}𝚫_\alpha ^i+\xi ^\alpha \mathbf{}\stackrel{~}{𝚫}_\alpha ^i$$
(1)
where $`𝝋^i`$ denotes a $`p_i`$-form field. The $`(p_i1)`$ and $`p_i`$ space-time forms, $`𝚫_\alpha ^i`$ and $`\stackrel{~}{𝚫}_\alpha ^i`$, are functions of the fields. Then the Noether current associated to a one parameter symmetry subgroup (1) is always given by:
$$𝑱_\xi =𝒅𝑼_\xi +𝑾_\xi $$
(2)
with the definition
$$𝑾_\xi :=\xi ^\alpha 𝚫_\alpha ^i\mathbf{}𝑬_i,$$
(3)
where the $`(Dp_i)`$-forms $`𝑬_i:=\frac{\delta 𝑳}{\delta 𝝋^i}`$ are the Euler-Lagrange equations associated to the field $`𝝋^i`$. We used equation (1) in the definition (3). Note that the Noether equation $`𝒅𝑱_\xi =\delta _\xi 𝝋^i\mathbf{}𝑬_i`$ follows from (2) and from the Noether identities due to the gauge symmetry, namely, $`𝒅𝑾_\xi =\delta _\xi 𝝋^i\mathbf{}𝑬_i`$. The current is associated to a one-parameter subgroup of “rigid” symmetries when $`\xi ^\alpha `$ can be globally taken as constant or at least has a canonical spacetime dependence as in the case of rotations, but the form (2) follows from local gauge invariance along that one parameter subgroup.
In our previous work , we emphasized the fact that the Noether method does not define the current $`𝑱_\xi `$ and its superpotential $`𝑼_\xi `$ unambiguously. This is the well-known fact that the Noether current is defined up to some exact form, namely $`𝑱_\xi 𝑱_\xi +𝒅𝒀`$. This exact term may contribute to the Noether charge when the space-time has some boundary with non-vanishing fields on it. A case by case prescription which depends on the boundary conditions is then needed in order to define $`𝑱_\xi `$. An attempt to give a general formula for $`𝑼_\xi `$ can give rise to incorrect results; see for instance our remarks in the case of supergravities . The same problem arises in the covariant symplectic phase space formalism. In that case, the ambiguity on the symplectic form can be fixed by some “covariant” criterion as in .
In , a way to solve the ambiguity in the Noether current was proposed. The main result was to give a formula for an arbitrary variation of the superpotential, namely:
$$__r\delta 𝑼_\xi =__r\delta 𝝋^i\mathbf{}\frac{𝑾_\xi }{𝒅𝝋^i}$$
(4)
We assume that our spacetime $``$ is bounded by a set of $`n`$ $`(D1)`$-dimensional time-like (or null) hypersurfaces denoted by $`_r`$, $`r=\{1,\mathrm{},n\}`$. Then, $`r`$ labels the connected time-like (or null) components of $``$. We denote by $`\mathrm{\Sigma }_t`$ a space-like Cauchy hypersurface at fixed time $`t`$. Then, the closed $`(D2)`$-dimensional manifold $`_r`$ of (4) is defined by $`_r=\mathrm{\Sigma }_t_r`$, for some $`r=\{1,\mathrm{},n\}`$ (and then $`\mathrm{\Sigma }_t=_{r=1}^n_r`$). In our gravitational example, we will choose $`_r`$ to be spatial infinity at fixed time $`_{\mathrm{}}`$.
It is important to recall that the formula (4) holds only if the theory has been rewritten in a first order formalism<sup>2</sup><sup>2</sup>2 About the possibility to construct a first order theory from a second order one preserving the symmetries, see ., in the sense that both, the equations of motion and the symmetries (1), depend at most on the first derivatives of the fields . Note that this definition does not imply that a $`1^{st}`$ order Lagrangian is linear in the first derivatives of the fields. The simplest example of a $`1^{st}`$ order theory quadratic in the first derivatives of the fields is the $`5`$-dimensional (abelian) Chern-Simons Lagrangian $`𝑳_{CS}=𝑨\mathbf{}𝒅𝑨\mathbf{}𝒅𝑨`$.
The formula for the superpotential (4) is non-ambiguous because it depends only on the equations of motion and on the functional form of the gauge symmetry (1) of the theory. It expresses an arbitrary variation of the superpotential to be integrated on $`_r`$ (for any $`r=\{1,\mathrm{},n\}`$) taking into account the chosen boundary conditions. No additional information on the behavior of the fields outside of $`_r`$ is needed for consistency. In fact, equation (4) can be justified using only the value of the fields (and their derivatives) at $`_r`$, no matter what happens in the bulk or on the other boundaries $`_s`$, $`sr`$. This is indeed expected from our experience: The ADM mass gives the total mass at spatial infinity, independently of the number of black holes inside the spacetime (and then independently of the boundary conditions used to describe their horizons). This has to be contrasted with the Hamiltonian Regge and Teitelboim procedure . There, the requirement of differentiability of the generators of first class constraints implies that the Hamiltonian version of equation (4) has to be satisfied on every $`_r`$. The same stronger condition is also needed in the so-called covariant phase space symplectic formalism for consistency (see for example ).
The purpose of this paper is to use formula (4) to compute the superpotentials associated to general relativity<sup>3</sup><sup>3</sup>3The examples of Yang-Mills, higher dimensional Chern-Simons and supergravity theories were respectively studied in and (see also for a review)., with or without cosmological constant, in any spacetime dimension ($`D3`$). As stated already in (using there a case by case prescription) we find that the superpotential associated to a “Dirichlet” boundary condition on the metric is the one proposed by Katz, Bic̆ák and Lynden-Bell . The superpotential associated to Dirichlet boundary conditions on the connection (ie Neumann condition on the metric), is one half (in four spacetime dimensions) the famous Komar superpotential<sup>4</sup><sup>4</sup>4This result was rediscovered in .. We use here the $`gl(D,\text{I}\mathrm{R})`$ first order formalism developed in . The motivation is that the Affine-$`gl(D,\text{I}\mathrm{R})`$ formalism generalizes both the Palatini (dear to relativists) and the tetrad-vielbein (needed for supergravities) formalisms. So a computation at the $`gl(D,\text{I}\mathrm{R})`$-level can be pulled back to either of these two well-known formalisms without much additional work.
We would like to insist on the simplicity and on the absence of any ambiguity or additional criteria of our present derivation of these gravitational superpotentials. In particular, the (general) covariance of our results is automatic and does not have to be required by hand.
We finally comment on the use of the KBL superpotential at null infinity.
For another approaches to compute superpotentials that do not emphasize the boundary conditions see .
## 2 The $`gl(D,\text{I}\mathrm{R})`$ formalism for gravity and the associated superpotentials
### 2.1 The $`gl(D,\text{I}\mathrm{R})`$ gravity
The basic idea of the $`gl(D,\text{I}\mathrm{R})`$ formalism is to unify the two known first order formulations of General Relativity, namely the Palatini and the vielbein (orthonormal frames) formulations, both will follow from partial extremisations of our action. Let us recall the results of .
In the Palatini case, the torsionless condition for the connection (namely $`\mathrm{\Gamma }_{\nu \rho }^\mu =\mathrm{\Gamma }_{\rho \nu }^\mu `$) is assumed from the beginning. The metric compatibility of the connection with the metric ($`_\mu g_{\nu \rho }=0`$) and the Einstein equations are derived from the variational principle.
On the other hand, the vielbein formulation assumes the metric compatibility condition between the flat Minkowski metric and the associated orthonormal connection, that is, $`D_\mu \eta ^{ab}=_\mu \eta ^{ab}+\omega _{\mu c}^a\eta ^{cb}+\omega _{\mu c}^b\eta ^{ac}=2\omega _\mu ^{(ab)}=0`$. The torsionless condition ($`D_{[\mu }\theta _{\nu ]}^a=0`$) and Einstein equations follow from the equations of motion of the basic fields.
The $`gl(D,\text{I}\mathrm{R})`$ first order formulation combines both ideas in a nice way: nothing is assumed from the beginning and the metric compatibility and torsionless condition are derived from the equations of motion of the connection of the linear frame bundle $`gl(D,\text{I}\mathrm{R})`$ (after fixing the Projective symmetry in the Einstein gauge). The Einstein equations are recovered as usual.
### The Lagrangian and the equations of motion
The Lagrangian of the $`gl(D,\text{I}\mathrm{R})`$ gravity is a D-form $`𝑳`$ (D $`3`$ is the spacetime dimension), it is a function of a linear 1-form connection $`𝝎_b^a`$ (for a Yang-Mills type $`gl(D,\text{I}\mathrm{R})`$), of the canonical 1-form $`𝜽^a`$ ($`\text{I}\mathrm{R}^D`$ valued) and of the metric $`g^{ab}`$ (which will be used to lift and lower the $`\text{I}\mathrm{R}^D`$-valued indices), as well as their first derivatives:
$$𝑳=\frac{1}{4\kappa ^2}𝑹_c^a\mathbf{}\sqrt{\left|g\right|}g^{cb}𝚺_{ab}$$
(5)
Where<sup>5</sup><sup>5</sup>5In our previous paper , we used a different notation, namely $`2k=4\kappa ^2=16\pi G`$. $`4\kappa ^2=16\pi G`$ and
$$𝚺_{a_1\mathrm{}a_r}:=\frac{1}{(Dr)!}ϵ_{a_1\mathrm{}a_rc_{r+1}\mathrm{}c_D}𝜽^{c_{r+1}}\mathbf{}\mathrm{}\mathbf{}𝜽^{c_D}$$
(6)
$`ϵ_{a_1\mathrm{}a_D}`$ being the Levi-Civita symbol ($`ϵ_{0\mathrm{}(D1)}=1`$).
Each field of the theory has a curvature,
$`𝑹_b^a`$ $`:=`$ $`𝒅𝝎_b^a+𝝎_c^a\mathbf{}𝝎_b^c`$ (7)
$`𝚯^a`$ $`:=`$ $`𝑫𝜽^a=𝒅𝜽^a+𝝎_b^a\mathbf{}𝜽^b`$ (8)
$`𝚵^{ab}`$ $`:=`$ $`𝑫g^{ab}=𝒅g^{ab}+𝝎^{ab}+𝝎^{ba}`$ (9)
called respectively the curvature 2-form, the torsion and the nonmetricity.
The Euler-Lagrange equations corresponding to (5) are given by:
$`{\displaystyle \frac{\delta 𝑳}{\delta g^{ab}}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{\left|g\right|}}{8\kappa ^2}}\left(𝑹_a^c\mathbf{}𝚺_{cb}+𝑹_b^c\mathbf{}𝚺_{ca}g_{ab}𝑹^{cd}\mathbf{}𝚺_{cd}\right)`$ (10)
$`{\displaystyle \frac{\delta 𝑳}{\delta 𝜽^a}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{\left|g\right|}}{4\kappa ^2}}𝑹^{bc}\mathbf{}𝚺_{bca}`$ (11)
$`{\displaystyle \frac{\delta 𝑳}{\delta 𝝎_b^a}}`$ $`=`$ $`{\displaystyle \frac{1}{4\kappa ^2}}𝑫\left(\sqrt{\left|g\right|}g^{bc}𝚺_{ac}\right)={\displaystyle \frac{\sqrt{\left|g\right|}}{4\kappa ^2}}\left(\mathbf{\Xi ̸}^{bc}\mathbf{}𝚺_{ac}+𝚯^c\mathbf{}𝚺_{aec}g^{eb}\right)`$ (12)
with $`\mathbf{\Xi ̸}^{ab}:=𝚵^{ab}g^{ab}\frac{𝚵}{2}`$ and $`𝚵:=𝚵^{ab}g_{ab}`$.
### The gauge symmetries
The Lagrangian (5) is invariant under three gauge symmetries:
1. 1) $`\mathrm{𝑇ℎ𝑒}\mathrm{𝑙𝑜𝑐𝑎𝑙}\mathrm{`}\mathrm{`}\mathrm{𝑓𝑟𝑎𝑚𝑒}\mathrm{𝑐ℎ𝑜𝑖𝑐𝑒}\mathrm{"}\mathrm{𝑓𝑟𝑒𝑒𝑑𝑜𝑚}`$, parametrized by an arbitrary infinitesimal local matrix of $`gl(D,\text{I}\mathrm{R})`$ namely, $`\lambda _b^a=\lambda _b^a(x)`$. The variations of the fields are:
$`\delta _\lambda 𝜽^a`$ $`=`$ $`\lambda _b^a𝜽^b`$
$`\delta _\lambda g^{ab}`$ $`=`$ $`\lambda _c^ag^{cb}+\lambda _c^bg^{ac}`$
$`\delta _\lambda 𝝎_b^a`$ $`=`$ $`𝑫\lambda _b^a=𝒅\lambda _b^a𝝎_c^a\lambda _b^c+𝝎_b^c\lambda _c^a`$ (13)
2. 2) $`\mathrm{𝑇ℎ𝑒}\mathrm{𝑑𝑖𝑓𝑓𝑒𝑜𝑚𝑜𝑟𝑝ℎ𝑖𝑠𝑚}\mathrm{𝑖𝑛𝑣𝑎𝑟𝑖𝑎𝑛𝑐𝑒}`$, parametrized by an arbitrary infinitesimal vector field $`\xi ^\rho =\xi ^\rho (x)`$:
$$\delta _\xi 𝜽^a=_\xi 𝜽^a$$
(14)
and so on for $`g^{ab}`$ and $`𝝎_b^a`$. Here $`_\xi `$ is the usual Lie derivative acting on forms and is given by $`_\xi =𝒅𝒊_\xi +𝒊_\xi 𝒅`$.
3. 3) $`\mathrm{𝑇ℎ𝑒}\mathrm{𝑃𝑟𝑜𝑗𝑒𝑐𝑡𝑖𝑣𝑒}\mathrm{𝑆𝑦𝑚𝑚𝑒𝑡𝑟𝑦}`$, parametrized by an arbitrary infinitesimal one-form $`𝜿=𝜿(x)`$:
$`\delta _\kappa 𝜽^a`$ $`=`$ $`\delta _\kappa g^{ab}=0`$
$`\delta _\kappa 𝝎_b^a`$ $`=`$ $`𝜿\delta _b^a`$ (15)
### Palatini and vielbein formalisms
As was shown in , the new feature here is that we need to fix the Projective symmetry (2) in what we called the Einstein gauge to recover the torsionless and metricity conditions from the equations of motion of the linear connection (12). The physics, namely Einstein equations, does not depend on this gauge choice.
Now, the Palatini formalism can be recovered after fixing all the $`gl(D,\text{I}\mathrm{R})`$ symmetry by the canonical choice $`\theta _\mu ^a=\delta _\mu ^a`$. On the other hand, the Cartan-Weyl (vielbein or orthonormal frames) formalism comes after choosing $`\theta _\mu ^a=e_\mu ^a`$, $`e_\mu ^a`$ being an orthonormal frame (i.e. $`g^{ab}=\eta ^{ab}`$, with $`\eta ^{ab}`$ the ordinary flat Minkowski metric). This last choice breaks the local $`gl(D,\text{I}\mathrm{R})`$ down to local<sup>6</sup><sup>6</sup>6In the Euclidian case, we fix $`g^{ab}=\delta ^{ab}`$ and the gauge group is broken to local $`so(D,\text{I}\mathrm{R})`$. $`so(D1,1,\text{I}\mathrm{R})`$.
The above gauge fixings can be generally rewritten as $`\theta _\mu ^a=\overline{\theta }_\mu ^a`$, with $`\overline{\theta }_\mu ^a(x)`$ being an arbitrary given frame. The residual gauge symmetry which preserves such a choice is a linear combination of a diffeomorphism and a $`gl(D,\text{I}\mathrm{R})`$ rotation such that:
$$_\xi \overline{𝜽}^a+\delta _\lambda \overline{𝜽}^a=0$$
(16)
The above equation gives the parameter $`\lambda _b^a`$ in terms of the diffeomorphism parameter $`\xi ^\rho `$ . Using now the identity $`_\xi \overline{𝜽}^a=𝑫\xi ^a+𝒊_\xi \overline{𝚯}^a𝒊_\xi 𝝎_b^a\overline{𝜽}^b`$ (where $`\xi ^a:=\xi ^\rho \theta _\rho ^a`$ and $`\overline{𝚯}^a`$ is the torsion (8) with $`𝜽^a=\overline{𝜽}^a`$ ) and the torsionless condition, equation (16) becomes:
$$\left(𝒊_\xi 𝝎_b^a\lambda _b^a\right)\overline{𝜽}^b=𝑫\xi ^a$$
(17)
In conclusion, after fixing the $`gl(D,\text{I}\mathrm{R})`$ gauge symmetry, the remaining symmetry is diffeomorphisms parametrized by $`\xi ^\rho `$ with simultaneous $`gl(D,\text{I}\mathrm{R})`$-rotations parametrized by $`\lambda _b^a[\xi ]`$ which satisfy (17) (or equivalently (16)).
Let us clarify the above point by a simple example: in the Palatini case ($`\theta _\mu ^a=\overline{\theta }_\mu ^a=\delta _\mu ^a`$), equation (16) (or (17)) simply gives $`\lambda _b^a=_b\xi ^a`$ (see also ). Using equations (13-14), we then check that the residual symmetry is the usual diffeomorphism symmetry:
$`\delta _{\xi ,\lambda }g^{ab}`$ $`=`$ $`_\xi g^{ab}+\lambda _c^ag^{cb}+\lambda _c^bg^{ac}`$
$`=`$ $`\xi ^c_cg^{ab}_c\xi ^ag^{cb}_c\xi ^bg^{ac}`$
$`\delta _{\xi ,\lambda }𝝎_b^a`$ $`=`$ $`_\xi 𝝎_b^a𝒅\lambda _b^a𝝎_d^a\lambda _b^d+𝝎_b^d\lambda _d^a`$
$`=`$ $`(\xi ^d_d\mathrm{\Gamma }_{cb}^a+_c\xi ^d\mathrm{\Gamma }_{db}^a+_c_b\xi ^a+\mathrm{\Gamma }_{cd}^a_b\xi ^d\mathrm{\Gamma }_{cb}^d_d\xi ^a)𝒅x^c`$
with $`𝝎_b^a:=\mathrm{\Gamma }_{cb}^a𝒅x^c`$.
### The Dirichlet and Neumann boundary conditions
The stationarity of the action constructed from the Hilbert Lagrangian (5) implies the following condition at spatial infinity<sup>7</sup><sup>7</sup>7For consistency, the variational principle is required to be satisfied at least near the boundary component where the superpotential and the conserved charges are computed .:
$$_\mathrm{\Sigma }_{\mathrm{}}\delta 𝝎_b^a\mathbf{}𝚺_a^b=0$$
(18)
where we used the compact notation:
$$𝚺_a^b:=\frac{\sqrt{\left|g\right|}}{4\kappa ^2}g^{bc}𝚺_{ac}$$
(19)
It is possible to add an appropriate total derivative $`𝒅`$$`𝑲`$ to the Lagrangian (5) in order to implement a given boundary condition and solve (18). The two generic cases are the following:
$`D:`$ $`{\displaystyle _\mathrm{\Sigma }_{\mathrm{}}}\delta 𝝎_b^a\mathbf{}\mathrm{\Delta }𝚺_a^b=0`$ (20)
$`N:`$ $`{\displaystyle _\mathrm{\Sigma }_{\mathrm{}}}\mathrm{\Delta }𝝎_b^a\mathbf{}\delta 𝚺_a^b=0`$ (21)
and the definitions,
$$\mathrm{\Delta }𝚺_a^b:=𝚺_a^b\overline{𝚺}_a^b,\mathrm{\Delta }𝝎_b^a:=𝝎_b^a\overline{𝝎}_b^a$$
(22)
with $`\overline{𝚺}_a^b`$ and $`\overline{𝝎}_b^a`$ our chosen asymptotic background fields<sup>8</sup><sup>8</sup>8The quantity $`\overline{𝚺}_a^b`$ is constructed with some chosen asymptotic background metric. For example, in the Palatini formalism, we use $`\theta _\mu ^a=\delta _\mu ^a`$ and $`g^{ab}=\overline{g}^{ab}`$ in definition (19). In the vielbein formalism, we fix at infinity the orthonormal frame by $`\theta _\mu ^a=\overline{e}_\mu ^a`$ in the gauge $`g^{ab}=\eta ^{ab}`$ (see above or ). whose arbitrary variation vanishes, $`\delta \overline{𝚺}_a^b=\delta \overline{𝝎}_b^a=0`$. Then, the equations (20-21) are obtained by adding respectively to the Lagrangian (5) the following surface terms $`𝒅(𝝎_b^a\mathbf{}\overline{𝚺}_a^b)`$ and $`𝒅(\mathrm{\Delta }𝝎_b^a\mathbf{}𝚺_a^b)`$.
We then recognize from (20-21) typical Dirichlet and Neumann boundary conditions: we respectively fix the asymptotic value of the metric (through $`\overline{𝚺}_a^b`$ ) or of the connection (through $`\overline{𝝎}_b^a`$). Any linear combination of (20-21) is also an appropriate boundary condition, compatible with a variational principle.
We will be mostly interested in the cases where the metric is asymptotically flat (or AdS) at spatial infinity. In that case, the appropriate boundary condition is given by (20), with $`\overline{𝚺}_a^b`$ constructed from the flat (or AdS) metric. A covariant way to check the vanishing of (20) is to use the compactification of Ashtekar and Romano for spatial infinity. For an asymptotically AdS space, we can use the usual compactification of Penrose (see also the recent ). We will not give the complete proof of these statements here. The case of null infinity is quite a bit more involved since neither (20) nor (21) are satisfied there . We shall elaborate on this in section 3.
### 2.2 The associated superpotentials
The superpotentials associated to the gauge symmetries (13-2) have been computed with the cascade equations in . The background contribution for the superpotential was however missing there. Let us emphasize that the background is only used at infinity, it is nothing but a covariant way to impose the required boundary conditions. Our purpose will now be to use formula (4) to get the complete result.
The total superpotential will receive contributions from the $`gl(D,\text{I}\mathrm{R})`$ (parametrized by $`\lambda _b^a`$) and diffeomorphism (parametrized by $`\xi ^\rho `$) gauge symmetries only<sup>9</sup><sup>9</sup>9The Projective symmetry (2) does not contribute to $`𝑾_\xi `$ (see definition (3) together with (1)). This is a consequence of the fact that no gauge field is associated to this symmetry. In other words, there is no field whose $`𝜿`$-transformation law is proportional to the derivative of the gauge parameter. Hence the associated current is identically zero .. Using the equations (10-2) in the definition (3) we obtain:
$$𝑾_{\xi ,\lambda }=𝒊_\xi 𝜽^a\frac{\delta 𝑳}{\delta 𝜽^a}+\left(𝒊_\xi 𝝎_b^a\lambda _b^a\right)\frac{\delta 𝑳}{\delta 𝝎_b^a}=𝑹_b^a\mathbf{}𝒊_\xi 𝚺_a^b\left(𝒊_\xi 𝝎_b^a\lambda _b^a\right)𝑫𝚺_a^b$$
(23)
where the shorthand notation (19) was used. And $`𝒊_\xi `$ denotes the interior product with the vector field $`\xi ^\rho `$.
Now from equation (4), the variation of the gravitational superpotential at some boundary $`_r`$ (and in particular at spatial infinity $`_{\mathrm{}}`$) is given by
$`\delta 𝑼_{\xi ,\lambda }`$ $`=`$ $`\delta 𝝎_b^a\mathbf{}𝒊_\xi 𝚺_a^b+\left(𝒊_\xi 𝝎_b^a\lambda _b^a\right)\delta 𝚺_a^b`$ (24)
$`=`$ $`\delta \left(\left(𝒊_\xi 𝝎_b^a\lambda _b^a\right)𝚺_a^b\right)𝒊_\xi \left(\delta 𝝎_b^a\mathbf{}𝚺_a^b\right)`$ (25)
$`=`$ $`\delta \left(D_b\xi ^a𝚺_a^b𝒊_\xi (𝝎_b^a\mathbf{}\overline{𝚺}_a^b)\right)𝒊_\xi \left(\delta 𝝎_b^a\mathbf{}\mathrm{\Delta }𝚺_a^b\right)`$ (26)
$`=`$ $`\delta \left(D_b\xi ^a𝚺_a^b𝒊_\xi \left(\mathrm{\Delta }𝝎_b^a\mathbf{}𝚺_a^b\right)\right)+𝒊_\xi \left(\mathrm{\Delta }𝝎_b^a\mathbf{}\delta 𝚺_a^b\right)`$ (27)
The first equation (24) follows from the criterion (4) applied to (23). The second one (25) follows from the first one after some simple algebraic manipulations assuming that the variation of the gauge parameters (namely $`\xi ^\rho `$ and $`\lambda _b^a`$) vanishes. The last two equations (26-27) are constructed from (25) such that the last term reproduces one of the two boundary conditions (20-21). We also used the result (17).
The next point is to integrate equation (26) or (27) using the Dirichlet or Neumann boundary condition (20-21). The equation (24) will be integrable iff
$$𝒊_\xi (\delta 𝝎_b^a\mathbf{}\delta 𝚺_a^b)=0$$
(28)
at spatial infinity. This integrability condition was derived apparently for the first time in within the covariant symplectic formalism (the term
$$\mathrm{\Omega }:=\delta 𝝎_b^a\mathbf{}\delta 𝚺_a^b$$
(29)
is nothing but the so-called pre-symplectic two-form). The correspondence between the covariant phase space formalism and equation (4) will be given in .
#### The (1/2) Komar superpotential
Let us first use the Neumann boundary condition on the metric (Dirichlet on the connection) given by (21) and integrate (24). If we assume that $`𝝎_b^a`$ approaches $`\overline{𝝎}_b^a`$ fast enough, the last two terms of (27) vanish. Then, up to some global constant, the superpotential is given by:
$$𝑼_\xi ^{\text{Ko}}=𝑫_b\xi ^a𝚺_a^b=\frac{\sqrt{\left|g\right|}}{4\kappa ^2}\mathbf{}(𝒅𝝃)$$
(30)
where $`𝝃`$ is the one-form associated to the vector field $`\xi ^\rho `$ (we also used the definition (19)).
We found the so-called Komar superpotential. However, the coefficient is not the usual one. If we compute the charge given by the above superpotential for the Schwarzschild black hole in $`D`$ spacetime dimensions, we will find $`\frac{D2}{D3}\times M`$ instead of $`M`$. The superpotential found by Komar was based on another one proposed by Møller . This author rescaled by hand the expression (30) to find the correct Schwarzschild mass for $`D=4`$.
Here, we are not allowed to change the natural normalization of (30), so this superpotential is not appropriate to compute the usual conserved charges at spatial infinity. This is not surprising because it must be derived using the boundary conditions (21) which are incompatible with asymptotic flatness (or AdS).
#### The KBL superpotential
Let us now use some Dirichlet boundary conditions on the metric (equation (20)) to integrate equation (24). In that case, the last term of equation (26) vanishes. Thus the superpotential is given by:
$$𝑼_\xi ^{\text{KBL}}=𝑼_\xi ^{\text{Ko}}𝒊_\xi (𝝎_b^a\mathbf{}\overline{𝚺}_a^b)C^t$$
(31)
Recall that equation (24) gives the superpotential up to some global constant $`C^t`$. A natural way to fix this constant is to require the vanishing of the superpotential when evaluated in the background metric $`\overline{g}^{ab}`$, that is $`𝑼_\xi ^{\text{KBL}}[\overline{g}]=0`$. We then find a $`D`$-dimensional version of the superpotential proposed by Katz, Bic̆ák and Lynden-Bell :
$`𝑼_\xi ^{\text{KBL}}`$ $`=`$ $`𝑼_\xi ^{\text{Ko}}𝑼_\xi ^{\text{Ko}}[\overline{g}]𝒊_\xi (\mathrm{\Delta }𝝎_b^a\mathbf{}\overline{𝚺}_a^b)`$ (32)
$`=`$ $`𝑼_\xi ^{\text{Ko}}𝑼_\xi ^{\text{Ko}}[\overline{g}]𝒊_\xi (\mathrm{\Delta }𝝎_b^a\mathbf{}𝚺_a^b)`$
where equation (20) (with $`\delta 𝝎_b^a=\mathrm{\Delta }𝝎_b^a`$) was again used in the second line.
It is straightforward to check that the Hodge-dual of the above $`(D2)`$-form is in components equal to (see for related calculation):
$$^{\text{KBL}}U_\xi ^{\mu \nu }=^{\text{Ko}}U_\xi ^{\mu \nu }^{\text{Ko}}\overline{U}_\xi ^{\mu \nu }+S^\mu \xi ^\nu S^\nu \xi ^\mu $$
(33)
with
$$^{\text{Ko}}U_\xi ^{\mu \nu }=\frac{\sqrt{\left|g\right|}}{4\kappa ^2}(^\mu \xi ^\nu ^\nu \xi ^\mu )\text{ and }S^\mu =\frac{\sqrt{\left|g\right|}}{4\kappa ^2}(\mathrm{\Delta }\mathrm{\Gamma }_{\rho \sigma }^\mu g^{\rho \sigma }\mathrm{\Delta }\mathrm{\Gamma }_{\rho \sigma }^\sigma g^{\mu \rho }).$$
(34)
In general, this superpotential<sup>10</sup><sup>10</sup>10The result (33) without the background contribution was derived in (and called the Katz superpotential ) using the cascade equations techniques. depends on the gauge parameter $`\xi ^\rho `$, which is not arbitrary but has to be compatible with the asymptotic boundary condition . In the case of Dirichlet boundary conditions (20), this parameter should be an asymptotic Killing vector, that is $`g_{\mu \nu }+\delta _\xi g_{\mu \nu }\overline{g}_{\mu \nu }`$ as $`r\mathrm{}`$.
Our derivation of the KBL superpotential is independent of the space-time dimension, of the first order theory used (Palatini or Cartan-Weyl), and of the presence of a cosmological constant<sup>11</sup><sup>11</sup>11If we add a cosmological constant to the Hilbert Lagrangian (5) by $`𝑳_\mathrm{\Lambda }=\mathrm{\Lambda }\sqrt{\left|g\right|}𝚺`$, the calculations leading to the result (33) remain almost unchanged.. It follows straightforwardly from equation (4) and from Dirichlet boundary conditions on the metric. This is to be contrasted with the more involved original derivation and our derivation in both with the Noether method where some prescription was needed in order to fix the ambiguity in the current $`𝑱_\xi `$. The comparison between our present derivation and the covariant phase space method used recently in will be given in . Moreover, the use of the Affine formalism for gravity gives directly the term proportional to $`_\mu \xi ^\rho `$ in the KBL superpotential which was added by hand in following some “covariant criterion”.
Some of the properties of the KBL superpotential, and the validity of equation (4) at null infinity are discussed in the last section.
## 3 The KBL superpotential at null infinity
As shown in the previous section, the KBL superpotential follows straightforwardly from the diffeomorphism invariance of general relativity, through equation (4) and Dirichlet boundary conditions (20). Moreover, it is the only superpotential which satisfies the following properties:
1. $``$ It is generally covariant and then can be computed in any coordinate system.
2. $``$ If the chosen coordinates are the Cartesian ones of an asymptotically flat (or AdS) spacetime, the KBL superpotential gives the ADM mass formula (or the AD mass ).
3. $``$ It gives the mass and angular momentum (and the Brown and Henneaux conformal charges for AdS<sub>3</sub>) with the right normalization in any spacetime dimensions $`D3`$. More generally, it can be used for any asymptotic Killing vector $`\xi ^\rho `$.
We just derived the KBL superpotential at spatial infinity. The important point is that it depends on some background metric $`\overline{g}_{\mu \nu }`$. This is crucial for its general covariance.
The asymptotic background metric is not a new object of spacetime. In fact, it fixes the boundary conditions, namely $`g_{\mu \nu }\overline{g}_{\mu \nu }`$ as $`r\mathrm{}`$. It is imposed naturally at the boundary of spacetime. It is however clear that in general it cannot be extended arbitrarily everywhere in the bulk, for instance flat space may have a different topology than our solution sector. The choice of such a $`\overline{g}_{\mu \nu }`$ everywhere is unnatural from the background-independent Einstein theory point of view. Moreover, there are many ways to define this background metric $`\overline{g}_{\mu \nu }`$ in the bulk such that its asymptotic value agrees with our chosen boundary conditions. It is not clear why one of these choices would be better than the others. If we cannot properly define $`\overline{g}_{\mu \nu }`$ everywhere then we will not be able to define $`{}_{}{}^{\text{KBL}}U_{\xi }^{\mu \nu }`$ in the bulk and interpret the KBL superpotential as an expression for a quasi-local mass. We believe the problem of quasi-local charges is ill posed and needs to be supplemented by specific boundary assumptions at the surface to be used for enveloping the physical object.
Suppose now that our spacetime is also flat at future (past) null infinity $`I^+`$ ($`I^{}`$). In that case, we can uniquely extend the background metric $`\overline{g}_{\mu \nu }`$ to that region. The KBL superpotential can then be covariantly defined also on $`I^+`$. In particular, we can integrate equation (2) on a piece of $`I^+`$, namely $`\mathrm{\Delta }`$ , bounded by spatial infinity $`i_0`$ and some time-dependent cross section $`𝒞`$ . The on-shell result is
$$_\mathrm{\Sigma }_{\mathrm{}}𝑼_\xi ^{\text{KBL}}=_𝒞𝑼_\xi ^{\text{KBL}}+_\mathrm{\Delta }J_\xi ^{\text{KBL}}$$
(35)
with $`J_\xi ^{\text{KBL}}=𝒅𝑼_\xi ^{\text{KBL}}`$ (on-shell).
It has been proved by Katz and Lerer that the first term in the rhs of (35) reproduces the Bondi mass , the Penrose linear momentum and the Penrose-Dray-Streubel angular momentum at null infinity (depending on the choice of the asymptotic $`\xi ^\rho `$). The second integral in the rhs of (35) is then nothing but the total amount of charge which crossed $`\mathrm{\Delta }`$. In the case where $`\xi ^\rho `$ is an asymptotic BMS translation , we recover the result pointed out in .
As will be shown in the forecoming paper , the equation (4) used to compute the superpotential is equivalent to the covariant symplectic phase space methods . In particular, it was shown by Wald and Zoupas that equation (4) is non integrable at null infinity. This is due in part to the fact that the use of equation (4) is justified only if the charge is conserved , that is, if the flux of the Noether current vanishes on the spacetime boundary $`_r`$ under consideration. This is not the case in general for the Bondi-type charges. The point is that the equation for the variation of the superpotential can be derived from the vanishing of the flux at infinity provided the set of boundary conditions is “Lagrangian” ie provided it leads to the vanishing of the symplectic form (29) there.
As a final comment, it seems quite natural that a formula like (35) should also exist on the horizon $``$ of a black hole. More precisely, if we consider $``$ as a boundary of spacetime, with some boundary conditions on it, we would hope to find an associated superpotential, and then, some conserved or Bondi-type charges depending on its dynamical behavior. The work of Ashtekar et al. goes precisely in that direction. So at least for isolated horizons the definition of a quasi-local mass on the horizon (even in the nonstationary case) should go through beyond the four dimensional case they have considered.
Acknowledgments. We are grateful to AEI for hospitality, and to A. Ashtekar, J. Bic̆ák and M. Henneaux for discussions.
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# Dirac operators and conformal invariants of tori in 3-space
## 1 Introduction
In this paper we show how to assign to any torus immersed into the three-space $`^3`$ or the unit three-sphere $`S^3`$ a complex curve such that the immersion is described by functions defined on this curve (a Riemann surface which is generically of infinite genus). We call this curve the spectrum of a torus (with a fixed conformal parameter). This spectrum has many interesting properties and, in particular, relates to the Willmore functional whose value is encoded in it.
The construction of the spectra for tori in $`^3`$ was briefly explained in . In this text we do that also for immersed tori in $`S^3`$.
Our conjecture that the spectrum of a torus in $`^3`$ is invariant under conformal transformations of $`^3`$ was proven modulo some analytic facts by Grinevich and Schmidt . In fact, their proof is rather physical which may be expected because the construction of the spectrum originates in soliton theory.
In this text we give a complete proof of the conformal invariance of the spectra for isothermic tori. This case already covers many interesting surfaces such as constant mean curvature tori and tori of revolution in $`^3`$.
Some spectral curves of finite genus already appeared in studies of harmonic tori in $`S^3`$ by Hitchin and constant mean curvature (CMC) tori in $`^3`$ by Pinkall and Sterling . It was shown that such tori are expressed in terms of algebraic functions corresponding to these complex curves . We show that for minimal tori in $`S^3`$ and CMC tori in $`^3`$ these spectral curves are particular cases of the general spectrum.
The general construction is based on the global Weierstrass representation of closed surfaces introduced in and a general construction of the Floquet(or Bloch) variety for a periodic differential operator. An existence of this variety is derived from the Keldysh theorem but an effective construction which gives more information about analytic behavior of this complex curve was proposed by Krichever in who used perturbation methods. It is as follows. Take an immersed torus with the induced metric $`e^{2\alpha }dzd\overline{z}`$ and consider differential operators
$$𝒟=\left(\begin{array}{cc}0& \\ \overline{}& 0\end{array}\right)+\left(\begin{array}{cc}U& 0\\ 0& U\end{array}\right)\text{with }U=\frac{1}{2}He^\alpha \text{ for a torus in }^3,$$
$$𝒟^S=\left(\begin{array}{cc}0& \\ \overline{}& 0\end{array}\right)+\left(\begin{array}{cc}V& 0\\ 0& \overline{V}\end{array}\right)\text{with }V=\frac{1}{2}(Hi)e^\alpha \text{ for a torus in }S^3,$$
where $`H`$ is the mean curvature.
Let $`\mathrm{\Lambda }`$ be the period lattice of a torus which means that the torus is an immersion of $`/\mathrm{\Lambda }`$ with a conformal parameter $`z`$ on it. Take a basis $`\gamma _1,\gamma _2`$ for $`\mathrm{\Lambda }`$ which is also considered as a basis for $`H_1(T^2)\mathrm{\Lambda }`$. Now consider all solutions $`\psi `$ to the equations
$$𝒟\psi =0\text{or}𝒟^S\psi =0$$
satisfying the following conditions
$$\psi (z+\gamma _j)=\mu _j\psi (z).$$
These are Floquet(–Bloch) functions and the pairs $`(\mu _1,\mu _2)`$ form a complex curve in $`^2`$. This is the Floquet zero-level spectrum of $`𝒟`$ and, by the definition, this is the spectrum of the immersed torus. The analytic properties of this curve are described by Pretheorem which is a modification of its analogs for two-dimensional scalar Schrödinger and heat operators proven in and it is clear that the proof of Pretheorem may be obtained by slight modifications of the reasonings of .
One of the most interesting properties of this construction is its relation to a conformal geometry and the Willmore functional which equals
$$4_{/\mathrm{\Lambda }}U^2𝑑xdy\text{or}4_{/\mathrm{\Lambda }}|V|^2𝑑xdy$$
for tori in $`^3`$ or $`S^3`$, where $`z=x+iy`$.
The global Weierstrass representation of closed surfaces represents any closed surface $`\mathrm{\Sigma }`$ in terms of a solution to the equation $`𝒟\psi =0`$ (a harmonic spinor) where $`\psi `$ takes values in some bundle over the constant curvature surface $`\mathrm{\Sigma }_0`$ which is is conformally equivalent to $`\mathrm{\Sigma }`$ (see Theorems 1–3 in 2.3) . The Willmore functional $`_\mathrm{\Sigma }(H^2K)𝑑\mu =4_{\mathrm{\Sigma }_0}U^2𝑑xdy2\pi \chi (\mathrm{\Sigma })`$ measures the $`L_2`$-norm of the potential $`U`$ of the surface. For small values of this functional the equation $`𝒟\psi =0`$ does not admit solutions which describe closed surfaces in $`^3`$ and therefore that explains physical meaning of lower bounds for the Willmore functional proposed by the Willmore conjecture and its generalizations. This gives a hint that the spectral properties of $`𝒟`$ have to have a geometric meaning. We shall discuss that in details in 4.4.
In it was established that the dimension of the kernel of $`𝒟`$ gives lower estimates for the Willmore functional for spheres. Actually we proved for spheres with one-dimensional potentials $`U`$ (examples of them are spheres of revolution but not only) and conjectured for all spheres the following inequality
$$_{\mathrm{\Sigma }_0}U^2𝑑xdy4\pi N^2\text{with }N=dim_{}\mathrm{ker}𝒟\text{.}$$
(1)
The proof of it is based on the inverse scattering problem for the one-dimensional Dirac operator and it also works for general Dirac operators with $`S^1`$-symmetry (i.e., with one-dimensional potentials) on special spinor bundles over the 2-sphere. This inequality can not be improved and the equality is achieved on “soliton spheres” .
Another treatment of the global representation belongs to Pedit and Pinkall who proposed to consider spinor $`^2`$-bundles introduced in as quaternionic line bundles and consider harmonic spinors as holomorphic quaternionic sections of such bundles. This enables them to apply ideas of algebraic geometry to surface theory and to generalize this representation for surfaces in $`^4`$ . Very recently they managed to relate (1) to the quaternionic analog of the Plücker formula and by that prove our conjecture, i.e., establish the inequality (1) for all spheres together with its generalizations for higher genus surfaces.
This paper is organized as follows.
In section 2 we recall the notion of the Weierstrass representation.
In section 3 we prove that the multipliers $`(\mu _1,\mu _2)`$ of Floquet functions form a spectral curve in $`^2`$ and discuss its analytic properties.
In section 4 we show how to assign such a spectrum to an immersed torus in $`^3`$, show the Willmore functional appear in this picture and discuss the modern state of the Willmore conjecture.
In section 5 we show that the spectra of CMC and isothermic tori are particular cases of the spectrum defined in section 4 and also prove our conjecture that the spectra of an isothermic torus and its dual surface coincide .
In section 6 we show how to assign a spectral curve to a torus in $`S^3`$ and prove that for a minimal torus in $`S^3`$ it coincides with a spectral curve defined by Hitchin for harmonic tori in $`S^3`$ .
In section 7 we prove that the spectrum of an isothermic torus in $`S^3`$ is invariant with respect to conformal transformations of $`\overline{}^3=^3\{\mathrm{}\}`$ preserving the torus in $`^3`$.
## 2 The Weierstrass representation
### 2.1 Basic equations of surface theory in $`^3`$
In this subsection we recall the main definitions and some well-known important facts from the classical surface theory.
Let $`𝒰`$ be a domain in $`^2`$, with coordinates $`(x^1,x^2)`$, regularly immersed into $`^3`$:
$$F:𝒰^3.$$
At every point $`p𝒰`$ the vectors
$$F_1=\frac{F}{x^1},F_2=\frac{F}{x^2},N=\frac{[F_1\times F_2]}{|F_1||F_2|}$$
form a linear basis $`\sigma =(F_1,F_2,N)^{}`$ for $`^3`$, where $`F_1`$ and $`F_2`$ are tangent vectors to the surface $`\mathrm{\Sigma }=F(𝒰)`$, and $`N`$ is a unit normal vector. The variables $`(x^1,x^2)`$ are local coordinates on $`\mathrm{\Sigma }`$ and the induced metric on it is
$$𝐈=g_{kl}dx^kdx^l\text{with }g_{kl}=F_k,F_l\text{ (the }\text{first fundamental form}\text{)}$$
where $`a,b=a_1b_1+a_2b_2+a_3b_3`$.
The derivatives of the basic vectors are expanded in $`\sigma `$ as
$$\frac{^2F}{x^kx^l}=\mathrm{\Gamma }_{kl}^j\frac{F}{x^j}+b_{kl}N,\frac{N}{x^k}=b_k^j\frac{F}{x^j}$$
(2)
where $`\mathrm{\Gamma }_{kl}^j`$ are the Christoffel symbols, $`\mathrm{𝐈𝐈}=b_{kl}dx^kdx^l`$ is the second fundamental form, and $`b_k^j=g^{jl}b_{lk}`$. The equations (2) are the Gauss–Weingarten derivation equations and have the form
$$\frac{\sigma }{x^1}=𝐔\sigma ,\frac{\sigma }{x^2}=𝐕\sigma ,$$
(3)
where $`U`$ and $`V`$ are $`(3\times 3)`$-matrices. The compatibility conditions for (3) are the Codazzi equations:
$$\frac{^2\sigma }{x^1x^2}\frac{^2\sigma }{x^2x^1}=\left(\frac{𝐔}{x^2}\frac{𝐕}{x^1}+[𝐔,𝐕]\right)\sigma =0,$$
which are equivalent to the zero-curvature equations
$$\frac{𝐔}{x^2}\frac{𝐕}{x^1}+[𝐔,𝐕]=0$$
(4)
for the connection $`(/x^1𝐔,/x^2𝐕)`$.
At every point $`p`$ of the surface the fundamental forms are diagonalized as
$$𝐈(p)=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\mathrm{𝐈𝐈}(p)=\left(\begin{array}{cc}k_1& 0\\ 0& k_2\end{array}\right)$$
and the principal curvatures $`k_1`$ and $`k_2`$ satisfy the equation $`det\left(b_{kl}kg_{kl}\right)=0`$, which divided by $`detg_{jk}`$ takes the form $`k^22Hk+K=0`$ where $`H`$ is the mean curvature and $`K`$ is the Gaussian curvature:
$$H=\frac{k_1+k_2}{2},K=k_1k_2.$$
A point $`p`$ is called an umbilic point if the principal curvatures coincide at $`p`$: $`k_1=k_2`$, which is equivalent to $`H^2K=0`$.
Let $`z=x^1+ix^2`$ be a conformal parameter on the surface, i. e., the first fundamental form is
$$𝐈=e^{2\alpha (z,\overline{z})}dzd\overline{z},$$
which means that
$$F_z,F_z=F_{\overline{z}},F_{\overline{z}}=0,F_z,F_{\overline{z}}=\frac{1}{2}e^{2\alpha }.$$
The family $`\stackrel{~}{\sigma }=(F_z,F_{\overline{z}},N)^{}`$ is a basis for $`^3`$ and the Gauss–Weingarten equations are written as
$$\stackrel{~}{\sigma }_z=\stackrel{~}{𝐔}\sigma ,\stackrel{~}{\sigma }_{\overline{z}}=\stackrel{~}{𝐕}\sigma ,$$
with
$$\stackrel{~}{𝐔}=\left(\begin{array}{ccc}2\alpha _z& 0& A\\ 0& 0& B\\ 2e^{2\alpha }B& 2e^{2\alpha }A& 0\end{array}\right),\stackrel{~}{𝐕}=\left(\begin{array}{ccc}0& 0& B\\ 0& 2\alpha _{\overline{z}}& \overline{A}\\ 2e^{2\alpha }\overline{A}& 2e^{2\alpha }B& 0\end{array}\right),$$
$`A=F_{zz},N`$, and $`B=F_{z\overline{z}},N`$. The second fundamental form equals
$$\mathrm{𝐈𝐈}=(2B+(A+\overline{A}))(dx^1)^2+2i(A\overline{A})dx^1dx^2+(2B(A+\overline{A}))(dx^2)^2,$$
and we have
$$H=2Be^{2\alpha },K=4(B^2A\overline{A})e^{4\alpha }.$$
Now the Codazzi equations $`\stackrel{~}{𝐔}_{\overline{z}}\stackrel{~}{𝐕}_z+[\stackrel{~}{𝐔},\stackrel{~}{𝐕}]=0`$ take the form
$$\alpha _{z\overline{z}}+e^{2\alpha }(B^2A\overline{A})=0,A_{\overline{z}}B_z+2\alpha _zB=0.$$
(5)
The first of them is the Gauss egregium theorem and another equation
$$A_{\overline{z}}=\frac{1}{2}H_ze^{2\alpha }$$
splits into two real-valued equations.
A quadratic differential $`\omega =Adz^2`$ is called the Hopf differential and has important geometrical properties. For instance, $`\omega `$ vanishes at a point if and only if this is an umbilic point.
It is said that the Gauss map of a surface $`\mathrm{\Sigma }=F(𝒰)`$
$$G:\mathrm{\Sigma }S^2,G(p)=N(p),$$
is harmonic if $`\mathrm{\Delta }G(p)=\lambda (p)N(p)`$, where $`\mathrm{\Delta }`$ is the Laplace–Beltrami operator: $`\mathrm{\Delta }=4e^{2\alpha }\overline{}`$. We have $`\mathrm{\Delta }F=2HN`$. Since $`N_{z\overline{z}}=2e^{2\alpha }(\overline{A}_zF_z+A_{\overline{z}}F_{\overline{z}}+(A\overline{A}+B^2)N)`$, we conclude that
the Gauss map $`G`$ is harmonic if and only if the Hopf differential $`\omega `$ is holomorphic, which, by (5), is equivalent to $`H_z=H_{\overline{z}}=0`$, i.e., $`H=\text{const}`$ and $`\mathrm{\Sigma }`$ is a constant mean curvature (CMC) surface .
There are two other important classes of surfaces:
1) a surface is called minimal if $`H=0`$, which is equivalent to $`F_{z\overline{z}}=0`$;
2) a surface is called isothermic if there is a conformal parameter on it such that $`\text{Im}A=0`$.
### 2.2 The local representation of a surface
In this subsection we follow .
Denote by $`𝒬`$ a quadric in $`^3`$ defined by the equation
$$Z_1^2+Z_2^2+Z_3^2=0,Z=(Z_1,Z_2,Z_3)^3.$$
For a conformal parameter $`z`$ on a surface $`\mathrm{\Sigma }`$ there is a mapping
$$f:𝒰\stackrel{F}{}\mathrm{\Sigma }𝒬\text{where }f(p)=F_z(p),$$
(6)
satisfying the conditions
$$\text{Im}\frac{f}{\overline{z}}=0.$$
(7)
It is clear that any mapping $`f:𝒰𝒬`$ satisfying (7) has the form $`f=_z\mathrm{\Phi }`$ for some real-valued function $`\mathrm{\Phi }`$ and therefore has the form (6) for some surface.
The set $`𝒬`$ is parameterized by $`(\phi _1,\phi _2)^2`$ as follows:
$$Z_1=\phi _1^2\phi _2^2,Z_2=i(\phi _1^2+\phi _2^2),Z_3=2\phi _1\phi _2.$$
For simplicity, renormalize $`\phi `$ as follows:
$$\psi _1=\sqrt{\frac{2}{i}}\phi _1,\psi _2=\sqrt{\frac{2}{i}}\overline{\phi }_2.$$
In terms of $`\psi `$ the equations (7) are
$$𝒟\psi =0$$
(8)
where
$$𝒟=\left(\begin{array}{cc}0& \\ \overline{}& 0\end{array}\right)+\left(\begin{array}{cc}U& 0\\ 0& U\end{array}\right)$$
(9)
is a Dirac operator with a real-valued potential $`U(z,\overline{z})`$.
Now, if $`F(p_0)=(x_0^1,x_0^2,x_0^3)^3`$, then the surface is described by the Weierstrass formulas
$$x^1(p)=x^1(p_0)+_{p_0}^p\left(\frac{i}{2}(\overline{\psi }_2^2+\psi _1^2)dz\frac{i}{2}(\psi _2^2+\overline{\psi }_1^2)d\overline{z}\right),$$
$$x^2(p)=x^2(p_0)+_{p_0}^p\left(\frac{1}{2}(\overline{\psi }_2^2\psi _1^2)dz+\frac{1}{2}(\psi _2^2\overline{\psi }_1^2)d\overline{z}\right),$$
(10)
$$x^3(p)=x^3(p_0)+_{p_0}^p\left(\psi _1\overline{\psi }_2dz+\overline{\psi }_1\psi _2\right),$$
which are just
$$F(p)=F(p_0)+_{p_0}^p(fdz+\overline{f}d\overline{z}).$$
These local formulas in different forms were known before (see, for instance, and comments in ) but this form was introduced by Konopelchenko who working with the formulas of Eisenhart elaborated them into a form most convenient for applications. He considered them for constructing some surfaces via solutions to (8) and defining soliton deformations of “induced surfaces” but as we see this gives a general local construction of surfaces.
By straightforward computations it is derived that
$$U=\frac{He^\alpha }{2},e^\alpha =|\psi _1|^2+|\psi _2|^2.$$
(11)
and we see that for $`H=0`$ the formulas (10) reduce to the classical formulas for minimal surfaces.
Compute $`N=e^\alpha (i(\overline{\psi }_1\overline{\psi }_2\psi _1\psi _2),(\overline{\psi }_1\overline{\psi }_2+\psi _1\psi _2),|\psi _2|^2|\psi _1|^2)`$ and taking (8) into account derive
$$A=F_{zz},N=\psi _{1z}\overline{\psi }_2\overline{\psi }_{2z}\psi _1,B=F_{z\overline{z}},N=Ue^\alpha .$$
(12)
The Gauss–Weingarten equations written in terms of $`\psi `$ describe the deformations of $`\psi `$ and the first half of them is just the Dirac equation (8). For obtaining another half of equations differentiate $`e^\alpha `$ by $`z`$:
$$\alpha _ze^\alpha =\overline{\psi }_1\psi _{1z}+\psi _2\overline{\psi }_{2z},$$
and, taking (12) into account, obtain
$$\psi _{1z}=\alpha _z\psi _1+Ae^\alpha \psi _2,\psi _{2\overline{z}}=\overline{A}e^\alpha \psi _1+\alpha _{\overline{z}}\psi _2.$$
Now the Gauss–Weingarten equations are written as
$$\left[\frac{}{z}\left(\begin{array}{cc}\alpha _z& Ae^\alpha \\ U& 0\end{array}\right)\right]\psi =\left[\frac{}{\overline{z}}\left(\begin{array}{cc}0& U\\ \overline{A}e^\alpha & \alpha _{\overline{z}}\end{array}\right)\right]\psi =0$$
(13)
and the compatibility conditions for them, the Codazzi equations, are
$$A_{\overline{z}}=(U_z\alpha _zU)e^\alpha ,\alpha _{z\overline{z}}+U^2A\overline{A}e^{2\alpha }=0.$$
(14)
In fact, the equation (8) is already the compatibility condition for an existence of a surface with the Gauss map given by $`\psi `$. The other half of the equations (13) follows from it.
In a paper by Friedrich this representation was explained by classical means of the theory of Dirac operators. We also would like to mention a paper by Matsutani where it was considered from the physical point of view.
### 2.3 The global Weierstrass representation
Here we follow where this global representation was introduced.
Any closed oriented surface $`\mathrm{\Sigma }`$ in $`^3`$ is conformally equivalent to a constant curvature surface $`\mathrm{\Sigma }_0`$ and a choice of a conformal parameter $`z`$ on $`\mathrm{\Sigma }`$ means that a conformal equivalence $`\mathrm{\Sigma }_0\mathrm{\Sigma }`$ is fixed.
To define a compact oriented surface globally via the formulas (10) we have to introduce fibre bundles over surfaces and Dirac operators on them.
Consider two cases:
1) Tori. Let $`\mathrm{\Sigma }`$ be a torus immersed into $`^3`$. Then it is conformally equivalent to a flat torus $`\mathrm{\Sigma }_0=/\mathrm{\Lambda }`$ and $`z`$ is a conformal parameter. The vector function $`\psi `$ is expanded to a section of a $`^2`$-fiber bundle over $`\mathrm{\Sigma }`$ defined by the monodromy rules
$$\psi (z+\gamma )=\epsilon (\gamma )\psi (z)$$
(15)
where $`\gamma \mathrm{\Lambda }`$ and $`\epsilon :\mathrm{\Lambda }\{\pm 1\}`$ is a character of $`\mathrm{\Lambda }`$, i. e., a homomorphism to $`\{\pm 1\}`$. The Dirac operator $`𝒟`$ acts on this bundle and
$$U(z+\gamma )=U(z).$$
(16)
Hence, we have
###### Theorem 1
$`(`$$`)`$ The formulas (15) and (16) define a $`^2`$ bundle over a flat torus $`\mathrm{\Sigma }_0`$. To any section $`\psi `$, of this bundle, satisfying the Dirac equation (8) corresponds a surface in $`^3`$ defined by the formulas (10) up to translations in $`^3`$.
2) Surfaces of genus $`g>1`$. Let $`\mathrm{\Sigma }_0`$ be a hyperbolic surface conformally equivalent to a surface $`\mathrm{\Sigma }`$ immersed into $`^3`$ and $`z`$ be a conformal parameter. The surface $`\mathrm{\Sigma }_0`$ is isometric to $`/\mathrm{\Lambda }`$, where $``$ is the Lobachevsky upper-half plane and $`\mathrm{\Lambda }`$ is a discrete subgroup of $`PSL(2,)`$.
Any element $`\gamma \mathrm{\Lambda }PSL(2,)`$ is represented by elements
$$\pm \left(\begin{array}{cc}a& b\\ c& d\end{array}\right),a,b,c,d,adbc=1.$$
The action on $``$ is
$$z\gamma (z)=\frac{az+b}{cz+d}.$$
Define over $`\mathrm{\Sigma }_0`$ a $`^2`$-bundle by the monodromy rules
$$\psi _1(\gamma (z))=(cz+d)\psi _1(z),\psi _2(\gamma (z))=(c\overline{z}+d)\psi _2(z).$$
(17)
The Dirac operator acts on this bundle and
$$U(\gamma (z))=|cz+d|^2U(z).$$
(18)
###### Theorem 2
$`(`$$`)`$ The formulas (17) and (18) define a $`^2`$-bundle over a hyperbolic surface $`\mathrm{\Sigma }_0`$. To any section $`\psi `$, of this bundle, satisfying the Dirac equation (8) corresponds a surface in $`^3`$ defined by the formulas (10) up to translations in $`^3`$.
Notice that $`\psi _1\sqrt{dz}`$ and $`\overline{\psi }_2\sqrt{dz}`$ are defined modulo $`\pm 1`$. In fact, they are spinors and therefore we shall call $`\psi _1`$ and $`\psi _2`$ also spinors. There are $`2^{2g}`$ such nonequivalent spinor bundles over $`\mathrm{\Sigma }_0`$ where $`g`$ is the genus of $`\mathrm{\Sigma }_0`$.
This representation of compact oriented surfaces in $`^3`$ via solutions of Dirac equations (i.e., harmonic spinors) in spinor bundles over constant curvature surfaces is called the Weierstrass representation of surfaces.
###### Theorem 3
$`(`$ for real-analytic surfaces, for $`C^3`$-regular surfaces$`)`$ Every smooth closed oriented surface in $`^3`$ has a Weierstrass representation.
We see from the direct construction in 2.2 that for a surface with a fixed conformal parameter such a representation is unique. This gives rise to the following definition.
###### Definition 1
Let $`(\mathrm{\Sigma },z)`$ be an immersed surface with a fixed conformal parameter. Then the potential $`U`$ of its Weierstrass representation is called the potential of a surface.
To any harmonic spinor on $`\mathrm{\Sigma }_0`$ with a potential $`U`$ there corresponds a surface whose Gauss map descends through $`\mathrm{\Sigma }_0`$. The criterion of closedness of such a surface is as follows.
###### Proposition 1
A surface represented by a harmonic spinor $`\psi `$ over a compact surface $`\mathrm{\Sigma }_0`$ is closed if and only if
$$_{\mathrm{\Sigma }_0}\overline{\psi }_1^2𝑑\overline{z}\omega =_{\mathrm{\Sigma }_0}\psi _2^2𝑑\overline{z}\omega =_{\mathrm{\Sigma }_0}\overline{\psi }_1\psi _2𝑑\overline{z}\omega =0$$
for any holomorphic differential on $`\mathrm{\Sigma }_0`$.
This proposition was proved by M. Schmidt for general tori (in this case $`\omega =\mathrm{const}dz`$) and by the author for higher genera surfaces.
One of the most important properties of this representation is the equality
$$4_{\mathrm{\Sigma }_0}U^2𝑑xdy=_\mathrm{\Sigma }H^2𝑑\mu $$
where $`d\mu `$ is the measure given by the induced metric on $`\mathrm{\Sigma }`$ . The functional
$$𝒲(\mathrm{\Sigma })=_\mathrm{\Sigma }(H^2K)𝑑\mu =_\mathrm{\Sigma }\left(\frac{k_1k_2}{2}\right)^2𝑑\mu $$
is called the Willmore functional. By the Gauss–Bonnet theorem, for closed oriented surfaces it equals
$$𝒲(\mathrm{\Sigma })=_\mathrm{\Sigma }H^2𝑑\mu 2\pi \chi (\mathrm{\Sigma })$$
where $`\chi (\mathrm{\Sigma })`$ is the Euler characteristic of the surface and, therefore, for tori $`𝒲=_\mathrm{\Sigma }H^2𝑑\mu `$.
There is a famous Willmore conjecture that for tori the Willmore functional is greater or equal than $`2\pi ^2`$. We shall discuss it in 4.4.
We see the main advantage of the global Weierstrass representation in using the spectral properties of $`𝒟`$ for study of conformal geometry of surfaces. The present paper is devoted to developing this idea for tori.
## 3 The Floquet spectrum of a periodic operator
### 3.1 Floquet functions and the spectral curve
Let $`L`$ be a differential operator acting on functions or vector functions on $`^n`$, whose coefficients are periodic with respect to translations by vectors from a lattice $`\mathrm{\Lambda }`$, i. e., $`\mathrm{\Lambda }`$-periodic. We assume that $`\mathrm{\Lambda }`$ has the maximal rank, which means that it is isomorphic to $`^n`$ and $`^n/^n`$ is a torus.
To any vector $`\gamma \mathrm{\Lambda }`$ there corresponds a translation operator $`T_\gamma `$:
$$T_\gamma f(x)f(x+\gamma ),$$
Since $`L`$ is $`\mathrm{\Lambda }`$-periodic, if $`Lf=\lambda f`$, then $`LT_\gamma f=\lambda T_\gamma f`$. Moreover
$$[T_\gamma ,L]=0$$
and there are joint eigenfunctions of these commuting operators. Such functions are called Floquet (or Bloch) functions. The rigorous definition is as follows.
###### Definition 2
A function $`f:^n`$ is called a Floquet function of a $`\mathrm{\Lambda }`$-periodic operator $`L`$ if
$$Lf=Ef\text{and}f(x+\gamma )=e^{2\pi ik,\gamma }f(x)$$
for $`\gamma \mathrm{\Lambda }`$. The quantities $`k_1,\mathrm{},k_n`$ are called the quasimomenta of $`f`$.
Any Floquet function defines the multiplier homomorphism $`\mu :\mathrm{\Lambda }`$:
$$f(x+\gamma )=\mu (\gamma )f(x).$$
Consider the case, when $`L=𝒟`$, the two-dimensional Dirac operator
$$𝒟=\left(\begin{array}{cc}0& \\ \overline{}& 0\end{array}\right)+\left(\begin{array}{cc}U& 0\\ 0& U\end{array}\right)$$
with a double-periodic continuous potential $`U(z)`$ where $`z=x^1+ix^2`$.
###### Theorem 4
There is an analytic set $`Q(𝒟)^3`$ of positive codimension such that there exists a Floquet function of $`𝒟`$ with the quasimomenta $`k=(k_1,k_2)`$ and the eigenvalue $`E`$ if and only if $`(k_1,k_2,E)Q(𝒟)`$.
Its intersection with the plane $`\lambda =0`$ is an analytic set $`Q_0(𝒟)`$ of complex dimension one.
Proof. Take a constant $`C`$ such that the operator $`𝒟+C`$ is invertible on $`L_2(T^2)=L_2(/\mathrm{\Lambda })`$. Then consider a polynomial operator pencil
$$A_{k,E}=1+\left(\begin{array}{cc}U(C+E)& \pi (k_2+ik_1)\\ \pi (k_2ik_1)& U(C+E)\end{array}\right)\left(\begin{array}{cc}C& \\ \overline{}& C\end{array}\right)^1.$$
Since for any function $`g`$ we have
$$e^{2\pi ik,x}[A_{k,E}(𝒟+C)]g=[𝒟E]e^{2\pi ik,x}g,$$
there is a Floquet function $`f(x)`$ with the quasimomenta $`(k_1,k_2)`$ and the eigenvalue $`E`$ if and only if there is a $`\mathrm{\Lambda }`$-periodic function $`g(x)`$ satisfying the equation
$$A_{k,E}[𝒟+C]g=0.$$
Such a solution exists if and only if there is a solution $`\phi L_2(T^2)`$ to the equation
$$A_{k,E}\phi =0.$$
(19)
If such a solution $`\phi `$ exists, then $`f=e^{2\pi ik,x}[𝒟+C]^1\phi `$ is the desired Floquet function. The operator pencil
$$1A_{k,E}$$
is polynomial in $`k_1,k_2`$, and $`E`$. Since $`U`$ is bounded, the multiplication operator
$$\times U:L_2(T^2)L_2(T^2)$$
is bounded and the pencil $`(1A_{k,E})`$ consists in compact operators on $`L_2(T^2)`$.
Now we apply the Keldysh theorem (or the polynomial Fredholm alternative) to it. This theorem reads that there is a regularized determinant $`\stackrel{~}{det}A_{k,E}`$ of this pencil analytic in $`k_1,k_2`$, and $`E`$ such that the equation (19) is solvable if and only if $`\stackrel{~}{det}A_{k,E}=0`$ . Now it remains to put
$$Q(𝒟)=\{\stackrel{~}{det}A_{k,E}=0\}.$$
In the same manner for the pencil $`(1A_{k,0})`$ we obtain a complex curve
$$Q_0(𝒟)=\{\stackrel{~}{det}A_{k,0}=0\}^2.$$
As shown by perturbation methods the codimensions of these sets are positive (see ). Also nontriviality of such determinants follows from their construction by Keldysh. This proves the theorem.
We applied this method to the operators $`\mathrm{\Delta }+u`$ and $`_t\mathrm{\Delta }`$ in 1985. Later it became known to us that Kuchment also proposed the same approach in and therefore our paper was not published and referred in as an unpublished paper. The theory of such determinants is developed in and one can show that in fact they are entire functions of $`k`$ and $`E`$.
We shall discuss another and more effective but technically difficult approach of Krichever in 3.3.
Recall that the dual lattice $`\mathrm{\Lambda }^{}^2`$ consists of vectors $`\gamma ^{}`$ such that $`\gamma ,\gamma ^{}=0`$ for any $`\gamma \mathrm{\Lambda }`$. The following proposition is evident.
###### Proposition 2
The sets $`Q(𝒟)`$ and $`Q_0(𝒟)`$ are invariant with respect to translations by vectors from $`\mathrm{\Lambda }^{}`$:
$$k_1k_1+\text{Re}\gamma ^{},k_2k_2+\text{Im}\gamma ^{}.$$
(20)
Indeed, this action preserves the multipliers.
In the sequel we shall confine to the set $`Q_0(𝒟)`$.
###### Definition 3
$`Q_0(𝒟)`$ is called the (zero) Floquet spectral data of $`𝒟`$. The genus of the normalization of $`Q_0(𝒟)/\mathrm{\Lambda }^{}`$ is called the spectral genus of $`𝒟`$.
Another two properties of the Floquet spectrum are easily derived in the manner usual for soliton theory.
###### Proposition 3
If $`U`$ is real-valued, then $`Q_0(𝒟)`$ is invariant under an antiholomorphic involution $`k\overline{k}`$.
Proof. If $`(\psi _1,\psi _2)^{}`$ is a Floquet function with the quasimomenta $`(k_1,k_2)`$, then $`(\overline{\psi }_2,\overline{\psi }_1)^{}`$ is a Floquet function with the quasimomenta $`(\overline{k}_1,\overline{k}_2)`$. This proves the proposition.
###### Proposition 4
$`Q_0(𝒟)`$ is invariant under a holomorphic involution $`kk`$.
Proof. Consider the pencil $`L_k=A_{k,0}(𝒟+C)`$. We have
$$\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\overline{L_k^{}}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)=L_k.$$
(21)
The index of a Fredholm operator $`A`$ is $`\mathrm{ind}A=dim\mathrm{ker}Adim\mathrm{ker}A^{}`$. Since $`𝒟`$ is selfadjoint, its index vanishes. The operators $`L_k`$ have the same principal terms as $`𝒟`$ and, by the index theorem, their indices also vanish. Hence if $`kQ_0(𝒟)`$, then $`dimL_k=dimL_k^{}>0`$ and the identity (21) implies that $`dimL_k^{}>0`$. Therefore $`(k)Q_0(𝒟)`$ and this proves the proposition.
Given a basis $`(\gamma _1,\gamma _2)`$ for $`\mathrm{\Lambda }`$, we have a mapping
$$:Q_0(𝒟)/\mathrm{\Lambda }^{}^2:(k)=(e^{2\pi ik,\gamma _1},e^{2\pi ik,\gamma _2}),$$
which maps quasimomenta into multipliers.
The submanifold $`(Q_0(𝒟)/\mathrm{\Lambda }^{})^2`$ is generically singular and its normalization is the normalization of $`Q_0(𝒟)/\mathrm{\Lambda }^{}`$.
###### Definition 4
A complex curve $`\mathrm{\Gamma }`$, which is the normalization of $`Q_0(𝒟)/\mathrm{\Lambda }^{}`$, is called the spectral curve of $`𝒟`$.
A normalization sometimes consists in unstucking double points: a pair of points of $`\mathrm{\Gamma }`$ corresponding to a double point of $`Q_0(𝒟)/\mathrm{\Lambda }^{}`$ is called a resonance pair.
The definition of $``$ depends on a choice of a basis for $`\mathrm{\Lambda }`$. Given another basis $`(\stackrel{~}{\gamma }_1,\stackrel{~}{\gamma }_2)`$ for $`\mathrm{\Lambda }`$ such that
$$\left(\begin{array}{c}\stackrel{~}{\gamma }_1\\ \stackrel{~}{\gamma }_2\end{array}\right)=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{c}\gamma _1\\ \gamma _2\end{array}\right),\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)SL(2,),$$
the multipliers $`(\mu _1,\mu _2)=(\mu (\gamma _1),\mu (\gamma _2))`$ and $`(\stackrel{~}{\mu }_1,\stackrel{~}{\mu }_2)=(\mu (\stackrel{~}{\gamma }_1),\mu (\stackrel{~}{\gamma }_2))`$ are related as follows
$$\stackrel{~}{\mu }_1=\mu _1^a\mu _2^b,\stackrel{~}{\mu }_2=\mu _1^c\mu _2^d,$$
(22)
and the sets of multipliers $`\{(\mu _1,\mu _2)\}`$ and $`\{(\stackrel{~}{\mu }_1,\stackrel{~}{\mu }_2)\}`$ with respect to different bases are biholomorphically equivalent.
We call the image of $``$ the (Floquet) spectrum of $`𝒟`$ (on the zero energy level, $`E=0`$). Given a basis of $`\mathrm{\Lambda }`$, this image is uniquely defined. In general, the spectral data are defined modulo the $`SL(2,)`$-action (22) and we say that the spectral data of two operators coincide if the $`SL(2,)`$-orbits of their spectral data coincide.
### 3.2 Examples of spectra
1) $`U=0`$.
Let $`\mathrm{\Lambda }=+i`$. The Floquet functions are parameterized by two planes:
$$\psi ^+=(e^{\lambda _+z},0),\psi ^{}=(0,e^{\lambda _{}\overline{z}}),$$
$`\mathrm{\Gamma }`$ is a union of these planes compactified by two points at infinities, and $`\psi `$ has exponential singularities at these points.
The quasimomenta of $`\psi ^+`$ are
$$k_1=\frac{\lambda _+}{2\pi i}+n_1,k_2=\frac{\lambda _+}{2\pi }+n_2,n_j,$$
(23)
and the quasimomenta of $`\psi ^{}`$ are
$$k_1=\frac{\lambda _{}}{2\pi i}+m_1,k_2=\frac{\lambda _{}}{2\pi }+m_2,m_j.$$
(24)
Hence
$$Q_0=\left(_{n_1,n_2}A_{n_1,n_2}\right)\left(_{m_1,m_2}B_{m_1,m_2}\right),$$
where $`A_{n_1,n_2}`$ and $`B_{m_1,m_2}`$ are planes described by (23) and (24).
The functions $`\psi ^+`$ and $`\psi ^{}`$ have the same multipliers at the points
$$\lambda _+^{m,n}=\pi (n+im),\lambda _{}^{m,n}=\pi (nim),m,n.$$
These are resonance pairs for this potential with $`\mathrm{\Lambda }=+i`$.
Considering the zero potential as $`\mathrm{\Lambda }`$-periodic with respect to a general lattice $`\gamma _1+\gamma _2`$, the Floquet functions are the same but the resonance pairs are
$$\lambda _+^{m,n}=\frac{2\pi i}{\overline{\gamma }_1\gamma _2\gamma _1\overline{\gamma }_2}(\overline{\gamma }_1n\overline{\gamma }_2m),\lambda _{}^{m,n}=\frac{2\pi i}{\overline{\gamma }_1\gamma _2\gamma _1\overline{\gamma }_2}(\gamma _1n\gamma _2m).$$
(25)
2) $`U=C=\text{const}0`$.
Assume for simplicity that $`\mathrm{\Lambda }=+i`$.
The Floquet functions are
$$\psi (z,\overline{z},\lambda )=(\mathrm{exp}\left(\lambda z\frac{C^2}{\lambda }\overline{z}\right),\frac{C}{\lambda }\mathrm{exp}\left(\lambda z\frac{C^2}{\lambda }\overline{z}\right))$$
where $`\lambda \mathrm{\Gamma }=^{}=\{0\}`$. Compactify $`\mathrm{\Gamma }`$ by the points $`0`$ and $`\mathrm{}`$ and define the Floquet function on $`\mathrm{\Gamma }`$ globally as
$$\psi (z,\overline{z},\lambda )=\frac{\lambda }{\lambda C}(\mathrm{exp}\left(\lambda z\frac{C^2}{\lambda }\overline{z}\right),\frac{C}{\lambda }\mathrm{exp}\left(\lambda z\frac{C^2}{\lambda }\overline{z}\right)).$$
It has the following asymptotics
$$\psi \left(\begin{array}{c}\mathrm{exp}(k_+z)\\ 0\end{array}\right)\text{as }\lambda \mathrm{},\psi \left(\begin{array}{c}0\\ \mathrm{exp}(k_{}\overline{z})\end{array}\right)\text{as }\lambda 0$$
(26)
with $`k_+=\lambda `$ and $`k_{}=C^2/\lambda `$. After the normalization $`\psi `$ gets a pole at $`\lambda =C`$.
The resonance pairs $`(\lambda ,\lambda ^{})`$ are
$$\lambda =\frac{q\overline{q}\pm \sqrt{(q\overline{q})^24C^2q\overline{q}}}{2\overline{q}},\lambda ^{}=\lambda q,q=\pi (n+im),m,n$$
and they are parameterized by $`q\pi ^2\{0\}`$.
3) $`U=U(x)`$ is a function of one variable.
Let $`U(x+T)=U(x)`$ where $`T`$ is the minimal period. Then the equation (8) for Floquet functions $`\psi (x,y)=\phi (x)e^{\lambda y}`$ is the Zakharov–Shabat system
$$L\phi =\left(\begin{array}{cc}U& \frac{1}{2}_x\\ \frac{1}{2}_x& U\end{array}\right)\left(\begin{array}{c}\phi _1\\ \phi _2\end{array}\right)=\left(\begin{array}{cc}0& \frac{i}{2}\lambda \\ \frac{i}{2}\lambda & 0\end{array}\right)\left(\begin{array}{c}\phi _1\\ \phi _2\end{array}\right)$$
(27)
and in terms of $`\eta _1=\phi _1+i\phi _2`$ and $`\eta _2=\phi _1i\phi _2`$ it is
$$(_x+2iU)\eta _1=i\lambda \eta _2,(_x2iU)\eta _2=i\lambda \eta _1.$$
We see that $`f=\eta _2`$ satisfies the equation
$$\left[_x^2+(2iU_x4U^2)\right]f=\nu f$$
where $`\nu =\lambda ^2`$. The transformation $`L\left[_x^2+(2iU_x4U^2)\right]`$ is called the Miura transformation. The same name is used for the transformation
$$\left[_x^2+(2iU_x4U^2)\right]\left[_x^2+(2iU_x4U^2)\right].$$
(28)
For any $`\lambda `$ take a two-dimensional space $`𝒱_\lambda `$ of solutions to (27) and consider the monodromy operator
$$\widehat{T}:𝒱_\lambda 𝒱_\lambda :\widehat{T}(\phi )(x)=\phi (x+T).$$
For any pair $`(\phi (x,\lambda ),\vartheta (x,\lambda ))`$ of solutions to (27) their Wronskian $`W(\vartheta ,\phi )=\vartheta _1(x)\phi _2(x)\vartheta _2(x)\phi _1(x)`$ is constant. Take the basis $`(c(x,\lambda ),s(x,\lambda ))`$ for $`𝒱_\lambda `$ normalized as
$$c(0,\lambda )=s_x(x,\lambda )=1,c_x(0,\lambda )=s(x,\lambda )=0.$$
As shown in this basis the entries of the matrix $`\widehat{T}`$ are entire functions of $`\lambda `$. Since $`W(c,s)`$ is constant, $`det\widehat{T}(\lambda )=1`$ and the characteristic equation for $`\widehat{T}(\lambda )`$ takes the form
$$k^2\text{Tr}\widehat{T}(\lambda )k+1=0.$$
(29)
If $`\widehat{T}(\lambda )`$ is diagonalized, then eigenfunctions of $`\widehat{T}`$ are the Floquet functions. The operator $`\widehat{T}(\lambda )`$ is not diagonalized if and only if $`\lambda `$ is a simple root of the equation
$$\text{Tr}^2\widehat{T}(\lambda )4=0,$$
(30)
which can have only simple and double roots. In this case the Jordan form for $`\widehat{T}(\lambda )`$ is a non-diagonal upper triangular matrix, and there is only one (up to multiple) eigenfunction of $`\widehat{T}(\lambda )`$ with this value of $`\lambda `$.
We conclude that the Floquet function is globally defined on the two-sheeted covering of $``$ branched at points $`\lambda _1,\mathrm{}`$ which are simple roots of (30) and this complex curve is exactly $`\mathrm{\Gamma }`$.
Resonance pairs of the spectrum are pairs of points which project into double roots of (30).
If there are finitely many simple roots of (30) $`L`$ is called finite gap. For finite gap operators $`\mathrm{\Gamma }`$ is compactified by two infinities and the Floquet functions are pasted in a meromorphic function on $`\mathrm{\Gamma }`$ with the asymptotics (26) at the infinities.
These analytic properties of the Floquet function for a one-dimensional Schrödinger operator are explained in and for the Dirac operator such results are obtained by using the Miura transformation or derived by the same reasonings straightforwardly.
### 3.3 Spectra via perturbation theory
Generically Floquet functions and the Floquet spectrum can not be found by solving ordinary differential equations as in examples in 3.2. For proving existence of the Floquet spectrum we use in 3.1 the Keldysh theorem. Another approach for finding this spectrum and describing it in rather efficient manner was proposed by Krichever who realized it for a two-dimensional Schrödinger operator and for the operator $`_y_x^2+U(x,y)`$ . It is based on perturbation theory.
The examples discussed above demonstrate how the spectrum deforms under a deformation of $`U`$. The main picture is as follows: deforming potential we deform double points on $`Q_0(𝒟)/\mathrm{\Lambda }^{}`$ into handles removing singularities. The norm of the deformation measures the “size” of such handles.
For the two-dimensional Dirac operator (9) this is not done until recently but, since it is clear that the method of works for this operator after a slight modification, we explain what is the expected picture:
###### Pretheorem 1
For a smooth potential $`U`$ the spectral curve $`\mathrm{\Gamma }`$ consists of two parts: $`M_0`$ and $`M_{\mathrm{}}`$ where
1) $`M_0`$ is a complex curve of finite genus whose boundary consists in a pair of circles;
2) $`M_{\mathrm{}}`$ is diffeomorphic to a union of the domains $`|\lambda _\pm |R`$ of the $`\lambda _\pm `$-planes for some $`R`$ with some resonance pairs $`(\lambda _+^{m,n},\lambda _{}^{m,n})`$ (25) “joined by handles” with decreasing sizes as $`m^2+n^2\mathrm{}`$;
3) $`M_0`$ and $`M_{\mathrm{}}`$ are pasted along their boundaries.
This “joining by a handle” means that some small disks $`|\lambda _+\lambda _+^{m,n}|<\epsilon _+^{m,n}`$ and $`|\lambda _{}\lambda _{}^{m,n}|<\epsilon _{}^{m,n}`$ are excluded and replaced by a cylinder pasted to their boundaries and $`\epsilon _\pm ^{m,n}0`$ as $`m^2+n^2\mathrm{}`$.
To any point of the subsets $`M_+,M_{}M_{\mathrm{}}`$, where
$$M_+=\{\{|\lambda _+|>R\}\left(_{m,n}\{|\lambda _+\lambda _+^{m,n}|\epsilon _+^{m,n}\}\right)\},$$
$$M_{}=\{\{|\lambda _{}|>R\}\left(_{m,n}\{|\lambda _{}\lambda _{}^{m,n}|\epsilon _{}^{m,n}\}\right)\},$$
corresponds a unique (up to multiple) Floquet function. These functions and their multipliers $`\mu (\gamma _j,\lambda _\pm )`$ asymptotically behave as in the case $`U=0`$, and, in particular,
$$\mu (\gamma _j,\lambda _+)=e^{\lambda _+\gamma _j}\left(1+O\left(\frac{1}{\lambda _+}\right)\right),\mu (\gamma _j,\lambda _{})=e^{\lambda _{}\overline{\gamma }_j}\left(1+O\left(\frac{1}{\lambda _{}}\right)\right)$$
as $`\lambda _\pm \mathrm{}`$.
If $`U`$ does not vanish identically, then $`\mathrm{\Gamma }`$ is irreducible.
A potential $`U`$ is finite gap if there is such a representation with no handles joining resonance points in $`M_{\mathrm{}}`$. In this case $`\mathrm{\Gamma }`$ is compactified by a pair of infinities $`\mathrm{}_\pm `$ to a Riemann surface of finite genus.
In fact this is a general description of the Floquet spectra of operators with periodic coefficients. This physical picture from was chosen by Feldman, Knörrer, and Trubowitz as a most convenient and reasonable definition of general (non-hyperelliptic) Riemann surfaces of infinite genus and they developed a nice theory of such surfaces for which analogs of many classical theorems on algebraic curves take place (see also their preprints published by ETH).
## 4 The spectrum of the Weierstrass representation
### 4.1 The spectral curve of an immersed torus
Let $`F:^3`$ be a conformal immersion of a plane whose Gauss map descends through a torus $`/\mathrm{\Lambda }`$, i. e., double periodic. We shall consider tori as a particular case of such planes, when the immersion is also double periodic.
For immersed plane with a periodic Gauss map we have the Weierstrass representation constructed in 2.2. The double periodic potential $`U(z)`$ of this representation is the potential of the surface (with fixed conformal parameter).
It is said that two surfaces (with fixed conformal parameters) $`F_1:^3`$ and $`F_2:^3`$ are isopotential if their potentials coincide.
###### Definition 5
The spectral curve $`\mathrm{\Gamma }`$ of the operator $`𝒟`$ with the potential $`U`$ is called the spectral curve of a surface and the spectral genus of $`𝒟`$ is called the spectral genus of a surface. Given a basis $`\gamma _1,\gamma _2`$ for $`\mathrm{\Lambda }`$, the image of the multiplier map
$$:Q_0(𝒟)/\mathrm{\Lambda }^{}^2:(k)=(e^{2\pi ik,\gamma _1},e^{2\pi ik,\gamma _2})$$
is called the spectrum of a surface.
Let us look how these spectral notions depend on a choice of a conformal parameter.
###### Proposition 5
The spectral curve and the spectral genus of a surface do not depend on a choice of a conformal parameter.
The spectrum of a surface depends on a choice of a basis for $`H_1(T^2)=H_1(/\mathrm{\Lambda })=\mathrm{\Lambda }`$ and for different bases $`(\stackrel{~}{\gamma }_1,\stackrel{~}{\gamma }_2)`$ and $`(\gamma _1,\gamma _2)`$ related by a $`SL(2,)`$-transformation
$$\left(\begin{array}{c}\stackrel{~}{\gamma }_1\\ \stackrel{~}{\gamma }_2\end{array}\right)=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{c}\gamma _1\\ \gamma _2\end{array}\right),$$
the spectra $`\{(\stackrel{~}{\mu }_1,\stackrel{~}{\mu }_2)\}`$ and $`\{(\mu _1,\mu _2)\}`$ are related as
$$\stackrel{~}{\mu }_1=\mu _1^a\mu _2^b,\stackrel{~}{\mu }_2=\mu _1^c\mu _2^d.$$
Proof. It remains to check that all these data are preserved by a passage from $`z`$ to $`w=t^2z`$ with $`t\{0\}`$. For these conformal parameters the surface is defined by functions $`\psi (z)=(\psi _1(z),\psi _2(z))^{}`$ and $`\stackrel{~}{\psi }(w)=(\stackrel{~}{\psi }_1(w),\stackrel{~}{\psi }_2(w))^{}`$, which are solutions to the equations
$$𝒟\psi =\left(\begin{array}{cc}U& _z\\ _{\overline{z}}& U\end{array}\right)\psi =0,\stackrel{~}{𝒟}\stackrel{~}{\psi }=\left(\begin{array}{cc}\stackrel{~}{U}& _w\\ _{\overline{w}}& \stackrel{~}{U}\end{array}\right)\stackrel{~}{\psi }=0.$$
We have $`dw=t^2dz,_z=\frac{1}{t^2}_w`$, and, since $`e^{2\stackrel{~}{\alpha }(w)}dwd\overline{w}=e^{2\alpha (z)}dzd\overline{z}`$ and $`\stackrel{~}{H}(w)=H(z)`$, the formulas (11) imply that $`\stackrel{~}{U}(w)=t\overline{t}U(z)`$ for $`w=t^2z`$. Therefore
$$\left(\begin{array}{cc}U& _z\\ _{\overline{z}}& U\end{array}\right)\left(\begin{array}{c}\phi _1\\ \phi _2\end{array}\right)=0\text{if and only if}\left(\begin{array}{cc}\stackrel{~}{U}& _w\\ _{\overline{w}}& \stackrel{~}{U}\end{array}\right)\left(\begin{array}{c}t\phi _1\\ \overline{t}\phi _2\end{array}\right)=0$$
and, since $`\phi =(\phi _1,\phi _2)^{}`$ and $`\stackrel{~}{\phi }=(t\phi _1,\overline{t}\phi _2)^{}`$ have the same multipliers.
The transformation of the spectra was already discussed in 3.1 and we just repeat here these formulas because now they appear in another situation.
The proposition is proved.
We say that two planes, which may convert to tori by immersions, are isospectral if
1) there are conformal parameters one them such that both of them are represented by mappings $`F_1:^3`$ and $`F_2:^3`$;
2) the corresponding potentials of Dirac operators are $`\mathrm{\Lambda }`$-periodic with the same lattice $`\mathrm{\Lambda }`$, and the spectra of these operators with respect to a fixed basis for $`\mathrm{\Lambda }`$ coincide.
### 4.2 On the Willmore functional
Given a Weierstrass representation of a torus $`\mathrm{\Sigma }`$, we have
$$𝒲(\mathrm{\Sigma })=4_\mathrm{\Pi }U^2𝑑xdy=2i_\mathrm{\Pi }U^2𝑑zd\overline{z},$$
(31)
(see ) and that shows that the Willmore functional measures the deviation of $`𝒟`$ from the Dirac operator with the zero potential and geometrically that means that it measures the deviation of a connection in a spinor bundle defined by $`𝒟`$ from the trivial connection.
A relation of this spectrum to the Willmore functional
$$𝒲(\mathrm{\Sigma })=_\mathrm{\Sigma }H^2𝑑\mu ,$$
was first established in where it was discussed for surfaces of revolution. In this case there is a direct construction of the Floquet spectrum which a hyperelliptic complex curve (see Example 3 in 3.2). When the curve is of finite genus there is a compactification of it by two infinities and the Floquet function $`\psi (z,P)`$ is defined on the compactification and is meromorphic outside these infinities. We shall discuss a general case using a construction of the spectrum by perturbation.
Let $`U`$ be a $`\mathrm{\Lambda }`$-periodic potential with $`\mathrm{\Lambda }=\gamma _1+\gamma _2`$ and let $`\mathrm{\Gamma }`$ be of finite genus and be compactified by two infinities $`\mathrm{}_\pm `$. In fact, as shown in Pretheorem, these infinities are inherited during the perturbation of $`U`$ from the compactification of the spectrum of the zero potential (see Example 1 in 3.2).
Take the Floquet function $`\psi (z,P)`$ meromorphic outside the infinities and with the asymptotics
$$\psi \left(\begin{array}{c}\mathrm{exp}(\lambda _+z)\\ 0\end{array}\right)\text{as }P\mathrm{}_\pm ,\psi \left(\begin{array}{c}0\\ \mathrm{exp}(\lambda _{}\overline{z})\end{array}\right)\text{as }P\mathrm{}_{}$$
(32)
where $`\lambda _\pm ^1`$ are local parameters near $`\mathrm{}_\pm `$.
The theory of finite gap integration gives a recipe for reconstructing $`U`$ from such asymptotic expansions. Let
$$\psi (z,\lambda _+)=\mathrm{exp}(\lambda _+z)\left(\left(\begin{array}{c}1\\ 0\end{array}\right)+\left(\begin{array}{c}\zeta _1\\ \zeta _2\end{array}\right)\frac{1}{\lambda _+}+O\left(\frac{1}{\lambda _+^2}\right)\right)\text{as }\lambda _+\mathrm{}\text{.}$$
Substitute this expansion into the Dirac equation and expand $`𝒟\psi =0`$ into the powers of $`\lambda _+`$. Every coefficient in this expansion equals zero. Take the first two of them:
$$U+\zeta _2=0,U\zeta _2\overline{}\zeta _1=0.$$
(33)
This gives a reconstruction formula for $`U`$ and also the identity
$$U^2=\overline{}\zeta _1.$$
Now let us show how the Willmore functional appears in this picture. By the perturbation theory, we expect that the spectrum asymptotically behaves as the spectrum of the zero potential and this leads to the following conclusion: there is a function $`W(\lambda _+)`$ defined near $`\mathrm{}_+`$ such that
1) $`W(\lambda _+)=C_1\lambda _+^1+O\left(\lambda _+^2\right)`$;
2) $`\psi (z,\lambda _+)=e^{\lambda _+z+W(\lambda _+)\overline{z}}\phi (z,\lambda _+)`$ where $`\phi (z,\lambda _+)`$ is $`\mathrm{\Lambda }`$-periodic.
The function $`W`$ measures the deviation of the Floquet spectrum from the spectrum of the zero potential. Indeed, the multipliers of $`\psi `$ are
$$(\mu (\gamma _1),\mu (\gamma _2))=(e^{\lambda _+\gamma _1+W(\lambda _+)\overline{\gamma }_1},e^{\lambda _+\gamma _2+W(\lambda _+)\overline{\gamma }_2})$$
and the multipliers of the Floquet function of the zero potential (which is considered as $`\mathrm{\Lambda }`$-periodic) are
$$(\mu _0(\gamma _1),\mu _0(\gamma _2))=(e^{\lambda _+\gamma _1},e^{\lambda _+\gamma _2}).$$
It is easy to notice that
$$\zeta _1=C_1\overline{z}+\text{( a }\mathrm{\Lambda }\text{-periodic function)}$$
and hence
$$_\mathrm{\Pi }U^2𝑑zd\overline{z}=C_1\text{Vol}\mathrm{\Pi }$$
where $`\mathrm{\Pi }`$ is a the parallelogram spanned by $`\gamma _1`$ and $`\gamma _2`$. Now using (31) we conclude that
$$𝒲(\mathrm{\Sigma })=4C_1(\gamma _1\overline{\gamma }_2\overline{\gamma }_1\gamma _2)=4C_1\text{Vol}\mathrm{\Pi }$$
(34)
where $`\mathrm{\Pi }`$ is the area of the fundamental domain of $`\mathrm{\Lambda }`$. This derivation of the formula (34) was exposed in . The analogous derivation for the area of minimal tori is given in .
Now consider the whole series
$$W(\lambda )=C_1\lambda ^1+C_2\lambda ^2+\mathrm{}.$$
Since an involution $`kk`$ inverts $`\lambda `$ and preserves the spectrum, we have $`C_{2k}=0`$ for $`k=1,2,\mathrm{}`$. The quantities $`C_{2k1}`$ for $`k2`$ depend on choices of a conformal parameter $`z`$ and a parameter $`\lambda `$ on the spectral curve. Given a parameter $`\lambda `$, the holomorphic differentials
$$𝒲_k=4(C_{2k1}\text{Vol}\mathrm{\Pi })dz^{2k2}$$
are geometric invariants. The first of them is the Willmore functional $`𝒲=𝒲_1`$.
If the conformal parameter is fixed, then $`C_{2k1}`$ are first integrals of the modified Novikov–Veselov equation and the question of what are their geometrical meanings was posed in .
On a surface of revolution there is a distinguished conformal parameter $`z=x+iy`$ where $`x`$ is a parameter on the rotating curve and $`y`$ is the angle of rotation, $`0y2\pi `$, and there is a distinguished parameter $`\lambda `$ which is the eigenvalue of the Miura transformation of the Dirac operator (see 3.2). For these parameters $`C_{2k1}`$ are the Kruskal–Miura integrals of the modified Korteweg–de Vries equations. They are geometric invariants of surfaces of revolution .
Critical points of the Willmore functional are Willmore surfaces. For higher functionals critical points were not studied but using trace formulas we showed in that for spheres of revolution higher invariants are not bounded both from above and below.
### 4.3 Surfaces in terms of theta functions
Assume that the spectrum $`\mathrm{\Gamma }`$ of a torus $`\mathrm{\Sigma }`$ has finite genus which equals $`g`$ and take two different points $`\mathrm{}_\pm `$ on $`\mathrm{\Gamma }`$ with local parameters $`\lambda _\pm ^1`$ near them such that $`\lambda _\pm ^1(\mathrm{}_\pm )=0`$. Then the theory of Baker–Akhieser functions (see also ) reads that for a generic effective divisor $`D`$, i. e. , a formal sum of points on $`\mathrm{\Gamma }`$, of degree $`g+1`$ ($`D=P_1+\mathrm{}+P_{g+1}`$) there is a unique function $`\psi (z,\overline{z},P)`$ such that
1) $`\psi `$ is meromorphic in $`P\mathrm{\Gamma }\{\mathrm{}_\pm \}`$ and has the asymptotics (32) as $`P\mathrm{}_\pm `$;
2) $`\psi `$ has poles only in $`D`$.
Then this function is constructed in terms of theta functions of a Riemann surface $`\mathrm{\Gamma }`$. From this function one can reconstruct a Dirac operator $`𝒟`$ by (33). Therefore to each point $`P\mathrm{\Gamma }\{\mathrm{}_\pm ,P_1,\mathrm{},P_{g+1}\}`$ there corresponds a surface in $`^3`$ constructed from $`\psi (P)`$ via (10).
In this event, the Riemann surface $`\mathrm{\Gamma }`$ parameterizes isopotential surfaces and each of these surfaces is described in terms of theta functions of $`\mathrm{\Gamma }`$. The detailed formulas are given in . For CMC tori in $`^3`$ such formulas were derived in .
For tori of infinite spectral genera one can apply the theory of theta functions on such surfaces developed by Feldman, Knörrer, and Trubowitz .
### 4.4 On Willmore surfaces and the Willmore conjecture
First the Willmore functional $`𝒲`$ appeared in the 20s in the papers by Blaschke and Thomsen . In this event it was called the conformal area and its extrema were called conformally minimal surfaces. Blaschke and Thomsen also established the main properties of this functional which are:
1) the Willmore functional is invariant with respect to conformal transformations of the ambient space: let $`\mathrm{\Sigma }^3`$ be a compact immersed oriented surface, $`z=x+iy`$ be a conformal parameter on it, and $`G:\overline{}^3\overline{}^3`$ be a conformal transformation which maps $`\mathrm{\Sigma }`$ into $`^3`$, then $`𝒲(\mathrm{\Sigma })=𝒲(G(\mathrm{\Sigma }))`$, and this follows from the conformal invariance of the quantity <sup>1</sup><sup>1</sup>1Pinkall indicated that he and Pedit recently proved that the quantity $`Ae^\alpha `$ is already conformally invariant.
$$(k_1k_2)^2d\mu =4(H^2K)d\mu =16|A|^2e^{2\alpha }dxdy$$
where $`e^{2\alpha }dzd\overline{z}`$ is the first fundamental form and $`Adz^2`$ is the Hopf differential of the surface;
2) if $`\mathrm{\Sigma }`$ is a minimal surface in $`S^3`$ and $`\pi :S^3\overline{}^3`$ is the stereographic projection which maps $`\mathrm{\Sigma }`$ into $`^3`$, then $`\pi (\mathrm{\Sigma })`$ is a conformally minimal (Willmore) surface.
Hence in difference with minimal surfaces there are compact immersed Willmore surfaces in $`^3`$.
All Willmore spheres were described by Bryant . The classification of Willmore tori is not complete until recently and we only mention the papers where the finite gap integration was applied to this problem. In the Dorfmeister–Pedit–Wu (DPW) method , was applied for constructing general Willmore surfaces.
The simplest example of a Willmore torus is the stereographic projection of the Clifford torus $`\{(x^1)^2+(x^2)^2=1/2,(x^3)^2+(x^4)^2=1/2\}S^3^4`$ into $`^3`$. In another way it may be obtained as a circle torus of revolution such that the ratio of the distance from the center of the circle to the axis of revolution and the radius of the circle equals $`\sqrt{2}`$. This torus in $`^3`$ is also called the Clifford torus.
Willmore conjectured that
the Willmore functional achieves its minimum for tori on the Clifford torus and its conformal transformations and this minimum equals $`2\pi ^2`$ (the Willmore conjecture)
and checked this conjecture for circle tori of revolution .
This conjecture implies the Hsiang–Lawson conjecture that the area of a minimal torus in $`S^3`$ is greater or equal than $`2\pi ^2`$, the area of the Clifford torus in $`S^3`$, but not equivalent to it since there are Willmore tori in $`^3`$ which are not images of minimal tori under the stereographic projection .
The Willmore conjecture is still open and there are some particular cases for which it was proved:
1) Langer and Singer proved it for tori of revolution and Hertrich-Jeromin and Pinkall generalized their result for channel tori which are the enveloping tori for one-parameter families of spheres ;
2) Li and Yau proved it for tori conformally equivalent to flat tori $`/\{+\tau \}`$ where $`\tau =a+ib,0a1/2,b>0`$ and $`\sqrt{1a^2}b1`$ and later Montiel and Ros improved the latter inequality to $`(a1/2)^2+(b1)^21/4`$ .
Simon proved that the minimum of the Willmore functional is achieved on a real-analytic torus .
For higher genera surfaces a generalization of the Willmore conjecture was proposed by Kusner . He conjectured that for such surfaces the Willmore functional attain its minima of the stereographic projections of some minimal surfaces in $`S^3`$ constructed by Lawson.
In (see also ) we conjectured that
for a fixed conformal classes of tori the minimum of $`𝒲`$ is attained on tori with minimal spectral genus.
If this conjecture is valid then we may reduce the Willmore conjecture to estimating $`𝒲`$ for Willmore tori of small spectral genera and this can be done by using soliton theory.
The global Weierstrass representation gives a physical explanation for lower bounds for $`𝒲`$: for small perturbations of the zero-potential $`U=0`$ the surfaces constructed from solutions to (8) via (10) do not convert into tori and the lower bound for $`𝒲`$, the squared $`L_2`$-norm of the perturbation, shows how large a perturbation has to be to convert planes into tori.
## 5 The spectra of integrable tori
### 5.1 Constant mean curvature tori
Let $`\mathrm{\Sigma }`$ be a CMC torus in $`^3`$, i.e., $`H=\text{const}`$, and let it be conformally equivalent to $`/\mathrm{\Lambda }`$ with $`\mathrm{\Lambda }=\gamma _1+\gamma _2`$.
As shown in 2.3 for CMC surfaces the Hopf differential $`\omega =Adz^2`$ is holomorphic and, since the space of quadratic holomorphic differentials on a torus is one-dimensional, we have $`Adz^2=\text{const}dz^2`$. This differential does not vanish because otherwise all points are umbilics which is impossible for tori in $`^3`$. Hence assume that
$$\omega =\frac{1}{2}dz^2,H=1$$
and this is achieved by rescaling $`z`$ and by a homothety in $`^3`$. Now the Codazzi equations in terms of $`u=2\alpha `$ read
$$u_{z\overline{z}}+\mathrm{sinh}u=0$$
(35)
which is the sinh-Gordon equation. The Codazzi equations (14) give a commutation representation for (35):
$$\left[\frac{}{z}\left(\begin{array}{cc}\alpha _z& \frac{\lambda ^2}{2}e^\alpha \\ \frac{1}{2}e^\alpha & 0\end{array}\right)\right]\psi =0,\left[\frac{}{\overline{z}}\left(\begin{array}{cc}0& \frac{1}{2}e^\alpha \\ \frac{\lambda ^2}{2}e^\alpha & \alpha _{\overline{z}}\end{array}\right)\right]\psi =0$$
(36)
where $`\lambda ^2=1`$ in (14).
Notice that
1) for any $`\lambda 0`$ the compatibility condition for (36) is (35);
2) if $`|\lambda |=1`$ then (36) are the Codazzi equations (14) for the surface defined by $`\psi (\lambda ,z,\overline{z})`$ via (10).
There is another representation of (35) which gives rise to the spectral curve of a CMC torus .
Consider the equation (35) as the compatibility condition for the system
$$\left[\frac{}{z}\frac{1}{2}\left(\begin{array}{cc}u_z& \lambda \\ \lambda & u_z\end{array}\right)\right]\phi =0,\left[\frac{}{\overline{z}}\frac{1}{2\lambda }\left(\begin{array}{cc}0& e^u\\ e^u& 0\end{array}\right)\right]\phi =0$$
(37)
which contains the linear problem
$$L\phi =_z\phi \frac{1}{2}\left(\begin{array}{cc}u_z& 0\\ 0& u_z\end{array}\right)\phi =\frac{1}{2}\left(\begin{array}{cc}0& \lambda \\ \lambda & 0\end{array}\right)\phi $$
for a general $`\mathrm{\Lambda }`$-periodic potential $`u`$. Since $`L`$ is a first order $`2\times 2`$-matrix operator, for every $`\lambda `$ the system (37) has a two-dimensional space $`𝒱_\lambda `$ of solutions and these spaces are invariant under the translation operators
$$\widehat{T}_j\phi (z)=\phi (z+\gamma _j),j=1,2.$$
Since $`\widehat{T}_1,\widehat{T}_2`$, and $`L`$ commute, they have common eigenvectors and these vectors are glued into a meromorphic function $`\mathrm{\Psi }(z,\overline{z},P)`$ on a two-sheeted covering
$$\mathrm{\Gamma }(L):P\mathrm{\Gamma }(L)\lambda ,$$
which ramifies at points where $`\widehat{T}_j`$ and $`L`$ are not diagonalized simultaneously.
To each point $`P\mathrm{\Gamma }(L)`$ corresponds a unique (up to a constant multiple) Floquet function $`\phi `$ with multipliers $`\mu (\gamma _1,P)`$ and $`\mu (\gamma _2,P)`$.
By the same reasonings as in the example 3 in 3.4 it is shown that there are the Floquet functions defined on a $`\mathrm{\Gamma }(L)`$ such that $`\mathrm{\Gamma }(L)`$ is compactified by four infinities $`\mathrm{}_\pm ^1,\mathrm{}_\pm ^2`$ such that $`\mathrm{}_\pm ^1`$ are mapped into $`\lambda =\mathrm{}`$ and $`\mathrm{}_\pm ^2`$ are mapped into $`\lambda =0`$ and there are asymptotics
$$\psi (z,P)\mathrm{exp}\left(\frac{\lambda z}{2}\right)\left(\begin{array}{c}1\\ \pm 1\end{array}\right)\text{as }P\mathrm{}_\pm ^1,$$
$$\psi (z,P)\mathrm{exp}\left(\frac{\overline{z}}{2\lambda }\right)\left(\begin{array}{c}1\\ \pm 1\end{array}\right)\text{as }P\mathrm{}_\pm ^2,$$
and therefore their multipliers tend to $`\mathrm{}`$ as $`\lambda 0,\mathrm{}`$ .
If $`\phi =(\phi _1,\phi _2)^{}`$ satisfies (37) for $`\lambda =\mu `$, then
$$\sigma (\phi )=(\phi _1,\phi _2)$$
(38)
satisfies (37) for $`\lambda =\mu `$ and this generates an involution $`\sigma :\mathrm{\Gamma }(L)\mathrm{\Gamma }(L)`$, descending to an involution of $`:\lambda \lambda `$. We also have $`\sigma (\mathrm{}_\pm ^1)=\mathrm{}_{}^1`$ and $`\sigma (\mathrm{}_\pm ^2)=\mathrm{}_{}^2`$. By (38), the immersion
$$:\mathrm{\Gamma }(L)^2:P(\mu (\gamma _1,P),\mu (\gamma _2,P))$$
descends through the quotient of $`\sigma `$, i.e., $`:\mathrm{\Gamma }(L)\mathrm{\Gamma }(L)/\sigma ^2`$.
###### Definition 6
The complex curve $`\mathrm{\Gamma }(L)/\sigma `$ is called the spectral curve of a CMC torus $`\mathrm{\Sigma }`$. It is said that $`(\mathrm{\Gamma }(L)/\sigma )`$ is the spectrum of this torus.
By straightforward computations, we obtain
###### Proposition 6
$`\phi =(\phi _1,\phi _2)^{}`$ satisfies (37) if and only if $`\psi =(\lambda \phi _2,e^\alpha \phi _1)^{}`$ satisfies (36).
This proposition implies that the Floquet functions of $`L`$ and $`𝒟`$ have the same multipliers:
###### Theorem 5
Given a CMC torus $`\mathrm{\Sigma }`$, the spectrum of the CMC torus form a component, of the spectrum of the surface as defined in 4.1, containing both asymptotic ends where $`\mu (\gamma _j)e^{\lambda _+\gamma _j}`$ and $`\mu (\gamma _j)e^{\lambda _{}\overline{\gamma }_j}`$ as $`\lambda _\pm \mathrm{}`$.
Therefore, the spectral curve of the CMC torus $`\mathrm{\Sigma }`$ is an irreducible component of the spectral curve of this surface as defined in 4.1.
Assuming that Pretheorem is valid, we have more:
the spectrum and the spectral curve of a CMC torus coincide with the spectrum and the spectral curve of this surface as defined in 4.1.
In fact, for this conclusion we need only one part of Pretheorem, which states that the spectral curve is irreducible for $`U0`$.
Notice that the genus of the spectral curve of a CMC torus is finite .
### 5.2 Isothermic tori
A surface is called isothermic if near every point there is a conformal parameter $`z=x+iy`$ such that the fundamental forms are
$$𝐈=e^{2\alpha }(dx^2+dy^2),\mathrm{𝐈𝐈}=e^{2\alpha }(k_1dx^2+k_2dy^2).$$
To each isothermic surface $`F:𝒰^3`$ there corresponds the dual isothermic surface $`F^{}:𝒰^3`$ defined up to translations by the formulas
$$F_z^{}=e^{2\alpha }F_{\overline{z}},F_{\overline{z}}^{}=e^{2\alpha }F_z.$$
The fundamental forms of the dual surface are
$$𝐈=e^{2\alpha }(dx^2+dy^2),\mathrm{𝐈𝐈}=k_1dx^2+k_2dy^2,$$
the Gauss maps of $`F`$ and $`F^{}`$ are antipodal: $`N=N^{}`$, and $`F^{}=F`$ (modulo translations).
It is obtained by straightforward computations that
###### Proposition 7
If an isothermic surface $`F`$ is represented via (10) by a vector function $`\psi =(\psi _1,\psi _2)^{}`$, then the dual surface $`F^{}`$ is represented via (10) by the function $`\psi ^{}=(ie^\alpha \psi _2,ie^\alpha \psi _1)`$.
The potentials of these surfaces are
$$U=\frac{k_1+k_2}{4}e^\alpha ;U^{}=\frac{k_2k_1}{4}e^\alpha $$
and the Hopf differentials $`Adz^2`$ and $`A^{}dz^2`$ are
$$A=\frac{k_1k_2}{4}e^{2\alpha },A^{}=\frac{k_1+k_2}{4}.$$
###### Corollary 1
Given an isothermic surface $`\mathrm{\Sigma }`$ and its Hopf differential $`Adz^2`$ and the metric $`e^{2\alpha }dzd\overline{z}`$, there is an equality
$$Ae^\alpha =\frac{k_1k_2}{4}e^\alpha =U^{},$$
where $`U^{}`$ is the potential of the dual surface.
The simplest examples of isothermic surfaces are surfaces of revolution and constant mean curvature surfaces.
Let $`\mathrm{\Sigma }`$ be an isothermic plane, which may convert into an immersed torus, whose Gauss map descends through $`/\mathrm{\Lambda }`$ with $`\mathrm{\Lambda }=\gamma _1+\gamma _2`$.
By Proposition 7, the Gauss–Weingarten equations (13) are written as
$$\psi _z=𝐔\psi ,\psi _{\overline{z}}=𝐕\psi ,$$
with
$$𝐔=\left(\begin{array}{cc}\alpha _z& \frac{k_1k_2}{4}e^\alpha \\ \frac{k_1+k_2}{4}e^\alpha & 0\end{array}\right),𝐕=\left(\begin{array}{cc}0& \frac{k_1+k_2}{4}e^\alpha \\ \frac{k_2k_1}{4}e^\alpha & \alpha _{\overline{z}}\end{array}\right),$$
and their compatibility conditions are
$$\alpha _{xx}+\alpha _{yy}+k_1k_2e^{2\alpha }=0,k_{2x}(k_1k_2)\alpha _x=k_{1y}+(k_1k_2)\alpha _y=0.$$
(39)
The equations (39) are also the compatibility conditions for linear problems with a spectral parameter:
$$\widehat{\phi }_z=\widehat{𝐔}(\lambda )\widehat{\phi },\widehat{\phi }_{\overline{z}}=\widehat{𝐕}(\lambda )\widehat{\phi },$$
(40)
where
$$\widehat{𝐔}(\lambda )=\left(\begin{array}{cc}𝐔& \lambda 𝐉^{}\\ \lambda 𝐉^+& 𝐔+\alpha _z𝐄\end{array}\right),\widehat{𝐕}(\lambda )=\left(\begin{array}{cc}𝐕& \lambda 𝐉^+\\ \lambda 𝐉^{}& 𝐕+\alpha _{\overline{z}}𝐄\end{array}\right),$$
$$𝐉^+=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),𝐉^{}=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right),𝐄=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
First such representation with a spectral parameter had been found in in terms of $`5\times 5`$-matrices and, by using the $`4`$-dimensional spinor representation of $`SO(5,)`$, had been written in terms of $`4\times 4`$-matrices in . Here we use another representation which is gauge equivalent to the latter one.
Now the reasonings for describing the spectrum of an isothermic torus are the same as for a one-dimensional Dirac operator (see Example 3 in 3.2) and CMC tori and we will only sketch them.
For every $`\lambda `$ the system (40) has a four-dimensional space $`𝒱_\lambda `$ of solutions. On each such a space the translation operators act
$$\widehat{T}_j\widehat{\phi }(z)=\widehat{\phi }(z+\gamma _j),j=1,2.$$
Since $`\widehat{T}_1,\widehat{T}_2`$, and $`L`$ commute, they have common eigenvectors and these vectors are glued into a meromorphic function $`\mathrm{\Phi }(z,\overline{z},P)`$ on a four-sheeted covering
$$\mathrm{\Gamma }(𝐔):P\mathrm{\Gamma }(𝐔)\lambda ,$$
which ramifies at points where $`\widehat{T}_j`$ and $`L`$ are not diagonalized simultaneously.
To each point $`P\mathrm{\Gamma }(𝐔)`$ there corresponds a unique (up to a constant multiple) Floquet function $`\widehat{\phi }`$ with multipliers $`\mu (\gamma _1,P)`$ and $`\mu (\gamma _2,P)`$. If $`\mathrm{\Gamma }(𝐔)`$ is of finite genus, it is compactified by four “infinities” and $`\widehat{\phi }`$ is normalized to a meromorphic function on $`\mathrm{\Gamma }(𝐔)`$ with exponential singularities at these “infinities”.
Notice that if $`\widehat{\phi }=(\widehat{\phi }_1,\widehat{\phi }_2,\widehat{\phi }_3,\widehat{\phi }_4)^{}`$ satisfies (40) for $`\lambda =\mu `$ then
$$\sigma (\widehat{\phi })=(\widehat{\phi }_1,\widehat{\phi }_2,\widehat{\phi }_3,\widehat{\phi }_4)^{}$$
(41)
satisfies (40) for $`\lambda =\mu `$ and this generates an involution $`\sigma :\mathrm{\Gamma }(𝐔)\mathrm{\Gamma }(𝐔)`$, which descends to an involution of $``$: $`\lambda \lambda `$.
By (41), the immersion
$$:\mathrm{\Gamma }(𝐔)^2:P(\mu (\gamma _1,P),\mu (\gamma _2,P))$$
descends through the quotient of $`\sigma `$, i. e., $`:\mathrm{\Gamma }(𝐔)\mathrm{\Gamma }(𝐔)/\sigma ^2`$.
###### Definition 7
The complex curve $`\mathrm{\Gamma }(𝐔)/\sigma `$ is called the spectral curve of an isothermic surface $`\mathrm{\Sigma }`$. It is said that $`(\mathrm{\Gamma }(𝐔)/\sigma ))`$ is the spectrum of this surface.
It is again checked by straightforward computations that
###### Proposition 8
If $`\widehat{\phi }`$ satisfies (40), then
1) $`\psi =(e^\alpha \phi _3,e^\alpha \phi _4)`$ satisfies the Dirac equation $`𝒟\psi =0`$ with the potential
$$U=\frac{k_1+k_2}{4}e^\alpha ;$$
2) $`\psi ^{}=(e^\alpha \phi _2,e^\alpha \phi _1)`$ satisfies the Dirac equation $`𝒟\psi ^{}=0`$ with the potential
$$U^{}=\frac{k_2k_1}{4}e^\alpha .$$
We see that the Floquet functions of $`(_z𝐔)`$ and the Dirac operators $`𝒟`$ with potentials $`U`$ and $`U^{}`$ have the same multipliers and conclude
###### Theorem 6
Given an isothermic surface $`\mathrm{\Sigma }`$, the spectral curve and the spectrum of this isothermic isothermic surface coincide with
1) a component, of the spectrum of $`\mathrm{\Sigma }`$ as defined in 4.1, containing both asymptotic ends where $`\mu (\gamma _j)e^{\lambda _+\gamma _j}`$ and $`\mu (\gamma _j)e^{\lambda _{}\overline{\gamma }_j}`$ as $`\lambda _\pm \mathrm{}`$;
2) a component, of the spectrum of the dual surface $`\mathrm{\Sigma }^{}`$ as defined in 4.1, containing both asymptotic ends where $`\mu (\gamma _j)e^{\lambda _+\gamma _j}`$ and $`\mu (\gamma _j)e^{\lambda _{}\overline{\gamma }_j}`$ as $`\lambda _\pm \mathrm{}`$.
The spectral curve of the isothermic surface is an irreducible component of the spectral curves of this surface and its dual as defined in 4.1.
Of course, speaking about irreducibility we exclude the case $`U0`$.
Now Pretheorem or, more precisely, its statement about irreducibility of the spectral curve implies that
the spectrum and the spectral curve of an isothermic surface coincide with the spectrum and the spectral curve of this surface and its dual as defined in 4.1.
We consider here a general case when surface may be an immersed plane but not only torus because usually the dual surface to an isothermic torus is not closed.
In we introduce a particular case of the conjecture on the isospectrality of an isothermic surface and its dual. We conjectured that for surfaces of revolution, for which the isospectrality is equivalent to coincidence of all Kruskal–Miura integrals. Theorem 6 proves the general conjecture modulo Pretheorem and, since for surfaces of revolution Pretheorem holds (see Example 3 in 3.2 and ), implies the following
###### Theorem 7
A torus of revolution $`\mathrm{\Sigma }`$ and its dual surface $`\mathrm{\Sigma }^{}`$ have the same values of the Kruskal–Miura invariants $`𝒲_k`$.
Tori of revolution also explain this passage from $`\mathrm{\Gamma }(𝐔)`$ to $`\mathrm{\Gamma }(𝐔)/\sigma `$. The spectra of one-dimensional Schrödinger operators related by the Miura transformation (28) coincide because they are both double covered by the spectrum of $`L`$ . Moreover the Floquet function $`\psi `$ of $`L`$ consists in two components $`\eta _1`$ and $`\eta _2`$ which are the Floquet functions of the Schrödinger operators. In the same manner the spectra of Dirac operators corresponding to an isothermic surface and its dual surface are double covered in by the spectrum of $`\widehat{L}=[_z\widehat{𝐔}(\lambda )]`$ for $`\lambda =0`$ and the Floquet function of $`\widehat{L}`$ consists in the Floquet functions $`\psi `$ and $`\psi ^{}`$ of the Dirac operators (Proposition 7).
Recall the formulas for the Kruskal–Miura invariants for surfaces of revolution. Let $`\mathrm{\Sigma }`$ be parameterized by $`x`$, a parameter on the rotating curve, and $`y`$, which is the angle of rotation and $`0y2\pi `$, and let $`z=x+iy`$ be a conformal parameter on $`\mathrm{\Sigma }`$. Let $`U`$ be the potential of its Weierstrass representation, which is periodic in $`x`$ for tori. It also may be defined for spheres of revolution and in this case it is fast decaying and defined on the whole real line .
The densities of the Kruskal–Miura integrals for the KdV equation are
$$R_1=q,R_{n+1}=R_{nx}\underset{k=1}{\overset{n1}{}}R_kR_{nk}.$$
Since $`R_{2n}`$ are the derivatives of fast decaying functions, only the integrals
$$H_n(q)=_0^TR_{2n1}𝑑x$$
do not vanish identically. Here the integration is taken over $`[0,T]`$, where $`T`$ is the minimal period, for tori and over $``$ for spheres. The simplest integrals are
$$H_1(q)=q𝑑x,H_2(q)=q^2𝑑x,H_3(q)=(2q^3q_x^2)𝑑x.$$
Put $`q=2iU_xU^2`$ and define the Kruskal–Miura invariants as
$$𝒦_l(U)=2\pi H_l(q).$$
Notice that $`𝒦_1`$ is the Willmore functional and others are multiples of $`𝒲_l`$.
## 6 Surfaces in the three-sphere
### 6.1 The Dirac equation for surfaces in the three-sphere
Let $`G`$ be the Lie group $`SU(2)`$ and $`𝒢`$ be its Lie algebra $`su(2)`$ identified with the tangent space $`T_eG`$ to $`G`$ at the unity $`e`$. This Lie algebra is spanned by
$$e_1=\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right),e_2=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),e_3=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)$$
(42)
which satisfy the commutation relations $`[e_j,e_k]=2\epsilon _{jkl}e_l`$. Take a biinvariant metric on $`G`$:
$$\xi ,\eta =\frac{1}{2}\text{Tr}(\xi \eta ),\xi ,\eta 𝒢=T_eG,$$
in which the basis $`\{e_1,e_2,e_3\}`$ is orthonormal and $`G`$ is isometric to the unit $`3`$-sphere in $`^4`$.
Let $`\mathrm{\Sigma }`$ be a surface immersed into $`G`$, let $`z=x+iy`$ be a conformal parameter on $`\mathrm{\Sigma }`$, let $`f:\mathrm{\Sigma }G`$ be the immersion, and let $`𝐈=e^{2\alpha }dzd\overline{z}`$ be the induced metric.
Take the pullback of $`TG`$ to a $`𝒢`$-bundle over $`\mathrm{\Sigma }`$: $`𝒢E=f^1(TG)\stackrel{\pi }{}\mathrm{\Sigma }`$ and the differential
$$d_𝒜:\mathrm{\Omega }^1(\mathrm{\Sigma };E)\mathrm{\Omega }^2(\mathrm{\Sigma };E),$$
acting on $`E`$-valued $`1`$-forms on $`\mathrm{\Sigma }`$ as follows. Let
$$\omega =fvdz+fv^{}d\overline{z}$$
where $`f:T_eGT_{f(p)}G`$ is the left translation by $`f(p)`$. Then
$$d_𝒜\omega =d_𝒜^{}\omega +d_𝒜^{\prime \prime }\omega $$
where
$$d_𝒜^{}\omega =d_𝒜^{}\left(fvdz\right)=f\left(\overline{}v\frac{1}{2}[f^1\overline{}f,v]\right)dzd\overline{z},$$
$$d_𝒜^{\prime \prime }\omega =d_𝒜^{\prime \prime }\left(fv^{}dz\right)=f\left(v^{}+\frac{1}{2}[f^1f,v^{}]\right)dzd\overline{z}.$$
By straightforward computations derive that
$$d_𝒜(df)=0.$$
(43)
Since $`dz=idz`$ and $`d\overline{z}=id\overline{z}`$, we have
$$d_𝒜(df)=f(i\overline{}(f^1f)+i(f^1\overline{}f))dzd\overline{z}.$$
By the definition of the mean curvature $`H`$, we have
$$d_𝒜(df)=f(e^{2\alpha }\tau (f))dxdy=\frac{i}{2}f(e^{2\alpha }\tau (f))dzd\overline{z}$$
where $`\tau (f)`$ is the tension vector and $`f\tau (f)=2HN`$ with $`N`$ the normal vector: $`f^1N=ie^{2\alpha }[f^1f,f^1\overline{}f]`$. Finally we derive
$$d_𝒜(df)=f\left(H[f^1f,f^1\overline{}f]\right)dzd\overline{z}.$$
(44)
The case $`H=0`$ is described by the harmonicity equation
$$d_𝒜(df)=0.$$
Put
$$df=f\left(\mathrm{\Psi }dz+\mathrm{\Psi }^{}d\overline{z}\right)$$
(45)
and rewrite (43) and (44) as
$$\overline{}\mathrm{\Psi }\mathrm{\Psi }^{}+[\mathrm{\Psi }^{},\mathrm{\Psi }]=0,$$
(46)
$$\overline{}\mathrm{\Psi }+\mathrm{\Psi }^{}=iH[\mathrm{\Psi }^{},\mathrm{\Psi }].$$
(47)
Since $`f^1f_x=a^je_j`$ and $`f^1f_y=b^ke_k`$ with $`a^j,b^k`$, we have
$$\mathrm{\Psi }=\underset{j=1}{\overset{3}{}}Z_je_j,\mathrm{\Psi }^{}=\underset{j=1}{\overset{3}{}}\overline{Z}_je_j$$
with $`Z_j=(a^jib^j)/2`$ and $`\mathrm{\Psi },\mathrm{\Psi }^{}su(2)`$. The induced metric is
$$e^{2\alpha }dzd\overline{z}=\frac{1}{2}\text{Tr}[(\mathrm{\Psi }dz+\mathrm{\Psi }^{}d\overline{z})^2]=$$
$$(Z_1^2+Z_2^2+Z_3^2)(dz)^2+2(|Z_1|^2+|Z_2|^2+|Z_3|^2)dzd\overline{z}+(\overline{Z}_1^2+\overline{Z}_2^2+\overline{Z}_3^2)(d\overline{z})^2$$
and we conclude that
$$|Z_1|^2+|Z_2|^2+|Z_3|^2=\frac{1}{2}e^{2\alpha },Z_1^2+Z_2^2+Z_3^2=0.$$
Representing solutions to the latter equation as in 2.2
$$Z_1=\frac{i}{2}(\overline{\psi }_2^2+\psi _1^2),Z_2=\frac{1}{2}(\overline{\psi }_2^2\psi _1^2),Z_3=\psi _1\overline{\psi }_2,$$
(48)
we derive that
$$e^\alpha =|\psi _1|^2+|\psi _2|^2.$$
Now rewriting (46) in terms of $`\psi _j`$ and expanding it in the basis $`\{e_j\}`$, show that (46) is equivalent to the system
$$\overline{}(\psi _1\overline{\psi }_2)(\overline{\psi }_1\psi _2)=i(|\psi _2|^4|\psi _1|^4),\overline{}(\psi _1^2)+(\psi _2^2)=2i\psi _1\psi _2e^\alpha .$$
Analogously show that (47) is equivalent to the system
$$\overline{}(\psi _1\overline{\psi }_2)+(\overline{\psi }_1\psi _2)=H(|\psi _2|^4|\psi _1|^4),\overline{}(\psi _1^2)(\psi _2^2)=2H\psi _1\psi _2e^\alpha .$$
Introduce $`V_1=\psi _2/\psi _1`$ and $`V_2=\overline{}\psi _1/\psi _2`$. It follows from (46) that $`\text{Re}V_1=\text{Re}V_2,\text{Im}V_1=e^\alpha /2`$, and $`\text{Im}V_2=e^\alpha /2`$, and (47) implies that $`\text{Re}V_1=\text{Re}V_2=He^\alpha /2`$.
Finally we derive that
###### Theorem 8
For any immersed surface $`\mathrm{\Sigma }`$ is $`S^3`$ the spinor field $`\psi `$ defined by (45) and (48) satisfies the Dirac equation
$$𝒟^S\psi =0$$
with
$$𝒟^S=\left(\begin{array}{cc}0& \\ \overline{}& 0\end{array}\right)+\left(\begin{array}{cc}V& 0\\ 0& \overline{V}\end{array}\right),V=\frac{1}{2}(Hi)(|\psi _1|^2+|\psi _2|^2).$$
(49)
This spinor field is unique by its construction and we say that $`\psi `$ is the generating spinor for a surface.
Notice that the converse is not always true as in the case of Theorems 1 and 2. Indeed, not to any solution of $`𝒟^S\psi =0`$ corresponds a surface: a solution related to a surface has to satisfy an additional condition
$$|\psi _1|^2+|\psi _2|^2=2\text{Im}V.$$
It is easy to check that if $`𝒟^S\psi =0`$, then $`𝒟^S\phi =0`$ with $`\phi =(\overline{\psi }_2,\overline{\psi }_1)^{}`$.
Let us write the complete system of the Gauss–Weingarten equations. Recall that the Hopf differential equals $`Adz^2=f_{zz},Ndz^2`$ and, since the metric is left-invariant, we have $`A=f^1f_{zz},f^1N`$. Now $`\mathrm{\Psi }=f^1f_z`$ and
$$f^1f_{zz}=\mathrm{\Psi }_z+\mathrm{\Psi }^2,$$
where
$$\mathrm{\Psi }=\left(\begin{array}{cc}iZ_1& Z_2+iZ_3\\ Z_2+iZ_3& iZ_1\end{array}\right),\mathrm{\Psi }^{}=\left(\begin{array}{cc}i\overline{Z}_1& \overline{Z}_2+i\overline{Z}_3\\ \overline{Z}_2+i\overline{Z}_3& i\overline{Z}_1\end{array}\right).$$
(50)
We have $`\mathrm{\Psi }^2=(Z_1^2+Z_2^2+Z_3^2)e_1`$ and, since $`z`$ is a conformal parameter, $`\mathrm{\Psi }^2=0`$. Therefore as in 2.2 the Hopf differential takes the same form
$$Adz^2=(\psi _{1z}\overline{\psi }_2\overline{\psi }_{2z}\psi _1)dz^2.$$
We also have
$$\alpha _ze^\alpha =\overline{\psi }_1\psi _{1z}+\psi _2\overline{\psi }_{2z}$$
and finally write down the Gauss–Weingarten equations for an immersed surface in $`S^3`$ as
$$\left[\frac{}{z}\left(\begin{array}{cc}\alpha _z& Ae^\alpha \\ V& 0\end{array}\right)\right]\psi =\left[\frac{}{\overline{z}}\left(\begin{array}{cc}0& \overline{V}\\ \overline{A}e^\alpha & \alpha _{\overline{z}}\end{array}\right)\right]\psi =0.$$
(51)
The compatibility conditions are the Codazzi equations
$$\alpha _{z\overline{z}}+|V|^2|A|^2e^{2\alpha }=0,A_{\overline{z}}=(\overline{V}_z\alpha _z\overline{V})e^\alpha .$$
(52)
Examples.
1) The Clifford torus. This torus in $`^4`$ is defined by the equations
$$(x^1)^2+(x^2)^2=(x^3)^2+(x^4)^2=\frac{1}{2}$$
where $`(x^1,\mathrm{},x^4)^4`$. It is immersed into $`SU(2)`$ by the formula
$$f(x,y)=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}e^{ix}& e^{iy}\\ e^{iy}& e^{ix}\end{array}\right)$$
and has a conformal type of the square torus, $`(x,y)^2/2\pi ^2`$. Compute that
$$\mathrm{\Psi }=\frac{1}{4}((1+i)e_1+(1i)\mathrm{sin}(xy)e_2+(1i)\mathrm{cos}(xy)e_3),$$
$$\psi _1=\sqrt{\frac{1i}{2}}\mathrm{sin}\left(\frac{xy}{2}\frac{\pi }{4}\right),\psi _2=\sqrt{\frac{1+i}{2}}\mathrm{cos}\left(\frac{xy}{2}\frac{\pi }{4}\right),$$
$$e^\alpha =\frac{1}{\sqrt{2}},V=\frac{i}{2\sqrt{2}},A=\frac{1}{4}.$$
2) Minimal tori in $`S^3`$. In this case we have $`V=ie^\alpha /2`$ and derive from (52) that $`A_{\overline{z}}=0`$. This means that the Hopf differential is holomorphic and as in the case of CMC tori in $`^3`$ we conclude that it is constant and by rescaling conformal parameter achieve $`A=1/2`$. The case $`A=0`$ is also excluded for tori: it is realized by the equatorial $`S^2`$-spheres in $`S^3`$ (complete CMC surfaces in $`^3`$ with $`A=0`$ are the round spheres). The first equation from (52) is
$$u_{z\overline{z}}+\mathrm{sinh}u=0,u=2\alpha .$$
For CMC tori in $`S^3`$ the Hopf differential is also holomorphic and they are described in the same manner.
Now it is clear that the analogs of Theorems 1, 2, and 3 hold for surfaces in $`S^3`$. We again have the same spinor bundles over constant curvature surfaces.
###### Definition 8
Given a torus, represented in $`S^3`$ via $`\psi `$ satisfying (51), the spectral curve $`\mathrm{\Gamma }^S`$ of the operator $`𝒟^S`$ with the potential (49) is called the spectral curve of the torus. Given in addition a basis $`\gamma _1`$,$`\gamma _2`$ for $`\mathrm{\Lambda }`$, the period lattice for $`U`$, the image of the multiplier map
$$:Q_0(𝒟^S)/\mathrm{\Lambda }^{}^2:(k)=(e^{2\pi ik,\gamma _1},e^{2\pi ik,\gamma _2})$$
is called the spectrum of the torus in $`S^3`$.
Example. Let $`\mathrm{\Sigma }`$ be the Clifford torus. Then $`V=i/2\sqrt{2}`$ and the Floquet eigenfunctions $`\psi =(\psi _1,\psi _2)^{}`$ satisfy the equation
$$\left(\overline{}+\frac{1}{8}\right)\psi _j=0,j=1,2.$$
We derive that the general Floquet function is
$$\psi (z,\overline{z},\lambda )=(e^{\lambda z\frac{1}{8\lambda }\overline{z}},\frac{i}{2\sqrt{2}\lambda }e^{\lambda z\frac{1}{8\lambda }\overline{z}})^{}$$
and find the spectrum as the image of the multiplier map
$$\lambda \{0\}(e^{2\pi (\lambda \frac{1}{8\lambda })},e^{2\pi i(\lambda +\frac{1}{8\lambda })}).$$
The spectral curve is the two-sphere, i.e., the punctured plane $`\{0\}`$ compactified by the points $`\lambda =0,\mathrm{}`$. This implies that the spectral genus of the Clifford torus equals zero.
We shall not discuss the spectra of tori in $`S^3`$ in detail but only mention that Pretheorem also has to hold for them.
### 6.2 The Hitchin system
Let us compare the previous computations with the Hitchin theory of harmonic tori in the $`3`$-sphere . For Riemannian manifolds $`N`$ and $`M`$ a mapping $`f:NM`$ is called harmonic if it satisfies the equations
$$d_𝒜(df)=d_𝒜(df)=0$$
where $`𝒜`$ is the pullback of the Levi-Civita connection on $`TM`$ and the Hodge operator $``$ is taken with respect to the metric on $`N`$. If $`f`$ is an immersion and the metric on $`N`$ is the induced metric, then $`f(N)`$ is a minimal submanifold.
Let $`N`$ be an immersed surface $`\mathrm{\Sigma }`$ with the induced metric and $`M=G`$ be Lie group with a biinvariant metric. We adopt the notation from 6.1.
The harmonicity equations take the form
$$d_𝒜(\mathrm{\Psi }dz+\mathrm{\Psi }^{}d\overline{z})=0,$$
$$d_𝒜((\mathrm{\Psi }dz+\mathrm{\Psi }^{}d\overline{z}))=id_𝒜(\mathrm{\Psi }dz\mathrm{\Psi }^{}d\overline{z})=0.$$
They describe minimal surfaces in $`S^3`$ and are rewritten as
$$\overline{}\mathrm{\Psi }\mathrm{\Psi }^{}+[\mathrm{\Psi }^{},\mathrm{\Psi }]=\overline{}\mathrm{\Psi }+\mathrm{\Psi }^{}=0.$$
(53)
Following put
$$\mathrm{\Phi }=\frac{1}{2}\mathrm{\Psi },\mathrm{\Phi }^{}=\frac{1}{2}\mathrm{\Psi }^{}$$
and rewrite (53) as the Hitchin system
$$d_𝒜^{\prime \prime }\mathrm{\Phi }=0,F_𝒜=d_𝒜^2=[\mathrm{\Phi },\mathrm{\Phi }^{}]=0,$$
(54)
where $`F_A`$ is the curvature of the connection
$$d_𝒜:\mathrm{\Omega }^p(\mathrm{\Sigma };f^1TG)\mathrm{\Omega }^{p+1}(\mathrm{\Sigma };f^1TG)$$
and the formula means that $`d_𝒜^2`$ coincides with the multiplication by $`F_𝒜`$. The system (54) describes general harmonic mappings of surfaces in $`S^3`$ (when the metric on the surface is not necessarily induced) in terms of a connection $`𝒜`$ associated to the harmonic map and the Higgs field $`\mathrm{\Phi }`$.
The equation $`d_𝒜df=0`$ is equivalent to
$$\overline{}\mathrm{\Psi }\mathrm{\Psi }^{}+[\mathrm{\Psi }^{},\mathrm{\Psi }]=0$$
and means that the connection $`𝒜=(+\mathrm{\Psi },\overline{}+\mathrm{\Psi }^{})`$ on $`f^1TG`$ is flat, which is evident from its construction. However the second of the equations (53) implies that this connection is extended to an analytic family of flat connections
$$𝒜_\lambda =(+\frac{1+\lambda ^1}{2}\mathrm{\Psi },\overline{}+\frac{1+\lambda }{2}\mathrm{\Psi }^{})$$
where $`𝒜=𝒜_1`$ and $`\lambda \{0\}`$. This commutation representation with a spectral parameter was found by Pohlmeyer for harmonic maps into $`SU(2)`$ and later developed by Mikhailov and Zakharov for the case, when the target space is not a group but a symmetric space $`S^2`$ . Finally these two papers gave rise to the “integrability” part of the modern theory of harmonic maps .
For harmonic tori Hitchin introduced spectral curves and showed that they are of finite genus . Their construction is as follows.
Let $`\mathrm{\Sigma }`$ be a harmonic torus in $`S^3`$. For any $`\lambda \{0\}`$ we have a flat $`Sl(2,)`$ connection. Fix a basis $`\{\gamma _1,\gamma _2\}`$ for $`H_1(\mathrm{\Sigma })`$. For $`\gamma _1`$ and $`\gamma _2`$ define matrices $`H(\lambda ),\stackrel{~}{H}(\lambda )SL(2,)`$ which describe the monodromies of $`𝒜_\lambda `$ along closed loops realizing $`\gamma _1`$ and $`\gamma _2`$. These matrices commute and have joint eigenfunctions $`\phi (\lambda ,\mu )`$ where $`\mu `$ is a root of the characteristic equation for $`H(\lambda )`$
$$\mu ^2\text{Tr}H(\lambda )+1=0$$
and therefore there is a Riemann surface on which the eigenvalues
$$\mu _{1,2}=\frac{1}{2}\left(\text{Tr}H(\lambda )\pm \sqrt{\text{Tr}^2H(\lambda )4}\right)$$
are defined. The complex curve $`\mathrm{\Gamma }`$, which is a two-sheeted covering of $`P^1`$, ramifying at the odd zeros of the function $`(\text{Tr}^2H(\lambda )4)`$ and at $`0`$ and $`\mathrm{}`$, is called the spectral curve of a harmonic torus in $`S^3`$.
On $`\mathrm{\Gamma }`$ the eigenvalues of $`H(\lambda )`$ paste into a single-valued function $`\mu `$ with singularities at $`0`$ and $`\mathrm{}`$. Moreover the joint eigenfunctions of $`H(\lambda )`$ and $`\stackrel{~}{H}(\lambda )`$ paste into a vector function $`\phi `$ meromorphic on $`\mathrm{\Gamma }\{0,\mathrm{}\}`$.
We shall show that in the case, when the harmonic tori is an immersed tori in $`S^3`$ with the induced metric, i. e., in the situation of 6.1, the spectral curve of Hitchin is the same as the spectrum of the torus as defined in 6.1.
Let $`f:\mathrm{\Sigma }S^3`$ be an immersion of a minimal torus and $`\mathrm{\Psi }=f^1f_z,\mathrm{\Psi }^{}=f^1f_{\overline{z}}`$. Let the surface be defined by a spinor $`\psi `$.
The Hitchin eigenfunction $`\phi (\lambda ,\mu )`$ satisfies the equations
$$\left[+\frac{1+\lambda }{2}\mathrm{\Psi }\right]\phi =\left[\overline{}+\frac{1+\lambda ^1}{2}\mathrm{\Psi }^{}\right]\phi =0.$$
Take the matrix
$$L=\left(\begin{array}{cc}\overline{a}& \overline{b}\\ b& a\end{array}\right)$$
(55)
with $`a=(i\overline{\psi }_1+\psi _2)/\sqrt{2},b=(i\psi _1+\overline{\psi }_2)/\sqrt{2}`$. By (48) and (50), compute that
$$L^1\mathrm{\Psi }L=e^\alpha \left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),L^1\mathrm{\Psi }^{}L=e^\alpha \left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).$$
We also have
$$L^1L_z=\left(\begin{array}{cc}\alpha _z& iV\\ iAe^\alpha & 0\end{array}\right),L^1L_{\overline{z}}=\left(\begin{array}{cc}0& i\overline{A}e^\alpha \\ i\overline{V}& \alpha _{\overline{z}}\end{array}\right).$$
The vector function $`L^1\phi `$ satisfies the equations
$$\left[+\left(\begin{array}{cc}\alpha _z& iV\\ iAe^\alpha & 0\end{array}\right)+\frac{1+\lambda }{2}e^\alpha \left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)\right]L^1\phi =0,$$
$$\left[\overline{}+\left(\begin{array}{cc}0& i\overline{A}e^\alpha \\ i\overline{V}& \alpha _{\overline{z}}\end{array}\right)+\frac{1+\lambda ^1}{2}e^\alpha \left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\right]L^1\phi =0.$$
These two equations are compatible only for minimal tori, which are described by the condition
$$V=\frac{ie^\alpha }{2}.$$
For $`\stackrel{~}{\phi }=e^\alpha L^1\phi `$ we derive that
$$\stackrel{~}{\phi }_1+\frac{\lambda }{2}e^\alpha \stackrel{~}{\phi }_2=0,\overline{}\stackrel{~}{\phi }_2\frac{1}{2\lambda }e^\alpha \stackrel{~}{\phi }_1=0.$$
Put $`\stackrel{~}{\psi }_1=i\lambda \stackrel{~}{\phi }_2,\stackrel{~}{\psi }_2=\stackrel{~}{\phi }_2`$ and notice that $`\stackrel{~}{\psi }`$ satisfies (49):
$$\left[\left(\begin{array}{cc}0& \\ \overline{}& 0\end{array}\right)+\left(\begin{array}{cc}V& 0\\ 0& \overline{V}\end{array}\right)\right]\stackrel{~}{\psi }=0\text{with }V=ie^\alpha /2.$$
As in the proofs of Theorems 5 and 6 we conclude
###### Theorem 9
Given a minimal torus $`f:\mathrm{\Sigma }S^3`$, the Hitchin eigenfunction $`\phi (\lambda ,\mu )`$ by the transformation
$$\left(\begin{array}{c}\phi _1\\ \phi _2\end{array}\right)\left(\begin{array}{c}\stackrel{~}{\psi }_1\\ \stackrel{~}{\psi }_2\end{array}\right)=e^\alpha \left(\begin{array}{cc}0& i\lambda \\ 1& 0\end{array}\right)L^1\left(\begin{array}{c}\phi _1\\ \phi _2\end{array}\right)$$
is mapped to a Floquet function $`\stackrel{~}{\psi }`$ of $`𝒟^S`$. There is the mapping of the eigenvalues $`\mu `$ and $`\stackrel{~}{\mu }`$ of $`\phi `$ with respect to the monodromy operators $`H(\lambda )`$ and $`\stackrel{~}{H}(\lambda )`$ to the multipliers of $`\stackrel{~}{\psi }`$:
$$(\mu ,\stackrel{~}{\mu })((1)^{\epsilon (\gamma _1)}\mu ,(1)^{\epsilon (\gamma _2)}\stackrel{~}{\mu })$$
(56)
where $`(1)^{\epsilon (\gamma _1)},(1)^{\epsilon (\gamma _2)}`$ are the multipliers of the spinor $`\psi `$ generating a minimal torus.
This mapping (56) establishes a biholomorphic equivalence between the Hitchin spectral curve of a minimal torus and the connected component of the spectrum of this torus as defined in 6.1. This connected component contains both asymptotic ends near which $`\stackrel{~}{\psi }(e^{\lambda _+}z,0)^{}`$ or $`\stackrel{~}{\psi }(0,e^\lambda _{}\overline{z})^{}`$.
If Pretheorem holds for $`𝒟^S`$, then the spectrum is irreducible and therefore (56) establishes a biholomorphic equivalence of both spectra.
## 7 Conformal invariance of the spectra of tori
### 7.1 The Möbius group
We consider $`^4`$ as the set of matrices
$$a=\left(\begin{array}{cc}x^4+ix^1& x^2+ix^3\\ x^2+ix^3& x^4ix^1\end{array}\right),x^1,x^2,x^3,x^4,$$
(57)
and consider $`^3`$ as a subset described by $`x^4=0`$. The unit sphere $`S^3=SU(2)`$ is defined by the equation
$$|x|=1.$$
Take the north pole $`P=(0,0,0,1)`$ and denote by $`\pi `$ the stereographic projection of $`S^3`$ to $`^3=\{x^4=0\}`$ from $`P`$:
$$\pi :a\frac{1}{1x^4}\left(\begin{array}{cc}ix^1& x^2+ix^3\\ x^2+ix^3& ix^1\end{array}\right)=(1+a)(1a)^1.$$
The inverse mapping is
$$\pi ^1:b(b1)(b+1)^1.$$
This mapping $`\pi `$ establishes a conformal equivalence between $`S^3`$ and $`^3`$ compactified by a point at infinity, i.e., by $`\pi (P)=\mathrm{}`$.
The group of conformal transformations of $`^3\mathrm{}`$ is isomorphic to $`O^+(4,1)`$, the subgroup of $`O(4,1)`$ formed by isochronic transformations. The geometric picture is as follows. Let $`^{1,4}`$ be a $`5`$-dimensional pseudo-Euclidean space with the metric
$$x,y_{1,4}=x^0y^0\underset{j=1}{\overset{4}{}}x^jy^j.$$
The $`4`$-dimensional hyperbolic space $`^4`$ is embedded into $`^{1,4}`$ as the upper half of a hyperboloid: $`x,x_{1,4}=1,x^0>0`$, with the metric on tangent vectors $`\xi ,\xi =\xi ,\xi _{1,4}`$. The group of isometries of $`^4`$ is $`O^+(4,1)`$ and it acts on $`S^3`$, the absolute of $`^4`$, by conformal transformations.
By the Liouville theorem, the group $`O^+(4,1)`$ of conformal transformations is generated by
1) isometries of $`^3`$;
2) inversions with centers in $`x_0^3`$: $`x\frac{xx_0}{|xx_0|^2}`$;
3) homotheties: $`x\lambda x`$, $`\lambda \{0\}`$.
Any conformal transformation of $`\overline{}^3=^3\{\mathrm{}\}`$ which preserves $`\mathrm{}`$ is a composition of isometries and homotheties. Notice that in terms of (57) the inversion of $`^3`$ centered at $`x_0`$ looks simply as $`x(x_0x)^1`$.
We think that
the spectrum of a torus in $`S^3`$ which is stereographically projected into a torus in $`^3`$ and the spectrum of this projection coincide.
We can not prove it now but would like to notice that this statement easily implies conformal invariance of both spectra: it is clear that the spectrum of a torus in $`^3`$ is invariant under translations of the torus and the spectrum of a torus in $`S^3`$ is invariant under rotations in $`S^3`$. However the stereographic projection converts rotations in $`S^3`$ into conformal transformations of $`^3`$ which together with translations and homotheties generate the conformal group $`O^+(4,1)`$. The same holds with translations of $`^3`$ whose compositions with the projection generate together with rotations the group of conformal transformations of $`S^3`$.
### 7.2 Conformal invariance of the spectra for isothermic tori in $`^3`$
###### Theorem 10
Let $`\mathrm{\Sigma }`$ be an isothermic torus in $`^3`$ and let $`F:\overline{}^3\overline{}^3`$ be a conformal transformation which maps $`\mathrm{\Sigma }`$ into a torus $`F(\mathrm{\Sigma })`$ lying in $`^3`$. Then the spectrum of an isothermic torus $`\mathrm{\Sigma }`$ coincides with the spectrum of $`F(\mathrm{\Sigma })`$.
Notice that the spectrum of an isothermic torus is defined as a component of the general spectrum which contains the asymptotic ends where $`\mu (\gamma _j)e^{\lambda _+\gamma _j}`$ and $`\mu (\gamma _j)e^{\lambda _j\overline{\gamma }_j}`$ as $`\lambda _\pm \mathrm{}`$. Pretheorem states that the general spectrum is irreducible and therefore coincides with this component.
Proof. By Theorem 6, the spectrum of an isothermic torus and its dual isothermic surface coincide. The potential of the dual surface equals
$$U^{}=\frac{k_2k_1}{4}e^\alpha $$
and, by the Blaschke theorem, the density of the Willmore functional
$$\left(\frac{k_2k_1}{2}\right)^2d\mu =4\left(U^{}\right)^2dxdy$$
is invariant under conformal transformations of $`\overline{}`$. As known, conformal transformations maps isothermic surfaces into isothermic ones.
Let $`z`$ be a conformal parameter on $`\mathrm{\Sigma }`$ which is mapped into a conformal parameter on $`F(\mathrm{\Sigma })`$ and $`V`$ be the potential of $`F(\mathrm{\Sigma })`$ with respect to this parameter. We see that, by the Blaschke theorem, $`V^2=(U^{})^2`$ and therefore $`V=\pm U^{}`$.
It is clear that the spectra of the Dirac operators whose potentials differs by sign coincide. Now we derive that the spectra of the isothermic tori $`\mathrm{\Sigma }`$ and $`F(\mathrm{\Sigma })`$ coincide with the spectrum of the isothermic surface with the potential $`U^{}`$. This proves the theorem.
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# Limits on the Mass of a Composite Higgs Boson aafootnote aFor a more complete discussion, see ref. 1 and references therein.
## 1 The Triviality of the Standard Higgs Model
This task is made easier, and is also motivated, by the fact that the standard one-doublet Higgs model does not strictly exist as a continuum field theory. This result is most easily illustrated in terms of the Wilson renormalization group.$`^\mathrm{?}`$ Any quantum field theory is defined using a regularization procedure which ameliorates the bad short-distance behavior of the theory. Following Wilson, we define the scalar sector of the standard model
$`_\mathrm{\Lambda }=`$ $`D^\mu \varphi ^{}D_\mu \varphi +m^2(\mathrm{\Lambda })\varphi ^{}\varphi +{\displaystyle \frac{\lambda (\mathrm{\Lambda })}{4}}(\varphi ^{}\varphi )^2`$
$`+{\displaystyle \frac{\eta (\mathrm{\Lambda })}{36\mathrm{\Lambda }^2}}(\varphi ^{}\varphi )^3+\mathrm{}`$
in terms of a fixed UV-cutoff $`\mathrm{\Lambda }`$. Here we have allowed for the possibility of terms of (engineering) dimension greater than four. While there are an infinite number of such terms, one representative term of this sort, $`(\varphi ^{}\varphi )^3`$, has been included explicitly for the purposes of illustration. Note that the coefficient of the higher dimension terms includes the appropriate number of powers of $`\mathrm{\Lambda }`$, the intrinsic scale at which the theory is defined.
Wilson observed that, for the purposes of describing experiments at some fixed low-energy scale $`E\mathrm{\Lambda }`$, it is be possible to trade a high-energy cutoff $`\mathrm{\Lambda }`$ for one that is slightly lower, $`\mathrm{\Lambda }^{}`$, so long as $`E\mathrm{\Lambda }^{}<\mathrm{\Lambda }`$. In order to keep low-energy measurements fixed, it will in general be necessary to redefine the values of the coupling constants that appear in the Lagrangian. Formally, this process is referred to as “integrating out” the (off-shell) intermediate states with $`\mathrm{\Lambda }^{}<k<\mathrm{\Lambda }`$. Keeping the low-energy properties fixed we find
$`_\mathrm{\Lambda }`$ $``$ $`_\mathrm{\Lambda }^{}`$
$`m^2(\mathrm{\Lambda })`$ $``$ $`m^2(\mathrm{\Lambda }^{})`$ (2)
$`\lambda (\mathrm{\Lambda })`$ $``$ $`\lambda (\mathrm{\Lambda }^{})`$
$`\eta (\mathrm{\Lambda })`$ $``$ $`\eta (\mathrm{\Lambda }^{}).`$
Wilson’s insight was to see that many properties of the theory can be summarized in terms of the evolution of these (generalized) couplings as we move to lower energies. Truncating the infinite-dimensional coupling constant space to the three couplings shown above, the behavior of the scalar sector of the standard model is illustrated in Figure 2. This figure illustrates a number of important features of scalar field theory. As we flow to the infrared, i.e. lower the effective cutoff, we find:
* $`\eta 0`$ — this is the modern interpretation of renormalizability. If $`m_H\mathrm{\Lambda }`$, the theory is drawn to the two-dimensional $`(m_H,\lambda )`$ subspace. Any theory, therefore, in which $`m_H\mathrm{\Lambda }`$ is close to a renormalizable theory with corrections suppressed by powers of $`\mathrm{\Lambda }`$.
* $`m^2\mathrm{}`$ — This is the naturalness/hierarchy problem. To maintain $`m_H𝒪(v)`$ we must adjust <sup>b</sup><sup>b</sup>bNothing we discuss here will address the hierarchy problem directly. the value of $`m_H`$ in the underlying theory to of order
$$\frac{\mathrm{\Delta }m^2(\mathrm{\Lambda })}{m^2(\mathrm{\Lambda })}\frac{v^2}{\mathrm{\Lambda }^2}.$$
(3)
* $`\lambda 0`$ — The coupling $`\lambda `$ has a positive $`\beta `$ function and, therefore, as we scale to low energies $`\lambda `$ tends to 0. If we try to take the “continuum” limit, $`\mathrm{\Lambda }+\mathrm{}`$, the theory becomes free or trivial.$`^\mathrm{?}`$
The triviality of the scalar sector of the standard one-doublet Higgs model implies that this theory is only an effective low-energy theory valid below some cut-off scale $`\mathrm{\Lambda }`$. Given a value of $`m_H^2=2\lambda (m_H)v^2`$, there is an upper bound on $`\mathrm{\Lambda }`$. An estimate of this bound can be obtained by integrating the one-loop $`\beta `$-function, which yields
$$\mathrm{\Lambda }\begin{array}{c}<\\ \end{array}m_H\mathrm{exp}\left(\frac{4\pi ^2v^2}{3m_H^2}\right).$$
(4)
For a light Higgs, the bound above is at uninterestingly high scales and the effects of the underlying dynamics can be too small to be phenomenologically relevant. For a Higgs mass of order a few hundred GeV, however, effects from the underlying physics can become important. I will refer to these theories generically as “composite Higgs” models.
Finally, while the estimate above is based on a perturbative analysis, nonperturbative investigations of $`\lambda \varphi ^4`$ theory on the lattice show the same behavior. This is illustrated in Figure 2.
## 2 $`T`$, $`S`$, and $`U`$ in Composite Higgs Models
In an $`SU(2)_W\times U(1)_Y`$ invariant scalar theory of a single doublet, all interactions of dimension less than or equal to four also respect a larger “custodial” symmetry $`^\mathrm{?}`$ which insures the tree-level relation $`\rho =M_W^2/M_Z^2\mathrm{cos}^2\theta _W1`$. The leading custodial-symmetry violating operator is of dimension six $`^\mathrm{?}`$ and involves four Higgs doublet fields $`\varphi `$. In general, the underlying theory does not respect the larger custodial symmetry, and we expect the interaction
$$\text{}\frac{b\kappa ^2}{2!\mathrm{\Lambda }^2}(\varphi ^{}\stackrel{}{D^\mu }\varphi )^2,$$
(5)
to appear in the low-energy effective theory. Here b is an unknown coefficient of $`𝒪(1)`$, and $`\kappa `$ measures size of couplings of the composite Higgs field. In a strongly-interacting theory, $`\kappa `$ is expected $`^\mathrm{?}`$ to be of $`𝒪(4\pi )`$.
Deviations in the low-energy theory from the standard model can be summarized in terms of the “oblique” parameters $`^{\mathrm{?},\mathrm{?}}`$ $`S`$, $`T`$, and $`U`$. The operator in eqn. 5 will give rise to a deviation ($`\mathrm{\Delta }\rho =\epsilon _1=\alpha T`$)
$$|\mathrm{\Delta }T|=|b|\kappa ^2\frac{v^2}{\alpha (M_Z)\mathrm{\Lambda }^2}\begin{array}{c}>\\ \end{array}\frac{|b|\kappa ^2v^2}{\alpha (M_Z^2)m_H^2}\mathrm{exp}\left(\frac{8\pi ^2v^2}{3m_H^2}\right),$$
(6)
where $`v246`$ GeV and we have used eqn. 4 to obtain the final inequality. The consequences of eqns. (4) and (6) are summarized in Figures 4 and 4. The larger $`m_H`$, the lower $`\mathrm{\Lambda }`$ and the larger the expected value of $`\mathrm{\Delta }T`$. Current limits imply $`|T|\stackrel{<}{}0.5`$, and hence $`\mathrm{\Lambda }\stackrel{>}{}4\mathrm{TeV}\kappa `$. (For $`\kappa 4\pi `$, $`m_H\begin{array}{c}<\\ \end{array}450`$ GeV.)
By contrast, the leading contribution to $`S`$ arises from
$$\text{}\frac{a}{2!\mathrm{\Lambda }^2}\left\{[D_\mu ,D_\nu ]\varphi \right\}^{}[D^\mu ,D^\nu ]\varphi .$$
(7)
This gives rise to ($`\epsilon _3=\alpha S/4\mathrm{sin}^2\theta _W`$)
$$\mathrm{\Delta }S=\frac{4\pi av^2}{\mathrm{\Lambda }^2}.$$
(8)
It is important to note that the size of contributions to $`\mathrm{\Delta }T`$ and $`\mathrm{\Delta }S`$ are very different
$$\frac{\mathrm{\Delta }S}{\mathrm{\Delta }T}=\frac{a}{b}\left(\frac{4\pi \alpha }{\kappa ^2}\right)=𝒪\left(\frac{10^1}{\kappa ^2}\right).$$
(9)
Even for $`\kappa 1`$, $`|\mathrm{\Delta }S||\mathrm{\Delta }T|`$.
Finally, contributions to $`U`$ ($`\epsilon _2=\frac{\alpha U}{4\mathrm{sin}^2\theta _W}`$), arise from
$$\frac{cg^2\kappa ^2}{\mathrm{\Lambda }^4}(\varphi ^{}W^{\mu \nu }\varphi )^2$$
(10)
and, being suppressed by $`\mathrm{\Lambda }^4`$, are typically much smaller than $`\mathrm{\Delta }T`$.
## 3 Limits on a Composite Higgs Boson
From triviality, we see that the Higgs model can only be an effective theory valid below some high-energy scale $`\mathrm{\Lambda }`$. As the Higgs becomes heavier, the scale $`\mathrm{\Lambda }`$ decreases. Hence, the expected size of contributions to $`T`$ grow, and are larger than the expected contribution to $`S`$ or $`U`$. The limits from precision electroweak data in $`(m_H,\mathrm{\Delta }T)`$ plane shown in Figure 6. We see that, for positive $`\mathrm{\Delta }T`$ at 95% CL, the allowed values of Higgs mass extend to well beyond 800 GeV. On the other hand, not all values can be realized consistent with the bound given in eqn. (4). As shown in figure 6, values of Higgs mass beyond approximately 500 GeV would likely require values of $`\mathrm{\Delta }T`$ much larger than allowed by current measurements.
I should emphasize that these estimates are based on dimensional arguments, and we am not arguing that it is impossible to construct a composite Higgs model consistent with precision electroweak tests with $`m_H`$ greater than 500 GeV. Rather, barring accidental cancellations in a theory without a custodial symmetry, contributions to $`\mathrm{\Delta }T`$ consistent with eqn. 4 are generally to be expected. Specific composite Higgs boson models are discussed in ref. 1, and the estimates given here are shown to apply.
These results may also be understood by considering limits in the $`(S,T)`$ plane for fixed $`(m_H,m_t)`$. In Figure 6, changes from the nominal standard model best fit ($`m_H=84`$ GeV) value of the Higgs mass are displayed as contributions to $`\mathrm{\Delta }S(m_H)`$ and $`\mathrm{\Delta }T(m_H)`$. Also shown are the 68% and 95% CL bounds on $`\mathrm{\Delta }S`$ and $`\mathrm{\Delta }T`$ consistent with current data. We see that, for $`m_H`$ greater than $`𝒪`$(200 GeV), a positive contribution to $`T`$ can bring the model within the allowed region.
At Run II of the Fermilab Tevatron, it may be possible to reduce the uncertainties in the top-quark and W-boson masses to $`\mathrm{\Delta }m_t=2`$ GeV and $`\mathrm{\Delta }M_W=30`$ MeV.$`^\mathrm{?}`$ Assuming that the measured values of $`m_t`$ and $`M_W`$ equal their current central values, such a reduction in uncertainties will result the limits in the $`(m_H,\mathrm{\Delta }T)`$ plane shown in Figure 7. Note that, despite reduced uncertainties, a Higgs mass of up to 500 GeV or so will still be allowed.
## 4 Conclusions
In conclusion, the triviality of the Standard Higgs model implies that it is at best a low-energy effective theory valid below a scale $`\mathrm{\Lambda }`$ characteristic of nontrivial underlying dynamics. As the Higgs mass increases, the upper bound on the scale $`\mathrm{\Lambda }`$ decreases. If the underlying dynamics does not respect a custodial symmetry, it will give rise to corrections to $`T`$ of order $`\kappa ^2v^2/\alpha \mathrm{\Lambda }^2`$, while the contributions to $`S`$ and $`U`$ are likely to be much smaller. For this reason, it is necessary to consider limits on a Higgs boson in the $`(m_H,\mathrm{\Delta }T)`$ plane. In doing so, we see that a Higgs mass larger than 200 GeV is consistent with precision electroweak tests if there is a positive $`\mathrm{\Delta }T`$. Absent a custodial symmetry, however, Higgs masses larger than $`500`$ GeV are unlikely: the scale of underlying physics is so low that $`\mathrm{\Delta }T`$ is likely to be too large.
## 5 Acknowledgments
I thank Bogdan Dobrescu, Nick Evans, Christian Hölbling, and E. H. Simmons for fruitful collaborations. This work was supported in part by the U.S. Department of Energy under grant DE-FG02-91ER40676.
## References
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# Orthogonality relations in Quantum Tomography
## 1 Introduction
Two fundamental restrictions limit the possibility of devising a state reconstruction method. On one hand, the quantum complementarity principle does not allow to recover the quantum state from measurements on a single system, unless we have some prior information on it. On the other hand, the no cloning theorem ensures that it is not possible to make exact copies of a quantum system, without having prior knowledge of its state. Hence, the only possibility for devising a state reconstruction procedure is to provide a measuring strategy that employs numerous identical (although unknown) copies of the system, so that different measurements may be performed on each of the copies.
The problem of state estimation resorts essentially to estimating arbitrary operators of a quantum system by using the result of measurements of a set of observables. If this set of observables is sufficient to give full knowledge of the system state, then we define it a quorum. Notice that, in general, a system may allow various, different quorums. Quantum tomography was born as a state reconstruction technique in the optical domain, and has recently been extended to a vast class of systems. By extension, we now denote as “Quantum Tomography” all unbiased quantum state reconstruction procedures, i.e. those procedures which are affected only by statistical errors that can be made arbitrarily small by increasing the number of measurements. Tomography makes use of the results of the quorum measurements in order to reconstruct the expectation value of arbitrary operators (even not observables) acting on the system Hilbert space.
The purpose of this work is to present in a formally familiar manner (employing the Dirac notation also on operator space) a constructive method to derive tomographic formulas for quantum systems, at least for finite dimensional Hilbert spaces. This is achieved by giving conditions to build quorums and to check whether a given set of operators is a quorum. In this way, we obtain an extension of the recently proposed group tomography , where similar conditions were derived for systems with an underlying group structure.
In Sect. 2 we give the definitions and the conditions to identify a quorum of operators by analyzing the space of operators of a system as a linear vector space. We derive a constructive algorithmic procedure to obtain tomographic formulas in the case of finite quorums. In Sect. 3 we give some examples of applications of the presented method in the domain of spin systems, where various different quorums are available .
## 2 General estimation
Consider the set of system operators, i.e. the Liouville space $`()`$. If we initially restrict ourselves to Hilbert-Schmidt operators in $`()`$, then this set is itself a Hilbert space of operators, with the scalar product
$`\widehat{A}|\widehat{B}\stackrel{\text{def}}{=}\text{Tr}\left[\widehat{A}^{}\widehat{B}\right].`$ (1)
It is then possible to employ all the properties of linear vector algebra, and to use the Dirac notation, by using the following definitions for bra and ket vectors:
$`\widehat{O}`$ $``$ $`|\widehat{O}`$
$`\text{Tr}[\widehat{O}^{}]`$ $``$ $`\widehat{O}|.`$ (2)
In this vision, quantum tomography consists of expressing the operator $`\widehat{A}`$ we want to evaluate as an expansion on the observables of the quorum as
$`|\widehat{A}={\displaystyle _𝒳}𝑑x|\widehat{C}(x)\widehat{B}(x)|\widehat{A},`$ (3)
where $`|\widehat{A}`$ is a generic operator in $`()`$, $`|\widehat{C}(x)`$ (with $`x𝒳`$) is the set of quorum observables ($`C(x)`$ is a generally complex function of a selfadjoint operator, hence it is observable in this sense), and the set $`\widehat{B}(x)|`$ is the dual of the quorum. In ordinary notation, Eq. (2) is the tomography identity, i.e.
$`\widehat{A}={\displaystyle _𝒳}𝑑x\text{Tr}\left[\widehat{B}(x)^{}\widehat{A}\right]\widehat{C}(x)`$ (4)
Notice that the extension of the theory to non-normalizable vectors in the operator Hilbert space is immediate: one only has to require the existence of the trace of Eq. (4). If, for example, $`\widehat{A}`$ is a trace–class operator, then we do not need to require $`\widehat{B}(x)`$ to be of Hilbert-Schmidt class, since it is sufficient to require $`\widehat{B}(x)`$ bounded. Through Eq. (3), the tomographic reconstruction procedure is immediately obtained. In fact, by measuring the observables $`|\widehat{C}(x)`$ of the quorum, we can(<sup>1</sup><sup>1</sup>1Eq. (5) is obtained by taking the expectation value of both members of Eq. (3) and by calculating the expectation value trace using the eigenvectors of the quorum observables $`|\widehat{C}(x)`$.) express the mean value of any operator $`\widehat{A}`$ in terms of the eigenvalues of $`|\widehat{C}(x)`$ as
$`\widehat{A}={\displaystyle _𝒳}𝑑x{\displaystyle \underset{m}{}}p(m,x)\lambda _m^{(x)}\text{Tr}[\widehat{B}^{}(x)\widehat{A}],`$ (5)
where $`p(m,x)`$ is the probability of obtaining the eigenvalue $`\lambda _m^{(x)}`$ when measuring the quorum observable $`\widehat{C}(x)`$.
Since we want Eq. (3) to be valid for a generic operator $`|\widehat{A}`$ in $`()`$ \[or also in a subspace of $`()`$\], then we must require that the $`|\widehat{C}(x)`$ constitute a spanning set for the operator (sub)space, with the set of $`\widehat{B}(x)|`$ acting as its dual. A spanning set is a generalized basis for a vector space: it is a complete set of vectors but it is not, in general, composed of linearly independent (or normalized) vectors. Define dual $`\widehat{B}(x)|`$ of the set $`|\widehat{C}(x)`$ as the set constructed so to have
$`\widehat{B}(x)|\widehat{C}(x^{})=\text{Tr}[\widehat{B}^{}(x)\widehat{C}(x^{})]=\delta (x,x^{})x,x^{}𝒳,`$ (6)
where $`\delta (x,x^{})`$ is a reproducing kernel for $`\widehat{B}(x)|`$, i.e.
$`{\displaystyle _𝒳}𝑑x\delta (x,x^{})\widehat{B}(x)|=\widehat{B}(x^{})|.`$ (7)
Since $`|\widehat{C}(x)`$ is a complete set, $`\delta (x,x^{})`$ is a reproducing kernel also for this set, i.e.
$`{\displaystyle _𝒳}𝑑x\delta (x,x^{})|\widehat{C}(x)=|\widehat{C}(x^{}).`$ (8)
From linear vector algebra we obtain the following four equivalent definitions of spanning set:
A set of vectors $`|\widehat{C}(x)`$ (with dual $`\widehat{B}(x)|`$) is a spanning set $``$
1. $`|\widehat{A}(),|\widehat{A}={\displaystyle _𝒳}𝑑x|\widehat{C}(x)\widehat{B}(x)|\widehat{A}`$, i.e. the tomographic identity, namely Eq. (3).
2. $`|\widehat{C}_n`$ is complete, i.e. (no nonzero element is orthogonal to $`|\widehat{C}(x)`$ $`x`$):
$`\widehat{A}|\widehat{C}(x)=\widehat{B}(x)|\widehat{A}=0x𝒳|\widehat{A}=0.`$ (9)
3. the following operatorial identity resolution applies,
$`{\displaystyle _𝒳}𝑑x|\widehat{C}(x)\widehat{B}(x)|=\widehat{\widehat{1}},`$ (10)
where $`\widehat{\widehat{1}}`$ is the identity super-operator, namely the operator acting on operators such that $`\widehat{\widehat{1}}[\widehat{A}]=\widehat{A}\widehat{A}()`$.
4. $`{\displaystyle _𝒳}𝑑x\widehat{A}|\widehat{C}(x)\widehat{B}(x)|\widehat{A}=\widehat{A}^2\stackrel{\text{def}}{=}\text{Tr}[\widehat{A}^{}\widehat{A}]|\widehat{A}().`$
In the usual notation, these equivalent definitions write as :
1. $`\widehat{A}={\displaystyle _𝒳}𝑑x\text{Tr}[\widehat{B}^{}(x)\widehat{A}]\widehat{C}(x)`$.
2. $`\text{Tr}[\widehat{B}^{}(x)\widehat{A}]=\text{Tr}[\widehat{A}^{}\widehat{C}(x)]=0x𝒳\widehat{A}=0`$.
3. $`{\displaystyle _𝒳}𝑑xi|\widehat{C}(x)|jk|\widehat{B}^{}(x)|l=\delta _{il}\delta _{jk}`$, where $`\{|n\}`$ is a basis for the system Hilbert space $``$.
4. $`{\displaystyle _𝒳}𝑑x\text{Tr}[\widehat{A}^{}\widehat{C}(x)]\text{Tr}[\widehat{B}^{}(x)\widehat{A}]=\text{Tr}[\widehat{A}^{}\widehat{A}]\widehat{A}()`$.
In order to obtain the dual set $`\widehat{B}(x)|`$ starting from a given set $`|\widehat{C}(x)`$, one in general has to solve the operatorial equation (6) that defines the quorum. For finite quorums, this resorts to a matrix inversion. An alternative procedure is now proposed. It derives from the Gram–Schmidt orthogonalization method , which allows to derive a basis starting from a complete set of vectors. Namely, one obtains a basis $`|y_k`$, given the complete set $`|C_k`$ (assume for simplicity that all $`|C_k`$ are non-zero and that in $`\{|C_k\}`$ there are no couples of proportional vectors), recursively defined as
$`\{\begin{array}{c}|y_0\frac{1}{N_0}|C_0\\ |y_k\frac{1}{N_k}\left(|C_k_{j=0}^{k1}|y_jy_j|C_k\right)\end{array},`$
where $`N_0|C_0`$ and $`N_k|C_k_{j=0}^{k1}|y_jy_j|C_k`$. Notice that in the recursion (2) one must take care of eliminating all the vectors $`|C_k`$ which are a linear combination of the $`|y_j`$ with $`j<k`$.
Write the identity resolution for the basis obtained with procedure (2), i.e.
$`\widehat{1}`$ $`=`$ $`{\displaystyle \underset{k=0}{}}|y_ky_k|`$ (11)
$`{\displaystyle \frac{|C_0}{N_0}}y_0|+{\displaystyle \underset{k=1}{}}{\displaystyle \frac{1}{N_k}}\left(|C_k{\displaystyle \underset{j=0}{\overset{k1}{}}}|y_jy_j|C_k\right)y_k|.`$
By using repeatedly Eq. (2) (expressing $`|y_j`$ of Eq. (11) in terms of the $`|C_n`$s) and by reorganizing the terms in the sums, we can find the dual set $`B_n|`$ as
$`B_0|={\displaystyle \frac{y_0|}{N_0}}{\displaystyle \frac{y_0|C_1y_1|}{N_0N_1}}+\left({\displaystyle \frac{y_0|C_2}{N_0N_2}}+{\displaystyle \frac{y_0|C_1y_1|C_2}{N_0N_1N_2}}\right)y_2|+\mathrm{}`$
$`B_1|={\displaystyle \frac{y_1|}{N_1}}{\displaystyle \frac{y_1|C_2y_2|}{N_1N_2}}+\left({\displaystyle \frac{y_1|C_3}{N_1N_3}}+{\displaystyle \frac{y_1|C_2y_2|C_3}{N_1N_2N_3}}\right)y_3|+\mathrm{}`$
$`\mathrm{}`$ (12)
Eq. (11) guarantees that it is possible to write
$`\widehat{1}={\displaystyle \underset{n}{}}|C_nB_n|,`$ (13)
which is just the definition iii \[i.e. Eq. (10)\] of spanning set.
Summarizing, we described a method for deriving tomographic formulas for arbitrary systems. One must start from a set of operators $`\widehat{C}_n`$ he would like to use as a quorum, and verify that such a set is complete, i.e. that no nonzero element of $`()`$ is orthogonal to all $`\widehat{C}_n`$:
$`\widehat{A}|\widehat{C}_n=\text{Tr}[\widehat{A}^{}\widehat{C}_n]=0n|\widehat{A}=0.`$ (14)
If the set is finite, then one can employ the orthogonalization procedure outlined previously to derive the dual set. If the set is infinite discrete or continuous, then one can only resort to finding appropriate solutions for Eq. (6). Once the dual is known, the tomographic identity (3) can be written explicitly. The reconstruction procedure, in terms of the probabilities of measurements of quorum observables, follows straightforwardly and yields Eq. (5), which allows to obtain arbitrary operator expectation values in terms of quorum outcome probabilities. Of course, one may think of similar procedures based on different orthogonalization algorithms.
Since no hypotheses were made on the structure of the system Hilbert space, the theory presented in this section is valid for any quantum system. In the following section we will give some example applications.
## 3 Example of application: Spin Tomography
Here we show an application of the theory presented in the previous section by rederiving the spin tomography , where various different quorums may be employed.
The simplest possible example is a spin $`s=\frac{1}{2}`$ system. In this case we expect that the Pauli matrix and the identity constitute a quorum (since any $`2\times 2`$ matrix can be written on such a basis). Take the quorum $`𝒬\stackrel{\text{def}}{=}\{\widehat{\sigma }_x,\widehat{\sigma }_y,\widehat{\sigma }_z,\widehat{1}\}`$: it is immediate to verify that it is complete. Since the quorum operators are orthogonal, i.e. $`\widehat{\sigma }_\alpha \widehat{\sigma }_\alpha ^{}=\widehat{1}\delta _{\alpha \alpha ^{}}`$ ($`\alpha ,\alpha ^{}=x,y,z`$), using the Gram-Schmidt procedure it is immediate to obtain the dual set as $`𝒞=\{\frac{1}{2}\widehat{\sigma }_x,\frac{1}{2}\widehat{\sigma }_y,\frac{1}{2}\widehat{\sigma }_z,\frac{1}{2}\widehat{1}\}`$. The expansion (3) of a matrix $`\widehat{A}`$ is, thus
$`|\widehat{A}={\displaystyle \frac{1}{2}}\left[{\displaystyle \underset{\alpha =x,y,z}{}}|\widehat{\sigma }_\alpha \widehat{\sigma }_\alpha ^{}|\widehat{A}+|\widehat{1}\widehat{1}|\widehat{A}\right],`$ (15)
which immediately yields the reconstruction procedure
$`\widehat{A}={\displaystyle \underset{m=\frac{1}{2}}{\overset{\frac{1}{2}}{}}}{\displaystyle \underset{\alpha =x,y,z}{}}p(m,\stackrel{}{n}_\alpha )m\text{Tr}\left[\widehat{A}\widehat{\sigma }_\alpha \right]+{\displaystyle \frac{1}{2}}\text{Tr}\left[\widehat{A}\right],`$ (16)
where $`p(m,\stackrel{}{n}_\alpha )`$ is the probability to obtain the eigenvalue $`m=\pm \frac{1}{2}`$ while measuring $`\stackrel{}{S}\stackrel{}{n}_\alpha `$. This equation allows the reconstruction of the expectation value of any spin $`s=\frac{1}{2}`$ operator $`\widehat{A}`$ from the measurement of the spin in the $`x,y,z`$ directions.
For an arbitrary spin $`s`$, a possible quorum is given by the spin component in all directions, i.e. the observable $`\stackrel{}{S}\stackrel{}{n}`$ ($`\stackrel{}{S}`$ being the spin operator and $`\stackrel{}{n}`$ a vector on the unit sphere). In order to find the dual $`\widehat{B}|`$, consider the exponential of the quorum, i.e. $`\widehat{D}(\psi ,\stackrel{}{n})=\mathrm{exp}(i\psi \stackrel{}{S}\stackrel{}{n})`$, which satisfies definition iii \[i.e. Eq. (10)\] of spanning set. In fact, $`\widehat{D}(\psi ,\stackrel{}{n})`$ constitutes a unitary irreducible representation of the group SU(2). The orthogonality relation between the matrix elements of the group representation $`D(g)`$ of dimension $`d`$ writes as
$`{\displaystyle _R}𝑑gD_{jr}(g)D_{tk}^{}(g)={\displaystyle \frac{V}{d}}\delta _{jk}\delta _{tr},`$ (17)
where $`dg`$ is the group Haar invariant measure, and $`V=_R𝑑g`$. For SU(2), with the $`2s+1`$ dimension unitary irreducible representation $`\widehat{D}(\psi ,\stackrel{}{n})`$, Haar’s invariant measure is $`\mathrm{sin}^2\frac{\psi }{2}\mathrm{sin}\vartheta d\vartheta d\phi d\psi `$, and $`V=4\pi ^2`$. Thus, the orthogonality relations in this case are given by
$`{\displaystyle \frac{2s+1}{4\pi ^2}}{\displaystyle _\mathrm{\Omega }}𝑑\stackrel{}{n}{\displaystyle _0^{2\pi }}𝑑\psi \mathrm{sin}^2{\displaystyle \frac{\psi }{2}}j|e^{i\psi \stackrel{}{n}\stackrel{}{S}}|rt|e^{i\psi \stackrel{}{n}\stackrel{}{S}}|k=\delta _{jk}\delta _{tr},`$ (18)
which is the the spanning set definition iii for the set of operators $`|\widehat{D}=\widehat{D}`$, with dual $`\widehat{D}^{}|=\text{Tr}[\widehat{D}^{}]`$.
Then, it is possible to write the spin tomography identity as
$`\widehat{A}={\displaystyle \frac{2s+1}{4\pi ^2}}{\displaystyle _\mathrm{\Omega }}𝑑\stackrel{}{n}{\displaystyle _0^{2\pi }}𝑑\psi \mathrm{sin}^2{\displaystyle \frac{\psi }{2}}\text{Tr}\left[\widehat{A}\widehat{D}^{}(\psi ,\stackrel{}{n})\right]\widehat{D}(\psi ,\stackrel{}{n}),`$ (19)
from which the following reconstruction procedure is derived
$`\widehat{A}={\displaystyle \frac{2s+1}{4\pi ^2}}{\displaystyle \underset{m=s}{\overset{s}{}}}{\displaystyle _\mathrm{\Omega }}𝑑\stackrel{}{n}p(m,\stackrel{}{n}){\displaystyle _0^{2\pi }}𝑑\psi \mathrm{sin}^2{\displaystyle \frac{\psi }{2}}\text{Tr}\left[\widehat{A}e^{i\psi \left(\stackrel{}{S}\stackrel{}{n}m\right)}\right],`$ (20)
where $`p(m,\stackrel{}{n})`$ is the probability of obtaining $`m`$ as the measurement result of $`\stackrel{}{S}\stackrel{}{n}`$. This equation allows the reconstruction of arbitrary spin $`s`$ expectation values $`\widehat{A}`$, from spin measurements in all directions $`\stackrel{}{n}`$.
Numerical simulations show that the two preceding quorums are (for spin $`s=\frac{1}{2}`$) equivalent, namely the same number of experimental measurement data yield the same results and the same statistical error bars, apart from statistical fluctuations. In Fig. 1 a Monte Carlo comparison of the two spin reconstruction strategies based on the two different quorums is given. Both reconstructions are applied to a coherent spin state, defined as $`|\alpha \stackrel{\text{def}}{=}\mathrm{exp}(\alpha S_+\alpha ^{}S_{})|s`$, where $`S_+,S_{}`$ are the spin lowering and raising operators and $`|s`$ is the eigenvector of $`S_z`$ relative to the minimum eigenvalue.
Weigert has shown that another spin $`s`$ quorum can be obtained by taking $`N_s\stackrel{\text{def}}{=}(2s+1)^2`$ arbitrary(<sup>2</sup><sup>2</sup>2 Actually the choice of the directions is not completely arbitrary, but “almost” any choice yields a complete set of operators in $`()`$.) directions $`\stackrel{}{n}_k`$ and measuring the observables $`\widehat{𝒬}_k\stackrel{\text{def}}{=}|\stackrel{}{n}_k\stackrel{}{n}_k|`$, which are the projectors for the eigenspace relative to the maximum eigenvalue $`s`$ of the observables $`\stackrel{}{S}\stackrel{}{n}_k`$. We define a dual $`\widehat{𝒬}_k|`$ for the $`|\widehat{𝒬}_k`$ by requiring
$`\widehat{𝒬}_k|\widehat{𝒬}_k^{}=\delta _{kk^{}},`$ (21)
i.e. Eq. (11) of , which is just the dual set definition (6). Condition (21) together with the completeness of the chosen quorum, guarantee that $`|\widehat{𝒬}_k`$ (with dual $`\widehat{𝒬}_k|`$) is a spanning set for $`()`$, thus allowing the tomographic identity
$`|\widehat{A}={\displaystyle \underset{k=1}{\overset{N_s}{}}}|\widehat{𝒬}_k\widehat{𝒬}_k|\widehat{A},`$ (22)
i.e. (using the notation of )
$`\widehat{A}={\displaystyle \underset{k=1}{\overset{N_s}{}}}\text{Tr}[\widehat{A}\widehat{𝒬}^k]\widehat{𝒬}_k,`$ (23)
where $`\widehat{𝒬}^k`$ is the dual operator of $`\widehat{𝒬}_k`$. The explicit form of the dual set $`\widehat{𝒬}^k`$ can be derived by a matrix inversion starting from Eq. (21) or by the Gram–Schmidt based procedure method given on page 2. The reconstruction procedure is, in this case,
$`\widehat{A}=s{\displaystyle \underset{k=1}{\overset{N_s}{}}}p(s,\stackrel{}{n}_k)\text{Tr}\left[\widehat{A}\widehat{𝒬}^k\right],`$ (24)
where $`p(s,\stackrel{}{n}_k)`$ is the probability of obtaining the maximum eigenvalue $`s`$, when measuring $`\stackrel{}{S}\stackrel{}{n}_k`$. This allows the reconstruction of arbitrary spin operators $`\widehat{A}`$ from measurements of the spin along $`N_s`$ fixed directions.
## 4 Conclusions
Recent Group Tomography gives a general framework that allows to derive all the state reconstruction procedures that employ quorums which exhibit a group symmetry. Here we extended these results to generic state reconstruction procedures. In fact, we have seen how it is possible to give a characterization of tomographic formulas in terms of linear vector algebra on the vectors of the Liouville space of the system.
A constructive method to derive new tomographic formulas has been proposed starting from the Gram–Schmidt orthogonalization procedure. At least in principle, it allows to calculate the quorum dual for the quantum systems that allow a discrete quorum. We have given some examples of the method in the spin domain, by re-obtaining all the known spin tomographies using linear vector algebra arguments. For the sake of illustrating the method, we limited our analysis to the description of spin systems, but all known tomographies can be analyzed in this framework . Moreover, one may expect to employ the presented procedures to uncover new tomographies for quantum systems for which state reconstruction procedures are not presently known.
This work has been partially supported by INFM through project PAIS-1999-TWIN.
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# The spectral properties of noncondensate particles in Bose-condensed atomic hydrogen
## Abstract
The strong spin-dipole relaxation, accompanying BEC in a gas of atomic hydrogen, determines the formation of a quasistationary state with a flux of particles in the energy space to the condensate. This state is characterized by a significant enhancement of the low-energy distribution of non-condensate particles, resulting in the growth of their spatial density in the trap. This growth leads to the anomalous reconstruction of the optical spectral properties of non-condensate particles.
The discovery of Bose-Einstein condensation (BEC) in a trapped dilute gas of atomic hydrogen has opened a new page in the studies of metastable gaseous systems at ultracold temperatures. For a long period, trapped atomic hydrogen remained a very promising system to achieve BEC (see, e.g.,). However, only after the involvement of the forced evaporative cooling mechanism typical for the works resulted in the BEC discovery in alkali metal vapors the achievement of BEC in atomic hydrogen became a reality.
The experiments reported in have brought out a set of specific features inherent in the hydrogen system. The anomalously low rate of a three-particle recombination allows one to reach the record density of a Bose condensate of $`n_c10^{16}`$cm<sup>-3</sup> with the total number of atoms $`N`$ by several orders of the magnitude greater than that in the experiments on alkali gases. The lightest mass of hydrogen atoms in combination with so high density determines the superfluid phase transition temperature $`T_c`$ exceeding that for the experiments with alkali gases by two orders of the magnitude.
One of most interesting features of the hydrogen experiments is impossibility to achieve a relatively high concentration of the condensate fraction. This result is a direct evidence for the phenomenon of the “burning of a condensate” related to a strong increase of the gas density at the condensation and, as a result, drastic enhancement of the spin relaxation rate (see also ). The low rate of the evaporation cooling due to an extremely small scattering length $`a0.65\stackrel{}{\mathrm{A}}`$ makes this phenomenon responsible for the kinetics of the formation of a condensate in a gas of atomic hydrogen. This led to the prediction that the concentration of the condensed fraction cannot exceed several percents. It is worthwhile to note that the spatial density of the condensate may be large.
In the experiments an unexpected phenomenon has been observed. The shift of the line of the 1s-2s transition for non-condensate particles is found to increase abruptly after the BEC transition. This is especially surprising since the condensate volume is only about $`1\%`$ of the volume of the thermal cloud and the estimations presented in could not explain the observed picture. The aim of the present paper is to reveal the nature of this interesting phenomenon.
Under conditions of the experiments the interparticle collision time $`\tau `$ is small compared with the lifetime of the system $`\tau _L`$. As a result, a quasistationary state sets in for $`\tau <t<\tau _L`$ . For $`T>T_c`$, this state is close to the equilibrium one. However, for $`T<T_c`$ the intense losses of the particles from the condensate due to spin relaxation determine the appearance of the quasistationary state with the compensating flux of particles into the region of low energies. As turns out, the distribution function for the non-condensate fraction of the gas in such state differs drastically from the equilibrium distribution function. An important feature of this “steady-flux” distribution is a sharp increase of the number of particles in the low-energy range of the spectrum. For the trap geometry, this results in a strong increase of the density of normal fraction and therefore causes the anomalous transition line shift for non-condensed atoms. Note that the attempt to explain the anomalous increase of the non-condensate density by introducing a rather artificial model of the formation of metastable dense droplets has recently been undertaken in .
We confine ourselves to the time scale of $`\tau t\tau _L`$ and suppose that the quasistationary regime sets in. In order to simplify the analysis we will consider a gas in a spherically symmetric harmonical potential with frequency $`\omega `$. We suppose that the condensate is in the so-called Thomas-Fermi regime at $`T<T_c`$, that is, $`n_c(0)U_0\mathrm{}\omega `$ where $`U_0=4\pi \mathrm{}^2a/m`$ and $`m`$ is the atom mass. This means that the condensate density is determined by the expression
$$n_c(r)=\frac{2\mu _cm\omega ^2r^2}{2U_0}$$
(1)
valid for $`r<r_\mu `$ where $`r_\mu =(2\mu _c/m\omega ^2)^{1/2}`$ and $`\mu _c=n_c(0)U_0`$ is the chemical potential of the condensate.
At $`T<T_c`$ the loss of the particles is determined mostly by the condensate
$$\dot{n}_c(r)=\frac{1}{2}L^{(2)}n_c^2(r).$$
(2)
Here $`L^{(2)}`$ is the spin relaxation coefficient determined for a normal gas and the factor of $`1/2`$ appears from the two-particle density correlator for a condensate. Note that, in spite of the spatially inhomogeneous loss of Bose-condensed particles, the condensate retains Thomas-Fermi density profile. The point is that the redistribution of particles in the condensate, induced by the mean-field interaction, is much faster than their loss. Indeed, the characteristic evolution time in the condensate is $`\tau =\mathrm{}/n_c(r)U_0`$, while the characteristic time of the particle loss is $`\tau _l=1/L^{(2)}n_c(0)`$. In the atomic hydrogen $`\tau _l/\tau 10^5`$*.*
With the condensate density distribution (1) we find for the total particle loss rate
$$Q=\frac{16\pi }{105}L^{(2)}n_c^2(0)r_\mu ^3.$$
(3)
Let us consider the kinetic equation for the normal fraction in a harmonic potential. We are interested in particles with energies $`ϵ>\mu \mathrm{}\omega `$. In this case we can disregard the discrete structure of the trap energy levels. Under the assumption of ergodicity the quantum Boltzmann equation acquires the form
$$\dot{n}(ϵ_1)=I_{coll}(ϵ_1,[n]).$$
(4)
where $`n(ϵ)`$ is the statistical average of occupation numbers. Within the energy interval where the occupation numbers are large
$`I_{coll}`$ $`=`$ $`W{\displaystyle _0^{\mathrm{}}}𝑑ϵ_2𝑑ϵ_3𝑑ϵ_4\chi (ϵ_1,ϵ_2,ϵ_3,ϵ_4)\delta (ϵ_1+ϵ_2ϵ_3ϵ_4)`$ (7)
$`\{n(ϵ_1)n(ϵ_2)[n(ϵ_3)+n(ϵ_4)]`$
$`n(ϵ_3)n(ϵ_4)[n(ϵ_1)+n(ϵ_2)]\}`$
where
$$\chi (ϵ_1,ϵ_2,ϵ_3,ϵ_4)=\left[\mathrm{min}(1,\frac{ϵ_2}{ϵ_1},\frac{ϵ_3}{ϵ_1},\frac{ϵ_4}{ϵ_1})\right]^2,$$
(8)
$$W=\frac{4ma^2}{\pi \mathrm{}^3}.$$
(9)
This kinetic equation differs from that for a homogeneous case (see, e.g., ) only by the exponent in expression (8) (see, e.g., ), which is characteristic for a harmonic potential. As we will see below, the increase of spatial density, accompanying BEC, comes from the energy region where $`n(ϵ)1`$. We took the latter condition into account explicitly, deriving Eq. (7).
The density in the energy space can be written as $`\rho _ϵ=g(ϵ)n(ϵ)`$ where $`g(ϵ)=ϵ^2/2(\mathrm{}\omega )^3`$ is the density of states in the harmonic potential. We can rewrite Eq. (4) in the form of the continuity equation for the energy space
$`{\displaystyle \frac{\rho _ϵ}{t}}={\displaystyle \frac{}{ϵ}}P(ϵ)`$
with the flux $`P(ϵ)`$
$$P(ϵ)=^ϵ𝑑ϵ^{}g(ϵ^{})I_{coll}(ϵ^{},[n]).$$
(10)
In the quasistationary regime $`P(ϵ)=Q=const`$. This means that the flux (10) has the same value for an arbitrary $`ϵ`$. Taking into account Eq. (7), this requirement can be fulfilled if the distribution function has the form
$$n(ϵ)=\frac{A}{\stackrel{~}{ϵ}^\gamma }$$
(11)
where $`\stackrel{~}{ϵ}=ϵ/\mathrm{}\omega `$. Let us substitute this distribution function into Eq. (7). Introducing dimensionless variables $`\xi _i=ϵ_i/ϵ_1`$, after simple transformations we find
$`I_{coll}=W(\mathrm{}\omega )^2A^3I_0(\gamma )\stackrel{~}{ϵ}_1^{23\gamma }`$
where
$`I_0(\gamma )`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑\xi _2𝑑\xi _3𝑑\xi _4\chi (\xi _2,\xi _3,\xi _4)\delta (1+\xi _2\xi _3\xi _4)`$ (13)
$`(\xi _4^\gamma +\xi _3^\gamma \xi _2^\gamma 1)(\xi _2\xi _3\xi _4)^\gamma `$
Inserting these relations into Eq. (10), we find for the flux
$`P(ϵ){\displaystyle \frac{I_0(\gamma )}{53\gamma }}ϵ^{53\gamma }.`$
The requirement of the independence of the flux on $`ϵ`$ leads to
$$\gamma =\frac{5}{3}$$
(14)
and simultaneously to $`I_0(\gamma )=0`$. The last result is fairly evident: $`I_{coll}`$ is equal to zero in a stationary case. At the same time the derivative $`I_0^{}(\gamma )`$ is finite and provides the permanent flux $`P=Q`$ to the low energies region. The obtained results are in correspondence with the general analysis of Zakharov (see and also ), showing that the equation $`I_{coll}=0`$ has a nontrivial solution relevant to a steady-state particle flux in the energy space.
For the relation between $`A`$ and $`Q`$, we obtain
$$A=C\left(\frac{Q}{W(\mathrm{}\omega )^2}\right)^{1/3}.$$
(15)
One can estimate the dimensionless numerical coefficient $`C`$ using the total number of non-condensate particles $`N_{nc}`$ if one assumes approximately that Eqs. (11),(14) remain valid up to the maximum energy.
It is worth to note that the numerical calculations of the BEC kinetics demonstrate the formation of the distribution of the form (11) (with $`\gamma >1`$) before the real growth of condensate starts (see, e. g., ).
Using $`n(ϵ)`$ from Eq. (11), one can find the spatial density distribution *$`n(r)`$* for the normal fraction. The direct numerical calculation shows that the main contribution to $`n(r)`$ comes from the energy levels $`\mu ϵ<T`$. These levels are weakly affected by the presence of the condensate (the volume occupied by the wavefunctions of a particle is much larger than the condensate volume) and by the interparticle interaction in the thermal cloud. This means that within a reasonable approximation we can neglect the renormalization of the levels and use oscillator vavefunctions *$`\psi _𝐩(𝐫)`$* where $`𝐩(n_r,l,m)`$. The energy region $`(ϵ\mu )\mu `$, where the modification of the level structure is noticeable, gives a negligible contribution to the spacial density, and we ignore this modification. Under the assumption of ergodicity
$$n(𝐫)=\underset{𝐩}{}|\psi _𝐩(𝐫)|^2n(ϵ_𝐩)$$
(16)
Since the spacial density of the thermal cloud is weakly sensitive to the low energies, at the presence of the condensate we simply truncate the sum (16) by the condition $`ϵ_𝐩\mu _c`$.
Since the concentration of the condensate fraction is confined to several percents, the main part of the particles at $`T<T_c`$ is in the non-condensate fraction and $`N_{nc}N`$. At $`T>T_c`$ the distribution function for low energies at the actual absence of the flux compensating losses has an equilibrium form $`n(ϵ)=T/(ϵ+|\mu |)`$. Comparing this distribution with Eqs. (11,14), one can see that the maximum density of the non-condensate fraction, determined in the trap mostly by the low-energy region in the presence of flux at $`T<T_c`$, can exceed significantly the equilibrium value of the density at $`T>T_c`$. Note that in the presence of condensate the thermal equilibrium is absent and therefore the temperature may be regarded only as a characteristic energy scale for the system.
For the Doppler-free two-photon excitation, the spectral distribution is determined by the red shift of the absorption line. The shift is caused by the change of the coupling (scattering length) of an excited atom compared with the atom in the ground state. In the quasiclassical approximation the shift is proportional to the local density of particles, coinciding with $`n(r)`$ beyond the condensate region. Thus, the appearance of a condensate, accompanied by a sharp growth of $`n(r)`$ for small $`r`$, causes a significant increase of the shift of the non-condensate spectral distribution.
For an optically thin sample, in the local density approximation the Doppler-free spectral distribution for the two-photon absorption is proportional to the density distribution $`G(n)`$. For the spherically symmetric configuration, this distribution reads
$$G(n)=4\pi r^2(n)\left|\frac{dr(n)}{dn}\right|n$$
(17)
where $`r(n)`$ is determined from condition $`n(r)=n`$. For the non-condensate fraction, $`n(r)`$ is determined by (16).
Below we present approximate quantitative estimations for the described picture on the basis of the experimental data . Since in the experiments the collision time $`\tau 10^3`$s is much less than the lifetime of the system $`\tau _L5`$s, the quasistationary regime sets in for times $`\tau t\tau _L`$ as explained above. The maximum condensate density is found to be $`n_c(0)5\times 10^{15}`$cm<sup>-3</sup>. We will consider a spherically symmetric harmonic trap with $`\omega =(\omega _{}^2\omega _z)^{1/3}=3.4\times 10^3`$Hz. Here $`\omega _{}`$ and $`\omega _z`$ are the frequencies of the cylindrical trap used in . For $`L^{(2)}=1.1\times 10^{15}`$cm<sup>3</sup>/s and $`a=0.65\stackrel{}{\mathrm{A}}`$ , we find $`Q=1.9\times 10^9`$s<sup>-1</sup> for the particle loss rate. In the Thomas-Fermi approximation, the number of particles in the condensate for the given parameters is $`N_c1.2\times 10^9`$. The largest uncertainty in the experimental data is related to the value of the relative condensate fraction $`\chi `$. Here we use some average value of $`\chi 0.07`$. Then the total number of particles is $`N1.7\times 10^{10}`$. The critical temperature of the BEC transition is $`T_c60\mu `$K in this case.
Let us assume that non-condensate particles are concentrated within the energy interval $`[\mu ,ϵ^{}]`$, having distribution (11). Comparing the energies of a system with the distribution (11) and with the equilibrium distribution at temperature $`T`$, we find $`ϵ^{}=4.7T_c(T/T_c)^4`$ for the same number of particles in the both cases. The parameter $`ϵ^{}\mu `$. The increase of the spatial density of non-condensed particles originates from the occupation of the energy interval where $`ϵ`$ is significantly smaller than $`ϵ^{}`$. The results are only weakly sensitive to the behavior of $`n(ϵ)`$ near the upper bound of the energy spectrum. Our approximate definition of $`ϵ^{}`$ is made with the aim to compare the results below and above $`T_c`$, assuming the same number of particles. Taking into account the estimated values of $`N_{nc}`$, $`Q`$ and $`W`$ from Eq.(9) we can determine the coefficients of $`A`$ and $`C`$ in (14,15). At $`T0.9T_c`$ we obtain $`A2.9\times 10^5`$ and $`C2.5`$.
Using distribution function (11) with the values found for $`A`$ and $`ϵ^{}`$, one can determine the spatial density distribution (16) for the non-condensate fraction.
The dependence $`n(r)`$ is shown in FIG.1 (curve 1). For comparison, we present the density profiles for two temperatures above $`T_c`$: 70$`\mu `$K (curve 2) and 120$`\mu `$K (curve 3). The length in the figure is given in units of the characteristic trap size $`l_0=\sqrt{\mathrm{}/m\omega }`$. The absence of condensate at $`T>T_c`$ allows one to neglect the spin relaxation losses and therefore to use the equilibrium distribution with the chemical potential fixing the same number of particles.
Knowing $`n(r)`$, we can find $`G(n)`$ (17) and therefore the spectral distribution of two-photon absorption. To make a comparison with the experimental results , we use the relation $`\mathrm{\Delta }\nu (n)=\eta n`$ for the shift of absorption line and, according to , assume that the coefficient $`\eta `$ has the same value both for the condensate and for the non-condensate fraction. Regardless of a specific value of the scattering length for the interaction between an excited atom and atoms in the ground state, one can use the relation $`\mathrm{\Delta }\nu (n)=\mathrm{\Delta }\nu _c(n/n_c(0))`$. Here $`\mathrm{\Delta }\nu _c`$ is the experimental value of the shift corresponding to the maximum condensate density $`n_c(0)`$.
The dependence $`I(\mathrm{\Delta }\nu )`$ for the density distributions presented in FIG.1 is shown in FIG.2, $`\mathrm{\Delta }\nu `$ being a shift of the laser frequency. Calculating $`I(\mathrm{\Delta }\nu )`$, we took into account the experimental linewidth of $`\delta \nu 3`$kHz (see ).
The curves plotted in FIG.2 reproduce the characteristic behavior and scale of the reconstruction observed in for the two-photon absorption spectrum of the normal fraction with the appearance of a condensate. The nature of this reconstruction lies in the formation of the special energy distribution of particles due to the presence of the quasistationary flux of particles into the small condensate region compensating the spin relaxation losses.
Note in conclusion that all the presented calculations have an estimational character due to adopted simplifications and uncertainty in a self-consistent set of experimental data.
One of the authors (Yu.K.) acknowledges the fruitful discussions with T. J. Greytak, D. Kleppner and L. Willmann. This work was supported by the Russian Foundation for Basic Research (Grants 98-02-16262 and 99-02-18024) and by INTAS (Grant INTAS-97-0972).
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# 1 Introduction
## 1 Introduction
In this note we study the geometrical setting of the super Krichever map in analogy to the standard non graded case . This map is an essential tool in the analysis of the algebraic geometric solutions to integrable systems of soliton type (see, e.g., ). Its super extension has already been introduced in , and studied in . The essential difference in our approach is to make full profit of the so–called Faà di Bruno approach to the KP theory and its super generalization , where the equivalence of this approach to the standard differential-operator approach to the Jacobian SKP is proved. It turns out, as in the classical case, that the Faà di Bruno recursion relation is (related to) the first cocycle condition for the hypercohomology group which controls the infinitesimal deformations of the spectral super line bundle together with its meromorphic sections.
The basic advantage of this approach is that it is directly related to the (Super)Grassmannian description of the hierarchy, and has an intrinsic geometrical meaning. In particular, we can avoid the difficult initial step of the introduction of the Baker–Akhiezer function, go on with the natural development of the geometrical construction, and recover the existence of the BA function at the end.
As in the classical case, the important technical tool is to construct a local universal deformation of the initial super line bundle, since the cocycle condition comes by considering point-wise the germ of such a deformation as a vector field on the base. This is a difficult point because we lack a sound definition of the Super Jacobian of a super curve $`𝒞`$. Indeed looking at the transition functions, one would say that this is the cohomology group $`H^1(𝒞,𝒪_0^\times )`$, where $`𝒪_0^\times `$ is the sheaf of units in the even part of the structure sheaf. Unfortunately, this set-up is not fully satisfactory because there are no naturally defined odd deformation directions. The way out we present in this paper is to work with the moduli space $`S_g\stackrel{~}{𝒞}`$ of effective superdivisors of degree $`g`$ or, which is the same, with the $`g`$–fold symmetric product of the dual curve . This is a supervariety with enough odd parameters over which we have a universal effective divisor, and we expect that, as in the classical case, the “Super Jacobian” will appear as a quotient of $`S_g\stackrel{~}{𝒞}`$.
The scheme of the paper is as follows: in Section 2 we briefly recall the Faà di Bruno recursion relations, their connection with the JSKP hierarchy, and with the Krichever map, referring to for more details. In Section 3 we give a brief resumè of the tools from deformation theory needed in the sequel. In Section 4 we construct the symmetric powers of the (dual) supercurve as a supervariety, and we prove the existence of a universal superdivisor. In the last Section we exploit the cohomological meaning of the Faà di Bruno recursion relations to insure that it gives a flow on the space of super Krichever data and, through the super Krichever map, the JSKP flow on the algebraic geometrical loci in the Super Grassmannian. Finally, in Appendix A we recall some basic definitions of the theory of super curves used in the paper.
## 2 The Jacobian super KP hierarchy
Let us start by fixing some notations. We denote by $`\mathrm{\Lambda }`$ a generic Grassmann algebra over $``$, $`B:=\mathrm{\Lambda }[[x,\phi ]]`$ is the $`\mathrm{\Lambda }`$–algebra of formal power series in the variables $`x`$ (even) and $`\phi `$ (odd) and $`𝒟:=_\phi +\phi _x`$. The ring of formal super pseudo–differential operators over $`X:=\mathrm{Spec}(B)`$ is the space of formal series
$$L:=\underset{j0}{}u_j𝒟^{nj},u_jB$$
endowed with the product induced by the super Leibniz rule
$$𝒟^kf=\underset{j0}{}(1)^{\overline{f}(kj)}\left[\genfrac{}{}{0.0pt}{}{k}{j}\right]f^{(j)}𝒟^{kj},$$
where $`\overline{f}`$ denotes the parity of $`f`$, $`f^{(j)}=𝒟^j(f)`$ and $`\left[\genfrac{}{}{0.0pt}{}{k}{j}\right]`$ is the super binomial coefficient .
Mulase and Rabin defined the Jacobian super KP hierarchy as the following set of evolutionary equations for the even dressing operator $`S:=1+_{j>0}s_j𝒟^j`$:
$$\{\begin{array}{c}_{t_{2k}}S:=(S𝒟^{2k}S^1)_{}S=(S_x^kS^1)_{}S\hfill \\ \\ _{t_{2k1}}S:=(S(𝒟^{2k1}\phi 𝒟^{2k})S^1)_{}S=(S_\phi _x^{k1}S^1)_{}S,\hfill \end{array}$$
where $`L_{}`$ is the pure pseudo–differential part of $`L`$ and the time $`t_k`$ has parity $`kmod2`$. One of the features which distinguishes this hierarchy from that of Manin and Radul is that for algebraic geometric solutions the equations describe super–commuting linear flows on the super Jacobian of a super curve. One way to approach this issue is to consider another description of the hierarchy, using the super Faà di Bruno polynomials instead of super pseudo–differential operators. We refer to for a detailed account and we only sketch what is relevant to the present discussion. Let $`V:=\mathrm{\Lambda }((z^1))\mathrm{\Lambda }((z^1))\theta `$ be the algebra of formal Laurent series in the even variable $`z^1`$ and odd variable $`\theta `$, let $`V_+:=\mathrm{\Lambda }[z,\theta ]`$, $`V_{}:=\mathrm{\Lambda }[[z^1,\theta ]]z^1`$ and let $`V_B:=V_\mathrm{\Lambda }B`$. The basic object of this formulation is the odd Faà di Bruno generator
$$\widehat{h}(z,\theta ;x,\phi ):=\theta +\phi z+O(z^1)V_B$$
where, abusing notations, we write $`O(z^1)`$ for an element of $`V_{}_\mathrm{\Lambda }B`$. Out of $`\widehat{h}`$ we construct iteratively the Faà di Bruno polynomials by
$$\{\begin{array}{c}\widehat{h}^{(0)}:=1\hfill \\ \\ \widehat{h}^{(k+1)}:=(𝒟+\widehat{h})\widehat{h}^{(k)}k\hfill \end{array}$$
(2.1)
and set $`W_B:=\mathrm{span}_B\{\widehat{h}^{(k)}:k\}`$. It is then easy to show that there exists a unique basis $`\{\widehat{H}^{(k)},k\}`$ of $`W_B`$, whose elements (called “super currents”) have the form
$$\widehat{H}^{(2k+p)}=\theta ^pz^k+O(z^1)$$
(2.2)
with $`p=0,1`$, in terms of which the equations of the Jacobian super KP hierarchy become
$$\frac{\widehat{h}}{t_k}=(1)^k𝒟\widehat{H}^{(k)}.$$
(2.3)
Since $`\widehat{H}^{(2)}=\widehat{h}^{(2)}`$ we have $`_{t_2}_x`$.
The study of these equations finds its most appropriate and natural setting in the concept of super universal Grassmannian $`\mathrm{SGr}_\mathrm{\Lambda }`$ defined as follows . The filtration $`\mathrm{}V_{j1}V_jV_{j+1}\mathrm{}V`$, where $`V_j=z^{j+1}V_{}`$, makes $`V`$ and its $`\mathrm{\Lambda }`$–submodule $`V_+`$ complete topological spaces and $`\mathrm{SGr}_\mathrm{\Lambda }:=\mathrm{SGr}_\mathrm{\Lambda }(V,V_+)`$ is the set of closed free $`\mathrm{\Lambda }`$–submodules $`W`$ of $`V`$ which are compatible with $`V_+`$ in the sense that the restriction $`\pi _W`$ of the projection $`\pi :VV_+`$ to $`W`$ is a Fredholm operator, i.e. its kernel (respectively cokernel) is a $`\mathrm{\Lambda }`$–submodule (respectively a quotient $`\mathrm{\Lambda }`$–module) of a finite rank free $`\mathrm{\Lambda }`$–module. As in the commutative setting, $`\mathrm{SGr}_\mathrm{\Lambda }`$ is the disjoint union of the denumerable set of its components $`\mathrm{SGr}_\mathrm{\Lambda }^{(i)}`$ labelled by the index $`i_W`$ of $`\pi _W`$, moreover each component acquires a structure of super scheme by means of projective limits. By definition, the space $`W_B`$ spanned by the $`\widehat{H}^{(k)}`$’s gives rise to a moving point of $`\mathrm{SGr}_\mathrm{\Lambda }`$ and the super currents evolve under JSKP along the equations of a dynamical system, the super central system , which gives vector fields on the Grassmannian. In particular one has that
$$_{t_j}\widehat{H}^{(k)}=(1)^{jk}_{t_k}\widehat{H}^{(j)}.$$
(2.4)
Whichever approach one takes, the link with algebraic geometric solutions is provided by the super Krichever map which associates a point $`W`$ of $`\mathrm{SGr}_\mathrm{\Lambda }`$ to the datum $`(𝒞,D,(z^1,\theta ),,\eta )`$ of
* a $`\mathrm{\Lambda }`$–super–curve $`𝒞=(C,𝒪_𝒞)`$ (see Appendix A),
* an irreducible divisor $`D`$ on $`𝒞`$ whose reduced support is a smooth point $`p_{\mathrm{}}C`$,
* local coordinates $`z^1`$ and $`\theta `$ in a neighbourhood $`U_0p_{\mathrm{}}`$,
* an invertible sheaf $``$ on $`𝒞`$ and
* a local trivialization $`\eta `$ of $``$ over $`U_0`$.
Let $`(\mathrm{}D)=lim_n\mathrm{}(nD)`$ be the sheaf of sections of $``$ with at most an arbitrary pole at $`D`$, then $`W=\eta (H^0(𝒞,(\mathrm{}D)))`$. Bergvelt and Rabin have shown in that the $`\mathrm{\Lambda }`$–module $`H^0(𝒞,(\mathrm{}D))`$ is free, so $`W`$ belongs indeed to $`\mathrm{SGr}_\mathrm{\Lambda }`$. We can invert the Krichever map on its image as explained in ; in particular, the ring of functions of $`𝒞`$ which are holomorphic on the open subset $`U_1:=C\{p_{\mathrm{}}\}`$ is the subalgebra $`𝒜_WV`$ of functions $`f`$ such that $`fWW`$. $`𝒜_W`$ is obviously graded.
As a consequence of equations (2.1) and (2.3), we recover the same picture in our approach:
###### Proposition 2.1 (Isospectrality)
Let $`\widehat{h}(x,\phi ,𝐭)`$ be a solution of the Jacobian super KP hierarchy and denote by $`W_T`$ the space generated by the corresponding super currents $`\widehat{H}^{(k)}(x,\phi ,𝐭)`$. For any specialization $`(x_0,\phi _0,𝐭_0)`$ of $`(x,\phi ,𝐭)`$ let $`𝒜_{(x_0,\phi _0,𝐭_0)}V`$ be the $`\mathrm{\Lambda }`$–algebra of functions that map by multiplication $`W_{T_0}`$ into itself. Then $`𝒜_{(x_0,\phi _0,𝐭_0)}`$ does not depend on $`(x_0,\phi _0,𝐭_0)`$. $`\mathrm{}`$
We limit ourselves to sketch the proof. We have to show that if $`f𝒜_{(x_0,\phi _0,𝐭_0)}`$, then $`fW_TW_T`$. Since $`1W_T`$, this is equivalent to showing that such an $`f`$ is in $`W_T`$, because $`f`$ supercommutes with $`𝒟+\widehat{h}`$. Since $`1W_{T_0}`$ as well, we can write
$$f=c_j\stackrel{~}{H^{(j)}}$$
where $`\stackrel{~}{H^{(j)}}`$ denote the specialization of $`\widehat{H}^{(j)}`$ at $`𝐭=𝐭_0`$. We have to prove that, calling
$$f^{}=c_j\widehat{H}^{(j)},$$
actually $`f^{}=f`$, that is that $`c_j\widehat{H}^{(j)}`$ is independent of the times $`t_k`$.
The identity:
$$f(𝒟+\widehat{h})^k=\underset{j0}{}(1)^{\frac{j(j+1)}{2}+k\overline{f}}\left[\genfrac{}{}{0pt}{}{k}{j}\right](𝒟+\widehat{h})^{kj}f^{(j)}$$
shows that $`𝒟^kf^{}W_T,k`$, so that $`𝒟f^{}=0`$. Similarly, one proves that $`_{t_k}f^{}=0k`$.
## 3 Deformation of super line bundles and of their sections
The meaning of the Isospectrality Lemma 2.1 is that when the solution is of algebraic geometric type the super curve $`𝒞`$ (also called the spectral curve) remains unaffected by the flows of the hierarchy. Indeed, it is also true that the divisor $`D`$ and the coordinates $`(z^1,\theta )`$ do not change, so the motion involves only the line bundle $``$ and its local trivialization $`\eta `$. Since our aim is to interpret geometrically the equations (2.1) and (2.3), which are of differential type, and the super Krichever map is defined in terms of sections of a super line bundle $``$, we have to understand how the sections change when we deform $``$.
###### Definition 3.1
Let $`𝒞`$ be a $`\mathrm{\Lambda }`$–super–curve, $``$ an invertible sheaf on $`𝒞`$, $`s`$ a global section of $``$ and $`(𝒳,x)`$ a pointed $`\mathrm{\Lambda }`$–super–scheme. An $`𝒳`$–family of invertible sheaves on $`𝒞`$ is an invertible sheaf $`_𝒳`$ over $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒳`$. A deformation of $`(,s)`$ over the pointed super–scheme $`(𝒳,x)`$ is a triple $`(_𝒳,\sigma ,\rho )`$ where
* $`_𝒳`$ is an $`𝒳`$-family of invertible sheaves on $`𝒞`$,
* $`\sigma `$ is a global section of $`_𝒳`$ and
* $`\rho `$ is an isomorphism $`\rho :\iota ^{}_𝒳`$, where $`\iota :𝒞𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒳`$ is the embedding identifying $`𝒞`$ with $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}\{x\}`$, such that $`\iota ^{}\sigma =\rho s`$.
Two deformations $`(_𝒳,\sigma ,\rho )`$ and $`(𝒩_𝒳,\tau ,\xi )`$ of $`(,s)`$ over $`(𝒳,x)`$ are isomorphic if and only if there exists an isomorphism of sheaves $`\eta :_𝒳𝒩_𝒳`$ compatible with $`\rho `$ and $`\xi `$ ($`\xi =\iota ^{}(\eta )\rho `$) and such that $`tau=\eta (\sigma )`$. The line bundle $`_𝒳|_{𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}\{x\}}`$ is sometimes called the central fibre of the deformation. Finally, an infinitesimal deformation of $`(,s)`$ is a deformation over the “one–point” $`\mathrm{\Lambda }`$–super-scheme
$$:=\mathrm{Spec}(\mathrm{\Lambda }[t,\epsilon ]/t^2,t\epsilon ),$$
where $`t`$ is even and $`\epsilon `$ is odd.
Let $`\{U_j\}_{jJ}`$ be a covering by open affine sub–super–schemes of $`𝒞`$ and denote by $`U_{j_1,\mathrm{},j_k}`$ the intersection $`_{l=1}^kU_{j_l}`$, by $`𝒪_{j_1,\mathrm{},j_k}`$ the super–commutative ring of sections of $`𝒪_𝒞`$ over $`U_{j_1,\mathrm{},j_k}`$ and by $`_{j_1,\mathrm{},j_k}`$ the $`𝒪_{j_1,\mathrm{},j_k}`$–module of sections of $``$ over $`U_{j_1,\mathrm{},j_k}`$. Finally, define
$$\{\begin{array}{c}𝒪_{j_1,\mathrm{},j_k}[t,\epsilon ]:=𝒪_{j_1,\mathrm{},j_k}_\mathrm{\Lambda }𝒪_{}\hfill \\ \\ _{j_1,\mathrm{},j_k}[t,\epsilon ]:=_{j_1,\mathrm{},j_k}_\mathrm{\Lambda }𝒪_{}\hfill \\ \\ U_{j_1,\mathrm{},j_k}[t,\epsilon ]:=\mathrm{Spec}(𝒪_{j_1,\mathrm{},j_k}[t,\epsilon ])=U_{j_1,\mathrm{},j_k}\times _{\mathrm{Spec}(\mathrm{\Lambda })}\hfill \end{array}.$$
Then, $`\{U_j[t,\epsilon ]\}_{jJ}`$ is an open affine covering of $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}`$ and the exact sequence of sheaves
$$\begin{array}{ccccccccc}0& & 𝒪_j& & 𝒪_j[t,\epsilon ]_0^\times & & 𝒪_{j,0}^\times & & 1\\ & & & & & & & & \\ & & f& & 1+tf_0+\epsilon f_1& & & & \end{array},$$
where $`f_0`$ and $`f_1`$ are the even and odd components of $`f`$, yields a group isomorphism $`\mathrm{Pic}(U_j[t,\epsilon ])\mathrm{Pic}(U_j)`$ due to the fact that the $`U_j`$’s are Stein (see , 1.3.8). Thus, if $`_{}`$ is an infinitesimal deformation of $``$ then
$$_{}|_{U_j[t,\epsilon ]}\left(|_{U_j}\right)[t,\epsilon ],$$
so it is described as the gluing of the last modules by means of a suitable isomorphism
$$G_{jk}:_{jk}[t,\epsilon ]\stackrel{}{}_{jk}[t,\epsilon ],$$
which in turn is given by the transition matrix
$$G_{jk}=\left(\begin{array}{ccc}g_{jk}& 0& 0\\ \delta _tg_{jk}& g_{jk}& 0\\ \delta _\epsilon g_{jk}& 0& g_{jk}\end{array}\right),$$
where we express an element $`\sigma _j_j[t,\epsilon ]`$, $`\sigma _j=f_j+t\delta _tf_j+\epsilon \delta _\epsilon f_j`$ as a column vector $`(f_j,\delta _tf_j,\delta _\epsilon f_j)^t`$, $`\delta _tg_{jk}𝒪_{jk,0}`$, $`\delta _\epsilon g_{jk}𝒪_{jk,1}`$ and $`g_{jk}`$ is the transition function of $``$. The cocycle condition for $`G_{jk}`$ implies that $`\{g_{jk}^1(\delta _tg_{jk}+\delta _\epsilon g_{jk})\}_{jk}`$ is a $`1`$-cocycle $`c_1`$ on $`𝒞`$ with values in $`𝒪_𝒞`$. Clearly, if we change $`c_1`$ by a coboundary we get an isomorphic infinitesimal deformation of the invertible sheaf $``$. Hence, the set of isomorphism classes of infinitesimal deformations of $``$ is isomorphic to $`H^1(𝒞,𝒪_𝒞)`$. If we have a deformation $`_𝒳`$ of $``$ over $`(𝒳,x)`$ and $`v:𝒳`$ is a ”tangent vector“ to $`𝒳`$ at $`x`$, then the pull–back of $`_𝒳`$ under $`id_𝒞\times v`$ is an infinitesimal deformation of $``$ and corresponds by the above argument to a class $`[c_1]H^1(𝒞,𝒪_𝒞)`$. This defines the Kodaira–Spencer map $`KS:T_x𝒳H^1(𝒞,𝒪_𝒞)`$ of the deformation.
Now we consider the deformation $`\sigma H^0(𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })},_{})`$ of $`sH^0(𝒞,)`$. Let us write the local expression of $`\sigma `$ as above: $`\sigma _j=f_j+t\delta _tf_j+\epsilon \delta _\epsilon f_j`$, where $`f_j`$ is the local function representing $`s`$. Then, the cocycle condition for $`\sigma `$ to be a global section reads
$$\{\begin{array}{c}g_{jk}^1\delta _tf_j=\delta _tf_k+g_{jk}^1\delta _tg_{jk}f_k\hfill \\ \\ g_{jk}^1\delta _\epsilon f_j=\delta _\epsilon f_k+g_{jk}^1\delta _\epsilon g_{jk}f_k\hfill \end{array}.$$
(3.1)
The meaning of these two equations is the following (see e.g. and the Appendix of ): the triple $`(\{U_j\}_j,\{g_{jk}^1(\delta _tg_{jk}+\delta _\epsilon g_{jk})\}_{jk},\{\delta _tf_j+\delta _\epsilon f_j\}_j)`$ gives rise to a class $`\gamma _1_s^1(𝒞,^{})`$ of the hyper–cohomology of the complex
$$^{}:0𝒪_𝒞\stackrel{s}{}0$$
of sheaves on $`𝒞`$. The set of isomorphism classes of infinitesimal deformations of $`(,s)`$ is isomorphic to $`_s^1(𝒞,^{})`$, and a corresponding Kodaira–Spencer map can be defined for any deformation.
Our goal is to show that the similarity between equation (3.1) and the second equation in (2.1) is not only formal, that is, we can interpret Eq.s (3.1) as the differential equations associated with a Kodaira–Spencer deformation of the spectral super line bundle together with its meromorphic sections,
$$\{\begin{array}{c}g_{jk}^1_tf_j=(_t+g_{jk}^1_tg_{jk})f_k\hfill \\ g_{jk}^1_\epsilon f_j=(_\epsilon +g_{jk}^1_\epsilon g_{jk})f_k\hfill \end{array}.$$
(3.2)
To achieve this we have first of all to construct a suitable family $`_𝒳`$ of line bundles on a $`\mathrm{\Lambda }`$–super–curve $`𝒞`$ and then to interpret the Faà di Bruno polynomials as local representatives of sections of $`_𝒳`$. These two steps will be taken in the next two Sections.
## 4 The universal relative positive super divisor
From now on we assume that $`𝒞`$ is a smooth super curve over $`\mathrm{\Lambda }`$. Since the points of the super universal Grassmannian associated with a solution of the Jacobian super KP hierarchy belong to the component of index $`0|0`$ we must require $`𝒞`$ to be a generic SKP curve and $``$ to have degree $`g`$ equal to the genus of $`𝒞`$.
###### Definition 4.1
(SKP curve ) A $`\mathrm{\Lambda }`$–super–curve $`𝒞=(C,𝒪_𝒞)`$ is called an SKP curve if its split structure sheaf $`𝒪_𝒞^{\mathrm{sp}}:=𝒪_𝒞_\mathrm{\Lambda }\mathrm{\Lambda }/𝔪`$ is of the form $`𝒪_𝒞^{\mathrm{rd}}|𝒮`$, where $`𝔪`$ is the maximal ideal of nilpotent elements of $`\mathrm{\Lambda }`$, $`𝒮`$ is an invertible $`𝒪_𝒞^{\mathrm{rd}}`$-module (a “reduced” line bundle) of degree zero and $`|`$ denotes a direct sum of free $`\mathrm{\Lambda }`$–modules, with on the left an evenly generated summand and on the right an odd one. If $`𝒮𝒪_𝒞^{red}`$ then $`𝒞`$ is called a generic SKP curve.
Let $`\stackrel{~}{𝒞}`$ be the dual super curve of $`𝒞`$, whose $`\mathrm{\Lambda }`$-points are the irreducible superdivisors of $`𝒞`$ (see Appendix A). Constructing a universal family of line bundles $`_𝒳`$ requires the construction of the super Picard scheme of $`𝒞`$ and the corresponding super Poincaré sheaf. However we can avoid this difficult step, since it suffices to produce the universal super divisor $`\mathrm{\Delta }^{(g)}`$ of degree $`g`$. In analogy to the commutative case, the central object we have to consider is the $`g`$–th symmetric product $`S_g\stackrel{~}{𝒞}`$ of the dual super curve $`\stackrel{~}{𝒞}`$, since $`\stackrel{~}{𝒞}`$ parameterizes irreducible positive super divisors on $`𝒞`$. Our discussion will follow closely that of , the only novelty being that we have to work over $`\mathrm{\Lambda }`$, instead over $``$.
Let $`𝒞^g:=𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}\mathrm{}\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒞`$ be the $`g`$–fold fibred product of $`𝒞`$ with itself over $`\mathrm{Spec}(\mathrm{\Lambda })`$. The symmetric group $`\mathrm{\Sigma }_g`$ of degree $`g`$ acts on $`𝒞^g`$ by
$$\begin{array}{cccc}\mathrm{\Sigma }_g\sigma :& C^g& & C^g\\ & & & \\ & (x_1,\mathrm{},x_g)& & (x_{\sigma (1)},\mathrm{},x_{\sigma (g)})\end{array}$$
and
$$\begin{array}{cccc}\sigma :& 𝒪_𝒞^{_\mathrm{\Lambda }g}& & 𝒪_𝒞^{_\mathrm{\Lambda }g}\\ & & & \\ & f_1_\mathrm{\Lambda }\mathrm{}_\mathrm{\Lambda }f_g& & \left(_{\genfrac{}{}{0.0pt}{}{j<k}{\sigma (j)>\sigma (k)}}(1)^{\overline{f}_{\sigma (j)}\overline{f}_{\sigma (k)}}\right)f_{\sigma (1)}\mathrm{}f_{\sigma (g)},\end{array}$$
(4.1)
where $`C`$ is the reduced curve associated with $`𝒞`$. We define the $`g`$–th symmetric product of $`𝒞`$ to be the ringed space
$$S_g𝒞:=(C^g/\mathrm{\Sigma }_g,(𝒪^{_\mathrm{\Lambda }g})^{\mathrm{\Sigma }_g}),$$
whose structure sheaf is the graded sheaf of invariants of $`𝒪^{_\mathrm{\Lambda }g}`$. Notice that, since $`\sigma `$ is an even map (i.e. it preserves degrees), the action above is the same as that in eq. (1) of . The form given above makes the proof of the following proposition quite immediate.
###### Proposition 4.1
The super space $`S_g𝒞`$ is a supermanifold over $`\mathrm{Spec}(\mathrm{\Lambda })`$ of dimension $`g|g`$.
Proof. It is well known that $`S_gC:=C^g/\mathrm{\Sigma }_g`$ is a smooth scheme, so we have to show that locally $`𝒪_{S_g𝒞}`$ is isomorphic to $`𝒪_{S_gC}\mathrm{\Lambda }[\theta _1,\mathrm{},\theta _g]`$. Obviously $`\mathrm{\Lambda }𝒪_{S_g𝒞}`$. By definition there exists an open covering $`\{U_j\}_{jJ}`$ of $`𝒞`$ such that $`𝒪_𝒞(U_j)\mathrm{\Lambda }_{}𝒪_C(U_j)[\theta _j]`$. Let $`p:C^gS_gC`$ be the natural projection of ordinary schemes. One has only to prove that if $`V`$ is an open affine subscheme of $`S_gC`$ such that $`𝒪_{𝒞^g}(U)`$ ($`U=p^1V`$) is isomorphic to $`\mathrm{\Lambda }_{}(𝒪_C(U)[\theta ])^_{}g`$, then $`𝒪_{S_g𝒞}(V)𝒪_{S_gC}(V)\mathrm{\Lambda }[\varsigma _1,\mathrm{},\varsigma _g]`$, for suitable odd coordinates $`\varsigma _1,\mathrm{},\varsigma _g`$. Now, $`\sigma \mathrm{\Sigma }_g`$ acts as the identity on the first factor $`\mathrm{\Lambda }`$: in fact we have
$`\sigma (\lambda _1f_1\mathrm{}\lambda _gf_g)=\sigma \left({\displaystyle \underset{j<k}{}}(1)^{\overline{f}_j\overline{\lambda }_k}(\lambda _1\mathrm{}\lambda _g)f_1\mathrm{}f_g\right)`$
$`=\left({\displaystyle \underset{\genfrac{}{}{0.0pt}{}{l<m}{\sigma (l)>\sigma (m)}}{}}(1)^{(\overline{f}_{\sigma (l)}+\overline{\lambda }_{\sigma (l)})(\overline{f}_{\sigma (m)}+\overline{\lambda }_{\sigma (m)})}\right)\lambda _{\sigma (1)}f_{\sigma (1)}\mathrm{}\lambda _{\sigma (g)}f_{\sigma (g)}`$
$`=\left({\displaystyle \underset{\genfrac{}{}{0.0pt}{}{l<m}{\sigma (l)>\sigma (m)}}{}}(1)^{\overline{f}_{\sigma (l)}\overline{\lambda }_{\sigma (m)}+\overline{f}_{\sigma (m)}\overline{\lambda }_{\sigma (l)}+\overline{f}_{\sigma (l)}\overline{f}_{\sigma (m)}}\right)\left({\displaystyle \underset{j<k}{}}(1)^{\overline{f}_{\sigma (j)}\overline{\lambda }_{\sigma (k)}}\right)\times `$
$`\lambda _1\mathrm{}\lambda _gf_{\sigma (1)}\mathrm{}f_{\sigma (g)}`$
$`=\left({\displaystyle \underset{\genfrac{}{}{0.0pt}{}{l<m}{\sigma (l)>\sigma (m)}}{}}(1)^{\overline{f}_{\sigma (l)}\overline{\lambda }_{\sigma (m)}+\overline{f}_{\sigma (m)}\overline{\lambda }_{\sigma (l)}}\right)\left({\displaystyle \underset{j<k}{}}(1)^{\overline{f}_{\sigma (j)}\overline{\lambda }_{\sigma (k)}}\right)\times `$
$`\lambda _1\mathrm{}\lambda _g\sigma (f_1\mathrm{}f_g)`$
and since
$$\left(\underset{\genfrac{}{}{0.0pt}{}{l<m}{\sigma (l)>\sigma (m)}}{}(1)^{\overline{f}_{\sigma (l)}\overline{\lambda }_{\sigma (m)}+\overline{f}_{\sigma (m)}\overline{\lambda }_{\sigma (l)}}\right)\left(\underset{j<k}{}(1)^{\overline{f}_{\sigma (j)}\overline{\lambda }_{\sigma (k)}}\right)=\underset{j<k}{}(1)^{\overline{f}_j\overline{\lambda }_k}$$
we get $`\sigma (\lambda _1\mathrm{}\lambda _gf_1\mathrm{}f_g)=\lambda _1\mathrm{}\lambda _g\sigma (f_1\mathrm{}f_g)`$. Therefore, it remains only to apply Theorem 1 of . If $`(z,\theta )`$ are graded local coordinates of $`𝒞`$ then a system of graded local coordinates for $`S_g𝒞`$ is given by $`(s_1,\mathrm{},s_g,\varsigma _1,\mathrm{},\varsigma _g)`$, where $`(s_1,\mathrm{},s_g)`$ are the (even) symmetric functions of $`z_j=1\mathrm{}z\mathrm{}1`$ (with $`z`$ in the $`j`$–th position), $`1jg`$ and $`(\varsigma _1,\mathrm{},\varsigma _g)`$ are the odd symmetric functions defined by $`\varsigma _j:=_{k=1}^g\theta _k\stackrel{~}{s}_{j1}^{(k)}`$, where $`\theta _k:=1\mathrm{}\theta \mathrm{}1`$ and $`\stackrel{~}{s}_j^{(k)}`$ is the $`j`$–th symmetric function of $`z_1,\mathrm{},z_{k1},z_{k+1},\mathrm{},z_g`$. $`\mathrm{}`$
To exploit this construction we give the following definition.
###### Definition 4.2
Let $`𝒳=(X,𝒪_𝒳)`$ be a super scheme over $`\mathrm{Spec}(\mathrm{\Lambda })`$. A positive relative super divisor of degree $`g`$ of $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒳𝒳`$ is a closed sub–super–scheme $`𝒵`$ of $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒳`$ of codimension $`1|0`$, defined by a homogeneous locally principal ideal $`𝒥`$ of $`𝒪_{𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒳}`$, such that $`𝒪_𝒵`$ is a locally free $`𝒪_𝒳`$–module of rank $`g|0`$ and its reduction (modulo nilpotents) $`Z`$ is a positive relative divisor of degree $`g`$ of $`C\times XX`$.
By definition, then, $`𝒵`$ is locally defined by an equation of type
$$f=z^g(a_1+\theta \alpha _1)z^{g1}+\mathrm{}+(1)^g(a_g+\theta \alpha _g)=0,$$
where $`f`$ is the local generator of $`𝒥`$ and the $`a_j`$’s (respectively the $`\alpha _j`$’s) are even (respectively odd) local functions on $`𝒳`$. Our aim is to show that the symmetric product $`S_g\stackrel{~}{𝒞}`$ is the parameter space for the universal relative super divisor of degree $`g`$, $`\mathrm{\Delta }^{(g)}`$, of $`𝒞`$. The universal relative super divisor of degree $`1`$ is simply the sub–super–scheme $`\mathrm{\Delta }^{(1)}`$ of $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}\stackrel{~}{𝒞}`$ locally defined by the equation
$$z_\mathrm{\Lambda }11_\mathrm{\Lambda }\stackrel{~}{z}\theta _\mathrm{\Lambda }\stackrel{~}{\rho }=0,$$
which we will write more compactly as $`z\stackrel{~}{z}\theta \stackrel{~}{\rho }=0`$, where $`(z,\theta )`$ are local coordinates of $`𝒞`$ and $`(\stackrel{~}{z},\stackrel{~}{\rho })`$ are the ”dual“ coordinates given in Appendix A, equation (A.2). Consider now the natural projections
$$\begin{array}{cccc}\pi _j:& 𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}\stackrel{~}{𝒞}^g& & 𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}\stackrel{~}{𝒞}\\ & (x,\stackrel{~}{x}_1,\mathrm{},\stackrel{~}{x}_g)& & (x,\stackrel{~}{x}_j)\end{array}$$
and define $`\stackrel{~}{\mathrm{\Delta }}_j:=\pi _j^1(\mathrm{\Delta }^{(1)})`$, $`\stackrel{~}{\mathrm{\Delta }}^{(g)}:=\stackrel{~}{\mathrm{\Delta }}_1+\mathrm{}+\stackrel{~}{\mathrm{\Delta }}_g`$. Since the local equation of $`\stackrel{~}{\mathrm{\Delta }}^{(g)}`$ is
$$\underset{j=1}{\overset{g}{}}(z\stackrel{~}{z}_j\theta \stackrel{~}{\rho }_j)=z^g(s_1+\theta \varsigma _1)z^{g1}+\mathrm{}+(1)^g(s_g+\theta \varsigma _g)=0,$$
where the $`s_j`$’s and the $`\varsigma _k`$’s are the symmetric functions of the $`\stackrel{~}{z}_m`$’s and $`\stackrel{~}{\rho }_n`$’s we introduced at the end of the proof of Proposition 4.1, the next lemma holds true.
###### Lemma 4.2
There exists a unique positive relative super divisor $`\mathrm{\Delta }^{(g)}`$ of degree $`g`$ of $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}S_g\stackrel{~}{𝒞}S_g\stackrel{~}{𝒞}`$ such that $`\pi ^{}(\mathrm{\Delta }^{(g)})=\stackrel{~}{\mathrm{\Delta }}^{(g)}`$, where $`\pi :𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}\stackrel{~}{𝒞}^g𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}S_g\stackrel{~}{𝒞}`$ is the natural projection. $`\mathrm{}`$
The most important result we need is Theorem 6 of , whose proof extends to the present situation.
###### Theorem 4.3
The pair $`(S_g\stackrel{~}{𝒞},\mathrm{\Delta }^{(g)})`$ represents the functor of relative positive super divisors of degree $`g`$ of $`𝒞`$, i.e. the natural map
$$\begin{array}{cccc}R:& \mathrm{Hom}(𝒳,S_g\stackrel{~}{𝒞})& & \mathrm{Div}_𝒳^g(𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒳)\\ & f& & (id\times f)^{}\mathrm{\Delta }^{(g)}\end{array}$$
is a functorial isomorphism for every $`\mathrm{\Lambda }`$-super scheme $`𝒳`$. $`\mathrm{}`$
## 5 The geometric super Faà di Bruno polynomials
The constructions of the previous Section allow us to define a canonical family of super line bundles together with an even section. For simplicity we call $`𝒳:=S_g\stackrel{~}{𝒞}`$ and we assume also that $`𝒞`$ is a generic SKP super curve of genus $`g`$. Then $`_𝒳:=𝒪_{𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒳}(\mathrm{\Delta }^{(g)})`$ is a $`𝒳`$–family of super line bundles on $`𝒞`$ and $`\mathrm{\Delta }^{(g)}`$ defines a section $`\sigma `$ of $`_𝒳`$. If we let $``$ be any non–special super line bundle on $`𝒞`$ (i.e. such that the reduced invertible sheaf $`^{\mathrm{rd}}`$ is non–special on $`C`$) of degree $`g`$ and we call $`s_{}`$ the unique (up to multiplication by a complex number) section that generates the even part of $`H^0(𝒞,)`$, then the divisor $`(s_{})`$ can be thought of as a $`\mathrm{Spec}(\mathrm{\Lambda })`$–family of positive relative super divisors of degree $`g`$ and the universality property of $`\mathrm{\Delta }^{(g)}`$ (Theorem 4.3) gives a unique map $`f_{}:\mathrm{Spec}(\mathrm{\Lambda })𝒳`$, i.e. a $`\mathrm{\Lambda }`$–point $`x`$ of $`𝒳`$, such that $`(s_{})=(id\times f_{})^{}\mathrm{\Delta }^{(g)}`$. In turn, this induces an isomorphism $`\rho _{}:(id\times f_{})^{}_𝒳`$ such that $`(id\times f_{})^{}\sigma =\rho _{}s_{}`$, so we can interpret the triple $`(𝒳,_𝒳,\sigma )`$ as a deformation of $`(,s_{})`$ for any non–special super line bundle $``$ on $`𝒞`$. Finally, if we put graded super coordinates $`𝐭=(t_1,\mathrm{},t_{2g})`$ on $`𝒳`$ ($`\overline{t}_j=jmod2`$) then the cocycle conditions (3.1) for the section $`\sigma `$ as a deformation of $`s_{𝐭_0}:=\sigma |_{𝒞\times \{x(𝐭_0)\}}`$, for any $`𝐭_0`$, become the differential equations
$$g_{jk}^1_{t_l}f_j=_{t_l}f_k+(g_{jk}^1_{t_l}g_{jk})f_k,$$
which are manifestly of the form of (2.1). To accomplish our goal of describing the algebraic geometric super Faà di Bruno polynomials, we have therefore only to choose a suitable open covering of $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒳`$ and to appropriately select two coordinates $`t_{2j}`$ and $`t_{2k+1}`$ and to call them $`x`$ and $`\phi `$ respectively.
As before select a non–special super line bundle $``$ of degree $`g`$ on $`𝒞`$. Let $`p_{\mathrm{}}C`$ be a reduced point of $`C`$ such that it is not Weierstrass for $`C`$ and the reduced section $`s_{}^{\mathrm{rd}}`$ does not vanish at $`p_{\mathrm{}}`$. Let $`U_0C`$ be an open neighbourhood of $`p_{\mathrm{}}`$ where we can define graded coordinates $`(z,\theta )`$ for $`𝒞`$ centered at $`p_{\mathrm{}}`$ (i.e. $`z(p_{\mathrm{}})=0`$) and let $`U_1:=C\{p_{\mathrm{}}\}`$. Then $`\{U_0,U_1\}`$ is a Stein open covering of $`𝒞`$. Since $`s_{}^{\mathrm{rd}}`$ does not vanish at $`p_{\mathrm{}}`$ the section $`s_{}`$ gives a local trivialization $`\eta `$ of $``$ on $`U_0`$ (suitably restricted). Then the quintuple $`(𝒞,D:=(z)|_{U_0},(z,\theta ),,\eta )`$ defines through the super Krichever map a point of $`\mathrm{SGr}_\mathrm{\Lambda }`$. Finally, let $`𝒱`$ be a Stein open neighbourhood of $`(s_{})𝒳`$ where the coordinates $`𝐭`$ are defined. The open subsets $`𝒰_0:=U_0\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒱`$ and $`𝒰_1:=U_1\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒱`$ define a Stein covering of $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒱`$ over which we can trivialize $`_𝒱:=_𝒳|_{𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒱}`$. Restricting $`𝒱`$ if necessary we can assume that $`\sigma `$ gives a local trivialization of $`_𝒱`$ over $`𝒰_0`$.
Now we move to the analytic category instead of the algebraic one. Let $`𝒩:=\pi ^{}𝒪_𝒞(gD)`$, where $`\pi `$ is now the projection of $`𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒱`$ to $`𝒞`$, and let $`\mu `$ be the pull back by $`\pi `$ of the section of $`𝒪_𝒞(gD)`$ which generates the even part of its module of global sections. Then $`\mu `$ gives a local trivialization of $`𝒩`$ over $`𝒰_1`$. Since $`_𝒱𝒩^1`$ has relative degree $`0`$ it follows that, restricting again $`𝒱`$ if necessary, it has a local analytic trivialization $`\nu `$ over $`𝒰_1`$ and $`\tau =\nu \mu `$ gives a trivialization of $`_𝒱`$ over $`U_1`$.
Summarizing, we have a trivialization $`(\sigma ,\tau )`$ of $`_𝒱`$ over $`(𝒰_0,𝒰_1)`$, with respect to which $`\sigma `$ is represented by the couple of functions $`(f_0=1,f_1)`$ and the transition function of $`_𝒱`$ is $`g_{10}=f_1/f_0`$. Let us define $`\widehat{H}^{(k)}:=_{t_k}\mathrm{log}g_{10}`$. Then these meromorphic functions on $`U_0`$ satisfy equation (2.4) and are therefore our candidates for the super currents of the hierarchy.
Observe that it is possible to choose the coordinates $`t_k`$ in such a way that (multiplying $`\tau `$ by the exponential of a suitable meromorphic function whose poles are only over $`\pi ^1(p_{\mathrm{}})`$) $`\widehat{H}^{(h)}`$ has the correct asymptotic behaviour (2.2) (here our coordinate $`z`$ is the inverse of the $`z`$ appearing there). Notice also that $`\widehat{H}^{(k)}|_{𝒰_0𝒰_1}`$ represents the class of $`H^1(𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒱,𝒪_{𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒱})`$ corresponding to the deformation of $``$ along $`t_k`$. Since the asymptotic behaviour of $`\widehat{H}^{(1)}`$ is $`\theta +O(z)`$ it follows that the first time $`t_1`$ does not deform $``$ at all. The super Jacobian $`\mathrm{Jac}(𝒞)`$ of $`𝒞`$ has dimension $`g|g1`$ while $`S_g\stackrel{~}{𝒞}`$ has dimension $`g|g`$ and maps surjectively to $`\mathrm{Pic}^g(𝒞)\mathrm{Jac}(𝒞`$), hence there is an odd direction in $`𝒱`$ which corresponds to trivial deformations of $``$, i.e. there indeed exists a coordinate like $`t_1`$.
In , Section 2.3 we have shown that the Faà di Bruno generator is computed by the formula $`\widehat{h}:=\widehat{H}^{(1)}|_{t_1+\phi }+\phi \widehat{H}^{(2)}|_{t_1+\phi }`$. The cocycle condition (3.1) can be interpreted also as saying that $`(_{t_k}f_0+\widehat{H}^{(k)}f_0,_{t_k}f_1)`$ is a section of $`_𝒱(\mathrm{}\pi ^{}D)`$ with pole of order $`k`$ at $`\pi ^{}D`$. Thus the super Faà di Bruno recurrence relation (2.1) corresponds to deformation along the non–integrable vector field $`𝒟`$ and the super Faà di Bruno polynomials $`\widehat{h}^{(k)}`$ are the local representatives on $`𝒰_0`$ of the meromorphic sections $`\sigma ^{(k)}`$ of $`_𝒱`$ obtained by iterative deformation of $`\sigma ^{(0)}:=\sigma `$ along $`𝒟`$. The form of $`\widehat{h}`$ implies also that the $`\sigma ^{(k)}`$’s form a basis of $`H^0(𝒞\times _{\mathrm{Spec}(\mathrm{\Lambda })}𝒱,_𝒱(\mathrm{}\pi ^{}D))`$ over $`𝒪_𝒱`$. We can restate the above discussion in the following proposition.
###### Proposition 5.1
The super Faà di Bruno recurrence relation is the cocycle condition for the hypercohomology group describing the deformations of the dynamical super line bundle $``$ on the spectral curve $`𝒞`$ and of its meromorphic sections which give rise to the super Krichever map. $`\mathrm{}`$
We end by remarking that equation (2.3) is an obvious consequence of the definition of $`\widehat{h}`$ and $`\widehat{H}^{(k)}`$.
### Acknowledgments
We thank D. Hernández Ruipérez and U. Bruzzo for useful discussions.
## Appendix A Super curves
In this Appendix we recall some facts concerning super curves, referring to for more details on supergeometry.
Let $`\mathrm{\Lambda }`$ be a Grassmann algebra over $``$. An algebraic super curve over $`\mathrm{\Lambda }`$, also called a $`\mathrm{\Lambda }`$–super–curve for brevity, is a proper irreducible superscheme $`𝒞\mathrm{Spec}(\mathrm{\Lambda })`$ over $`\mathrm{Spec}(\mathrm{\Lambda })`$ with fibre dimension $`1|1`$ and whose underlying reduced scheme is a proper irreducible algebraic curve over $``$. Throughout this paper we assume $`𝒞`$ to be a supermanifold, so it is given by a pair $`(C,𝒪_𝒞)`$ where $`C`$ is a topological space and $`𝒪_𝒞=𝒪_{𝒞,0}𝒪_{𝒞,1}`$ is a sheaf of super–commutative $`\mathrm{\Lambda }`$–algebras on $`C`$ such that
* $`(C,𝒪_C:=𝒪_𝒞^{\mathrm{rd}}=𝒪_𝒞/𝒥_𝒞)`$ is a smooth irreducible proper algebraic curve over $``$, where $`𝒥_𝒞`$ is the ideal sheaf $`𝒪_{𝒞,1}+𝒪_{𝒞,1}^2`$,
* there exists an open covering $`\{U_j\}_{jJ}`$ of $`C`$ and odd elements $`\theta _j𝒪_𝒞(U_j)`$ such that
$$𝒪_𝒞(U_j)𝒪_C(U_j)_{}\mathrm{\Lambda }[\theta _j].$$
A $`\mathrm{\Lambda }`$–point of $`𝒞`$ is a map $`\mathrm{Spec}(\mathrm{\Lambda })𝒞`$ whose composition with the projection $`𝒞\mathrm{Spec}(\mathrm{\Lambda })`$ is the identity morphism.
An invertible sheaf $``$ on $`𝒞`$ is a locally free evenly generated $`𝒪_𝒞`$–module of rank $`1|0`$; it is the sheaf of sections of a super line bundle that, abusing notations, we still call $``$. We can find a suitable open covering $`\{U_j\}_{iJ}`$ of $`𝒞`$ over the elements of which $``$ is trivial. Then the super line bundle is completely described in terms of its (even invertible) transition functions $`g_{jk}\mathrm{\Gamma }(U_jU_k,𝒪_{𝒞,0}^\times )`$ satisfying the usual cocycle conditions. The set of isomorphism classes of super line bundles on $`𝒞`$ is therefore $`H^1(𝒞,𝒪_{𝒞,0}^\times )`$ and tensor product of invertible sheaves (or, equivalently, multiplication in $`𝒪_{𝒞,0}^\times `$) gives it a group structure under which it is called the Picard group $`\mathrm{Pic}(𝒞)`$ of $`𝒞`$.
Another way to describe an invertible sheaf is by means of super (Cartier) divisors. A super divisor on $`𝒞`$ is a collection $`D:=\{(U_j,f_j)\}_{jJ}`$ of even non–zero rational functions $`f_j`$ defined, up to even invertible regular functions, on the open subsets $`U_j`$ of a covering of $`𝒞`$, and agreeing in the intersections $`U_jU_k`$ up to an element of $`𝒪_{𝒞,0}^\times (U_jU_k)`$: i.e., $`D`$ is a section of $`Rat_{𝒞,0}^\times /𝒪_{𝒞,0}^\times `$, where $`Rat_𝒞`$ is the sheaf of rational functions on $`𝒞`$. With the super divisor $`D`$ one associates the invertible subsheaf $`𝒪_𝒞(D)Rat_𝒞`$ whose local sections over $`U_j`$ span the module $`f_j^1𝒪_𝒞(U_j)`$ and the transition functions of the corresponding super line bundle are $`g_{jk}=f_jf_k^1`$. We have the exact sequence
$$0𝒪_{𝒞,0}^\times Rat_{𝒞,0}^\times Rat_{𝒞,0}^\times /𝒪_{𝒞,0}^\times 0$$
and a super divisor $`D`$ is called principal if it is the image of a global non–zero even rational function $`f`$, in which case we write $`D=(f)`$. Of course, the invertible sheaf associated with a principal divisor is trivial and vice versa. $`D`$ is called effective (or positive) if $`f_j`$ is regular for every $`j`$, and irreducible if $`f_j=z_j\stackrel{~}{z}_j\theta _j\stackrel{~}{\theta }_j`$, where $`\stackrel{~}{z}_j,\stackrel{~}{\theta }_j\mathrm{\Lambda }`$.
A useful concept associated with irreducible super divisors is the dual super curve $`\stackrel{~}{𝒞}`$ of $`𝒞`$, which we briefly review (see for more details). Let $`\underset{¯}{𝒞}=(C,𝒪_{\underset{¯}{𝒞}})`$ be the $`N=2`$ super curve whose reduced algebraic curve is again $`C`$ and whose structure sheaf is the only super conformal extension of $`Ber_𝒞`$ by $`𝒪_𝒞`$
$$0𝒪_𝒞𝒪_{\underset{¯}{𝒞}}Ber_𝒞0.$$
(A.1)
Here $`Ber_𝒞`$ is the dualizing sheaf of $`𝒞`$, whose transition functions $`g_{jk}`$ are the Berezinians of the (super) Jacobian matrices of the coordinate transformations between $`U_j`$ and $`U_k`$, and the super conformal property means that the local form $`\omega _j:=\text{d}z_j\text{d}\theta _j\rho _j`$ is globally defined up to a scalar factor, where $`(z_j,\theta _j,\rho _j)`$ are graded local coordinates on $`\underset{¯}{𝒞}`$ adapted to $`𝒞`$ (i.e. $`(z_j,\theta _j)`$ are coordinates on $`𝒞`$). The kernel of $`\omega _j`$ is generated by $`𝒟_j:=_{\rho _j}`$ and $`\stackrel{~}{𝒟}_j:=_{\theta _j}+\rho _j_{z_j}`$ and one can easily convince himself that $`𝒟_j`$ represents locally the map $`𝒪_{\underset{¯}{𝒞}}Ber_𝒞`$, thus the structure sheaf $`𝒪_𝒞`$ of $`𝒞`$ is the kernel of $`𝒟`$.
Introducing the new coordinates
$$\stackrel{~}{z}_j:=z_j\theta _j\rho _j,\stackrel{~}{\theta }_j:=\theta _j,\stackrel{~}{\rho }_j:=\rho _j$$
(A.2)
on $`\underset{¯}{𝒞}`$, the two operators above become $`𝒟_j=_{\stackrel{~}{\rho }_j}+\stackrel{~}{\theta }_j_{\stackrel{~}{z}_j}`$ and $`\stackrel{~}{𝒟}_j=_{\stackrel{~}{\theta }_j}`$ respectively, so the kernel of $`\stackrel{~}{𝒟}_j`$ consists of functions of $`\stackrel{~}{z}_j`$ and $`\stackrel{~}{\rho }_j`$. One shows that this makes sense globally obtaining therefore a new exact sequence
$$0𝒪_{\stackrel{~}{𝒞}}𝒪_{\underset{¯}{𝒞}}\stackrel{\stackrel{~}{𝒟}}{}𝒬0,$$
where $`𝒪_{\stackrel{~}{𝒞}}`$ is the structure sheaf of a $`1|1`$ $`\mathrm{\Lambda }`$–super–curve $`\stackrel{~}{𝒞}`$ which is called the dual super curve of $`𝒞`$, moreover $`𝒬Ber_{\stackrel{~}{𝒞}}`$ and $`\stackrel{~}{\stackrel{~}{𝒞}}𝒞`$, which explains the terminology. The interesting fact is that the $`\mathrm{\Lambda }`$–points of $`\stackrel{~}{𝒞}`$ correspond to the irreducible divisors of $`𝒞`$.
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# 1 Introduction
## 1 Introduction
Although both the Casimir effect of quantum theory and the existence of symmetry-breaking condensates in both the strong and the electroweak sectors of the Standard Model indicate that empty space is a dynamical medium that ‘ought to’ have a large mass density, gravity, which couples universally to mass, does not reveal it. This is the problem of the cosmological term . Various mechanisms have been proposed to address this problem, but so far none has won wide acceptance. This situation is especially challenging for string theory, and its conjectured non-perturbative definition M theory, in as much as string theory is proposed as a fully specified dynamical theory of gravity.
An interesting approach to the solution of the cosmological term problem is the proposal that it is relaxed by jumps (saltation) associated with some rather exotic dynamics. There is an important conceptual advantage to having the relaxation connected to some — necessarily exotic — dynamics that responds only to a source taking the form of an effective cosmological term. For if the dynamics responds to several influences, it is difficult to see how a particularly simple value for one partial determinant of its behavior can become overwhelmingly preferred. A model for this logic is the advantage one gains from Peccei-Quinn symmetry in relieving the $`\theta `$-problem of QCD. Indeed, that symmetry produces an exotic dynamics (the axion field) which responds only to a source having the form of an effective $`\theta `$ term. Note that this logic also applies, in connection with continuous relaxation of scalar fields, to self-interactions (including kinetic terms), as has been quantified by Weinberg .
Saltation in the effective cosmological term has been considered in the context of stepwise false vacuum decay in a quasi-periodic staircase potential by Abbott , or alternatively by Brown and Teitelboim (BT) through nucleation of fundamental membrane degrees of freedom . Since the membrane formulation is close in spirit to the ones that appear to arise naturally in string theory, we shall phrase our discussion in its framework.
The essential ingredients in the BT model are a fundamental membrane degree of freedom (in modern language, a 2-brane) and a 4-form gauge field strength $`F_4`$ (deriving from a 3-form potential). The world-volumes of the membranes couple to the 3-form potential. The field strength $`F_4`$ has no other couplings and no local dynamics in the four-dimensional spacetime, but its expectation value contributes to the effective cosmological term. The presence of an expectation value for $`F_4`$ induces the nucleation of the 2-brane to which it couples, in a manner analogous to the Schwinger mechanism for electric field decay through nucleation of $`e^+e^{}`$ pairs. When a membrane is nucleated, say as a spherical shell, the effective value of the cosmological term on the inside differs from its previous value (the value on the outside) by an amount proportional to the coupling constant of the membrane. If a membrane of the correct sign is nucleated, the contribution of the 4-form to the effective value of the cosmological term will be reduced. For a suitable value of $`F_4`$, this contribution can cancel that arising from the combined effects of all other degrees of freedom in the theory. If the steps between adjacent false vacua are sufficiently small, and if there is a reason why vacuum decay stops or slows down dramatically as $`\mathrm{\Lambda }_{\mathrm{eff}}0`$, this mechanism could in principle relax a large microscopic cosmological term (arising from all sources other than the membrane-$`F`$ field dynamics) to a value within observational bounds.
In their pioneering work, Abbott, and Brown and Teitelboim, postulated dynamical entities ad hoc, with no richer context to enhance their credibility or connect them to other problems and facts of physics. On the other hand, developments in string theory in the past few years have emphasized the importance of extended objects — branes — of various types. It therefore seems appropriate to consider whether these branes might allow improved mechanisms for saltatory relaxation of the cosmological term. We believe that this is the case, for several reasons:
* A wide variety of membranes play a fundamental role in string theory. There are corresponding gauge fields, which naturally couple (only) to these extended objects.
* Since the theory is naturally formulated in higher dimensions, the couplings of these membranes as seen in four dimensions are not fixed and quantized, but rather are determined in terms of the fundamental (fixed, quantized) couplings together with properties of the extra-dimensional compactification.
* Similarly, the effective tension as seen in four dimensions depends on the properties of the extra-dimensional compactification.
* There are significant and, in suitable cases, exponentially large, density of states factors associated with semi-classical brane processes.
* Small tension, which may be favored for dynamical reasons, and large density of states factors make possible rapid relaxation of the cosmological constant.
We should stress that regardless of whether membrane nucleation is the final solution, or perhaps an ingredient, in solving the riddle of the cosmological constant, nucleation processes involving extended objects generically occur in string theory. Studying these processes is likely to give some insight into the question of vacuum selection and into the question of a background independent formulation of string theory.
In the following section, we review the BT formalism for relaxation of the cosmological constant through brane nucleation and the primary obstacles encountered. In the next two sections, we discuss features of string theory that may alleviate these difficulties: in Sec. 3 we describe the possibility of small charge densities and tensions arising from compactification, and in Sec. 4 we note the relevance of exponentially large density of states factors. Motivated by these features, we describe two possible scenarios for the cosmological constant in Sec. 5. In Sec. 6, we discuss a number of outstanding issues and summarize.
## 2 Relaxation Dynamics
### 2.1 Basic mechanism
We now recall the basic dynamics of the BT mechanism , which we have been able to express in a somewhat simplified fashion. For other recent discussions of membrane nucleation, see .
Consider gravity in $`D=4`$ spacetime dimensions with a 2-brane $`X^\alpha `$ coupled to a 3-form gauge potential $`A_3`$. The Minkowski action is
$`S_M`$ $`=`$ $`\tau _2{\displaystyle d^3\xi \sqrt{detg_{ab}}}+{\displaystyle \frac{\rho _2}{6}}{\displaystyle d^3\xi A_{\alpha \beta \gamma }\frac{X^\alpha }{\xi ^a}\frac{X^\beta }{\xi ^b}\frac{X^\gamma }{\xi ^c}\epsilon ^{abc}}`$ (1)
$`{\displaystyle \frac{1}{48}}{\displaystyle d^4x\sqrt{g}F_{\alpha \beta \gamma \delta }F^{\alpha \beta \gamma \delta }}+{\displaystyle \frac{1}{6}}{\displaystyle d^4x_\alpha \left[\sqrt{g}F^{\alpha \beta \gamma \delta }A_{\beta \gamma \delta }\right]}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle d^4x\sqrt{g}M^2(R2\lambda )}M^2{\displaystyle d^3x\sqrt{h}K},`$
where the $`\xi ^a`$ parameterize the membrane world-volume, and $`g_{ab}=_aX^\alpha _bX_\alpha `$ is the induced world-volume metric. The surface integral is over spacetime boundaries with $`h`$ and $`K`$ the induced metric and extrinsic curvature, respectively. This term and the total derivative integral ensure that the action has well-defined functional derivatives with respect to the metric and gauge field. An important point is that in four dimensions the 4-form field strength contains no independent propagating degrees of freedom, its value, up to a constant, being fully determined by the background of sources charged with respect to $`A_{\alpha \beta \gamma }`$.
The parameters entering this action, and their mass dimensions in $`D=4`$, are
$$\begin{array}{ccc}\text{2-brane tension}\hfill & \tau _2& 3\hfill \\ \text{2-brane charge density}\hfill & \rho _2& 2\hfill \\ \text{(Reduced) Planck mass}\hfill & M& 1\hfill \\ \text{Bare cosmological constant}\hfill & \lambda & 2.\hfill \end{array}$$
(2)
Numerically, $`M=(8\pi G)^{1/2}=2.4\times 10^{18}\mathrm{GeV}`$. Here, in agreement with BT, we use the canonical — positive energy — sign for the $`F^2`$ term. In Sec. 3, we shall see that this is appropriate for the branes relevant to us.
Rotating to Euclidean space, we find
$`S_E`$ $`=`$ $`\tau _2{\displaystyle d^3\xi \sqrt{detg_{ab}}}+{\displaystyle \frac{\rho _2}{6}}{\displaystyle d^3\xi A_{\alpha \beta \gamma }\frac{X^\alpha }{\xi ^a}\frac{X^\beta }{\xi ^b}\frac{X^\gamma }{\xi ^c}\epsilon ^{abc}}`$ (3)
$`{\displaystyle \frac{1}{48}}{\displaystyle d^4x\sqrt{g}F_{\alpha \beta \gamma \delta }F^{\alpha \beta \gamma \delta }}+{\displaystyle \frac{1}{6}}{\displaystyle d^4x_\alpha \left[\sqrt{g}F^{\alpha \beta \gamma \delta }A_{\beta \gamma \delta }\right]}`$
$`+{\displaystyle d^4x\sqrt{g}\frac{1}{2}M^2(R+2\lambda )}+M^2{\displaystyle d^3x\sqrt{h}K}.`$
The sign of the $`F^2`$ term in Eq. (3) depends on the Euclideanization procedure. Here, following Ref. , we make the conventional rotations $`x^0ix^0`$, $`X^0iX^0`$ for time-like quantities, but take $`A^{0\mu _2\mathrm{}\mu _{D1}}A^{0\mu _2\mathrm{}\mu _{D1}}`$ and $`A^{\mu _1\mathrm{}\mu _{D1}}iA^{\mu _1\mathrm{}\mu _{D1}}`$ for the space-like components so that the field strength $`F`$ is invariant. Alternatively, one may adopt the prescription $`A^{0\mu _2\mathrm{}\mu _{D1}}iA^{0\mu _2\mathrm{}\mu _{D1}}`$ and keep the space-like components invariant. In this case, the sign of the $`F^2`$ term in Eq. (3) changes. However, the field strength is not invariant under this prescription; taking $`F_4`$ in Eq. (3) to be pure imaginary leaves the following analysis unchanged.
The instanton solution is a membrane that divides space into two regions, an outside $`O`$ and an inside $`I`$. In each region, the field strength is a constant
$$F_{O,I}^{\alpha \beta \gamma \delta }=\frac{c_{O,I}}{\sqrt{g}}\epsilon ^{\alpha \beta \gamma \delta },$$
(4)
and the field strengths are matched across the membrane boundary via
$$c_I=c_O\rho _2.$$
(5)
The effective cosmological terms are
$$\mathrm{\Lambda }_{O,I}=\lambda +\frac{1}{2M^2}c_{O,I}^2,$$
(6)
where the field strength contribution follows from Einstein’s equations. Alternatively, it may be ‘read off’ from the action if one is careful to include the on-shell contribution from the total derivative term, which is double in magnitude and opposite in sign relative to the usual $`F^2`$ term. From Eq. (6), it is clear that if the bare cosmological term is to be canceled, it must be negative, and we therefore assume $`\lambda <0`$.
The tunneling probability is $`Pe^B`$, where the bounce action for this false vacuum decay is
$$B=\{\begin{array}{cc}\mathrm{},\hfill & \mathrm{if}\mathrm{\Lambda }_\mathrm{O},\alpha _\mathrm{O}<0\hfill \\ 12\pi ^2M^2\left[\frac{1}{\mathrm{\Lambda }_O}(1b\alpha _O)\frac{1}{\mathrm{\Lambda }_I}(1b\alpha _I)\right],\hfill & \mathrm{otherwise}.\hfill \end{array}$$
(7)
Here the bubble radius, defined so that the area of the bubble slice when continued back to Minkowski signature is $`4\pi b^2`$, is
$$b=\frac{1}{\sqrt{\frac{1}{3}\mathrm{\Lambda }_O+\alpha _O^2}}=\frac{1}{\sqrt{\frac{1}{3}\mathrm{\Lambda }_I+\alpha _I^2}},$$
(8)
and
$`\alpha _{\genfrac{}{}{0pt}{}{O}{I}}`$ $`=`$ $`{\displaystyle \frac{1}{3xM}}\left[c_O\left({\displaystyle \frac{1}{2}}\pm {\displaystyle \frac{3}{4}}x^2\right)\rho _2\right]`$ (9)
$`\mathrm{\Lambda }_I`$ $`=`$ $`\mathrm{\Lambda }_O+{\displaystyle \frac{1}{2M^2}}(\rho _2^22\rho _2c_O)`$ (10)
$`x`$ $`=`$ $`{\displaystyle \frac{\tau _2}{\rho _2M}}.`$ (11)
In the following sections, we will often work in Planck units with $`M1`$.
### 2.2 Naturalness
The value of the cosmological term at present is
$$\mathrm{\Lambda }_{\mathrm{obs}}M^2\stackrel{<}{}(2\times 10^3\mathrm{eV})^410^{120}M^4.$$
(12)
Present data prefer a non-zero, positive value; for recent reviews, see . One virtue of saltatory relaxation is that, in contrast to mechanisms involving symmetries or continuous relaxation, a small non-zero value may emerge naturally as a consequence of a non-vanishing jump size.
For the observed value to be natural in the framework of brane nucleation, the spacing between allowed values of the effective cosmological term, near the observed value, cannot be much larger. This translates into the condition
$$\rho _2\stackrel{<}{}\frac{\mathrm{\Lambda }_{\mathrm{obs}}}{|\lambda |^{1/2}}.$$
(13)
This is an extremely stringent condition on the microphysics, even for plausible $`|\lambda |1`$, since the observed cosmological constant is so small. For example, even if the bare cosmological constant is generated only at the TeV scale through low-energy supersymmetry breaking so that $`|\lambda |10^{60}`$, one still requires $`\rho _2\stackrel{<}{}10^{90}`$, which translates into an associated mass scale of $`10^{18}\mathrm{eV}`$. The requirement of such a small coupling is, at best, unsettling, and one might hope for an explanation in some more fundamental framework. This is especially true in string theory, where, in the absence of such an explanation, one naively expects couplings to be of order the Planck scale.
As an aside, note that throughout this study, we consider evolution of the cosmological term down the staircase of values allowed by brane nucleation with fixed charge density $`\rho _2`$, and thus with essentially fixed step size. However, very small fractional changes in $`\rho _2`$, of order
$$\frac{\delta \rho _2}{\rho _2}\stackrel{<}{}\frac{\mathrm{\Lambda }_{\mathrm{obs}}}{|\lambda |},$$
(14)
assuming $`\rho _2\sqrt{|\lambda |}`$, are sufficient to bring the observed value of the effective cosmological term into range by distorting the size of the steps near $`\mathrm{\Lambda }=0`$. Another possibility is that the entire staircase moves up or down by some suitably small amount. Although we will not attempt to exploit these features here, they may play an important role in future work, since, as we shall emphasize below, $`\rho _2`$ is in principle a dynamical variable. It depends, in particular, on the expectation values of the string theory compactification moduli, including the dilaton.
### 2.3 Absolute Stability
The cosmological term must not only relax to within its observational bounds, but it must also stop evolving once it reaches this interval. For de Sitter space, additional bubble nucleations are always possible. However, for $`\mathrm{\Lambda }_O<0`$, transitions can take place only when $`\alpha _O>0`$. It is not hard to show that this constraint, along with the condition $`\frac{1}{3}\mathrm{\Lambda }_O+\alpha _O^2>0`$, implies that further lowering of the effective cosmological term will not occur beyond the first anti-de Sitter step if the tension is sufficiently large. This remarkable result is closely related to the Coleman-de Luccia gravitational suppression of false vacuum decay. BT hypothesized that our universe is at the endpoint of such an evolution.
A sufficient condition to ensure absolute vacuum stability is $`\tau _2^2>\frac{4}{3}\rho _2c_O`$ or, in terms of the tension to charge density ratio,
$$x=\frac{\tau _2}{\rho _2}\stackrel{>}{}\sqrt{\frac{|\lambda |}{\mathrm{\Lambda }_{\mathrm{obs}}}}.$$
(15)
Thus, even for the smallest plausible $`|\lambda |`$, a large hierarchy between tension and charge density is required.
More problematic still, BT showed that, upon combining the stability and naturalness conditions, the time required to reach the endpoint is excessively large, so large that even the very slow inflation that occurs in the penultimate vacuum would leave the universe entirely devoid of matter and energy.
## 3 Saltation in String Theory: Tension and Charge Density
### 3.1 Framework
In its long wavelength approximation, M theory supports BPS M2-branes and M5-branes . The M2-branes couple electrically to the 3-index gauge potential of $`11`$-dimensional supergravity, which we denote by $`A_3`$. We define a dual gauge potential $`A_6`$ in the standard way:
$`F=dA_3=F_7=dA_6.`$ (16)
The M5-branes couple magnetically to $`A_3`$, or directly to $`A_6`$. In terms of the 11-dimensional Planck scale $`l_p`$, M2-branes have a tension $`T_2(l_p)^3`$ while M5-branes have a tension $`T_5(l_p)^6`$.
The simplest way to arrive at a $`4`$-dimensional world is via compactification on a $`7`$-dimensional internal space $``$. We obtain 2-branes from either the fundamental M2-branes of M theory, or by wrapping M5-branes on a 3-cycle $`a_3`$ of the internal space $``$. Our spectrum of 4-forms in spacetime comes about in the following way: let $`\omega ^i`$ be a basis for $`H^3(,Z)`$. We require integer forms so that the resulting 4-form fields satisfy Dirac quantization. We can then expand $`F_7`$ in this basis,
$$F_7=\stackrel{~}{F}_i\omega ^i.$$
(17)
Note that we get many 4-forms from a generic compactification of this kind. There is a natural geometric picture associated to this expansion. To each form $`\omega ^i`$, we can associate a dual $`3`$-cycle on which we wrap an M5-brane to obtain a 2-brane.
After setting the fermions to zero, the 11-dimensional metric satisfies the equation of motion
$$R_{ab}\frac{1}{2}g_{ab}R+\frac{1}{6}\left(F_{acde}F_b^{cde}\frac{1}{8}g_{ab}F_{cdef}F^{cdef}\right)=0.$$
(18)
The sign of the four-form contribution is fixed by 11-dimensional supergravity. For both kinds of effective 2-brane described above, the contribution of $`F`$ (or $`\stackrel{~}{F}`$) to the effective cosmological constant is positive. We therefore require a negative bare cosmological constant $`\lambda `$, as assumed in section 2.
If $``$ is circle-fibered then we can reduce M theory on the fiber to obtain various string compactifications. For example, if we take $``$ to have the form $`S^1\times K`$ then for a small $`S^1`$ of size $`R_{11}`$, M theory reduces to type IIA string theory on $`K`$ . The string scale $`\alpha ^{}`$ and string coupling $`g_s`$ are given in terms of $`l_p`$ and $`R_{11}`$,
$$\alpha ^{}=l_p^3/R_{11},g_s=(R_{11}/l_p)^{3/2}.$$
(19)
Type IIA string theory contains an NS-NS (Neveu-Schwarz) 3-form field strength $`H_3`$ which measures fundamental string charge. The magnetic dual $`H_7=H_3`$ measures the charge of NS 5-branes. These 5-branes have a tension proportional $`1/g_s^2`$ and so are very heavy at weak string coupling. In addition, there are Ramond-Ramond (RR) field strengths $`F_2`$ and $`F_4`$ (and a non-dynamical $`F_{10}`$ ) together with their 8-form and 6-form (and 0-form) magnetic duals. These $`(p+2)`$-form RR field strengths couple electrically to dynamical Dirichlet $`p`$-branes through the world-volume action
$$d^{p+1}\xi A_{p+1}.$$
(20)
In type IIA, we therefore have the additional possibility of wrapped D6-branes and D8-branes giving rise to effective 2-branes. In (9+1) dimensions, the low-energy string and D$`p`$-brane effective action is
$`S_{10}`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _{10}^2}}{\displaystyle d^{10}x\sqrt{G}\left[e^{2\varphi }\left(R+4(\varphi )^2\right)\frac{1}{2n!}F_n^2+\mathrm{}\right]}`$ (21)
$`T_p{\displaystyle d^{p+1}\xi e^\varphi \sqrt{detG_{ab}}}+\rho _p{\displaystyle d^{p+1}\xi A_{p+1}}+\mathrm{},`$
where $`\varphi `$ is the dilaton, $`G_{ab}`$ is the induced metric on the brane, and only the relevant bosonic terms are displayed. The tension and charge density of type II D$`p`$-branes is
$$T_p^2=\rho _p^2=\frac{\pi }{\kappa _{10}^2}(4\pi ^2\alpha ^{})^{3p}.$$
(22)
This may be written more conveniently by using the relation $`\kappa _{10}^2=2^6\pi ^7\alpha ^4`$ and defining the string length $`\mathrm{}_s`$ in terms of the fundamental string tension through $`T_{\mathrm{F1}}=1/(2\pi \alpha ^{})=2\pi /\mathrm{}_s^2`$. With these conventions,
$$T_p=\rho _p=\frac{2\pi }{\mathrm{}_s^{p+1}}.$$
(23)
(Note that there is no factor of $`1/g_s`$ in these expressions because of the conventions employed in Eq. (21).)
The 4-dimensional physical brane tensions and charge densities are proportional to $`T_p`$ and $`\rho _p`$. As we discuss in the next subsections, the exact relations depend on the complicated dynamics of compactification which are not well understood. For this reason, when we consider scenarios in Sec. 5, we will take a 4-dimensional effective action approach, and simply assume suitable values for the 2-brane tension and charge density. In the rest of this section, however, we explore the possible effects of compactification on these quantities.
There are also other ways of arriving at a 4-dimensional world that do not fit into the above framework. For example, compactifying F theory on a Calabi-Yau 4-fold naturally gives a large class of N=1 compactifications. The 4-fold must be elliptically-fibered with a section. If $`B`$ denotes the 6-dimensional base of the fibration then over each point of $`B`$, the structure of the elliptic fibration specifies a torus. By F theory compactified on the 4-fold, we mean type IIB string theory compactified on $`B`$. However, the complexified string coupling constant $`\tau `$ is not constant but varies over $`B`$ in a manner determined by the shape of the torus fiber. Therefore, the string coupling constant is typically not a tunable modulus for these compactifications.
To determine, the spectrum of 2-branes for an F theory compactification, it is natural to use the duality between M theory on $`T^2`$ and the type IIB string. From this duality, we see that 2-branes arise only from M5-branes. An M5-brane wrapping the elliptic fiber and a 1-cycle in $`B`$ would look like a D3-brane wrapping the 1-cycle in $`B`$. There are additional possibilities. For example, an M5-brane with a leg on the elliptic fiber and 3 legs on $`B`$ can give rise to a 2-brane. From the IIB perspective, it would appear to be a combination of $`(p,q)`$ 5-branes wrapping a 3-cycle of $`B`$. In backgrounds such as these with N=1 supersymmetry, we also need to worry about the spacetime superpotential. However, our subsequent discussion will be largely classical. Our aim is to be as simple as possible in order to isolate the aspects of compactifications that are most relevant for the brane nucleation mechanism.
We also note that there is a much broader class of compactifications that involve background fluxes. These compactifications, which are actually generic string compactifications, typically warp spacetime . These backgrounds can potentially give rise to the kind of novel infra-red physics that could make saltatory relaxation a viable mechanism for reducing the cosmological constant.
### 3.2 Direct descent
One might suppose that M2-branes are made to order for saltatory relaxation of the cosmological term. However, as we shall see, they do not have the right properties, at least when taken in their straightforward form. The wrapped M5-branes appear more promising.
Let us first consider the case of a D2-brane in 10 dimensions that descends directly to a 2-brane in 4 dimensions. Our discussion will be in the context of 10-dimensional type IIA supergravity compactified on a Calabi-Yau manifold $`K`$ with string-frame volume $`V_6`$. This compactification preserves N=2 supersymmetry. The physical effective tensions and charge densities are then determined in 4-dimensional Einstein frame, where the gravitational action takes the conventional Einstein-Hilbert form. This frame follows after performing the Weyl rescaling $`g_{\mu \nu }e^{2\varphi }g_{\mu \nu }`$ on the 4-dimensional metric $`g_{\mu \nu }`$. The 4-dimensional action is then
$`S_4`$ $`=`$ $`{\displaystyle \frac{V_6}{2\kappa _{10}^2}}{\displaystyle d^4x\sqrt{g}\left[R2(\varphi )^2\frac{1}{2n!}e^{2(n2)\varphi }F_n^2+\mathrm{}\right]}`$ (24)
$`T_p{\displaystyle d^{p+1}\xi e^{p\varphi }\sqrt{detg_{ab}}}+\rho _p{\displaystyle d^{p+1}\xi A_{p+1}}+\mathrm{}.`$
The 4-dimensional effective tension is $`\tau _p|_{4\mathrm{D},\mathrm{eff}}=T_pg_s^p`$, where $`g_s=e^\varphi `$ is the string coupling. To determine the effective charge density, note that the $`F_n^2`$ kinetic terms are not canonically normalized. Upon restoring the normalization, we find a 4-dimensional effective charge density of $`\rho _p|_{4\mathrm{D},\mathrm{eff}}=\sqrt{2}\rho _pg_s^p/M`$, where the reduced Planck mass is fixed by $`M^2=V_6/\kappa _{10}^2`$. The tension to charge density ratio, in Planck units, is then $`x=1/\sqrt{2}`$, as we expect for BPS branes.
For the 2-brane case of interest,
$$\rho _2|_{4\mathrm{D},\mathrm{eff}}=\frac{2\sqrt{2}\pi g_s^2}{M\mathrm{}_s^3}.$$
(25)
Clearly, to obtain a sufficiently small charge density, we must have extreme values for $`\mathrm{}_s`$ and/or $`g_s`$. For example, for the canonical choice of string scale $`\mathrm{}_s(10^{17}\mathrm{GeV})^1`$, we find that a charge density of $`\rho _2\stackrel{<}{}10^{90}`$ requires $`g_s\stackrel{<}{}10^{44}`$. Alternatively we could take $`\mathrm{}_s`$ to be a larger length scale. These non-canonical cases include the ‘large extra dimension’ scenario with $`g_s1`$ and some number of sub-millimeter dimensions. However, given the success of quantum field theory at colliders such as LEP and the Tevatron, $`\mathrm{}_s\stackrel{<}{}(\mathrm{TeV})^1`$ at the very best. Although an improvement, this still requires a tiny string coupling $`g_s\stackrel{<}{}10^{23}`$ to generate a sufficiently small $`\rho _2`$. It is possible that this string coupling is unrelated to the gauge couplings of the Standard Model; for example, the Standard Model may arise from some non-perturbative sector of string theory. However, if $`g_s`$ is related to the Standard Model gauge couplings in a straightforward way, it is difficult to understand how to accommodate the Standard Model in such an extreme corner of string theory moduli space, much less gauge coupling unification with $`\alpha _{\mathrm{unif}}1/25`$ .
### 3.3 Degenerating cycles
A more promising alternative to direct descent branes are branes wrapped on homology cycles. These branes become tensionless when the volume of these cycles approaches zero classically, as for conifolds . For compactifications with $`N2`$ supersymmetry, these branes give rise to BPS states and so quantum corrections do not change this conclusion.
Specifically, if a $`p`$-brane of tension $`T_p`$ wraps a $`k`$-cycle $`a_k`$ of the compactification manifold, where $`kp`$ and the volume of $`a_k`$ is $`\mathrm{Vol}(a_k)`$, then the result in the effective 4-dimensional theory is a $`(pk)`$-brane of tension $`\tau _{(pk)}T_p\mathrm{Vol}(a_k)`$. In particular, for an effective 2-brane in 4 dimensions coming from the wrapping of a D$`p`$-brane on a cycle $`a_{p2}`$, the Einstein frame effective tension is
$$\tau _2|_{4\mathrm{D},\mathrm{eff}}=T_p\mathrm{Vol}(a_{p2})g_s^2.$$
(26)
If $`\mathrm{Vol}(a_{p2})`$ approaches zero, i.e., $`a_{p2}`$ is a degenerating cycle, a nearly tensionless object exists in the 4-dimensional theory. The analogous formula for the tension of a wrapped NS 5-brane just differs by a factor of $`g_s`$. This is consistent with the higher-dimensional quantization rules for brane properties.
What determines the charge density? As we saw in the case of the direct reduction of an M2-brane to a D2-brane, the kinetic terms for the $`3`$-form gauge-field become moduli-dependent. The term,
$$d^{p+1}\xi A_{p+1},$$
is reduced on integer classes and so contains no moduli dependence. For purposes of determining the scaling of the charge density with the moduli, we only need the behavior of the moduli space metric near the singularity. At the level of classical geometry, this is determined by the behavior of the Weil-Petersson metric. For N=2 compactifications, this metric behavior is not changed by quantum corrections, although it is certainly renormalized in situations with less supersymmetry. Let us take the case of the conifold where an $`S^3`$ shrinks to zero size in $`K`$. We choose coordinates for the moduli space so that the singularity is located at $`Z=0`$. The tension of a wrapped brane goes like,
$$\tau _2|_{4\mathrm{D},\mathrm{eff}}|Z|.$$
However, the moduli space metric scales like $`\mathrm{ln}(Z)`$ and so the charge density behaves like ,
$$\rho _2|_{4\mathrm{D},\mathrm{eff}}1/\sqrt{|\mathrm{ln}(Z)|}.$$
Thus although the tensions can be vanishingly small, the charge densities generated are only slightly smaller than in the direct descent case. We would therefore need to stabilize $`Z`$ at extremely small values in order to obtain an acceptably small charge density.<sup>2</sup><sup>2</sup>2We thank J. Polchinski for correcting an error in an earlier version of this paper, and J. Maldacena for discussions.
The log-singularity in the metric has a simple space-time interpretation as arising from integrating out a massless hypermultiplet . We desire a singularity at finite distance in the moduli space whose behavior is worse than that of a conifold. It is not hard to see that in the context of N=2 compactifications, such singularities cannot arise by simply tuning vector multiplet moduli. For example, one can try generalizations of the conifold with multiple simple nodes. These singularities are all at finite distance in the moduli space. However, in each case and in general, it appears that gravity can be decoupled. The resulting theory of vector multiplets has a moduli space metric with singularities than metrically cannot be worse than logarithmic by standard non-renormalization theorems.
Too little is currently known about the hypermultiplet moduli space to decide whether a singularity with the right metric behavior exists. Some results on classifying singularities at finite distance have appeared in . However, when we relax the condition of N=2 supersymmetry and consider the far larger class of N=1 compactifications, it seems far more likely that a sufficiently bad singularity exists (classically). This is quite exciting since it connects the possibility of novel infra-red physics with a mechanism for reducing the cosmological constant. A recent F theory example with unusual infra-red physics arising from a bad singularity appeared in . This kind of example certainly has the right qualitative features, and it actually seems quite hard to rule out the existence of a compactification with the features we desire. Clearly more work along these lines is needed.
We have yet to discuss supersymmetry breaking. One natural mechanism worth mentioning in this context breaks supersymmetry by turning on $`RR`$ and $`NS`$ fluxes on $`K`$. For example, we can give (3+1)-Lorentz invariant expectation values to the RR fields $`F_n`$ by taking,
$$F_2=v_2^i\alpha _i^{(2)},F_4=v_4^i\alpha _i^{(4)}+v_4\epsilon ^{(4)}.$$
The $`\alpha _i`$ are harmonic forms on $`K`$, and $`\epsilon ^{(4)}`$ is the spacetime volume element. If we take the 10-dimensional Poincaré duals of these expectation values $`v_{(4,2)}^i`$, we then find expectation values of the general form $`F=v_{(4,2)}^i\epsilon ^{(4)}\stackrel{~}{\omega }^{(2,4)}`$, which couple magnetically to the D$`p`$-branes. Let us denote the expectation values collectively by $`v_a`$.
Generically, when $`v_a0`$, supersymmetry is broken , and a scalar potential is generated that depends on the Calabi-Yau moduli $`t_\alpha `$ and the expectation values $`v_a`$. What is interesting to us, however, is that the potential naturally tends to drive us to singular compactifications. In some examples involving just RR fluxes, two types of critical point for the scalar potential have been found : either the Calabi-Yau runs away to infinite volume, or the theory is driven to conifold-like points where homology cycles of $`K`$ degenerate (classically approach zero volume). The minima tend to either break supersymmetry completely or restore the full N=2. Moreover, since the configurations with $`v_a0`$ are not iso-potential, there are dynamical processes whereby the values of the $`v_a`$’s change. The $`v_a`$ ‘discharge’ by the nucleation and expansion of charged membranes, the D-branes of string theory.
Classically, the vacuum expectation values $`v_a`$ of the RR fields can vary continuously. Quantum mechanically, they satisfy quantization rules that follow from the quantization of brane properties, which are standard results in string theory . These brane properties are fixed, given information about the compactification and the parameters of our low-energy effective Lagrangian.
There is an additional issue concerning this means of supersymmetry breaking that deserves mention. In the context of M theory on $`K\times S^1`$, a class of 2-branes arise from M5-branes wrapping 3-cycles. In type IIA, these 2-branes arise from either D4-branes wrapping 2-cycles of the compactification space $`K`$, or from NS 5-branes wrapping 3-cycles of $`K`$. The two cases are quite different. In the case of the wrapped D4-branes, we must also consider D2-branes which can also wrap the shrinking 2-cycle. These wrapped branes give rise to particle states with mass $`\tau _0\mathrm{Vol}(a_2)/\mathrm{}_s^30`$. With multiple collapsed nodes, we will find an interacting gauge theory in $`3+1`$-dimensions. In this case, it seems possible that the wrapped D4-branes will correspond to collective excitations of the gauge theory. This is not an issue, however, for the wrapped NS 5-branes since there are no BPS D3-branes in type IIA string theory. These wrapped NS-branes give rise to inherently non-perturbative stringy excitations.
The phenomenon of classically vanishing tensions arising from branes wrapped on degenerating cycles is intriguing. It is also worth noting that the classical phenomenon of true degeneration and corresponding tensionless, or massless, states is typically not realized in the full quantum theory. Instead the effective (3+1)-dimensional 2-brane will have a dynamically generated non-perturbative tension, which may be exponentially small. It is known that in some cases a tension is generated from the dynamics of the would-be massless particle states arising from a D2-brane wrapping the 2-cycle. These particle states realize a non-Abelian gauge theory, presumably in a sector hidden with respect to the Standard Model, whose low-energy non-perturbative dynamics can break supersymmetry. (This sector is conceptually distinct from hidden sectors postulated to provide supersymmetry breaking for the supersymmetric Standard Model.) This leads to a potential for the volume modulus of the cycle, which stabilizes it at a scale $`\mathrm{\Lambda }\mathrm{exp}(8\pi ^2/b_1g_{YM}^2)M`$ . Once the cycle is stabilized at this small scale, membranes wrapping this cycle have a tension that is also proportional to $`\mathrm{\Lambda }`$ and thus can be very small, even for $`g_{YM}1`$. The precise conditions under which very small tensions occur deserve much further study.
Some years ago, one of us made the numerical joke
$$\mathrm{\Lambda }_{\mathrm{obs}}M^2(e^{\pi /\alpha _{\mathrm{unif}}}M)^4,$$
(27)
with $`\alpha _{\mathrm{unif}}1/25`$. Our present considerations suggest an intimate connection between non-perturbative effects and the value of the observed cosmological term, which could conceivably lead to a relation of just this form.
## 4 Saltation in String Theory: Density of States
### 4.1 Multi-bounce properties
An essentially new feature is introduced by multi-bounce solutions arising from coincident branes. Such coincident branes support low-energy internal degrees of freedom. Thus there are potentially large density of states factors accompanying their nucleation. Calculations performed in the context of checking duality between type I and heterotic SO(32) string theory demonstrate that D-branes do make contributions that can be interpreted semi-classically as incorporating degeneracy factors reflecting the non-Abelian structure of coincident D-branes. Another aspect of this is that many coincident branes with large total charge can be described in appropriate limits as ‘black’ objects, similar to black holes, with event horizons, and with associated Bekenstein-Hawking (BH) entropy .
Consider first the case of $`k`$ coincident D3-branes. Such a configuration possesses a U($`k`$) super Yang-Mills (SYM) gauge theory on its world-volume. In the limit where the interactions with the bulk string theory are weak, and where the temperature (or excitation energy) of the SYM is small, one can compute the entropy of this system.
When the effective SYM gauge coupling $`g_{\mathrm{eff}}^2kg_{\mathrm{YM}}^2`$ is small, the entropy of the gas of massless gauge bosons and their superpartners at temperature $`T`$ is simply
$$S_3=\frac{2\pi ^2}{3}k^2VT^3,$$
(28)
where $`V`$ is the spatial volume of the 3-branes. What happens when the effective coupling is large? In this case one can use the type IIB supergravity solution describing the $`k`$ coincident 3-branes. These classical solutions with the asymptotic geometry and quantum numbers appropriate for $`k`$ coincident 3-branes contain a non-extremality parameter upon which their masses and horizon areas depend. If we associate the area of the horizon with BH entropy, we can derive a temperature by taking an appropriate derivative. By this procedure, the supergravity picture yields the strong coupling form of the entropy. In this case, the entropy agrees with the preceding weak coupling formula up to a numerical prefactor 3/4, which is then a prediction for the strong coupling behavior of the theory.
For $`k`$ coincident D2-branes the UV theory (in the decoupling limit) is again a weakly-coupled U($`k`$) SYM theory, so the UV entropy again scales as $`k^2`$. However we are interested in the IR entropy since, as we will argue in the next subsection, the physically motivated temperatures are the ambient de Sitter temperatures which are small (vanishingly small as $`\mathrm{\Lambda }_O\mathrm{\Lambda }_{\mathrm{obs}}`$). In the IR the SYM theory on the (2+1)-dimensional world-volume becomes strongly coupled, so one must switch over to the supergravity description. As shown in Ref. , the theory flows in the far IR to that of the M2 brane with BH entropy inferred from the horizon area given by
$$S_2k^{3/2}AT^2,$$
(29)
with $`A`$ the 2-brane area. This strongly suggests that such strongly-coupled brane configurations support $`𝒪(k^{3/2})`$ light degrees of freedom, though the physical nature of these degrees of freedom remains somewhat mysterious. Thus the probability for a semi-classical process involving $`k`$ coincident D2-branes is multiplied by a density of states factor of the form $`N_k\mathrm{exp}\left(k^{3/2}AT^2\right)`$ in the IR limit $`T0`$.
The branes of interest to us are effective 2-branes arising from the wrapping of, say, a D4-brane on a 2-cycle, $`a_2`$, or in the M theory picture an M5-brane wrapped on $`S^1\times a_2`$. Forgetting for a moment about the wrapping on a cycle, $`k`$ coincident D4-branes have a (4+1)-dimensional SYM theory which now flows to a free theory in the IR, so the entropy would scale as $`k^2V_4T^4`$ (the M5-branes have a rather unusual, and microscopically not fully understood, scaling $`Sk^3`$). In the case of the wrapped branes, the excitations in the wrapped directions are massive and not excited at the low temperatures we consider, so the entropy scales as $`SAT^2`$ with $`A`$ being the area of the 2-brane in our extended (3+1)-dimensional spacetime. Such behavior is correct until the temperature $`T`$ falls below the inverse linear size $`1/b`$ of the nucleated brane, below which point the density of states factor just counts the number of zero modes. From the preceding discussion, we expect the scaling of the entropy with $`k`$ in this case to lie between $`k^2`$ and $`k^{3/2}`$. A more exact treatment requires an additional analysis of the way in which the exponent scales as we approach the IR, as the physically motivated temperature is small but non-zero, and it is incorrect to scale infinitely far into the IR. In the case of a degenerating cycle, however, since the near degeneration implies that the world-volume theory of the effective 2-brane in (3+1) dimensions has a bare coupling that is large, we regard the lower value $`\beta =3/2`$ as being more likely.
In any case, if we accept this reasoning, the probability for a semi-classical process involving $`k`$ coincident 2-branes must be multiplied by an appropriate density of states factor of the form
$$N_k\mathrm{exp}\left(k^\beta AT^2\right),$$
(30)
for an appropriate temperature $`T`$, with an exponent $`\beta `$ that likely lies between $`\beta =2`$ and $`\beta =3/2`$. (Although we will not utilize it here, there might also be the possibility of $`k^3`$ scaling in the M5-brane limit.) As we will show below, a larger $`\beta `$ exponent implies more complete saltation, so we will adopt the more ‘conservative’ value of $`\beta =3/2`$ as our canonical choice.
### 4.2 Temperature: ambiguity, black hole analogy
The only temperatures intrinsic to our scenario are the de Sitter temperatures
$$T_{O,I}=\frac{H_{O,I}}{2\pi }=\frac{1}{2\pi }\sqrt{\frac{\mathrm{\Lambda }_{O,I}}{3}}.$$
(31)
Ambient ordinary matter might supply a much higher temperature (see Sec. 6), but the branes are in very poor thermal contact with it. (Ambient D-matter, i.e., the light particles mentioned above, might supply a better coupled temperature bath, but we shall not pursue this possibility.)
If the initial cosmological term is much larger than the change brought about by the $`k`$-bounce,
$$\mathrm{\Lambda }_Ok\rho _2c_O,$$
(32)
then the de Sitter temperatures before and after nucleation are almost identical, $`T_OT_I`$, and we may use either one in calculating the density of states factor.
In the case that a given transition produces large changes in the effective cosmological constant, an ambiguity arises. One possibility is that the temperature scale for tunneling from a highly curved (high temperature) de Sitter space to a less curved de Sitter space (or even to a flat or anti-de Sitter space) is substantially set by the high temperature. In this case one would take
$$TT_O$$
(33)
in the density of states factor of Eq. (30). We will consider the dynamics of this possibility in Sec. 5.1.
However, when the change in the nominal de Sitter temperature is comparable to the temperature itself, the thermal description of the tunneling process is internally inconsistent. A similar situation has been encountered before, in black hole physics . The problem arises in its most acute form for near-extremal holes, as the temperature approaches zero. If one uses the temperature of the initial hole, one finds a significant probability for radiating a quantum that will take the hole past extremality to a naked singularity. A more refined analysis shows that radiation is not thermal with regard to the initial temperature, and in particular that radiation beyond extremality is forbidden.
If we make an analogy between maximally homogeneous cosmologies and black holes based on their temperatures, then de Sitter spaces correspond to ordinary holes, flat space corresponds to an extremal hole, and anti-de Sitter spaces to naked singularities. This analogy suggests, in view of the previous paragraph, that we should not consider finite temperature branes that mediate transitions from de Sitter to anti-de Sitter spaces. A crude working hypothesis, which interpolates smoothly to this suggestion, is that in the density of states factor of Eq. (30), we should, instead of $`TT_O`$, employ the geometric mean of the de Sitter temperatures
$$T\sqrt{T_OT_I}.$$
(34)
The dynamics of this possibility is explored in Sec. 5.2.
### 4.3 A catastrophe for $`TT_O`$?
If the analogy to black hole results holds, tunneling to (and through) anti-de Sitter space is excluded. On the other hand, if the temperature for transitions is effectively set by $`T_O`$, multi-bounce transitions with arbitrarily large $`k`$ are possible and must be considered.
The action for $`k`$-bounce tunneling is given by Eqs. (7)–(11), with $`\rho _2k\rho _2`$ and $`\tau _2k\tau _2`$. The action is monotonically increasing as $`k`$ increases, and has the limiting behavior
$$B=\frac{24\pi ^2M^2}{\mathrm{\Lambda }_O},k\mathrm{}.$$
(35)
The limit $`k\mathrm{}`$ therefore appears problematic, since $`PN_ke^B`$ approaches a constant for large $`k`$. Note, however, that as $`k\mathrm{}`$, the bubble size $`b\frac{1}{k}0`$. Thus, the apparent ‘instability’ is toward creation of highly curved branes. We do not expect the action of Eq. (1) to be valid in the regime where the brane curvature far exceeds the brane tension.<sup>3</sup><sup>3</sup>3We thank R. Sundrum for emphasizing this point to us. Presumably, a calculation of the tunneling probability with the degrees of freedom appropriate for this regime will be well-behaved.
## 5 Scenarios
Let us now gather the pieces and attempt to envisage how — and whether! — they may be assembled into a complete scenario.
The cosmological constant evolves from some initial value through multi-bounce transitions. The probability for such transitions is
$$Pe^De^B.$$
(36)
The bounce action $`B`$ is given by Eqs. (7)–(11), with $`\rho _2k\rho _2`$ and $`\tau _2k\tau _2`$, where $`k`$ is the bounce number. The density of states prefactor is specified by $`D=k^\beta AT^2`$, where $`A=4\pi b^2`$ is the 2-brane area, and $`T`$ is the temperature. For concreteness, we will assume the low ‘conservative’ value of $`\beta =3/2`$, but the qualitative features of the following analysis hold more generally, for example, for $`\beta =2`$ or larger.
As noted above, the temperature $`T`$ is not under good theoretical control. We will therefore explore both of the broad alternatives mentioned previously.
### 5.1 $`TT_O`$, 1-step relaxation
We first consider the possibility that the temperature is given by the scale of the initial (outside) de Sitter temperature, so $`TT_O\sqrt{\mathrm{\Lambda }_O}`$. The tunneling probability from a given background configuration is then fixed in terms of the initial effective cosmological constant $`\mathrm{\Lambda }_O`$, the initial field strength $`c_O`$, and the 2-brane charge density and tension, parameterized by $`\rho _2`$ and $`x\tau _2/\rho _2`$.
We begin with some bare cosmological constant $`\lambda <0`$. Assume that the initial field strength gives a similar contribution to the effective cosmological constant, so $`\mathrm{\Lambda }_O,c_O^2|\lambda |`$. We assume also a very small charge density $`\rho _2\mathrm{\Lambda }_{\mathrm{obs}}/\sqrt{|\lambda |}`$, consistent with the naturalness condition discussed in Sec. 2.2, and $`x1`$.
With such initial conditions and brane properties, the maximal bounce action is $`B1/|\lambda |`$, while the degeneracy factor may be as large as $`D\lambda ^2/\mathrm{\Lambda }_{\mathrm{obs}}^2`$. Recall that $`\mathrm{\Lambda }_{\mathrm{obs}}10^{120}`$, while $`\lambda `$ is plausibly in the range of 1 (for Planck scale cosmological constants) to $`10^{60}`$ (for low-scale supersymmetry breaking at the TeV scale). Thus, the degeneracy enhancement overpowers the bounce action suppression, and tunneling proceeds rapidly.
It is not difficult to show that $`D`$ is maximized for $`k\rho _2c_O`$, i.e., for field strength step sizes of the right order to neutralize the field strength contribution to the effective cosmological constant. Indeed the most probable tunneling events nucleate bubbles of deep anti-de Sitter space. Such events produce small, short-lived interior universes, so the meaning of ‘probable’ in this context must be carefully qualified. Among universes that live a long time and even remotely resemble ours, the exponentially most favored are those closest to having zero effective cosmological term.
This scenario invokes a form of the anthropic principle. It is a uniquely weak one, however, in the following sense. Anthropic bounds on the cosmological term are highly asymmetrical . For positive cosmological terms, the formation of sufficiently large gravitational condensations requires cosmological terms below $`100`$ in units of $`\rho _c`$, the critical density. For negative cosmological terms, the lifetime of the universe requires cosmological terms roughly above $`1`$. Thus if the spacing between allowed near-zero saltatory values of the cosmological term in units of $`\rho _c`$ is, say, 3 and allows the values $`\mathrm{},2,1,4,\mathrm{}`$, then among values that can be experienced by sentient observers, 1 is by far the most likely.
Now suppose such a transition to nearly flat space has occurred. As discussed in Sec. 2.3, absolute stability is possible only for anti-de Sitter spaces, and only for extremely large $`x`$. However, absolute stability is not required on empirical grounds. We need only require that the effective value of the cosmological term is at present stable on cosmological time scales. For a starting effective cosmological constant of $`\mathrm{\Lambda }_O\mathrm{\Lambda }_{\mathrm{obs}}`$, the degeneracy factor is highly suppressed by small $`T`$, and the bounce action is dominant. The stability of the vacuum is then determined solely by $`B`$.
Since the bounce action can become truly infinite for slightly anti-de Sitter spaces, we anticipate, by continuity, that it can become very large even for slightly de Sitter spaces. The least suppressed transition (and so most dangerous from the point of view of vacuum instability) is that mediated by $`k=1`$. For small $`\rho _2`$ and small $`\mathrm{\Lambda }_O`$, the $`k=1`$ bounce action is
$$B\frac{27\pi ^2x^4}{2}\left[\frac{\rho _2}{c_O^3}\frac{\rho _2^2}{c_O^2\mathrm{\Lambda }_O}\right],\rho _2,\sqrt{\mathrm{\Lambda }_O}c_O.$$
(37)
Neglecting numerical factors, we find
$$Bx^4\mathrm{\Lambda }_{\mathrm{obs}}/\lambda ^2.$$
(38)
If we now take $`|\lambda |1`$, then the action is very small, and there is no effective stability.
On the other hand, if supersymmetry is broken at a low scale, then we expect $`|\lambda |1`$. Let us inquire when $`B\stackrel{>}{}1`$. This translates into
$$|\lambda _{\mathrm{halting}}|\stackrel{<}{}x^2\sqrt{\mathrm{\Lambda }_{\mathrm{obs}}},$$
(39)
or, restoring the mass units,
$$|\lambda _{\mathrm{halting}}|M^2\stackrel{<}{}x^2(2\times 10^3\mathrm{eV})^2(2.4\times 10^{18}\mathrm{GeV})^2x^2(2\mathrm{TeV})^4.$$
(40)
In models where supersymmetry breaking is transmitted with little suppression to the Standard Model fields (or is even present in the observable sector itself), it is reasonable to expect the supersymmetry breaking scale to set the scale of the bare cosmological constant, so
$$M_{\mathrm{weak}}^4M_{\mathrm{SUSY}}^4|\lambda _{\mathrm{halting}}|M^2.$$
(41)
Given the present experimental lower bounds on the supersymmetry breaking scale, we find then that the stability of the vacuum in this scenario requires low scale supersymmetry breaking, and relates the cosmological constant, Planck, and weak scales according to
$$M_{\mathrm{weak}}^2(10^3\mathrm{eV})(M_{\mathrm{Planck}}),$$
(42)
in accord with observation.
In a more careful analysis, one may require $`B1`$ for stability. However, the required supersymmetry breaking scale will not differ significantly from the above estimates, as $`B`$ goes as the inverse 8th power of the energy scale appearing in $`|\lambda |M^2`$. Finally, note that we have assumed the qualitative validity of the bounce action Eq. (7)–(11) throughout this section. For large $`\lambda 10^{60}`$, this may not be appropriate, as the brane curvature may be much greater than the brane tension. However, we have checked that in the case of greatest interest to us with $`\lambda 10^{60}`$, the brane curvature never greatly exceeds the brane tension, and the analysis above is under reasonable control.
### 5.2 $`T(T_IT_O)^{1/2}`$, multi-step relaxation
Motivated by the black hole analogy, we now consider an effective temperature that is the geometric mean of the initial and final de Sitter temperatures. In this case, tunneling to (and through) anti-de Sitter space is forbidden by fiat. However, the requirements of rapid tunneling to the observed cosmological constant and its stability are non-trivial constraints, and we now investigate their implications.
As in the previous scenario, we consider initial conditions $`c_O^2,\mathrm{\Lambda }_O|\lambda |`$. Now, however, the density of states factor $`D`$ is typically maximized for $`\mathrm{\Lambda }_I`$ within an order of magnitude of $`\mathrm{\Lambda }_O`$. To see this, a very rough estimate may be obtained by neglecting the bubble radius dependence on $`k`$ and approximating $`\mathrm{\Lambda }_O\mathrm{\Lambda }_I=(2\rho _2c_O\rho _2^2)/2\rho _2c_O`$. We then have $`D\sqrt{\mathrm{\Lambda }_O\mathrm{\Lambda }_I}(\mathrm{\Lambda }_O\mathrm{\Lambda }_I)^{3/2}`$, which is maximized for $`\mathrm{\Lambda }_I=\frac{1}{4}\mathrm{\Lambda }_O`$.
For $`\mathrm{\Lambda }_I\mathrm{\Lambda }_O`$,
$$D_{\mathrm{max}}\frac{\mathrm{\Lambda }_O^{5/2}}{|\lambda |^{7/4}\rho _2^{3/2}}.$$
(43)
It is not hard to verify that this degeneracy factor dominates the bounce suppression when $`\mathrm{\Lambda }_O|\lambda |`$. Thus, initially the effective cosmological constant tunnels rapidly as in the previous scenario, but in contrast to the previous case, the cosmological constant relaxes through several steps, with values roughly following a geometric series.
The effective cosmological constant will relax as described until $`\mathrm{\Lambda }_Oc_O^2`$, when Eq. (37) holds. At this point, the condition that tunneling continue is the requirement $`D_{\mathrm{max}}\stackrel{>}{}B`$, or, since $`B\mathrm{\Lambda }_O/\lambda ^2`$,
$$\mathrm{\Lambda }_O^{3/2}\stackrel{>}{}|\lambda |^{1/4}\rho _2^{3/2}.$$
(44)
For vanishing $`\rho _2`$, tunneling may continue to arbitrarily small $`\mathrm{\Lambda }_O`$. However, if we require stability from $`B\stackrel{>}{}1`$, we find, from Eq. (37),
$$\rho _2\stackrel{>}{}c_O^3|\lambda |^{3/2},$$
(45)
so $`\rho _2`$ cannot be arbitrarily small. Combining Eqs. (44) and (45), we find that tunneling stops when
$$\mathrm{\Lambda }_O\stackrel{>}{}|\lambda |^{4/3}.$$
(46)
Thus, even for $`|\lambda |10^{60}`$, although the effective cosmological constant is reduced by a factor of $`10^{20}`$, one membrane cannot suppress it to the observed value.
In general, however, it is important to note that several different 2-branes with various fundamental charge densities may be expected to arise. Suppose that another brane begins nucleating as the first membrane reaches its endpoint. The initial conditions for this new membrane are identical to those for the first brane, except that now the role of the initial bare cosmological constant is played by $`\mathrm{\Lambda }_O|\lambda |^{4/3}`$. For appropriate charge densities, $`n`$ branes may reduce the cosmological constant to $`|\lambda |^{(4/3)^n}`$. For $`|\lambda |10^{60}`$, three branes are sufficient to reduce the cosmological constant to its observed value.
So far we have considered only the ‘conservative’ $`\beta =3/2`$ case. For larger values of $`\beta `$ more complete relaxations of the cosmological term are possible. For general $`\beta `$, a single membrane may relax the cosmological constant to $`\mathrm{\Lambda }_O\stackrel{>}{}|\lambda |^{2\beta ^1}`$. Thus, even for the $`\beta =2`$ case, only two stages are required. Note also that in these multi-brane scenarios, in principle quite complex dynamics can arise, with periods of slow relaxation interspersed with more rapid changes.
## 6 Summary and Discussion
On very general grounds, it is appealing to think that relaxation of the cosmological term might be associated with very special degrees of freedom that have no conventional couplings to matter, and no conventional kinetic energy, but respond only to 3+1 dimensionally uniform form fields and, of course, to gravity. Several difficulties in such an approach must be addressed. The phenomenologically required energy scale is very small and not easily manufactured out of conventional energy scales. In addition, in any reasonable scenario, the cosmological constant must relax sufficiently quickly from high scales, but must be stable on cosmological time scales at its present value.
String theory provides a promising microscopic framework for such a mechanism. The necessary degrees of freedom are naturally supplied by string/M theory branes, and the dependence of brane properties on compactification appears, in principle, to be capable of producing a very small scale. We have also identified a candidate mechanism, the enhancement of multi-step jumps due to large density of states factors, which typically leads to large tunneling probabilities. Finally, as we have seen, the absolute stability of certain anti-de Sitter universes implies that near flat universes may be also be sufficiently stable.
We have considered two representative scenarios differing in the treatment of the effective temperature entering the density of states factor. In the simplest scenario, with $`TT_O`$, the exponentially most probable transition, excluding extremely short-lived universes, is to universes that are most nearly flat. By requiring that this new vacuum be sufficiently stable, we derived a non-trivial constraint for a mechanism of this kind. This constraint provides an intriguing relationship, between the supersymmetry breaking scale and the geometric mean of the present-day effective cosmological constant and Planck scales:
$$M_{\mathrm{SUSY}}^2\stackrel{<}{}(10^3\mathrm{eV})(M_{\mathrm{Planck}}).$$
(47)
Large supersymmetry breaking scales are thereby excluded, and the largest possible scale is plausibly though not necessarily satisfied in Nature.
In reaching this relation, we have assumed the bare cosmological constant to be of order the supersymmetry breaking scale. Indeed, while the bare cosmological term appears as a free parameter in supergravity, in string/M theory, phenomenologically interesting models with unbroken supersymmetry have zero cosmological term. In this context, it is therefore reasonable to expect that the relevant scale for the bare cosmological term is indeed the supersymmetry breaking scale. Many models of supersymmetry breaking invoke ‘hidden sectors’ with a characteristic mass scale much larger than the TeV scale. Unless the hidden sector contribution to the bare cosmological term is somehow suppressed to this TeV scale, in our simplest scenario the vacuum will be unstable.
Alternatively, an analogy with black holes suggests a richer dynamics, in which flat space plays a distinguished role and tunneling to anti-de Sitter space is forbidden. Within this circle of ideas, and in contrast to the previous scenario, we found that the cosmological constant relaxes through a several jumps, roughly following a geometric series. The constraint of stability limits the range over which the cosmological constant may be relaxed by any given membrane. However, two or more types of branes with radically different scales may relax the cosmological constant to within observational bounds, and appeal to the anthropic principle may be avoided.
Our work so far is very seriously incomplete, in that we have not attempted to incorporate it into a realistic (e.g., Friedmann-Robertson-Walker) cosmological model including matter. Thus, in particular, we have not addressed the dynamics of relaxation following a phase transition. There is a potential problem here, since, if after relaxation to zero effective cosmological term a later matter phase transition drives it negative, recovery may be difficult. We note also that the existence of very light membrane degrees of freedom may have a variety of observational and experimental consequences. We reserve discussion of these issues for a future publication.
We also require, for our dynamics, compactification schemes that produce the desired brane properties. Most model building in string/M theory has been based, implicitly or explicitly, on the paradigm of minimizing a potential. This has always been problematic within a quantum theory of gravity, but it was not clear what could replace it. The saltatory mechanism suggests a different principle, based on the dynamics of relaxation of the cosmological term. Whether this principle, or any other, is powerful enough to select uniquely a vacuum as complex as the one we observe remains to be seen.
## 7 Acknowledgments
We are very grateful to Luis Alvarez-Gaume, Paul Aspinwall, Raphael Bousso, Amit Giveon, Alex Kusenko, Juan Maldacena, Emil Martinec, Peter Mayr, Greg Moore, Yaron Oz, Joe Polchinski, Raman Sundrum, and Claudio Teitelboim for discussions. The work of JLF is supported in part by the Department of Energy under contract DE–FG02–90ER40542 and through the generosity of Frank and Peggy Taplin. JMR wishes to thank the Alfred P. Sloan Foundation for the award of a Fellowship, and the US Department of Energy for an Outstanding Junior Investigator Award. The work of SS is supported in part by the William Keck Foundation and by NSF grant No. PHY–9513835. The work of FW is supported in part by the Department of Energy under contract DE–FG02–90ER40542 and by the National Science Foundation under grant PHY–9513835.
## Note
As this manuscript was being completed, we learned of independent work by Bousso and Polchinski proposing quite a different scenario for fixing the cosmological term in string theory through a generalization of the Brown-Teitelboim mechanism. They do not utilize degenerating cycles nor enhanced density of states factors, and instead invoke the anthropic principle in an essential way. We thank them for conversations regarding our respective approaches.
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# A Note on Regularized Shannon’s Sampling Formulae
## I Introduction
In previous work, one of the present authors proposed a discrete singular convolution (DSC) algorithm for computer realization of singular convolutions involving singular kernels of delta type, Abel type and Hilbert type. One of illustrations for the algorithm was Shannon’s sampling formulae which plays an important role in the approximation of the delta distribution and generalized derivatives. However, in practical computations, the truncation error of Shannon’s sampling formulae is substantial. A regularization technique was used to construct a regularized Shannon’s sampling formulae, which was found to be extremely accurate and robust for resolving various challenging dynamical problems, such as the homoclinic orbit excitation of the Sine-Gordon equation, the Navier-Stokes flow in complex geometries, shock capturing of the inviscid Burgers’ equation, molecular quantum system described by the Schrödinger equation and nonlinear pattern formation of the Cahn-Hilliard equation.
The objective of the present note is to provide a theoretical analysis for the previous excellent numerical results. Rigorous error estimations of the regularized Shannon’s sampling formulae are given for their applications to interpolations and derivatives of a function.
## II Main result
Theorem. Let $`f`$ be a function $`fL^2(R)C^s(R)`$ and bandlimited to $`B`$, $`(B<\frac{\pi }{\mathrm{\Delta }},\mathrm{\Delta }`$ is the grid spacing). For a fixed $`tR`$ and $`\sigma >0`$, denote $`g(x)=f(x)H_k(\frac{tx}{\sqrt{2}\sigma })`$, where $`H_k(x)`$ is the $`k`$th order Hermite polynomial. If $`g(x)`$ satisfies
$$g^{}(x)g(x)\frac{(xt)}{\sigma ^2}$$
(1)
for $`xt+(M_11)\mathrm{\Delta }`$, and
$$g^{}(x)g(x)\frac{(xt)}{\sigma ^2}$$
(2)
for $`xtM_2\mathrm{\Delta }`$, where $`M_1,M_2𝒩`$, then for any $`s𝒵^+`$
$`f^{(s)}(t){\displaystyle \underset{n=\frac{t}{\mathrm{\Delta }}M_2}{\overset{\frac{t}{\mathrm{\Delta }}+M_1}{}}}f(n\mathrm{\Delta })\left[{\displaystyle \frac{\mathrm{sin}\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}\mathrm{exp}({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}})\right]^{(s)}_{L^2(R)}`$ (3)
$`\sqrt{3}[{\displaystyle \frac{f^{(s)}(t)_{L^2(R)}}{2\pi \sigma (\frac{\pi }{\mathrm{\Delta }}B)\mathrm{exp}(\frac{\sigma ^2(\frac{\pi }{\mathrm{\Delta }}B)^2}{2})}}`$ (4)
$`+{\displaystyle \frac{f(t)_{L^2(R)}_{i+j+k=s}\frac{s!\pi ^{i1}H_k(\frac{M_1\mathrm{\Delta }}{\sqrt{2}\sigma })}{i!k!\mathrm{\Delta }^{i1}(\sqrt{2}\sigma )^k((M_11)\mathrm{\Delta })^{j+1}}}{\mathrm{exp}(\frac{(M_1\mathrm{\Delta })^2}{2\sigma ^2})}}`$ (5)
$`+{\displaystyle \frac{f(t)_{L^2(R)}_{i+j+k=s}\frac{s!\pi ^{i1}H_k(\frac{M_2\mathrm{\Delta }}{\sqrt{2}\sigma })}{i!k!\mathrm{\Delta }^{i1}(\sqrt{2}\sigma )^k(M_2\mathrm{\Delta })^{j+1}}}{\mathrm{exp}(\frac{(M_2\mathrm{\Delta })^2}{2\sigma ^2})}}],`$ (6)
where superscript, $`(s)`$, denotes the $`s`$th order derivative.
## III Proof
### A Separation of the error
The error breaks naturally into a few components. Denote
$`E(t)=f^{(s)}(t){\displaystyle \underset{n=\frac{t}{\mathrm{\Delta }}M_2}{\overset{n=\frac{t}{\mathrm{\Delta }}+M_1}{}}}f(n\mathrm{\Delta })\left[{\displaystyle \frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}\mathrm{exp}\left({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}}\right)\right]^{(s)}`$ (7)
$`E_1(t)={\displaystyle \underset{n=\mathrm{}}{\overset{n=+\mathrm{}}{}}}f(n\mathrm{\Delta })\left[{\displaystyle \frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}{\displaystyle \frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}\mathrm{exp}\left({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}}\right)\right]^{(s)}`$ (8)
$`E_2(t)={\displaystyle \underset{n\frac{t}{\mathrm{\Delta }}+M_1}{}}f(n\mathrm{\Delta })\left[{\displaystyle \frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}\mathrm{exp}\left({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}}\right)\right]^{(s)}`$ (9)
$`E_3(t)={\displaystyle \underset{n\frac{t}{\mathrm{\Delta }}M_2}{}}f(n\mathrm{\Delta })\left[{\displaystyle \frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}\mathrm{exp}\left({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}}\right)\right]^{(s)}.`$ (10)
Here, $`E_1(t)`$ is regularization error. $`E_2(t)`$ and $`E_3(t)`$ are truncation errors. From Shannon’s sampling theorem
$$f(t)=\underset{n=\mathrm{}}{\overset{n=+\mathrm{}}{}}f(n\mathrm{\Delta })\frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })},$$
(11)
which can be differentiated term by term, the total error can be written as a sum of three components
$$E(t)=\underset{i=1}{\overset{3}{}}E_i(t).$$
(12)
The corresponding error norms satisfy
$$E(t)_{L^2(R)}\sqrt{3}\underset{i=1}{\overset{3}{}}E_i(t)_{L^2(R)}.$$
(13)
### B Estimation of $`E_1(t)`$
Let $`f\widehat{(}\omega )`$ be the Fourier transform of $`f(x)`$, and $`f\widehat{(}\omega )=_Rf(x)\mathrm{exp}(ix\omega )dx`$. Since
$$\left[\frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}\right]\widehat{(}\omega )=\mathrm{\Delta }\mathrm{exp}(in\mathrm{\Delta }\omega )\chi _{[\frac{\pi }{\mathrm{\Delta }},\frac{\pi }{\mathrm{\Delta }}]}(w)$$
(14)
and
$$\left[\mathrm{exp}(\frac{(tn\mathrm{\Delta })^2}{2\sigma ^2})\right]\widehat{(}\omega )=\sqrt{2\pi }\sigma \mathrm{exp}(in\mathrm{\Delta }\omega \frac{\sigma ^2\omega ^2}{2}),$$
(15)
one writes
$`\left[{\displaystyle \frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}\right]\widehat{(}\omega )\left[\mathrm{exp}({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}})\right]\widehat{(}\omega )`$ (16)
$`={\displaystyle _R}\mathrm{\Delta }\sqrt{2\pi }\sigma \mathrm{exp}(in\mathrm{\Delta }(\omega \theta ))\chi _{[\theta \frac{\pi }{\mathrm{\Delta }},\theta +\frac{\pi }{\mathrm{\Delta }}]}(\omega )\mathrm{exp}(in\mathrm{\Delta }\theta {\displaystyle \frac{\sigma ^2\theta ^2}{2}})𝑑\theta `$ (17)
$`=\mathrm{\Delta }\sqrt{2\pi }\sigma \mathrm{exp}(in\mathrm{\Delta }\omega ){\displaystyle _{\theta \frac{\pi }{\mathrm{\Delta }}}^{\theta +\frac{\pi }{\mathrm{\Delta }}}}\mathrm{exp}({\displaystyle \frac{\sigma ^2\theta ^2}{2}})𝑑\theta .`$ (18)
From Eq. (16)
$`\left[{\displaystyle \underset{n=\mathrm{}}{\overset{n=+\mathrm{}}{}}}f(n\mathrm{\Delta }){\displaystyle \frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}\mathrm{exp}({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}})\right]\widehat{(}\omega )`$ (19)
$`={\displaystyle \underset{n=\mathrm{}}{\overset{n=+\mathrm{}}{}}}f(n\mathrm{\Delta }){\displaystyle \frac{1}{2\pi }}[{\displaystyle \frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}]\widehat{(}\omega )[\mathrm{exp}({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}})]\widehat{(}\omega )`$ (20)
$`={\displaystyle \underset{n=\mathrm{}}{\overset{n=+\mathrm{}}{}}}f(n\mathrm{\Delta })\mathrm{\Delta }\mathrm{exp}(in\mathrm{\Delta }\omega ){\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle _{\frac{\sigma (\omega \frac{\pi }{\mathrm{\Delta }})}{\sqrt{2}}}^{\frac{\sigma (\omega +\frac{\pi }{\mathrm{\Delta }})}{\sqrt{2}}}}\mathrm{exp}(t^2)𝑑t.`$ (21)
Since function $`f`$ satisfies
$$f\widehat{(}\omega )L^2[B,B]L^2[\frac{\pi }{\mathrm{\Delta }},\frac{\pi }{\mathrm{\Delta }}],$$
(22)
it has a Fourier series expansion
$$f\widehat{(}\omega )=_{n=\mathrm{}}^{\mathrm{}}c_n\mathrm{exp}(in\mathrm{\Delta }\omega ),$$
(23)
where coefficients is given by
$$c_n=\frac{\mathrm{\Delta }}{2\pi }_{\frac{\pi }{\mathrm{\Delta }}}^{\frac{\pi }{\mathrm{\Delta }}}f\widehat{(}\omega )\mathrm{exp}(in\mathrm{\Delta }\omega )d\omega =\mathrm{\Delta }f(n\mathrm{\Delta }).$$
(24)
Equivalently, $`f\widehat{(}\omega )`$ can be written
$$f\widehat{(}\omega )=f\widehat{(}\omega )\chi _{[B,B]}(\omega )=_{n=\mathrm{}}^{\mathrm{}}\mathrm{\Delta }f(n\mathrm{\Delta })\mathrm{exp}(in\mathrm{\Delta }\omega )\chi _{[B,B]}(\omega ).$$
(25)
Denote
$$\epsilon (\omega )=\chi _{[B,B]}(\omega )\frac{1}{\sqrt{\pi }}_{\frac{\sigma (\omega \frac{\pi }{\mathrm{\Delta }})}{\sqrt{2}}}^{\frac{\sigma (\omega +\frac{\pi }{\mathrm{\Delta }})}{\sqrt{2}}}\mathrm{exp}(t^2)𝑑t,$$
(26)
then combining Eqs. (8), (21), (23) and (25), one has
$$E_1\widehat{(}\omega )=(i\omega )^sf\widehat{(}\omega )\epsilon (\omega ).$$
(27)
For $`\omega [B,B]`$, $`\epsilon (\omega )`$ can be evaluated as
$`\epsilon (\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}\left[{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{exp}(t^2)𝑑t{\displaystyle _{\frac{\sigma (\omega \frac{\pi }{\mathrm{\Delta }})}{\sqrt{2}}}^{\frac{\sigma (\omega +\frac{\pi }{\mathrm{\Delta }})}{\sqrt{2}}}}\mathrm{exp}(t^2)𝑑t\right]`$ (28)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}\left[{\displaystyle _{\frac{\sigma (\frac{\pi }{\mathrm{\Delta }}\omega )}{\sqrt{2}}}^{\mathrm{}}}\mathrm{exp}(t^2)𝑑t+{\displaystyle _{\frac{\sigma (\omega +\frac{\pi }{\mathrm{\Delta }})}{\sqrt{2}}}^{\mathrm{}}}\mathrm{exp}(t^2)𝑑t\right].`$ (29)
Moreover, for $`x0`$, the following inequality is valid
$$\frac{1}{x+\sqrt{x^2+2}}\mathrm{exp}(x^2)_x^{\mathrm{}}\mathrm{exp}(t^2)𝑑t\frac{1}{x+\sqrt{x^2+\frac{4}{\pi }}}.$$
(30)
Therefore, the estimation for $`\epsilon (\omega )`$ is obtained as
$`\epsilon (\omega )`$ $``$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}\left({\displaystyle \frac{\mathrm{exp}\left(\frac{\sigma ^2(\frac{\pi }{\mathrm{\Delta }}\omega )^2}{2}\right)}{\sqrt{2}\sigma (\frac{\pi }{\mathrm{\Delta }}\omega )}}+{\displaystyle \frac{\mathrm{exp}\left(\frac{\sigma ^2(\frac{\pi }{\mathrm{\Delta }}+\omega )^2}{2}\right)}{\sqrt{2}\sigma (\frac{\pi }{\mathrm{\Delta }}+\omega )}}\right)`$ (31)
$``$ $`{\displaystyle \frac{1}{\sigma (\frac{\pi }{\mathrm{\Delta }}B)\mathrm{exp}\left(\frac{\sigma ^2(\frac{\pi }{\mathrm{\Delta }}B)^2}{2}\right)}}.`$ (32)
It follows from Eqs. (27) and (32) that
$`E_1\widehat{(}\omega ){\displaystyle \frac{f\widehat{(}\omega )(i\omega )^s}{\sigma (\frac{\pi }{\mathrm{\Delta }}B)\mathrm{exp}(\frac{\sigma ^2(\frac{\pi }{\mathrm{\Delta }}B)^2}{2})}}.`$ (33)
From the Parseval identity, one has
$`E_1(t)_{L^2(R)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}E_1\widehat{(}\omega )_{L^2(R)}`$ (34)
$``$ $`{\displaystyle \frac{f^{(s)}(t)_{L^2(R)}}{2\pi \sigma (\frac{\pi }{\mathrm{\Delta }}B)\mathrm{exp}\left(\frac{\sigma ^2(\frac{\pi }{\mathrm{\Delta }}B)^2}{2}\right)}}.`$ (35)
### C Estimation of $`E_2(t)`$
Differentiations can be written
$`E_2(t)={\displaystyle \underset{n\frac{t}{\mathrm{\Delta }}+M_1}{}}f(n\mathrm{\Delta })[{\displaystyle \underset{i+j+k=s}{}}{\displaystyle \frac{s!}{i!k!}}({\displaystyle \frac{\pi }{\mathrm{\Delta }}})^{i1}\mathrm{sin}({\displaystyle \frac{\pi }{\mathrm{\Delta }}}tn\pi +{\displaystyle \frac{\pi i}{2}})`$ (36)
$`{\displaystyle \frac{(1)^j}{(tn\mathrm{\Delta })^{j+1}}}{\displaystyle \frac{(1)^k}{(\sqrt{2}\sigma )^k}}H_k({\displaystyle \frac{tn\mathrm{\Delta }}{\sqrt{2}\sigma }})\mathrm{exp}({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}})],`$ (37)
where $`H_k(x)`$ is the Hermite polynomial
$$\mathrm{exp}(x^2)H_k(x)=(1)^k(\frac{d}{dx})^k\mathrm{exp}(x^2).$$
(38)
Let $`l=n\frac{t}{\mathrm{\Delta }}`$, where $`x`$ is the integral part of $`x`$ and $`x=xx`$, then
$`E_2(t)`$ $`=`$ $`{\displaystyle \underset{lM_1}{}}f(t+l\mathrm{\Delta }{\displaystyle \frac{t}{\mathrm{\Delta }}}\mathrm{\Delta })[{\displaystyle \underset{i+j+k=s}{}}{\displaystyle \frac{s!}{i!k!}}({\displaystyle \frac{\pi }{\mathrm{\Delta }}})^{i1}(1)^{j+k+l}`$ (41)
$`\mathrm{sin}\left({\displaystyle \frac{t}{\mathrm{\Delta }}}\pi +{\displaystyle \frac{\pi i}{2}}\right){\displaystyle \frac{1}{(\sqrt{2}\sigma )^k\mathrm{\Delta }^{j+1}(l+\frac{t}{\mathrm{\Delta }})^{j+1}}}`$
$`\left[H_k\left({\displaystyle \frac{l\mathrm{\Delta }+\frac{t}{\mathrm{\Delta }}\mathrm{\Delta }}{\sqrt{2}\sigma }}\right)\right]\mathrm{exp}\left({\displaystyle \frac{(l\mathrm{\Delta }\frac{t}{\mathrm{\Delta }}\mathrm{\Delta })^2}{2\sigma ^2}}\right)`$
$`=`$ $`{\displaystyle \underset{i+j+k=s}{}}{\displaystyle \underset{lM_1}{}}F_k(l)s_{i,j,k}(l),`$ (42)
where
$`F_k(l)=`$ $`f(t+l\mathrm{\Delta }{\displaystyle \frac{t}{\mathrm{\Delta }}}\mathrm{\Delta })H_k\left({\displaystyle \frac{l\mathrm{\Delta }+\frac{t}{\mathrm{\Delta }}\mathrm{\Delta }}{\sqrt{2}\sigma }}\right)`$ (44)
$`\mathrm{exp}\left({\displaystyle \frac{(l\mathrm{\Delta }\frac{t}{\mathrm{\Delta }}\mathrm{\Delta })^2}{2\sigma ^2}}\right)`$
$`s_{i,j,k}(l)=`$ $`(1)^l{\displaystyle \frac{s!\pi ^{i1}(1)^{j+k}\mathrm{sin}(\frac{t}{\mathrm{\Delta }}\pi +\frac{\pi i}{2})}{i!k!\mathrm{\Delta }^{i+j}(l+\frac{t}{\mathrm{\Delta }})^{j+1}(\sqrt{2}\sigma )^k}}.`$ (45)
Two simple lemmas are required.
Lemma 1 (Abel’s inequality). For two sequences $`\{a_n\},\{b_n\},b_1b_2\mathrm{}b_n,a_n,b_nR`$, set
$`s_k`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}a_i`$ (46)
$`m`$ $`=`$ $`\underset{1kn}{\mathrm{min}}s_k`$ (47)
$`M`$ $`=`$ $`\underset{1kn}{\mathrm{max}}s_k,`$ (48)
then
$$mb_1\underset{i=1}{\overset{n}{}}a_ib_iMb_1.$$
(49)
Lemma 2. As all the notations unchanged, set $`g(x)=f(x)H_k(\frac{tx}{\sqrt{2}\sigma })`$. If $`g(x)`$ satisfies
$$g^{}(x)g(x)\frac{(xt)}{\sigma ^2},$$
(50)
whenever $`xt+(M_11)\mathrm{\Delta }`$, then $`\{F_k(l)\}_{l𝒩}`$ is a decreasing sequence.
The proof is obvious by taking the first order derivative.
Let denote
$$S_{i,j,k}(N)=\underset{lM_1}{\overset{M_1+N}{}}s_{i,j,k}(l).$$
(51)
It is estimated that
$`|S_{i,j,k}(N)|`$ $`=`$ $`\left|{\displaystyle \underset{lM_1}{\overset{M_1+N}{}}}(1)^l{\displaystyle \frac{s!\pi ^{i1}\mathrm{sin}(\frac{t}{\mathrm{\Delta }}\pi +\frac{\pi i}{2})(1)^{k+j}}{i!k!\mathrm{\Delta }^{i+j}(\sqrt{2}\sigma )^k}}{\displaystyle \frac{1}{(l+\frac{t}{\mathrm{\Delta }})^{j+1}}}\right|`$ (52)
$`=`$ $`\left|{\displaystyle \frac{s!\pi ^{i1}\mathrm{sin}(\frac{t}{\mathrm{\Delta }}\pi +\frac{\pi i}{2})(1)^{k+j}}{i!k!\mathrm{\Delta }^{i+j}(\sqrt{2}\sigma )^k}}{\displaystyle \underset{lM_1}{\overset{M_1+N}{}}}(1)^l{\displaystyle \frac{1}{(l+\frac{t}{\mathrm{\Delta }})^{j+1}}}\right|`$ (53)
$``$ $`{\displaystyle \frac{s!\pi ^{i1}}{i!k!\mathrm{\Delta }^{i+j}(\sqrt{2}\sigma )^k}}{\displaystyle \frac{1}{(M_11)^{j+1}}}.`$ (54)
Then from $`(28)`$ and $`(37)`$, and by using lemma $`1`$ and lemma $`2`$, one has
$`E_2(t)f(t+M_1\mathrm{\Delta })\mathrm{exp}\left({\displaystyle \frac{(M_1\mathrm{\Delta })^2}{2\sigma ^2}}\right)`$ (55)
$`\times {\displaystyle \underset{i+j+k=s}{}}{\displaystyle \frac{s!\pi ^{i1}H_k(\frac{M_1\mathrm{\Delta }}{\sqrt{2}\sigma })}{i!k!\mathrm{\Delta }^{i1}(\sqrt{2}\sigma )^k\left((M_11)\mathrm{\Delta }\right)^{j+1}}}.`$ (56)
This gives rise to
$$E_2(t)_{L^2(R)}\frac{f(t)_{L^2(R)}_{i+j+k=s}\frac{s!\pi ^{i1}H_k(\frac{M_1\mathrm{\Delta }}{\sqrt{2}\sigma })}{i!k!\mathrm{\Delta }^{i1}(\sqrt{2}\sigma )^k\left((M_11)\mathrm{\Delta }\right)^{j+1}}}{\mathrm{exp}\left(\frac{(M_1\mathrm{\Delta })^2}{2\sigma ^2}\right)}.$$
(57)
### D Estimation of $`E_3(t)`$
A result like lemma 2 is required.
Lemma 3. Notations are the same as before. Denote $`g(x)=f(x)H_k(\frac{tx}{\sqrt{2}\sigma })`$, if $`g(x)`$ satisfies
$$g^{}(x)g(x)\frac{(xt)}{\sigma ^2},$$
(58)
whenever $`xtM_2\mathrm{\Delta }`$, then $`\{F_k(l)\}_{l𝒩}`$ is an increasing sequence.
The proof is also direct. Therefore, by the same treatment as that in the previous subsection, we obtain
$$E_3(t)_{L^2(R)}\frac{f(t)_{L^2(R)}_{i+j+k=s}\frac{s!\pi ^{i1}H_k(\frac{M_2\mathrm{\Delta }}{\sqrt{2}\sigma })}{i!k!\mathrm{\Delta }^{i1}(\sqrt{2}\sigma )^k(M_2\mathrm{\Delta })^{j+1}}}{\mathrm{exp}\left(\frac{(M_2\mathrm{\Delta })^2}{2\sigma ^2}\right)}.$$
(59)
### E The end of the proof
By combining Eqs. (13), (34), (57) and (59), one obtains Eq. (6).
## IV Discussion
*Remark 1*. For $`s=0`$, one has
$`f(t){\displaystyle \underset{n=\frac{t}{\mathrm{\Delta }}M_2}{\overset{\frac{t}{\mathrm{\Delta }}+M_1}{}}}f(n\mathrm{\Delta }){\displaystyle \frac{\mathrm{sin}\left(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })\right)}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}}\mathrm{exp}\left({\displaystyle \frac{(tn\mathrm{\Delta })^2}{2\sigma ^2}}\right)_{L^2(R)}`$ (60)
$`\sqrt{3}f(t)_{L^2(R)}\{{\displaystyle \frac{1}{2\pi \sigma (\frac{\pi }{\mathrm{\Delta }}B)\mathrm{exp}\left(\frac{\sigma ^2(\frac{\pi }{\mathrm{\Delta }}B)^2}{2}\right)}}`$ (61)
$`+{\displaystyle \frac{1}{(M_11)\pi \mathrm{exp}\left(\frac{(M_1\mathrm{\Delta })^2}{2\sigma ^2}\right)}}+{\displaystyle \frac{1}{M_2\pi \mathrm{exp}\left(\frac{(M_2\mathrm{\Delta })^2}{2\sigma ^2}\right)}}\}.`$ (62)
This is a rigorous error statement for the formulae widely used in the aforementioned numerical computations. Roughly speaking, if $`\mathrm{exp}(\frac{x^2}{2})=10^\eta `$, then $`\eta =\frac{x^2}{2ln10}`$, so the error is
$`\epsilon (r,B,\mathrm{\Delta },M)=\sqrt{3}f(t)_{L^2(R)}({\displaystyle \frac{1}{2\pi r(\pi B\mathrm{\Delta })10^{\frac{r^2(\pi B\mathrm{\Delta })^2}{2ln10}}}}`$ (63)
$`+{\displaystyle \frac{1}{(M_11)\mathrm{\Delta }10^{\frac{(M_11)^2}{2r^2ln10}}}}+{\displaystyle \frac{1}{M_2\mathrm{\Delta }10^{\frac{(M_2)^2}{2r^2ln10}}}}),`$ (64)
where $`r=\frac{\sigma }{\mathrm{\Delta }}`$. One may choose $`r,B,\mathrm{\Delta }`$ and $`M`$ appropriately to attain desired accuracy. Assume all non-exponential quantities are combined to give unit, and $`M=M_11=M_2`$, one has
$$r(\pi B\mathrm{\Delta })>\sqrt{\eta 2\mathrm{ln}10}$$
(65)
and
$$\frac{M}{r}>\sqrt{\eta 2\mathrm{ln}10},$$
(66)
where $`\eta `$ is the desired order of accuracy. There are some general rules for attaining high accuracy. These are discussed from two different arguments.
1.) For a given function $`f(x)`$ with a known bandlimit $`B`$, other parameters, $`\mathrm{\Delta },r`$ and $`M`$, are to be chosen appropriately to achieve a desired accuracy order $`\eta `$:
(i) From Eqs. (65) and (66) one has $`B\mathrm{\Delta }\pi \frac{\sqrt{\eta 2\mathrm{ln}10}}{r}`$. For fixed $`r`$, the higher the frequency bandlimit $`B`$ is, the smaller $`\mathrm{\Delta }`$ should be, which means the more grid points in the computational domain. When $`\mathrm{\Delta }`$ varies from $`0`$ to $`\frac{\pi }{B}`$, $`r`$ changes from $`\frac{\sqrt{\eta 2\mathrm{ln}10}}{\pi }`$ to $`+\mathrm{}`$, therefore for sufficiently small $`\mathrm{\Delta }`$, $`r`$ is near $`\frac{\sqrt{\eta 2\mathrm{ln}10}}{\pi }`$.
(ii) No matter how many grid points are in the computational domain, $`r`$ and $`M`$ cannot be too small. Equations (65) and (66) indicate $`r>\frac{\sqrt{\eta 2\mathrm{ln}10}}{\pi }`$ and $`M>r\sqrt{\eta 2\mathrm{ln}10}`$. If $`M`$ and $`r`$ are less than the minimal requirements, the accuracy deteriorates quickly. On the other hand, if sufficiently large $`r`$ and $`M`$ are used, say, $`M=30`$ and $`r=3.5`$, high approximation accuracy can be achieved.
2.) In practical computations, such as in solving a partial differential equation, the function $`f`$ and its $`B`$ are unknown. In this case, $`\mathrm{\Delta }`$ is selected a priori. Then $`r`$ and $`M`$ are to be chosen properly for achieving a desired accuracy order $`\eta `$:
(i) For a given grid spacing, $`\mathrm{\Delta }`$, and accuracy requirement $`\eta `$, $`r`$ value determines frequency bandlimit $`B`$ which can be reached. Then the set of functions $`f`$ which are almost bandlimited to $`B`$ can be accurately approximated (where ‘almost bandlimited to $`B`$’ means the function $`f`$ is not necessarily bandlimited but its Fourier amplitude outside $`|B|`$ is much smaller than the given error $`10^\eta `$). The choice of $`M`$ should be consistent with $`r`$ for a given accuracy requirement. In general, small $`r`$ and $`M`$ values lead to an accurate approximation for low frequency component of a function of interest. But the prediction of a high frequency component will not be accurate in such a case.
(ii) For a given grid spacing $`\mathrm{\Delta }`$ and $`r`$ value, the larger $`M`$ is, the higher bandlimit $`B`$ can be reached.
(iii) To improve computational efficiency with a given $`\mathrm{\Delta }`$, $`B`$ shall be very close to $`\frac{\pi }{\mathrm{\Delta }}`$. However, to maintain certain approximation accuracy, $`r`$ has to be sufficiently large, which implies that $`M`$ has to be very large too. This in turn results in low efficiency (It takes $`M\mathrm{}`$ to maintain the accuracy if one samples at the Nyquist rate).
(iv) If $`\mathrm{\Delta }`$, $`M`$ and $`\eta `$ are chosen, then $`r`$ is also fixed. For example, to achieve the machine precision $`10^\eta 10^{15}`$, Eq. (65) estimates $`r>2.8`$. If this is achieved by using $`M=33`$, then Eq. (66) estimates $`r<4`$. In fact, $`M30`$ and $`2.8<r<4.0`$ are the parameter regions found from an earlier numerical test and were used in many applications.
*Remark 2*. A comparison between the truncation errors of Shannon’s sampling formulae and the regularized Shannon’s sampling formulae is in order. Reference estimates that the expression
$$(T_Nf)(t)=f(t)\underset{n=N}{\overset{n=+N}{}}f(n\mathrm{\Delta })\frac{\mathrm{sin}(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta }))}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}$$
(67)
has error of
$$|T_N(t)|\frac{\sqrt{2}}{\pi }\sqrt{E}\left|\mathrm{sin}(\frac{\pi t}{\mathrm{\Delta }})\right|\sqrt{\frac{N\mathrm{\Delta }}{(N\mathrm{\Delta }^2t^2)}},$$
(68)
where $`t<N\mathrm{\Delta }`$, and $`E`$ is the total ‘energy’ of the function given by
$$E=_{\frac{\pi }{\mathrm{\Delta }}}^{\frac{\pi }{\mathrm{\Delta }}}|f\widehat{(}w)|^2dw.$$
(69)
This is not directly comparable with our error estimate because our sampling is centered around a point of interest, $`x`$. Let consider a truncation error of the form
$$(E_Mf)(t)=f(t)\underset{n=\frac{t}{\mathrm{\Delta }}M}{\overset{n=\frac{t}{\mathrm{\Delta }}+M}{}}f(n\mathrm{\Delta })\frac{\mathrm{sin}(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta }))}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}.$$
(70)
In Appendix A, it is shown that in a finite computational domain, the $`L^2`$ norm of $`(E_Mf)(t)`$ has the order of $`f(t)_{L^2(R)}\sqrt{\frac{1}{M\mathrm{\Delta }}}`$, which is much larger than the truncation error of the regularized Shannon’s formulae. On the other hand, to achieve the same accuracy, the regularized formulae requires much fewer computational grids.
*Remark 3*. Discussions for the higher order derivatives can be presented in a similar manner as those of Remarks 1 and 2. In fact, previous work of solving partial differential equations. involved such derivatives, and results are consistent with the present theorem. Detailed comparison is omitted.
*Remark 4*. In many practical applications, such as in solving partial differential equations, error estimations and discussions in other spaces are often required. Moreover, in real computations, the computational domain is always limited to a finite interval, such as $`[a,b]`$. Therefore, the norm $`f_{L^2(R)}`$ in Eqs. (57) and (59) are required to be changed into $`f_{L^2(a,b)}`$, which can be evaluated by integrations along $`[a+M_1\mathrm{\Delta },b+M_1\mathrm{\Delta }]`$ and $`[aM_2\mathrm{\Delta },bM_2\mathrm{\Delta }]`$ respectively. Therefore various $`L^p(1p2)`$ error estimates of $`E(t)`$ can be derived accordingly. If we know the size of $`L^p`$ norm $`(1p2)`$ of the function of interest, then we can deduce from the theorem the critical values, $`r`$ and $`M`$, to achieve desired accuracy.
## A Truncation error of Shannon’s sampling formulae
Lemma 4. In the computational domain $`[a,b]`$, the error expression
$$(E_Mf)(t)=f(t)\underset{n=\frac{t}{\mathrm{\Delta }}M}{\overset{n=\frac{t}{\mathrm{\Delta }}+M}{}}f(n\mathrm{\Delta })\frac{\mathrm{sin}(\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta }))}{\frac{\pi }{\mathrm{\Delta }}(tn\mathrm{\Delta })}$$
(A1)
satisfies the estimate
$$(E_Mf)(t)_{L^2(a,b)}\frac{2f(t)_{L^2(R)}}{\sqrt{(M2)\mathrm{\Delta }}}.$$
(A2)
Proof. Let denote
$$f_M(t)=\frac{1}{\pi }\mathrm{sin}\left(\frac{\pi t}{\mathrm{\Delta }}\right)\underset{n=\frac{t}{\mathrm{\Delta }}M}{\overset{n=\frac{t}{\mathrm{\Delta }}+M}{}}\frac{f(n\mathrm{\Delta })(1)^n}{(\frac{t}{\mathrm{\Delta }}n)}.$$
(A3)
By using Schwartz’s inequality, one has
$`\left({\displaystyle \underset{n=\frac{t}{\mathrm{\Delta }}+M}{\overset{n=+\mathrm{}}{}}}{\displaystyle \frac{f(n\mathrm{\Delta })(1)^n}{(\frac{t}{\mathrm{\Delta }}n)}}\right)^2`$ (A4)
$`f(t)_{L^2(R)}^2{\displaystyle \underset{n=\frac{t}{\mathrm{\Delta }}+M}{\overset{n=+\mathrm{}}{}}}{\displaystyle \frac{1}{(tn\mathrm{\Delta })^2}}`$ (A5)
$`{\displaystyle \frac{1}{\mathrm{\Delta }^2}}f(t)_{L^2(R)}^2{\displaystyle \underset{lM}{}}{\displaystyle \frac{1}{(l1)^2}}`$ (A6)
$`{\displaystyle \frac{1}{\mathrm{\Delta }^2}}f(t)_{L^2(R)}^2{\displaystyle _{M2}^+\mathrm{}}{\displaystyle \frac{dx}{x^2}}`$ (A7)
$`={\displaystyle \frac{f(t)_{L^2(R)}^2}{(M2)\mathrm{\Delta }^2}}.`$ (A8)
Similarly one obtains
$$\left(\underset{n=\mathrm{}}{\overset{n=\frac{t}{\mathrm{\Delta }}M}{}}\frac{f(n\mathrm{\Delta })(1)^n}{(\frac{t}{\mathrm{\Delta }}n)}\right)^2\frac{f(t)_{L^2(R)}^2}{(M1)\mathrm{\Delta }^2}.$$
(A9)
By combining Eqs. (A3), (A8) and (A9), one finishes the proof.
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# A gentle introduction to the foundations of classical electrodynamics: The meaning of the excitations (𝒟,ℋ) and the field strengths (𝐸,𝐵)
## I Introduction
In Cologne, we teach a course on Theoretical Physics II (electrodynamics) to students of physics in their fourth semester. For several years, we have been using for that purpose the calculus of exterior differential forms, see , because we believe that this is the appropriate formalism: It is based on objects which possess a clear operational interpretation, it elucidates the fundamental structure of Maxwell’s equations and their mutual interrelationship, and it invites a 4-dimensional representation appropriate for special and general relativity theory (i.e., including gravity, see ).
Our experimental colleagues are somewhat skeptical; and not only them. Therefore we were invited to give, within 90 minutes, a sort of popular survey of electrodynamics in exterior calculus to the members of one of our experimental institutes (group of H. Micklitz). The present article is a worked-out version of this talk. We believe that it could also be useful for other universities.
Subsequent to the talk we had given, we found the highly interesting and historically oriented article of Roche on “$`B`$ and $`H`$, the intensity vectors of magnetism…”. Therein, the corresponding work of Bamberg and Sternberg , Bopp , Ingarden and Jamiołkowski , Kovetz , Post , Sommerfeld , and Truesdell and Toupin , to drop just a few names, was neglected yielding a picture of $`H`$ and $`B`$ which looks to us as being not up of date; one should also compare in this context the letter of Chambers and the book of Roche , in particular its Chapter 9. Below we will suggest answers to some of Roche’s questions.
Moreover, “…any system that gives $`E`$ and $`B`$ different units, when they are related through a relativistic transformation, is on the far side of sanity” is an apodictic statement of Fitch . In the sequel, we will prove that we are on the far side of sanity: The absolute dimension of $`E`$ turns out to be magnetic flux/time and that of $`B`$ magnetic flux, see Sec. IV.
According to the audience we want to address, we will skip all mathematical details and take recourse to plausibility considerations. In order to make the paper self-contained, we present though a brief summary of exterior calculus in the Appendix. A good reference to the mathematics underlying our presentation is the book of Frankel , see also and . For the experimental side of our subject we refer to Bergmann-Schaefer .
As a preview, let us collect essential information about the electromagnetic field in Table I. The explanations will follow below.
Table I. The electromagnetic field
| Field | name | math. | independent | related | reflec- | absolute |
| --- | --- | --- | --- | --- | --- | --- |
| | | object | components | to | tion | dimension |
| $`𝒟`$ | electric | odd | $`𝒟_{23},𝒟_{31},𝒟_{12}`$ | area | $`𝒟`$ | $`q=`$ electric |
| | excitation | 2-form | | | | charge |
| $``$ | magnetic | odd | $`_1,_2,_3`$ | line | $``$ | $`q/t`$ |
| | excitation | 1-form | | | | |
| $`E`$ | electric | even | $`E_1,E_2,E_3`$ | line | $`E`$ | $`\mathrm{\Phi }_0/t`$ |
| | field strength | 1-form | | | | |
| $`B`$ | magnetic | even | $`B_{23},B_{31},B_{12}`$ | area | $`B`$ | $`\mathrm{\Phi }_0=`$ mag- |
| | field strength | 2-form | | | | netic flux |
It was Maxwell himself who advised us to be very careful in assigning a certain physical quantity to a mathematical object. As it turns out, the mathematical images of $`𝒟,,E,B`$ are all different from each other. This is well encoded in Schouten’s images of the electromagnetic field in Fig.1.
## II Electric charge conservation
The conservation of electric charge was already recognized as fundamental law during the time of Franklin (around 1750) well before Coulomb discovered his force law in 1785. Nowadays, when one can catch single electrons and single protons and their antiparticles in traps and can count them individually (see, e.g., Dehmelt , Paul , Devoret et al. , and Lafarge et al. ), we are more sure than ever that electric charge conservation is a valid fundamental law of nature. Therefore matter carries as a primary quality something called electric charge which only occurs in positive or negative units of an elementary charge $`e`$ (or, in the case of quarks, of $`1/3`$th of it) and which can be counted. Thus it is justified to introduce the physical dimension of charge $`q`$ as a new and independent concept. Ideally one should measure a charge in units of $`e/3`$. However, for practical reasons, the SI-unit C (coulomb) is used in laboratory physics.
Let us start with the 3-dimensional Euclidean space in which we mark a 3-dimensional domain $`V`$. Hereafter, the local coordinates in this space will be denoted by $`x^a`$ and the time by $`t`$, with the basis vectors $`e_a:=_a`$ and $`a,b,\mathrm{}=1,2,3`$, see Fig.2. The total charge in the domain $`V`$ is given by the integral
$$Q=\underset{V}{}\rho ,$$
(1)
where the electric charge density $`\rho `$ is the 3-form $`\rho =\frac{1}{3!}\rho _{abc}dx^adx^bdx^c=\rho _{123}dx^1dx^2dx^3`$. Here summation is understood over the indices $`a,b,c`$ and $`\rho _{abc}=\rho _{bac}=\rho _{bca}=\mathrm{}`$, i.e., the components $`\rho _{abc}`$ of the charge density 3-form $`\rho `$ are antisymmetric under the exchange of two indices, leaving only one independent component $`\rho _{123}`$. The wedge $``$ denotes the (anticommutative) exterior product of two forms, and $`dx^1dx^2dx^3`$ represents the volume “element”. For our present purpose it is enough to know, for more details see the Appendix, that a 3-form (a $`p`$-form) is an object that, if integrated over a 3-dimensional ($`p`$-dimensional) domain, yields a scalar quantity, here the charge $`Q`$.
The dimension of $`Q`$ is $`[Q]=q`$. Since an integral (a summation after all) cannot change the dimension, the dimension of the charge density 3-form and its components are, respectively, $`[\rho ]=q`$, and $`[\rho _{abc}]=q/\mathrm{}^3`$, with $`\mathrm{}=`$ length.
In general, the elementary charges are not at rest. The electric current $`J`$ flowing across a 2-dimensional surface $`S`$ is given by the integral
$$J=\underset{S}{}j,$$
(2)
see Fig.3. Accordingly, the electric current density $`j`$ turns out to be a 2-form: $`j=\frac{1}{2!}j_{ab}dx^adx^b=j_{12}dx^1dx^2+j_{13}dx^1dx^3+j_{23}dx^2dx^3`$, with $`j_{ab}=j_{ba}`$. If $`t=`$ time, then the dimensions of the current integral and the current 2-form and its components are $`[J]=[j]=q/t`$ and $`[j_{ab}]=q/(t\mathrm{}^2)`$, respectively.
If we use the abbreviation $`_t:=/t`$, the global electric charge conservation can be expressed as
$$_t\underset{V}{}\rho +\underset{V}{}j=0(\mathrm{Axiom}\mathrm{\hspace{0.25em}\hspace{0.25em}1}),$$
(3)
where the surface integral is evaluated over the (closed and 2-dimensional) boundary of $`V`$, which we denote by $`V`$, see Fig.2. The change per time interval of the charge contained in $`V`$ is balanced by the flow of the electric current through the boundary of the domain.
The closed surface integral $`_Vj`$ can be transformed into a volume integral $`_V𝑑j`$ by Stokes’s theorem (34). Here $`d`$ denotes the exterior derivative that increases the rank of a form by one, i.e. $`dj`$ is a 3-form. Thus (3) translates into
$$\underset{V}{}\left(_t\rho +dj\right)=0.$$
(4)
Since this is valid for an arbitrary domain, we arrive at the local form of electric charge conservation,
$$dj+_t\rho =0.$$
(5)
## III Excitations
Since the charge density $`\rho `$ is a 3-form, its exterior derivative vanishes: $`d\rho =0`$. Then, by a theorem of de Rham, it follows that $`\rho `$ can be derived from a “potential”:
$$d\rho =0\rho =d𝒟.$$
(6)
In this way one finds the electric excitation 2-form $`𝒟`$. Its absolute dimension is $`[𝒟]=[\rho ]=q`$, furthermore, for the components, $`[𝒟_{ab}]=[𝒟]/\mathrm{}^2=q/\mathrm{}^2`$.
Substituting (6)<sub>2</sub> into charge conservation (5) and once again using the de Rham theorem, we find another “potential” for the current density:
$$d\left(j+_t𝒟\right)=0j+_t𝒟=d.$$
(7)
The magnetic excitation $``$ turns out to be a 1-form, see Fig.4. Its dimension is $`[]=[j]=q/t,[_a]=q/(t\mathrm{})`$.
Consequently, the excitations $`(𝒟,)`$ are potentials of the sources $`(\rho ,j)`$. All these (additive) quantities (How much?) are described by odd differential forms.
In this way, we find the inhomogeneous Maxwell equations (the Gauss law and the Oersted-Ampère law):
$`d𝒟`$ $`=`$ $`\rho ,`$ (8)
$`d_t𝒟`$ $`=`$ $`j.`$ (9)
Electric charge conservation is valid in microphysics. Therefore the corresponding Maxwell equations (8) and (9) are valid on the same “microphysical” level as the notions of charge density $`\rho `$ and current density $`j`$. And with them the excitations $`𝒟`$ and $``$ are microphysical quantities of the same type likewise – in contrast to what is stated in most textbooks.
Before we ever talked about forces on charges, charge conservation alone gave us the inhomogeneous Maxwell equations including the appropriate dimensions for the excitations $`𝒟`$ and $``$.
Under the assumption that $`𝒟`$ vanishes inside an ideal electric conductor, one can get rid of the indeterminacy of $`𝒟`$ which is inherent in the definition of the excitation as a “potential” of the charge density, and we can measure $`𝒟`$ by means of two identical conducting plates (“Maxwellian double plates”) which touch each other and which are separated in the $`𝒟`$-field to be measured. The charge on one plate is then measured. Analogously, $``$ can be measured by the Gauss compensation method or by a superconductor and the Meissner effect ($`B=0=0`$). Accordingly, the excitations do have a direct operational significance.
## IV Field strengths
So far, conserved charge was the notion at center stage. Now energy enters the scene, which opens the way for introducing the electromagnetic field strengths. Whereas the excitations $`(𝒟,)`$ are linked to (and measured by) the charge and the current $`(\rho ,j)`$, the electric and magnetic field strengths are usually introduced as forces acting on unit charges at rest or in motion, respectively.
Let us consider a point particle with electric charge $`e`$ and velocity vector $`v`$. The force $`F`$ acting on it is a 1-form since its (1-dimensional) line integral yields a scalar, namely the energy. Thus $`F`$ carries the absolute dimension of an energy or of action over time: $`[F]=\mathrm{energy}=h/t`$, where $`h`$ denotes the dimension of an action. Accordingly, the local components $`F_a`$ of the force $`F=F_adx^a`$ possess the dimension $`[F_a]=h/(t\mathrm{})=`$ force.
In an electromagnetic field, the motion of a test particle is governed by the Lorentz force:
$$F=e(EvB)(\mathrm{Axiom}\mathrm{\hspace{0.17em}2}).$$
(10)
The symbol $``$ denotes the interior product of a vector (here the velocity vector) with a $`p`$-form. It decreases the rank of a form by 1 (see the Appendix), and since $`vB`$ is to be a 1-form, then $`B`$ is a 2-form. Newly introduced by (10) are the electric field strength 1-form $`E`$ and the magnetic field strength 2-form $`B`$. They are both intensities (How strong?).
The dimension of the velocity is $`[v]=1/t`$. With the decomposition $`v=v^a_a`$, we find for its components $`[v^a]=\mathrm{}/t`$. Then it is straightforward to read off from (10) the absolute dimension of the electric field strength $`[E]=h/(tq)=\varphi _0/t`$, with $`\varphi _0:=h/q`$. For its components we have $`[E_a]=\varphi _0/(t\mathrm{})`$. Analogously, for the dimension of the magnetic field strength we find $`[B]=(h/t)/(q/t)=h/q=\varphi _0`$ and $`[B_{ab}]=\varphi _0/\mathrm{}^2`$, respectively. The field $`B`$ carries the dimension of a magnetic flux $`\varphi _0`$. In superconductors, magnetic flux can come in quantized flux tubes, so-called fluxoids, underlining the importance of the notion of magnetic flux.
The definition (10) of the field strengths makes sense only if the charge of the test particle is conserved. In other words, axiom 2 presupposes axiom 1 and should not be seen as a stand alone pillar of electrodynamics.
## V Magnetic flux conservation
Taking into account the rank (as exterior forms) of the field strengths, the only integral we can build up from $`E`$ and $`B`$, respectively, are line integrals and surface integrals $`_{\mathrm{line}}E`$ and $`_{\mathrm{surface}}B`$. Apart from a factor $`t`$, the dimensions are equal. Hence, from a dimensional point of view, it seems sensible to postulate the conservation theorem (see Fig.5),
$$_t\underset{S}{}B+\underset{S}{}E=0(\mathrm{Axiom}\mathrm{\hspace{0.17em}3}).$$
(11)
Magnetic flux conservation (11) has to be seen as an analog of electric charge conservation (3). Magnetic flux, though, is related to a 2-dimensional surface whereas electric charge is related to a 3-dimensional volume. Thus the integration domains in the conservation theorems (11) and (3) differ by one spatial dimension always.
Axiom 3 gains immediate evidence from the dynamics of an Abrikosov flux line lattice in a superconductor. There the quantized flux lines can be counted, they don’t vanish nor are they created from nothing, rather they move in and out crossing the boundary $`S`$ of the surface $`S`$ under consideration.
Again, by means of Stokes’s theorem (34), we can transform the boundary integral: $`_SE=_S𝑑E`$. Taking into account the arbitrariness of the surface $`S`$, we recover Faraday’s induction law
$$dE+_tB=0,$$
(12)
which is experimentally very well established. We differentiate Faraday’s law by means of $`d`$ and find $`_t(dB)=0`$. Since an integration constant other than zero is senseless (recall the relativity principle), we have
$$dB=0.$$
(13)
The homogeneous Maxwell equations (12) and (13) (Faraday’s induction law and the closure of the magnetic field strength) nearly complete the construction of the theory.
We find the $`3+3`$ time evolution equations (9) and (12) for the electromagnetic field ($`𝒟,`$; $`E,B`$), i.e., for $`6+6`$ components. Before we can find solutions of these equations, we have to reduce the number of variables to 6, i.e., we have to cut them in half. Such a reduction is achieved by Axiom 4.
## VI The Maxwell-Lorentz spacetime relations
In Axiom 4 we assume linear, isotropic, and centrosymmetric relations between the (additive) quantities and the intensities :
$$𝒟=\epsilon _0^{}E\mathrm{and}=\frac{1}{\mu _0}^{}B(\mathrm{Axiom}\mathrm{\hspace{0.17em}4}).$$
(14)
The proportionality coefficients $`\epsilon _0,\mu _0`$ encode all the essential information about the electric and magnetic properties of spacetime. The Hodge star operator $``$ is needed, since we have to map a 1-form into a 2-form and vice versa or, more generally, a $`k`$-form into a $`(3k)`$-form. Then the operator $``$ in (14) has the dimension of a length or its reciprocal. Note that the Hodge star depends on the metric of our Euclidean space, see the Appendix. Recalling the dimensions of the excitations and the field strengths, we find the dimensions of the electric constant $`\epsilon _0`$ and the magnetic constant $`\mu _0`$ as
$$[\epsilon _0]=\frac{qt}{\varphi _0\mathrm{}}=\frac{1}{c\mathrm{\Omega }_0}\mathrm{and}[\mu _0]=\frac{\varphi _0t}{q\mathrm{}}=\frac{\mathrm{\Omega }_0}{c},$$
(15)
respectively. They are also called vacuum permittivity and vacuum permeability, see the new CODATA report . Here we define $`\mathrm{\Omega }_0:=\mathrm{\Phi }_0/q=h/q^2`$ and velocity $`c:=\mathrm{}/t`$. Dimensionwise, it is clearly visible that
$$\left[\frac{1}{\sqrt{\epsilon _0\mu _0}}\right]=c\mathrm{and}\left[\sqrt{\frac{\mu _0}{\epsilon _0}}\right]=\mathrm{\Omega }_0.$$
(16)
Obviously, the velocity $`c`$ and the resistance $`\mathrm{\Omega }_0`$ are constants of nature, the velocity of light $`c`$ being a universal one, whereas $`\mathrm{\Omega }_0`$, the characteristic impedance (or wave resistance) of the vacuum , seemingly refers only to electromagnetic properties of spacetime. Note that $`1/\mathrm{\Omega }_0`$ plays the role of the coupling constant of the electromagnetic field which enters as a factor into the free field Maxwell Lagrangian.
The Maxwell equations (8)-(9) and (12)-(13), together with the Maxwell-Lorentz spacetime (or aether) relations (14), constitute the foundations of classical electrodynamics. These laws, in the classical domain, are assumed to be of universal validity. Only if vacuum polarization effects of quantum electrodynamics are taken care of or hypothetical nonlocal terms should emerge from huge accelerations, Axiom 4 can pick up corrections yielding a nonlinear law (Heisenberg-Euler electrodynamics, see ) or a nonlocal law (Volterra-Mashhoon electrodynamics, see ), respectively. In this sense, the Maxwell equations are “more universal” than the Maxwell-Lorentz spacetime relations. The latter ones are not completely untouchable. We may consider (14) as constitutive relations for spacetime itself.
## VII SI-units
The fundamental dimensions in the SI-system for mechanics and electrodynamics are $`(\mathrm{},t,M,q/t)`$, with $`M`$ as mass. And for each of those a unit was defined. However, since action – we denote its dimension by $`h`$ – is a relativistic invariant quantity and since the electric charge is more fundamental than the electric current, we rather choose as the basic units
$$(\mathrm{},t,h,q),$$
(17)
see Schouten and Post . Thus, instead of the kilogram and the ampere, we choose joule$`\times `$second (or weber$`\times `$coulomb) and the coulomb:
$$(m,s,Wb\times C,C).$$
(18)
Numerically, in the SI-system, one puts (for historical reasons)
$$\mu _0=4\pi \times 10^7\frac{Wbs}{Cm}(\mathrm{magnetic}\mathrm{constant}).$$
(19)
Then measurements à la Weber-Kohlrausch yield
$$\epsilon _0=8.85\times 10^{12}\frac{Cs}{Wbm}(\mathrm{electric}\mathrm{constant}).$$
(20)
The SI-units of the electromagnetic field are collected in Table II.
Table II. SI-units of the electromagnetic field
| Field | SI-unit of field | SI-unit of components of field |
| --- | --- | --- |
| $`𝒟`$ | $`C`$ | $`C/m^2`$ |
| $``$ | $`A=C/s`$ | $`A/m=C/(sm)(\mathrm{oersted})`$ |
| $`E`$ | $`Wb/s=V`$ | $`V/m=Wb/(sm)`$ |
| $`B`$ | $`Wb`$ | $`Wb/m^2=\mathrm{tesla}=T(\mathrm{gauss})`$ |
## VIII Electrodynamics in matter
“It should be needless to remark that while from the mathematical standpoint a constitutive equation is a postulate or a definition, the first guide is physical experience, perhaps fortified by experimental data.” C. Truesdell and R.A. Toupin (1960)
In a great number of the texts on electrodynamics the electric and magnetic properties of media are described following the macroscopic averaging scheme of Lorentz (1916). However, this formalism has a number of serious limitations, see the relevant criticism of Hirst , e.g.. An appropriate modern presentation of this subject has been given in the textbook of Kovetz .
Here we follow a consistent microscopic approach to the electrodynamics in media, cf. . The total charge or current density is the sum of the two contributions originating “from the inside” and “from the outside” of the medium:
$$\rho =\rho ^{\mathrm{mat}}+\rho ^{\mathrm{ext}},j=j^{\mathrm{mat}}+j^{\mathrm{ext}}.$$
(21)
Hereafter, the bound charge in matter is denoted by mat and the external charge by ext. The same notational scheme will also be applied to the excitations $`𝒟`$ and $``$. Bound charge and bound current are inherent characteristics of matter determined by the medium itself. They vanish outside the medium. In contrast, external charge and external current in general do not vanish outside matter. They can be prepared for a specific purpose (such as the scattering of a current of particles by a medium or a study of the reaction of a medium in response to a prescribed configuration of charges and currents).
We assume that the charge bound by matter fulfills the usual conservation law:
$$dj^{\mathrm{mat}}+_t\rho ^{\mathrm{mat}}=0.$$
(22)
Taking into account (5), this means that there is no physical exchange (or conversion) between the bound and the external charges. As a consequence of this assumption, we can repeat the arguments of Sec.III that will give rise to the corresponding excitations $`𝒟^{\mathrm{mat}}`$ and $`^{\mathrm{mat}}`$ as “potentials” for the bound charge and the bound current. The conventional names for these newly introduced excitations are polarization $`P`$ and magnetization $`M`$, i.e.,
$$𝒟^{\mathrm{mat}}P,^{\mathrm{mat}}M.$$
(23)
The minus sign is conventional, see Kovetz . Thus, in analogy to the inhomogeneous Maxwell equations, we find
$$dP=\rho ^{\mathrm{mat}},dM+_tP=j^{\mathrm{mat}}.$$
(24)
The identifications (23) are only true up to an exact form. However, if we require $`𝒟^{\mathrm{mat}}=0`$ for $`E=0`$ and $`^{\mathrm{mat}}=0`$ for $`B=0`$, as we will do in (28), uniqueness is guaranteed.
The Maxwell equations are linear partial differential equations. Therefore we can define
$$𝒟^{\mathrm{ext}}:=𝒟𝒟^{\mathrm{mat}}=𝒟+P,^{\mathrm{ext}}:=^{\mathrm{mat}}=M.$$
(25)
The external excitations $`(𝒟^{\mathrm{ext}},^{\mathrm{ext}})`$ can be understood as auxiliary quantities. In terms of these quantities, the inhomogeneous Maxwell equations for matter finally read:
$`d𝒟^{\mathrm{ext}}`$ $`=`$ $`\rho ^{\mathrm{ext}},𝒟^{\mathrm{ext}}=\epsilon _0^{}E+P[E,B],`$ (26)
$`d^{\mathrm{ext}}_t𝒟^{\mathrm{ext}}`$ $`=`$ $`j^{\mathrm{ext}},^{\mathrm{ext}}={\displaystyle \frac{1}{\mu _0}}^{}BM[B,E].`$ (27)
Here the polarization $`P[E,B]`$ is a functional of the electromagnetic field strengths $`E`$ and $`B`$. In general, it can depend also on the temperature $`T`$, and possibly of other thermodynamic variables specifying the material continuum under consideration; similar remarks apply to the magnetization $`M[B,E]`$. The system (26)<sub>1</sub> and (27)<sub>1</sub> looks similar to (8) and (9). However, the former equations refer only to the external fields and sources. The homogeneous Maxwell equations (12) and (13) remain valid in their original form.
In the simplest cases, we have the linear constitutive laws
$$P=\epsilon _0\chi _\mathrm{E}^{}E,M=\frac{1}{\mu _0}\chi _\mathrm{B}^{}B,$$
(28)
with the electric and magnetic susceptibilities $`(\chi _\mathrm{E},\chi _\mathrm{B})`$. With material constants
$$\epsilon :=\epsilon _0(1+\chi _\mathrm{E}),\mu :=\frac{\mu _0}{1\chi _\mathrm{B}},$$
(29)
one can rewrite the material laws (28) as
$$D^{\mathrm{ext}}=\epsilon ^{}E,H^{\mathrm{ext}}=\frac{1}{\mu }^{}B.$$
(30)
For the discussion of the concrete applications of the developed microscopic theory in modern condensed matter physics, we refer to the review of Hirst .
## IX Conclusion
The Maxwell equations
$`d𝒟`$ $`=`$ $`\rho ,d_t𝒟=j,`$ (31)
$`dB`$ $`=`$ $`0,dE+_tB=0,`$ (32)
are the cornerstones of any classical theory of electromagnetism. As an expression of charge and flux conservation, they carry a high degree of plausibility as well as solid experimental support. The Maxwell equations in this form remain valid in an accelerated reference frame and in a gravitational field likewise, without any change.
The Maxwell-Lorentz spacetime relations
$$𝒟=\frac{1}{c\mathrm{\Omega }_0}^{}E\mathrm{and}=\frac{c}{\mathrm{\Omega }_0}^{}B$$
(33)
are necessary for developing the Maxwellian system into a predictive physical theory. They depend, via the star operator, on the metric of space and are, accordingly, influenced by the gravitational field. They are valid in very good approximation, but there are a few exceptions known (if the Casimir effect is to be described, e.g.).
For the description of matter, the sources $`(\rho ,j)`$ and the excitations $`(𝒟,)`$ have to be split suitably in order to derive, from the microscopic equations (31) and (32), appropriate macroscopic expressions.
Summing up, we can give an answer to one of the central questions posed by Roche : The need for the different notations and different dimensions and units for the excitation $``$ and the field strength $`B`$ (and, similarly, for $`𝒟`$ and $`E`$) is well substantiated by the very different geometrical properties and physical origins of these fields, see Table I and Fig. 1. Even in vacuum, these differences do not disappear.
## Acknowledgments
We are grateful to H. Micklitz (Cologne) for arranging this lecture. One of the authors (FWH) would like to thank W. Raith (Berlin-Bielefeld) for an extended exchange of letters on the fundamental structure of Maxwell’s theory. Moreover, he is grateful to R.G. Chambers (Bristol), A. Kovetz (Tel Aviv), J. Roche (Oxford), and S. Scheidl (Cologne) for most useful remarks.
## Appendix: The ABC of exterior calculus
The formalism of exterior differential forms is widely used in different domains of mathematics and theoretical physics. In particular, in electromagnetic theory, exterior calculus offers an elegant and transparent framework for the introduction of the basic notions and for the formulation of the corresponding laws. Here, we will give a very elementary description of the objects and operations used above.
We will confine ourselves only to the case of a 3-dimensional space. Let be given the set of local coordinates $`x^a=\{x^1,x^2,x^3\}`$; hereafter Latin indices $`a,b,\mathrm{}`$ will run over $`1,2,3`$. Then the vectors $`e_a=\{_1,_2,_3\}`$ will serve as a basis of the tangent vector space at every point. The symbol $`dx^a`$ denotes the set of linear 1-forms dual to the coordinate vector basis, $`dx^a(e_b)=\delta _b^a`$. An arbitrary $`k`$-form can be described, in local coordinates, by its components: $`\phi =\phi _adx^a=\phi _1dx^1+\phi _2dx^2+\phi _3dx^3`$ for 1-forms and $`\omega =\frac{1}{2}\omega _{ab}dx^adx^b=\omega _{12}dx^1dx^2+\omega _{23}dx^2dx^3+\omega _{31}dx^3dx^1`$ for 2-forms. Any 3-form has a single nontrivial component, $`\eta =\eta _{123}dx^1dx^2dx^3`$. When $`\eta `$ is smoothly defined on the whole space, it is called a volume form. Zero-forms are just the ordinary differentiable functions.
It is often stated that the exterior product$``$” generalizes the vector product. However, one should be careful with such statements, because the vector product in the standard 3-dimensional analysis is, strictly speaking, a superposition of the wedge product and of the Hodge duality operator. Thus, the vector product necessarily involves the metric on the manifold. In contrast, the exterior product is a pre-metric operation, although it resembles the vector product. For example, the exterior product of the two 1-forms $`\omega `$ and $`\phi `$ with the components $`\omega _a`$ and $`\phi _a`$ yields a 2-form $`\omega \phi `$ with the local components $`\{(\omega _2\phi _3\omega _3\phi _2),(\omega _3\phi _1\omega _1\phi _3),(\omega _1\phi _2\omega _2\phi _1)\}`$.
The exterior differential $`d`$ increases the rank of a form by 1. It is most easily described in local coordinates, see Table III. Thus, $`d`$ naturally generalizes the “grad” operator which leads from a scalar to a vector and, at the same time, it represents a pre-metric extension of the “curl” operator. The exterior differential is a nilpotent operator, i.e., $`dd=0`$.
Table III. Operators acting on an exterior form
| | $`k`$-form $`\omega =\frac{1}{k!}\omega _{a_1\mathrm{}a_k}dx^{a_1}\mathrm{}dx^{a_k}`$, with $`k=0,1,2,3`$ |
| --- | --- |
| $`d`$ | $`(k+1)`$-form $`d\omega =\frac{1}{(k+1)!}\left(_{[a_1}\omega _{a_1\mathrm{}a_{k+1}]}\right)dx^{a_1}\mathrm{}dx^{a_{k+1}}`$ |
| $``$ | $`(k1)`$-form $`v\omega =\frac{1}{(k1)!}v^a\omega _{aa_1\mathrm{}a_{k1}}dx^{a_1}\mathrm{}dx^{a_{k1}}`$ |
| | $`(3k)`$-form $`{}_{}{}^{}\omega =\frac{1}{k!}\omega _{a_1\mathrm{}a_k}g^{a_1b_1}\mathrm{}g^{a_kb_k}e_{b_k}\mathrm{}e_{b_1}\eta `$ |
Complementary to $`d`$, one can define an operation which decreases the rank of a form by 1. This is the interior product of a vector with a $`k`$-form. Given the vector $`v`$ with the components $`v^a`$, the interior product with the coframe 1-form yields $`vdx^a=v^a`$, which is a sort of a projection along $`v`$. By linearity, the interior product of $`v`$ with a $`k`$-form is defined as described in Table III.
The Hodge dual operator maps $`k`$-forms into $`(3k)`$-forms. Its introduction necessarily requires the metric which assigns a real number $`g(u,v)=g(v,u)`$ to every two vectors $`u`$ and $`v`$. In local coordinates, the components of the metric tensor are determined as the values of the scalar product of the basis vectors, $`g_{ab}:=g(e_a,e_b)`$. This matrix is positive definite. The metric introduces a natural volume 3-form $`\eta :=\sqrt{detg_{ab}}dx^1dx^2dx^3`$ which underlies the definition of the Hodge operator . The general expression is displayed in Table III. Explicitly the Hodge dual of the coframe 1-form reads, for example: $`{}_{}{}^{}dx^a=\sqrt{detg_{ab}}(g^{a1}dx^2dx^3+g^{a2}dx^3dx^1+g^{a3}dx^1dx^2)`$, where $`g^{ab}`$ is inverse to $`g_{ab}`$.
The notions of the odd and even exterior forms are closely related to the orientation of the manifold. In simple terms, these two types of forms are distinguished by their different behavior with respect to a reflection (i.e., a change of orientation): an even (odd) form does not change (changes) sign under a reflection transformation. These properties of odd and even forms are crucial in the integration theory, see, e.g., .
For a $`k`$-form an integral over a $`k`$-dimensional subspace is defined. For example, a $`1`$-form can be integrated over a curve, a $`2`$-form over a 2-surface, and a volume 3-form over the whole 3-dimensional space. We will not enter into the details here, limiting ourselves to the formulation of Stokes’s theorem which occupies a central place in integration theory:
$$\underset{C}{}\omega =\underset{C}{}𝑑\omega .$$
(34)
Here $`\omega `$ is an arbitrary $`k`$-form, and $`C`$ is an arbitrary $`(k+1)`$-dimensional (hyper)surface with the boundary $`C`$.
For a deeper and a more rigorous introduction into exterior calculus, see, e.g., .
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# Untitled Document
hep-th/0005129 CALT-68-2272 CITUSC/00-023
Space-Time Noncommutative Field Theories And Unitarity
Jaume Gomis and Thomas Mehen
Department of Physics
California Institute of Technology
Pasadena, CA 91125
and
Caltech-USC Center for Theoretical Physics
University of Southern California
Los Angeles, CA 90089
gomis, mehen@theory.caltech.edu
We study the perturbative unitarity of noncommutative scalar field theories. Field theories with space-time noncommutativity do not have a unitary S-matrix. Field theories with only space noncommutativity are perturbatively unitary. This can be understood from string theory, since space noncommutative field theories describe a low energy limit of string theory in a background magnetic field. On the other hand, there is no regime in which space-time noncommutative field theory is an appropriate description of string theory. Whenever space-time noncommutative field theory becomes relevant massive open string states cannot be neglected.
May 2000
1. Introduction
Noncommutative field theories are constructed from conventional (commutative) field theories by replacing in the Lagrangian the usual multiplication of fields with the $``$-product of fields. The $``$-product is defined in terms of a real antisymmetric matrix $`\theta ^{\mu \nu }`$ that parameterizes the noncommutativity of Minkowski space-time<sup>1</sup> Throughout the paper we will use the $`(+,,\mathrm{},)`$ convention for the signature of space-time.
$$[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }\mu ,\nu =0,\mathrm{},D1.$$
The $``$-product of two fields $`\varphi _1(x)`$ and $`\varphi _2(x)`$ is given by
$$\varphi _1(x)\varphi _2(x)=e^{\frac{i}{2}\theta ^{\mu \nu }\frac{}{\alpha ^\mu }\frac{}{\beta ^\nu }}\varphi _1(x+\alpha )\varphi _2(x+\beta )|_{\alpha =\beta =0}.$$
The noncommutativity in (1.1) gives rise to a space-time uncertainty relation
$$\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu \frac{1}{2}|\theta ^{\mu \nu }|$$
which leads to number of unusual phenomena such as the mixing of the ultraviolet with the infrared as well as apparent acausal behavior -.
These field theories are non-local and this nonlocality has important consequences for the dynamics -. The structure of the product in $`(1.1)`$ leads to terms in the action with an infinite number of derivatives of fields which casts some doubts on the unitarity of noncommutative field theories. In this paper we will check the unitarity of scalar noncommutative field theories at the one loop level and show that theories with $`\theta ^{0i}=0`$ are unitary while theories with $`\theta ^{0i}0`$ are not unitary.
Noncommutative field theories with space noncommutativity (that is $`\theta ^{0i}=0`$) have an elegant embedding in string theory . They describe the low energy excitations of a D-brane in the presence of a background magnetic field<sup>2</sup> We will sometimes refer to these theories as magnetic theories.. In this limit the relevant description of the dynamics is in terms of the noncommutative field theory of the massless open strings. Both the massive open strings and the closed strings decouple<sup>3</sup> In - one-loop string theory amplitudes were shown to exhibit this decoupling. See also -.. The consistent truncation of the full unitary string theory to field theory with space noncommutativity leads one to suspect that these field theories are unitary. Moreover, these field theories are nonlocal in space but are local in time. Therefore, a Hamiltonian can be constructed and it gives rise to unitary time evolution of noncommutative magnetic field theories.
Theories with space-time noncommutativity<sup>4</sup> Likewise, we will sometimes refer to these theories as electric theories. (that is $`\theta ^{0i}0`$) have an infinite number of time derivatives of fields in the Lagrangian and are nonlocal in time. The commutator in (1.1) leads to noncommutativity of the time coordinate. Noncommutativity of the time coordinate and the corresponding nonlocality in time results in theories where it is far from clear whether the usual framework of quantum mechanics makes sense. As such, noncommutative field theories with space-time noncommutativity are excellent laboratories in which to test the possible breakdown of the conventional notion of time or the familiar framework of quantum mechanics in string theory at the Planck scale<sup>5</sup> In recent years, several examples of the breakdown of the conventional notion of space at very short distances have been found in string theory such as in topology changing transitions . It seems, therefore, imperative to address the issue of possible breakdown of time.. In this paper we will test in these exotic field theories one of the basic principles of quantum mechanics, the existence of a unitary S-matrix. We explicitly show that several one loop amplitudes in noncommutative scalar electric field theory are not unitary which demonstrates that noncommutative field theories with space-time noncommutativity clash with quantum mechanics.
This field theory result meshes very nicely with string theory expectations. $`\theta ^{0i}0`$ is obtained by studying string theory in the presence of a background electric field (recent work in this direction has recently appeared in -, see also ). There are three important parameters that characterize the open strings : $`\alpha ^{}`$, the metric $`G_{\mu \nu }`$ and the noncommutativity matrix $`\theta ^{\mu \nu }`$. One must also keep in mind that there is an upper critical value on the magnitude of the background electric field $`E_c`$ -, beyond this value the string vacuum becomes unstable. It can be shown using the relation between these open string parameters with the closed string metric and background electric field that it is impossible to take a consistent limit of string theory in which $`\theta ^{\mu \nu }`$ and $`G_{\mu \nu }`$ are kept fixed while $`\alpha ^{}0`$. Therefore, unlike the case of strings in a background magnetic field, it is impossible to find a limit of string theory in which one is left only with a noncommutative field theory with fixed background metric $`G_{\mu \nu }`$ and space-time noncommutativity parameter $`\theta ^{0i}`$. It is possible to find a limit of string theory with nonvanishing $`\theta ^{0i}`$ in which the closed strings decouple. However, $`\theta ^{0i}\alpha ^{}`$ making it impossible to decouple massive open string states and keep $`\theta ^{0i}`$ finite. Thus, there is no sense in which the electric field theories give an approximate description of a limit of string theory. The lack of decoupling of the massless open string modes from the massive ones gives very strong indication that the noncommutative field theory truncation to the massless modes is not unitary. This is consistent with what we find from our field theory analysis.
The paper is organized as follows. In section $`2`$ we compute several one loop amplitudes in noncommutative scalar field theory and show that Feynman diagrams of space-time noncommutative theories do not satisfy the usual cutting rules and the S-matrix does not satisfy unitarity constraints. We also show that one loop amplitudes are unitary in the presence of only space noncommutativity and the Feynman diagrams satisfy the cutting rules. We conclude in section $`3`$ with a discussion of our results and their relation to limits of string theory in electromagnetic field backgrounds.
2. Unitarity of Noncommutative Scalar Field Theory
In this section we examine one loop diagrams of noncommutative $`\varphi ^3`$ and $`\varphi ^4`$ theories to see if they satisfy constraints from unitarity. For on-shell matrix elements unitarity implies that
$$2\mathrm{Im}M_{ab}=\underset{n}{}M_{an}M_{nb}$$
where $`M_{ab}`$ is the transition matrix element between states $`a`$ and $`b`$. The sum over intermediate states on the right hand side includes phase space integrations for each particle in $`n`$. Quantum field theories actually satisfy more restrictive relations called generalized unitarity relations or cutting rules. These state that the imaginary part of a Feynman diagram can be obtained by the following procedure: First, “cut” the diagram by drawing a line through virtual lines such that the graph is severed in two. Next, wherever the cut intersects a virtual line, place that virtual particle on-shell by replacing the propagator with a delta function:
$$\frac{1}{p^2m^2+iϵ}2\pi i\delta (p^2m^2).$$
Summing over all cuts yields the imaginary part of the Feynman diagram. Cutting rules are a generalization of (2.1) to Feynman diagrams. Unitarity of the S-matrix (2.1) follows from the cutting rules<sup>6</sup> This assumes of course that the poles of the propagators correspond to physical states. In gauge theories unphysical states can propagate in loops and one must demonstrate that these states decouple from the physical S-matrix. This will not be a concern for the scalar theories considered in this paper.. Note that the cutting rules are more restrictive than the constraint of unitarity since they apply to off-shell Green’s functions as well as S-matrix elements.
We will first show that the two-point function of the noncommutative $`\varphi ^3`$ theory does not obey the usual cutting rules when there is space-time noncommutativity ($`\theta ^{0i}0)`$. In the case of space noncommutativity $`(\theta ^{0i}=0,\theta ^{ij}0)`$ the cutting rules are satisfied. Next, we consider $`22`$ scattering in noncommutative $`\varphi ^4`$ theory. The S-matrix is nonunitary at one-loop if $`\theta ^{0i}0`$, but is unitary if the noncommutativity is only in the spatial directions.
It is somewhat surprising that Feynman diagrams of space-time noncommutative theories do not obey the usual cutting rules. Since the Feynman rules for the vertices of noncommutative theories are manifestly real functions of momenta, one would expect that Feynman graphs could only develop a branch cut when internal lines go on-shell. This would imply that the imaginary parts of Feynman diagrams would be given by the same cutting rules as ordinary commutative field theories. The resolution of this puzzle requires an examination of the high energy behavior of the oscillatory factors that typically arise in these theories. We will find that a necessary condition for one-loop Feynman integrals to converge in these theories is that the following inner product
$$pp=p^\mu \theta _{\mu \nu }^2p^\nu ,$$
be positive definite. This inner product is positive definite when $`p_\mu `$ and $`\theta _{\mu \nu }`$ are analytically continued to Euclidean space. In Minkowski space $`pp`$ can be negative if $`\theta ^{0i}0`$. Feynman integrals can be defined via analytic continuation, but the resulting amplitudes will develop branch cuts when $`pp<0`$. These additional branch cuts are responsible for the failure of cutting rules and unitarity in noncommutative theories with space-time noncommutativity.
For $`pp=0`$, the S-matrix does not suffer from lack of unitarity, but is ill-defined because of infrared divergences. $`pp=0`$ is possible whether the noncommutativity is space-time or space-space. Obviously an outstanding problem in noncommutative field theory is to construct the infrared safe observables of the theory. This may require all order resummation of infrared divergent terms in the perturbative series. We will not atttempt to address this issue in this paper, and focus only on perturbative unitarity constraints for matrix elements which do not suffer from infrared singularities.
2.1. Noncommutative $`\varphi ^3`$ two-point function
Fig. 1: Generalized unitarity relation for $`\varphi ^3`$ two-point function.
The cutting rule for the noncommutative $`\varphi ^3`$ theory two-point function at lowest order is displayed in fig. 1. The propagators of fields in noncommutative field theories are identical to those of commutative field theory. The Feynman rule for the vertex in this theory is
$$i\lambda \mathrm{cos}\left(\frac{kq}{2}\right),kq=k_\mu \theta ^{\mu \nu }q_\nu .$$
where $`k`$ and $`q`$ are any two of the momenta flowing into the vertex. Because of conservation of momentum and the antisymmetry of $`\theta _{\mu \nu }`$ it does not matter which two momenta are chosen. The amplitude for the one loop diagram appearing in fig. 1 is:
$$\begin{array}{cc}\hfill iM=\frac{\lambda ^2}{2}\frac{d^Dl}{(2\pi )^D}\frac{1+\mathrm{cos}(pl)}{2}\frac{1}{l^2m^2+iϵ}\frac{1}{(l+p)^2m^2+iϵ},& \end{array}$$
while the expression for the right hand side of fig. 1 is
$$|M|^2=\frac{\lambda ^2}{2}\frac{1}{(2\pi )^{D2}}\frac{d^{D1}k}{2k_0}\frac{d^{D1}q}{2q_0}\delta ^D(pkq)\frac{1+\mathrm{cos}(pk)}{2}.$$
The mass of the $`\varphi `$ quanta is $`m`$ and $`p`$ denotes the external momentum which is not required to be on-shell. In both (2.1) and (2.1) we have used the identity $`\mathrm{cos}^2x=(1+\mathrm{cos}(2x))/2`$ to separate the planar and nonplanar contributions -. We will focus on the nonplanar terms since it is obvious that the planar parts satisfy unitarity constraints.
First we compute the one loop graph. We combine denominators using Feynman parameters then represent propagators via Schwinger parameters to obtain
$$M=\frac{\lambda ^2}{8}\frac{d^Dl_E}{(2\pi )^D}_0^1dx_0^{\mathrm{}}d\alpha \alpha (\mathrm{exp}(\alpha (l_E^2+x(1x)p_E^2+m^2iϵ)+il_Ep_E)+c.c.).$$
We have performed the usual analytic continuation $`l^0=il_E^0,p^0=ip_E^0`$. The subscript $`E`$ denotes Euclidean momenta. In addition, if there is space-time noncommutativity we must analytically continue $`\theta ^{0i}i\theta ^{0i}`$. In the string theory realization this can be easily undertood since $`\theta ^{0i}`$ is related to a background electric field. This continuation leaves the Moyal phase invariant. Otherwise the phases appearing in (2.1) , $`\mathrm{exp}(\pm il_Ep_E)`$, which render the integral finite, become $`\mathrm{exp}(\pm l_Ep_E)`$ and the integral is no longer convergent. Integrating over the loop momentum $`l_E`$ gives
$$M=\frac{\lambda ^2}{4}\frac{1}{(4\pi )^{D/2}}_0^1𝑑x_0^{\mathrm{}}𝑑\alpha \alpha ^{1D/2}\mathrm{exp}\left(\alpha (x(1x)p_E^2+m^2iϵ)\frac{p_Ep_E}{4\alpha }\right).$$
We will now evaluate this integral for $`D=3`$ and $`D=4`$ space-time dimensions and analytically continue back the answer to Minkowski space. The amplitudes are given by
$$M_{D=3}=\frac{\lambda ^2}{32\pi }_0^1𝑑x\frac{\mathrm{exp}(\sqrt{pp(m^2p^2x(1x)iϵ})}{\sqrt{m^2p^2x(1x)iϵ})},$$
and
$$M_{D=4}=\frac{\lambda ^2}{32\pi ^2}_0^1𝑑xK_0(\sqrt{pp(m^2p^2x(1x)iϵ)}),$$
where $`K_0`$ is a modified Bessel function. A crucial point to note is that the $`\alpha `$ integral is convergent only if $`p_Ep_E>0`$. For Euclidean momenta this is always true<sup>7</sup> We will stay away from the region where $`pp=0`$ where infrared singularities appear.. On the other hand, $`pp`$ need not be positive definite in Minkowski space when space-time is noncommuting. Specifically, let us choose $`\theta ^{01}=\theta ^{10}=\mathrm{\Theta }_E,\theta ^{23}=\theta ^{32}=\mathrm{\Theta }_B`$ with all other components vanishing. Then
$$pp=\mathrm{\Theta }_E^2(p_0^2p_1^2)+\mathrm{\Theta }_B^2(p_2^2+p_3^2).$$
Therefore, in the case of only space noncommutativity $`pp`$ is positive definite but for space-time noncommutativity $`pp`$ can be negative. This fact has very important consequences in the unitarity analysis.
We will now proceed to verify that the generalized unitarity relation (2.1) is satisfied for magnetic theories and violated for electric field theories. First we compute the imaginary part of the Feynman diagram when $`p^2>0`$ and $`pp>0`$. It is then easy to show that
$$\begin{array}{cc}\hfill \mathrm{Im}M_{D=3}=\frac{\lambda ^2}{32\pi }_{(1\gamma )/2}^{(1+\gamma )/2}𝑑x\frac{\mathrm{cos}(\sqrt{pp}\sqrt{m^2+p^2x(1x)})}{\sqrt{m^2+p^2x(1x)}}& \\ \hfill =\frac{\lambda ^2}{32\sqrt{p^2}}J_0\left(\frac{\gamma \sqrt{p^2pp}}{2}\right)& \end{array},$$
where $`\gamma =\sqrt{14m^2/p^2}`$.
Using the fact that $`\text{Im}K_0(ix)=\frac{\pi }{2}J_0(x)`$, where $`J_0`$ is a Bessel function, one obtains for $`D=4`$ space-time dimensions
$$\begin{array}{cc}\hfill \mathrm{Im}M_{D=4}=\frac{\lambda ^2}{64\pi }_{(1\gamma )/2}^{(1+\gamma )/2}𝑑xJ_0(\sqrt{pp}\sqrt{m^2+p^2x(1x)})& \\ \hfill =\frac{\lambda ^2}{32\pi }\frac{\mathrm{sin}(\gamma \sqrt{p^2pp}/2)}{\sqrt{p^2pp}}& \end{array}.$$
We will now evaluate the sum over final states (2.1). The integrals evaluate to
$$|M_{D=3}|^2=\frac{\lambda ^2}{4}\frac{1}{8\pi \sqrt{p^2}}_0^{2\pi }𝑑\theta \mathrm{cos}(pk)=\frac{\lambda ^2}{16\sqrt{p^2}}J_0\left(\frac{\gamma \sqrt{p^2pp}}{2}\right)$$
and for $`D=4`$ space-time (2.1) gives
$$|M_{D=4}|^2=\frac{\lambda ^2}{4}\frac{\gamma }{32\pi ^2}𝑑\mathrm{\Omega }\mathrm{cos}(pk)=\frac{\lambda ^2}{16\pi }\frac{\mathrm{sin}(\gamma \sqrt{p^2pp}/2)}{\sqrt{p^2pp}}.$$
We see that for $`pp>0`$ the generalized unitarity relation (2.1) is satisfied.
We will now consider the case $`pp<0`$. From (2.1) it follows that this configuration of momenta can only exist in the presence of space-time noncommutativity. Moreover $`p^2`$ must be negative so it corresponds to space-like momentum. Then
$$\mathrm{Im}M_{D=3}=\frac{\lambda ^2}{32\pi }_0^1𝑑x\frac{\mathrm{sin}(\sqrt{|pp|(m^2+|p^2|x(1x)})}{\sqrt{m^2+|p^2|x(1x)})}$$
and
$$\mathrm{Im}M_{D=4}=\frac{\lambda ^2}{64\pi }_0^1𝑑xJ_0(\sqrt{|pp|(m^2+|p^2|x(1x))}).$$
which are obviously nonzero. However, the right hand side of the equation in fig. 1 is zero because energy-momentum conservation (2.1) forbids a particle with space-like momenta to decay into two massive on-shell particles. Therefore, when $`pp<0`$, the generalized unitarity relation (2.1) is violated.
Summarizing, we have shown that field theories with space-time noncomutativity violate the equation in fig. 1 and that field theories with space noncommutativity satisfy it for arbitrary momenta.
2.2. Noncommutative $`\varphi ^4`$ Scattering Amplitude
Next we consider the $`22`$ scattering amplitude in noncommutative $`\varphi ^4`$ theory. The Feynman rule for the 4-point vertex in this theory is
$$i\frac{\lambda }{3}\left(\mathrm{cos}\left(\frac{p_1p_2}{2}\right)\mathrm{cos}\left(\frac{p_3p_4}{2}\right)+\mathrm{cos}\left(\frac{p_1p_3}{2}\right)\mathrm{cos}\left(\frac{p_2p_4}{2}\right)+\mathrm{cos}\left(\frac{p_1p_4}{2}\right)\mathrm{cos}\left(\frac{p_2p_3}{2}\right)\right)$$
where the $`p_i`$ are momenta entering the vertex.
Fig. 2: One loop diagrams for $`22`$ scattering
The one loop contribution to the $`22`$ scattering amplitude are shown in fig. 2. Evaluating these graphs leads to rather complicated expressions which involve integrals over modified Bessel functions, but these simplify greatly if we expand the expressions in powers of $`\theta ^{\mu \nu }`$. The optical theorem (2.1) for this S-matrix element has to be true term by term in a power series in $`\theta ^{\mu \nu }`$. The leading contribution in $`\theta `$ to the right hand side of (2.1) is the same as commutative $`\varphi ^4`$ theory, so
$$\underset{n}{}M_{p_1+p_2n}M_{np_3+p_4}^{}=\frac{\gamma \lambda ^2}{16\pi }\mathrm{\Theta }(s4m^2)$$
where $`\mathrm{\Theta }(x)`$ is a step function, $`\gamma =\sqrt{14m^2/s}`$ and $`s=(p_1+p_2)^2=(p_3+p_4)^2`$.
The leading contribution to the one loop scattering amplitude is
$$\begin{array}{cc}& M_{p_1+p_2p_3+p_4}=\hfill \\ & \frac{\lambda ^2}{2(4\pi )^2}(_0^1dx[\mathrm{ln}(1\frac{s}{m^2}x(1x))+(st,su)]+\frac{2}{3}\mathrm{ln}\left(\frac{m^2}{\mu ^2}\right)+\mathrm{const}.\hfill \\ & +\frac{1}{3}(\underset{i=1}{\overset{4}{}}\mathrm{ln}(m^2p_ip_i)+\frac{1}{3}\mathrm{ln}(m^2p_{12}p_{12})+\mathrm{ln}(m^2p_{13}p_{13})+\mathrm{ln}(m^2p_{14}p_{14})))\hfill \end{array}.$$
Here we have defined $`t=(p_1p_3)^2,u=(p_1p_4)^2,p_{12}=p_1+p_2,p_{13}=p_1p_3`$, and $`p_{14}=p_1p_4`$. The first line in (2.1) includes the contributions present in ordinary $`\varphi ^4`$ theory, while the second includes logarithms that are unique to the noncommutative theory. The first logarithm has an imaginary piece when $`s>4m^2`$, corresponding to threshold production of two $`\varphi `$ particles, which gives the precise contribution so that (2.1) is satisfied to leading order in $`\theta ^{\mu \nu }`$ with the right hand side of (2.1) given by formula (2.1). The noncommutative logarithms $`\mathrm{ln}(m^2p_ip_i)`$ and $`\mathrm{ln}(m^2p_{12}p_{12})`$ do not contribute an imaginary piece because $`p_i,p_{12}`$ are time-like, and hence $`p_ip_i,p_{12}p_{12}`$ are positive definite. However, $`p_{13}`$ and $`p_{14}`$ are space-like and therefore $`p_{13}p_{13}`$ and $`p_{14}p_{14}`$ can be negative if there is space-time noncommutativity. In this case these logarithms have imaginary parts which violate the unitarity relation (2.1) . Therefore, theories with space-time noncommutativity do not have a unitary S-matrix. Moreover, since $`pp>`$ is always positive for space noncommutative theories there are no new imaginary parts and the optical theorem is satisfied.
3. Discussion
In this paper we have investigated the unitarity of noncommutative scalar field theories. The results we have obtained have a natural interpretation in string theory. We have shown that field theories with space noncommutativity appear to have perturbatively unitarity S-matrix elements and satisfy the generalized unitarity relations of field theory Green’s functions. On the other hand, theories with space-time noncommutativity do not have a unitary S-matrix and do not satisfy the cutting rules for Feynman diagrams.
We have done calculations for noncommutative scalar field theories. Even though we have not checked the unitarity of noncommutative gauge theories we have strong reasons to believe that the same results still hold, that is the magnetic theories are unitary while the electric theories are not. This is because the structure of Feynman integrals is the same in gauge and scalar theories. Both have oscillating phases in loop integrations. After analytically continuing momenta and $`\theta ^{0i}`$ to Euclidean spacetime, then performing loop integrals one encounters integrals of the form: <sup>8</sup> We are not including integration over Feynman parameters which are not needed for the argument.
$$A𝑑\alpha \alpha ^{1D/2}\mathrm{exp}^{\frac{pp}{\alpha }},$$
where $`p`$ denotes some external momentum. In the Euclidean theory $`pp0`$ so that $`1/pp`$ regulates (3.1) and acts like an ultraviolet cutoff which renders the integral finite. In the theory with only space noncommutativity, the Minkowski expression for $`pp>0`$ is never negative. The only possible singularity arises when $`\theta ^{\mu \nu }p_\nu =0`$ which leads to characteristic infrared singularities of noncommutative field theories -. On the other hand when $`\theta ^{0i}0`$, the Minkowski expression for $`pp`$ can be positive or negative so that when one analytically continues the Euclidean answer, one finds branch cuts in the Feynman diagrams for Minkowski $`pp<0`$. It is the presence of these extra branch cuts in the loop diagrams of field theories with space-time noncommutativity that are responsible for the failure of the cutting rules and lack of a unitary S-matrix.
The fact that the magnetic gauge theories are unitary can be easily understood from string theory. They provide the appropiate effective description of a low energy limit of string theory in the presence of a background magnetic field . In this limit, all the massive open string states and the closed strings decouple and the relevant degrees of freedom for the description of the dynamics are the massless open strings. One can build up the effective action for these modes from string theory and show that they are described by noncommutative field theory. Therefore, we expect the field theory to be unitary since string theory in this limit can be appropriately described in terms of noncommutative field theory, without the need of adding any further degrees of freedom. This is indeed what we have found from our field theory analysis.
Field theories with space-time noncommutativity should appear from studying string theory in the presence of a background electric field. This follows from constructing the effective action of open strings in this vacuum. One might expect, based on analogy with a background magnetic field, that there is a similar limit of the string dynamics which is described just by the electric field theory. However, electric fields behave differently than magnetic fields in that they lead to pair production of strings and these destabilize the vacuum if the background electric field exceeds the upper critical value $`E_c`$ . Consider for simplicity the electric field to be aligned in the $`x^1`$ direction and the metric to be diagonal in the $`x^0,x^1`$ plane with each metric component given by $`g`$. Reality of the brane action requires that the background electric field on the brane satisfy
$$EE_c,\text{where}E_c=\frac{g}{2\pi \alpha ^{}}.$$
The open strings see a diagonal metric along the $`x^0,x^1`$ plane given by $`G`$ and a noncommutative parameter $`\theta ^{0i}=\theta `$. In terms of the metric $`g`$ and the background electric field, these parameters are related by the following formula
$$\alpha ^{}G^1=\frac{1}{2\pi }\frac{E}{E_c}\theta .$$
In order to obtain a field theory of only the massless modes one has to go to the point particle limit $`\alpha ^{}0`$. From formula (3.1) this implies that, for finite $`G`$, that the noncommativity parameter must vanish. Therefore, if one wants a truncation of the full string theory to the theory of only the massless open string modes this can be done but the description of these modes is given by conventional field theory and not noncommutative field theory. Thus, we expect the conventional field theory description to be unitary and it is. Clearly, in order to have a finite noncommutativity parameter $`\theta `$, $`\alpha ^{}`$ must be kept finite. This is a string theory and not a field theory. Moreover, since $`\theta \alpha ^{}`$ there is no scattering process in this string theory which is accurately described only by noncommutative field theory. For scattering processes involving massless open strings of characteristic energy $`E\theta ^{1/2}`$ conventional field theory is a proper description. Noncommutative field theory becomes relevant for energies of the order of $`E\theta ^{1/2}`$. However, since $`\theta \alpha ^{}`$, the energy scale where noncommutativity becomes relevant is precisely the energy scale at which the massive open string states cannot be neglected. Thus, there is no regime in which space-time noncommutative field theory is an appropiate description of string theory. Whenever space-time noncommutative field theory becomes relevant massive open string states cannot be neglected. This gives a very strong indication that noncommutative field theories with space-time noncommutativity are not unitary, since they do not have all the relevant degrees of freedom necessary for an approximate description of a unitary string theory. This is what we found from our field theoretic analysis.
Recently, it has been noticed by several groups that it is possible to define a limit of string theory in a background electric field in which the full tower of open string states decouple from the closed strings . In this limit the background electric field is sent to its critical value (see for previous analysis of this limit). There is a very simple way of showing that indeed closed strings decouple in this limit. Quantization of open strings with the modified boundary conditions due to the electric field lead to familiar looking mode expansions for the light-cone coordinates $`X^\pm `$. In the limit that $`EE_c`$, the waves on the string for the $`X^\pm `$ directions become chiral, that is, they are either purely right moving or left moving waves. Therefore, in this limit, it is impossible for such an open string to become a closed string, since closed strings require waves which are left moving and right moving.
There is still a lot to learn about noncommutative theories, both field theories with space noncommutativity as well as the recently discovered decoupled open string theories with space time noncommutativity . The infrared divergences in the magnetic theories for $`pp=0`$ certainly need to better understood within a field theory framework. These theories, as they stand, have no infrared safe observables and the S-matrix is ill-defined. Perhaps nonperturbative input will be required to address this problem.
Acknowledgments
We would like to thank Steve Giddings, Djordje Minic, Hirosi Ooguri, Mark Wise and Edward Witten for helpful discussions and Hirosi Ooguri for carefully reading of the manuscript. J.G. and T.M. are supported in part by the DOE under grant no. DE-FG03-92-ER 40701.
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# Time evolution of the Partridge-Barton Model
## 1 Introduction
Early in life we perceive that everything around us, inanimate objects, animals and human beings undergo a variety of changes that accompany the passage of time. Everything suffers a progressive deterioration with time. This phenomenon is called ageing or senescence and it is characterized by a decline in the physical capabilities of the individuals. Several theories (see and references therein) have been suggested to explain why there is senescence, when it occurs and what are the biological processes responsible for it. Usually, these theories are divided into three classes: biochemical, evolutionary and telomeric. The first invokes damages on DNA, cells, tissues and organs and connect senescence with imperfections of the biochemical processes. One kind of this biochemical imperfection is the presence of free radicals which can cause death of the cells or may even lead to cancer . The evolutionary theory , on the other hand, explains the senescence as a competitive result of the reproductive rate, mutation, heredity and natural selection. In the telomere hypothesis , senescence depends on the cumulative number of cell divisions. The replication of a normal cell is followed by a telomeric shortening. This acts as a counting mechanism which controls the number of divisions.
Evolutionary theories of ageing are hypothetico-deductive in character, not inductive. They do not contain any specific genetic parameter, but only physiological factors and constraints imposed by the environment. There are two kinds: the optimality theory and the mutational theory. In the optimality theory , senescence is a result of searching an optimal life history where survival late in life is sacrificed for the sake of early reproduction. For the mutational theory , on the other hand, ageing is a process which comes from a balance between Darwinian selection and accumulation of mutations. The natural selection efficiency to remove harmful alleles in a population depends on when in the lifespan they come to express. Alleles responsible for lethal diseases that express late in life, escape from the natural selection and accumulate in the population, provoking senescence. Nevertheless, if the natural selection is too strong then deleterious mutations might not accumulate in the population and the eternal youth could be reached. An evolutionary model with such characteristics was recently studied and solved by Onody and de Medeiros .
A simple evolutionary model of ageing is the Partridge-Barton model. It was introduced to illustrate the optimality theories of ageing. Its principal feature is the inclusion of the antagonistic balance mechanism . This mechanism arises out from processes which enhance the lifespan early in life, but have deleterious effects latter.
In this work, we find an exact solution for the whole dynamics of the Partridge-Barton model. When only deleterious somatic mutations and pleiotropy are present the time evolution of the model can be formulated in a matricial form. Explicit analytic expressions can be written for the mean survival probabilities and the growth rate. For large time $`t`$, the system behavior is dominated by the matrix largest eigenvalue. The existent integrals can be solved by the saddle point approximation, allowing us to determine precisely the steady state values of the survival probabilities. A time expansion for the population’s mean age shows that it converges to a constant value according to a $`t^1`$ power law, a result which was first obtained by Ray . All the results were confirmed by some Monte Carlo simulations that we performed.
## 2 The Partridge-Barton Model
In the Partridge-Barton model there are only three ages. The population consists of babies $`(age=0)`$, juveniles $`(age=1)`$ and adults $`(age=2)`$. The survival probabilities from infancy to juvenile is $`J_1`$ and from juvenile to adulthood is $`J_2`$. Reproduction is permitted only to juveniles and adults, with rates $`m_1`$ and $`m_2`$, respectively. Babies don’t have offsprings and adults are eliminated from the population after reproduction.
The population grows at a steady rate $`r`$. The Malthusian growth exponent $`r`$ is related to the other parameters of the model through a discrete version of the Euler-Lotka equation
$$m_1J_1e^r+m_2J_1J_2e^{2r}=1.$$
(1)
The antagonistic pleiotropy arises when the same gene is responsible for multiple effects. For example, genes enhancing early survival by promotion of bone hardening might reduce later survival by promoting arterial hardening. Partridge and Barton implemented the basic ideas of the antagonistic pleiotropy by adopting the constraint, $`J_1+J_2^x=1`$, between the survival probabilities $`J_1`$ and $`J_2`$. The parameter $`x`$ is a real positive number whose value depends on the kind of population we are dealing with. The pleiotropic condition states that it is impossible to sustain simultaneously both high juvenile and adult survivals. For the particular case in which $`m_1=m_2=1`$ and $`x=4`$, Partridge and Barton found $`J_1=0.935`$ and $`J_2=0.505`$ as the values which maximize the growth rate $`r`$.
Also the action of deleterious or helpful mutations can be added to the model. Using Monte Carlo simulations, Stauffer studied the case in which the pleitropic constraint $`J_1+J_2^4=1`$ is accompanied by random somatic mutations. His results clearly show that the survival probabilities $`J_1`$ and $`J_2`$ move rapidly to stationary values with $`J_1>J_2`$. This fact means that the model exhibits senescence, in the sense that the adult survival is lower than the juvenile. In the absence of mutations, $`J_1`$ and $`J_2`$ tend towards $`0.935`$ and $`0.505`$ in accord with the Partridge-Barton conclusions. However, it is not clear how the system drives itself towards these optimal values.
## 3 Analytical Solution
In this section we obtain the exact time solution of the Partridge-Barton model in the presence of pleiotropy and somatic mutations.
Let $`N_i(J_i,t)`$ be the number of individuals at age $`i`$ ($`i=0,1,2`$) with survival probability between $`J_i`$ and $`J_i+dJ_i`$ at time $`t`$. We choose, as initial condition, a population with the profile
$$N_i(J_i,0)=N_0\delta _{i,0},$$
(2)
that is, in $`t=0`$ there are only $`N_0`$ babies with the survival probabilities $`J_0`$ uniformly distributed in the interval $`[0,1]`$.
At time $`t`$, all babies are equally submitted to somatic and deleterious mutations with strength $`\alpha `$ ($`\alpha <1`$). Their survival probabilities $`J_0`$ are changed to $`J_1=\alpha J_0`$. Subsequently, all these babies pass through natural selection in a such way that, on average, the number of juveniles with survival probability $`J_1`$ at the instant $`t+1`$ is given by
$$N_1(J_1,t+1)=J_1N_0(J_0,t).$$
(3)
Since the mutation is restricted to be somatic, each one of the $`N_1(J_1,t+1)`$ juveniles will give birth to exactly $`m_1`$ offspring with survival probability $`J_0`$.
Now, the probability with which a juvenile will reach adulthood must take into account the antagonistic pleiotropy and the somatic deleterious mutations. As pleiotropy is not affected by the somatic mutations, a juvenile with survival probability $`J_1`$ (formerly, a baby with survival probability $`J_0`$) will change its survival probability to $`(1J_0)^{1/x}`$, where $`x`$ is a real positive number and a measurement of the pleiotropic constraint. Under the action of a deleterious somatic mutation, described by a parameter $`\beta `$ ($`\beta <1`$, fixed), the new survival probability can be written as $`J_2=\beta (1J_0)^{1/x}`$. Submitting all juveniles to natural selection we get, on average, the number of adults with survival probability $`J_2`$ which is given by
$$N_2(J_2,t+1)=J_2N_1(J_1,t).$$
(4)
Each one of these adults will generate $`m_2`$ descendants with survival probability $`J_0`$ since the mutations are not inherited.
In general, the number of babies with survival probability $`J_0`$ is given by
$$N_0(J_0,t)=m_1N_1(J_1,t)+m_2N_2(J_2,t),fort1$$
(5)
where $`J_1=\alpha J_0`$ and $`J_2=\beta (1J_0)^{1/x}`$. If we substitute equation (5) into (3) we can write the following recursive matricial equation
$$\left(\begin{array}{c}N_1(J_1,t+1)\\ N_2(J_2,t+1)\end{array}\right)=A\left(\begin{array}{c}N_1(J_1,t)\\ N_2(J_2,t)\end{array}\right),$$
(6)
where $`A`$ is the matrix
$`A=\left(\begin{array}{cc}m_1J_1& m_2J_1\\ J_2& 0\end{array}\right).`$
Iterating the equation above and using the initial condition, we get for $`t2`$
$$\left(\begin{array}{c}N_1(J_1,t)\\ N_2(J_2,t)\end{array}\right)=J_1N_0(J_0,0)A^{t2}\left(\begin{array}{c}m_1J_1\\ J_2\end{array}\right),$$
(8)
with $`A^0`$ meaning the identity matrix.
The complete dynamics of the Partridge-Barton model can be obtained by diagonalizing the matrix $`A`$. We have, explicitly (for $`t2`$)
$`N_1(J_1,t)`$ $`=`$ $`{\displaystyle \frac{J_1N_0(J_0,0)}{\sqrt{m_1^2J_1^2+4m_2J_1J_2}}}[m_1J_1(\lambda _+^{t1}\lambda _{}^{t1})`$ (9)
$`+m_2J_1J_2(\lambda _+^{t2}\lambda _{}^{t2})],`$
$`N_2(J_2,t)`$ $`=`$ $`{\displaystyle \frac{J_1N_0(J_0,0)}{\sqrt{m_1^2J_1^2+4m_2J_1J_2}}}[m_1J_1J_2(\lambda _+^{t2}\lambda _{}^{t2})`$ (10)
$`+m_2J_1J_2^2(\lambda _+^{t3}\lambda _{}^{t3})],`$
where
$$\lambda _\pm =\frac{m_1J_1\pm \sqrt{m_1^2J_1^2+4m_2J_1J_2}}{2}$$
(11)
are the eigenvalues of the matrix $`A`$, $`J_1=\alpha J_0`$ and $`J_2=\beta (1J_0)^{1/x}`$. Let us point out that the time evolution of the babies distribution $`N_0(J_0,t)`$, can be calculated using equations (5), (8) and (9). Having the expressions above, we can determine the evolution of many other quantities like the total number of inviduals at age $`i`$ $`N_i(t)=_0^1N_i(J_i,t)𝑑J_i`$ or their mean survival probabilities $`J_i(t)=\frac{_0^1J_iN_i(J_i,t)𝑑J_i}{_0^1N_i(J_i,t)𝑑J_i}`$. The given input parameters are the initial population $`N_0`$, the birth rates ($`m_1`$ and $`m_2`$), the mutation strengths ($`\alpha `$ and $`\beta `$) and the pleiotropic constraint ($`x`$).
## 4 Asymptotic Limit
Before taking the asymptotic limit, we observe that $`\lambda _+`$ is the largest eigenvalue for all possible values of the input parameters. Once these parameters are fixed and $`J_2=\beta (1J_1/\alpha )^{1/x}`$, $`\lambda _+`$ is in the last instance a function of $`J_1`$. From the equation ($`8`$) we have, asymptotically
$$N_1(J_1,t)e^{tln[\lambda _+(J_1)]}.$$
(12)
By integrating in $`J_1`$ the expression above, we can get the total number of juveniles $`N_1(t)`$. It is convenient to change the integration variable $`J_1`$ for a new variable $`y`$ (a monotonically increasing function of $`J_1`$), $`y=cot(\pi J_1)`$, such that
$$N_1(t)_{\mathrm{}}^{\mathrm{}}\frac{e^{tln[\lambda _+(y)]}}{\pi (1+y^2)}𝑑y.$$
(13)
For large time $`t`$, this integral can be evaluated by the saddle point approximation. We thus obtain
$$N_1(t)=A(\stackrel{~}{y})\frac{e^{tln[\lambda _+(y)]}}{\sqrt{t}},$$
(14)
where $`\stackrel{~}{y}`$ is the value which maximize the eigenvalue $`\lambda _+`$ and $`A(\stackrel{~}{y})=\sqrt{\frac{\pi }{\frac{1}{2\lambda _+}\frac{d^2\lambda _+}{dy^2}|_{y=\stackrel{~}{y}}}}`$.
In the original paper of Partridge and Barton, the optimization process was achieved by a direct (and not well explained) maximization of the growth rate. Here, in our formalism, it is a simple and a natural consequence of taking the asymptotic time limit in the exact evolving equations. Further, the growth rate or the Malthusian exponent is simply given by $`ln[\lambda _+(\stackrel{~}{y})]`$.
To have deepest insight in the dynamics, let us determine the probability density $`P_1(J_1,t)`$ of finding a juvenile at time $`t`$ with survival probability between $`J_1`$ and $`J_1+dJ_1`$. It is given by
$$P_1(J_1,t)=\frac{N_1(J_1,t)}{_0^1N_1(J_1,t)𝑑J_1}=\frac{N_1(J_1,t)}{N_1(t)}\sqrt{t}e^{tln[\frac{\lambda _+(J_1)}{\lambda _+(\stackrel{~}{J_1})}]}$$
(15)
where we have used equations ($`11`$) and ($`13`$) and $`\stackrel{~}{y}=cot(\pi \stackrel{~}{J_1})`$. Clearly, at the asymptotic limit, the distribution probability $`P_1(J_1,t\mathrm{})`$ approaches the Dirac delta function $`\delta (J_1\stackrel{~}{J_1})`$ and the mean survival probability at age $`1`$, is simply given by $`J_1=\stackrel{~}{J_1}`$. Similar results can be obtained for the ages $`0`$ and $`2`$. Another interesting quantity which can be calculated is the population mean age $`A(t)`$ defined as $`A(t)=\frac{_{i=0}^2iN_i(t)}{_{i=0}^2N_i(t)}`$. It is straightforward to show that
$$A(t)=\frac{\gamma +2}{\gamma (1+m_1)+(1+m_2)}+\left\{\frac{2\gamma (1+m_1)+\gamma (1+m_2)}{2[\gamma (1+m_1)+(1+m_2)]^2}\right\}t^1+O(t^2)$$
(16)
where $`\gamma =\frac{\lambda _+(\stackrel{~}{J_1})}{\stackrel{~}{J_2}}`$ with $`\stackrel{~}{J_2}=\beta (1\stackrel{~}{J_1}/\alpha )^{1/x}`$. So we rederive, in a quite simple way, the power law decayment first found by Ray .
## 5 Discussion
We solved exactly in this paper the Partridge-Barton model under the action of arbitrary pleiotropic constraints and deleterious somatic mutations. Through a matricial formalism we were able to predict the complete time evolution of the population. We derived analytic expressions for the time dependence of the mean survival probabilities and the Malthusian exponent. Since for large time $`t`$ the system behavior is controlled by the largest eigenvalue, it was possible to obtain the steady state values of the survival probabilities and to demonstrate, in a simple way, that the population mean age has a power law $`t^1`$ decayment to its final constant value.
For comparison with our analytical results, we also performed some Monte Carlo simulations. In these simulations, the natural selection is implemented by discarding any individual with survival probability smaller than a random number (generated from a uniform distribution). The deleterious somatic mutations and the antagonic pleiotropy can be easily incorporated into the computer program. More difficult is to avoid an explosion of the computer’s memory due to the unlimited growth of the population. To take this problem into account, we resort to the Verhulst factor which is commonly used in such circumstances.
In Figure 1 we put together the analytical solution and the Monte Carlo result. The exact solution was plotted by inserting equations ($`8`$, $`9`$ and $`10`$) into the expressions for the mean survival probabilities $`J_i(t)`$ and by integrating them using the software Maple . We conclude that the Monte Carlo simulations confirm very well the theoretical results.
Finally, let us to point out that, unfortunately, the technique developed here cannot be applied to the case in which mutations are hereditary. The main reason for this come from the fact that the equation ($`5`$) is no longer valid.
## 6 Acknowledgements
We acknowledge CNPq (Conselho Nacional de Desenvolvimento Científico e Tecnológico) for the financial support.
FIGURE CAPTION
Figure 1. The continuous lines correspond to the analytical solutions and the square symbols to the Monte Carlo simulations. We used $`\alpha =0.82`$, $`\beta =0.67`$ , $`x=4`$, $`m_1=m_2=1`$ and $`N_0=4000`$. The steady state values are $`\stackrel{~}{J_1}=0.77`$ and $`\stackrel{~}{J_2}=0.33`$. There is senescence, i. e., $`J_2<J_1`$.
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# CIRRUS SPECTRA OF LOW SURFACE BRIGHTNESS REGIONS
## 1. INTRODUCTION
At wavelengths longer than $`100`$$`\mu `$m the far-infrared emission of the Galaxy is dominated by thermal emission from large dust particles that are in thermal equilibrium with the interstellar radiation field. Dust temperatures determined from DIRBE measurements are typically 16–19 K (e.g. ??; ??) in rough agreement with model predictions (e.g. ??).
Most previous FIR dust temperature determinations have been based on COBE observations. These have, however, poor spatial resolution and in low surface brightness regions the sensitivity of the detectors becomes a limiting factor so that only an average dust spectrum over a large area can be determined.
We use ISOPHOT observations to study the spectrum of galactic cirrus in regions selected for their low surface brightness. Because of the better sensitivity and angular resolution of the ISOPHOT instrument it is possible to determine the cirrus spectrum in these fields based on the cirrus intensity variations. In this way constant or slowly varying components, e.g. Zodiacal light, can be eliminated.
The determination of the cirrus spectrum is important also for extragalactic studies where the galactic foreground emission must be separated. In future work the same fields will be used to estimate the FIR component of the cosmic background radiation.
We start by presenting the observations used in this study. After that calibration differences between ISOPHOT and DIRBE observations are discussed. Calibration issues are obviously essential for the determination of the spectra and the dust temperatures. Finally, the method used to derive the cirrus spectra and the derived dust temperatures are presented. A $`\nu ^2`$ dust emissivity law will be assumed throughout this paper.
## 2. OBSERVATIONS
The cirrus spectra were determined in six fields that had been observed at 90 $`\mu `$m, 150 $`\mu `$m and 180 $`\mu `$m with the ISOPHOT instrument (??) aboard the ISO satellite (??). The observation consist of PHT22 raster maps. Detailed description of the observations is given in Table 1. Data reduction was done using the ISOPHOT Interactive Analysis software package (PIA) (??) version 8.0. In the following the analysis is based on measurements reduced to the AAP level.
Flatfielding was done outside PIA using custom routines that also removed slow drifts of the detector pixels relative to each other. Good flat fielding is not critical for the determination of the spectra but it will reduce the noise in the subsequent steps.
## 3. COMPARISON WITH DIRBE SURFACE BRIGHTNESS VALUES
The results presented in this paper are based on the calibration performed with Fine Calibration Source (FCS) measurements. The absolute calibration of ISOPHOT observations has been estimated to be better than 30% (??). ?? recently reported for both C100 and C200 accuracies better than 20% relative to DIRBE. These results apply to fields with surface brightness above $``$5 MJy sr<sup>-1</sup> while several of our C200 maps are below 4 MJy sr<sup>-1</sup>. It is therefore useful to compare our observations with DIRBE data especially as the calibration accuracy is essential for the determination of the cirrus spectra.
The DIRBE and ISOPHOT calibration was studied in the case of the fields NGPS, NGPN, EBL22, and EBL26 (see Table 1). Three methods were used to perform the comparison:
1. Comparison of absolute surface brightness levels. The ISOPHOT maps were compared with DIRBE ZSMA (Zodi-Subtracted Mission Average) data to which Zodiacal light was added according to the model given by ??). Using the DIRBE weekly maps observed with the same solar elongation as the ISO data did not change the results significantly. Because of the lower noise we chose to use the ZSMA data. Both DIRBE ZSMA and ISO were colour corrected assuming a $`\nu ^2B_\nu `$ spectrum with $`T_{\mathrm{dust}}`$=18 K. The zodiacal light estimates were colour corrected assuming black body spectrum with $`T_{\mathrm{dust}}`$=270 K (??; ??). Corresponding curves fitted to the DIRBE ZSMA 100 $`\mu `$m, 140 $`\mu `$m and 240 $`\mu `$m data and the zodiacal light estimates were used to derive the DIRBE surface brightness estimates interpolated to the wavelength of the ISO observations. The average ISO flux density was compared with the derived DIRBE values calculated as the average over ISO map and weighted by the DIRBE beam. The result is the surface brightness ratio $`S`$(DIRBE)/$`S`$(ISOPHOT).
2. Direct comparison of the surface brightness variations in ISOPHOT and DIRBE ZSMA maps. All data were first colour corrected for a spectrum $`\nu ^2B(\nu )`$ with $`T_{\mathrm{dust}}`$ =18.0 K. For each DIRBE pixel the data were interpolated to the wavelength of the ISOPHOT observations and the corresponding ISOPHOT surface brightness estimate was calculated as a weighted average over the DIRBE beam. Linear fit was done to these points to derive the slope between the surface brightness values $`k=\mathrm{\Delta }S`$(DIRBE)/$`\mathrm{\Delta }S`$(ISOPHOT).
3. Comparison via IRAS. Linear relations were established between the IRAS 100 $`\mu `$m data and the ISO observations and between the IRAS data and the DIRBE ZSMA data interpolated to the wavelength of the ISO observations using a fitted curve $`\nu ^2B_\nu `$($`T`$=18K). For each IRAS ISSA map pixel inside the field the corresponding ISO surface brightness was calculated as an average weighted with a Gaussian with FWHM$``$5 arcmin. The relation between IRAS and DIRBE surface brightnesses was determined with a similar method but over an area with radius $``$2 degrees where for each DIRBE pixel the average IRAS value was calculated by weighting with the DIRBE beam. The slopes of these two linear relations $`\mathrm{\Delta }S(\mathrm{DIRBE})`$/$`\mathrm{\Delta }S(\mathrm{IRAS}(100\mu \mathrm{m}))`$ $`\mathrm{\Delta }S(\mathrm{IRAS}(100\mu \mathrm{m}))`$/$`\mathrm{\Delta }S(\mathrm{ISOPHOT})`$ give the ratio between DIRBE and ISO scales, $`k=\mathrm{\Delta }S`$(DIRBE)/$`\mathrm{\Delta }S`$(ISOPHOT).
Method (1) suffers from the small size of the ISOPHOT maps which means that for each DIRBE pixel a large portion of the flux comes from outside the area mapped with ISOPHOT.
Method (2) is even more affected by the limited map sizes and the large size of the DIRBE beam since we must determine the relation between brightness variations. However, NGPN and NGPS form together a 3.1 degrees long strip where the surface brightness gradient is along the longer side of the map. Here the fact that the ISO map is narrow compared with the DIRBE beam should not lead to significant errors. On the other hand, in EBL26 the surface brightness drops quickly outside the bright region in the northern part of the map. The surface brightness variation is not resolved by DIRBE beam and the values $`\mathrm{\Delta }(\mathrm{DIRBE})`$/$`\mathrm{\Delta }(\mathrm{ISOPHOT})`$ obtained would be underestimated. In EBL22 and NGPN the method does not work because of the lack of sufficient surface brightness variations.
Methods (2) and (3) are not affected by the presence of zodiacal light. The last method uses IRAS ISSA maps as an intermediate step and is therefore much less affected by the poor resolution of the DIRBE data. It allows also the use of DIRBE data over a much larger area and thereby reduces the effect of the large noise present in 140 $`\mu `$m and 240 $`\mu `$m DIRBE observations. In these low surface brightness regions the error estimates of 140 $`\mu `$m and 240 $`\mu `$m DIRBE measurements exceed 50%. On the other hand, the method is based on the assumption that the ratio of the emission at 100 $`\mu `$m and at other wavelengths remains constant.
The results are given in Table 2. It can be seen that DIRBE and ISOPHOT surface brightness values agree typically to within 30%. However, at 90 $`\mu `$m the ISOPHOT surface brightness tends to be above the value predicted from DIRBE observations while at longer wavelengths the situation is reversed. The extrapolation of the DIRBE values to 90 $`\mu `$m using the modified black body fitted to three longer wavelength may not be entirely justified and may lead to underestimation of the surface brightness. Also, in these faintest regions the accuracy of the DIRBE 140 $`\mu `$m and 240 $`\mu `$m is poor and a reliable comparison with ISOPHOT C200 values is difficult. The modified black body curves fitted to DIRBE data are determined mostly by the 100 $`\mu `$m DIRBE value which has significantly smaller error estimates than the two longer wavelengths, and the assumed dust temperature, $`T`$=18.0 K. Increasing the temperature by one degree would increase $`k`$ at 90 $`\mu `$m by some 5% and decrease it at longer wavelengths by some 15%.
In spite of uncertainties the results suggest that the dust temperatures derived from ISOPHOT observations are, because of calibration differences, higher than those based on DIRBE data. Using PIA versions prior to 8.0 the differences between ISOPHOT and DIRBE surface brightnesses were typically larger by some 20%.
## 4. CIRRUS SPECTRA
As the first step all data were colour corrected in order to derive monochromatic surface brightness values at the reference wavelengths of the filters. Correction was done for spectrum $`\nu ^2B_\nu (T=18K)`$ which approximates the expected cirrus spectrum. A large fraction of the observed surface brightness values is due to Zodiacal light. Since we will study only the relative surface brightness variations the Zodiacal light is eliminated from the analysis and it does not affect the required colour correction.
In the calculations the 180 $`\mu `$m maps are used as the reference. For each 180 $`\mu `$m pixel the corresponding surface brightness value at another wavelength is calculated as a weighted average. The spatial weighting is done with a Gaussian with FWHM$``$100 arcsec and the individual measurements are also weighted according to their error estimates. The calculated values are plotted against the 180 $`\mu `$m surface brightness and a straight line is fitted to the points with a least squares algorithm that takes into account the error estimates on both axes. The procedure is repeated for all maps and the slopes of the fitted lines provide the emission spectrum i.e. for each observed wavelength the relative intensity with respect to the 180 $`\mu `$m emission.
Figure 2 shows as an example the relations between 180 $`\mu `$m and 90 $`\mu `$m and 180 $`\mu `$m and 150 $`\mu `$m surface brightness values in the field EBL26. The fitted least square lines are also shown. Zodiacal light will appear merely as an offset and does not affect the slopes that determine the cirrus spectrum. Furthermore, zodiacal light is weak at 180 $`\mu `$m and we are looking only for variations correlated with 180 $`\mu `$m emission.
Galactic cirrus is the dominant source of brightness fluctuations and therefore the spectra obtained characterize the galactic cirrus emission. However, in low surface brightness regions there may be a significant contribution from faint extragalactic sources. The DIRBE (??) and FIRAS (??) experiments indicated a FIR cosmic infrared background flux of 1 MJy sr<sup>-1</sup> between 100 and 240 $`\mu `$m. In fields like EBL22 this would mean that at 180 $`\mu `$m about one third of the surface brightness is due to extragalactic sources. Since the cirrus spectra are determined from the brightness variations we must consider the brightness fluctuations caused by these sources in relation to the cirrus fluctuations.
The cirrus power spectrum is proportional to $`B_0/k^3`$, where $`B_0`$ is the cirrus brightness and $`k`$ the spatial frequency (??). In other words, cirrus fluctuations decrease rapidly as we move to fainter cirrus regions and smaller spatial scales. ?? have reported a detection of the extragalactic background fluctuations at 175 $`\mu `$m in the Marano 1 field that has cirrus surface brightness comparable to our fields. Even after removal of detected sources they found at scales below 10 arcmin fluctuations in excess of the expected cirrus contribution. This was interpreted to be caused by point sources below the detection limit.
The larger dimension of our maps is typically 1.5 degrees and most of the surface brightness variations can be expected to be due to cirrus. This is particularly clear in fields NGPS and EBL26 where there is a clear intensity gradient along the mapped strip consistent with the larger scale cirrus distribution visible in IRAS maps. On the other hand, the field EBL22 is remarkably flat and it is conceivable that in the absence of large cirrus structures the extragalactic point sources could affect the derived spectrum. The small dynamic range means, however, that in this case the accuracy of the derived spectrum is not very good. The fields VCN and VCS were observed close the bright galaxy NGC 5907 and the VCS field actually extends over the centre of the galaxy. Pixels that were clearly affected by the emission from NGC 5907 were first discarded but it is still possible that the derived spectra are affected by emission of the galaxy. Because of the small field size ($``$10 arcmin) the VCN spectrum could be affected by faint extragalactic sources.
The spectra obtained for the six fields are presented in Figure 3. The error bars include the error estimates obtained from the fitting procedure and a further 10% error representing the error in the relative calibration of the observations made with different filters. Modified Planck curves, $`\nu ^2B_\nu (T)`$, were fitted to the data and the dust temperatures obtained are shown in the figure.
## 5. DISCUSSION
We have determined cirrus dust temperatures in six fields with low surface brightness. The derived dust temperatures, $``$18-20 K, are higher than the typical temperatures estimated for galactic dust from DIRBE and FIRAS experiments assuming the same $`\nu ^2`$ emissivity law. However, the values are still within the range found in these studies. ?? found an average temperature of 17.5 K while higher temperatures were mostly associated with star forming regions. Temperatures exceeding 20 K were seen only towards the galactic centre and some prominent star forming regions although the scatter in the temperatures is more than 1 K at all galactic latitudes.
Our fields are all at high galactic latitudes (see Table 1). Based on COBE measurements at 100–240 $`\mu `$m ?? derived average dust temperatures of $`T=17.5`$ K for regions at $`|b|>`$30 degrees. Earlier FIRAS estimates (??) also suggested similar or lower values. However, based on the DIRBE 140 $`\mu `$m to 240 $`\mu `$m ratios ??) reported colour temperatures of $`T`$19 K for regions with $`|b|>`$45 degrees.
In Section 3. we found systematic differences between ISOPHOT and DIRBE calibrations. Compared with DIRBE the ISOPHOT surface brightnesses were found to be higher at 90 $`\mu `$m and lower at longer wavelengths. According to Table 2 adopting the DIRBE calibration would increase our 90 $`\mu `$m data by up to 30% relative to the longer wavelengths. This would reduce the estimated dust temperature by $``$1.5 K. Therefore the temperature difference between DIRBE and our results can be explained entirely by the difference in calibration.
Our temperature determinations are based on data at three wavelengths: 90 $`\mu `$m, 150 $`\mu `$m and 240 $`\mu `$m. The emission at wavelengths $`\lambda >`$100 $`\mu `$m can be explained by classical grains in equilibrium with the interstellar radiation field. Smaller grains that are transiently heated to higher temperatures are responsible for the emission at shorter wavelengths. Based on both observations and theoretical calculations the contribution from transietly heated particles is significant already at 60 $`\mu `$m and may extend even to slightly longer wavelengths (??; ??; ??). The ISOPHOT 90 $`\mu `$m filter is wide enough to pick up some radiation from wavelenths down to $``$60 $`\mu `$m and our observations may be affected to a small degree by this other grain population. Our temperature estimates do not correspond exactly to the equilibrium temperature of large dust particles but the bias towards higher temperatures is expected to be small.
In the low surface brightness regions the accuracy of DIRBE measurements is not very good and one can obtain only the average spectrum of a very large area. Because of the better resolution and sensitivity of ISOPHOT we have been able to study the cirrus spectrum in separate, small regions. Adopting the ISOPHOT calibration the accuracy of the temperature estimates is $``$1 K. For example, the 2 K difference between the adjacent fields NGPS and NGPN is likely to be true.
## ACKNOWLEDGMENTS
The work was supported by the Academy of Finland Grant no. 1011055 and funding from DLR and MPG in Germany.
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# Covariant Quark-Representation of Composite Meson Systems and Chiral Symmetry
## 1 Introduction
There are the two contrasting view points of composite quark-antiquark mesons: The one is non-relativistic, based on the approximate symmetry of $`LS`$-coupling in the non-relativistic quark model (NRQM); while the other is relativistic, based on the dynamically broken chiral symmetry typically displayed in the Nambu Jona-Lasinio (NJL) model. The $`\pi `$-meson (or $`\pi `$-nonet) is now widely believed to have a dual nature of non-relativistic particle with $`(L,S)=(0,0)`$ and also of relativistic particle as a Nambu-Goldston boson with $`J^P=0^{}`$ in the case of spontaneous breaking of chiral symmetry. However, no successful attempts to unify the above two view points have been yet proposed. On the other hand we have developed the covariant oscillator quark model (COQM) for many years as a covariant extension of NRQM, which is based on the boosted $`LS`$-coupling scheme. The meson wave functions (WF) in COQM are tensors in the $`\stackrel{~}{U}(4)O(3,1)`$ space and reduce at the rest frame to those in the $`SU(2)_{\mathrm{spin}}O(3)_{\mathrm{orbit}}`$ space in NRQM. The COQM has been successful especially in treating the $`Q\overline{Q}`$ meson system and the ($`q,\overline{Q}`$) or ($`Q,\overline{q}`$) meson system, leading, respectively, to a satisfactory understanding of radiative transitions and to the same weak form factor relations as in the heavy quark effective theory (HQET). However, in COQM no consideration on chiral symmetry has been given and it is not able to explain the dual nature of $`\pi `$ meson.
The purpose of this paper is to get rid of this defect in COQM and is to give a unified view point of the two contrasting ones of the composite meson systems, extending the tensors of WF from the restricted ones necessary only in the boosted $`LS`$ coupling scheme to the general ones in the $`\stackrel{~}{U}(4)O(3,1)`$ space, which are required for taking into account chiral symmetry.
## 2 Covariant Framework for Describing Composite Mesons
For meson WF described by $`\mathrm{\Phi }_A{}_{}{}^{B}(x_1,x_2)`$ ($`x_1,x_2`$ denoting the space-time coordinate and $`A=(\alpha ,a)(B=(\beta ,b))`$ denoting the Dirac spinor and flavor indices of constituent quark (anti-quark)) we set up the bilocal Yukawa equation
$`[{\displaystyle \frac{^2}{X_\mu ^2}}^2(x_\mu ,{\displaystyle \frac{}{x_\mu }})]\mathrm{\Phi }_A{}_{}{}^{B}(X,x)=0`$ (2.1)
($`X(x)`$ denoting the center of mass (CM) (relative) coordinate of meson), where the $`^2`$ is squared mass operator including only a central, Dirac-spinor-independent<sup>\*)</sup><sup>\*)</sup>\*) In the boosted $`LS`$-coupling scheme the squared mass operator $`^2`$ was assumed to be only Pauli-spinor independent, while in the present scheme it is assumed more generally to be Dirac-spinor independent. confining potential. The WF is separated into the plane wave describing CM motion and the (Fierz-component) internal WF as
$`\mathrm{\Phi }_A{}_{}{}^{B}(x_1,x_2)={\displaystyle \underset{𝐏_n,n}{}}(e^{iP_nX}\mathrm{\Psi }_{n,A}{}_{}{}^{(+)B}(x,P_n)+e^{iP_nX}\mathrm{\Psi }_{n,A}{}_{}{}^{()B}(x,P_n)),`$ (2.2)
where the Fierz components $`\mathrm{\Psi }_n^{(\pm )}`$ are eigenfunctions of $`^2`$ as
$`^2(x_\mu ,{\displaystyle \frac{}{x_\mu }},P_n)\mathrm{\Psi }_n^{(\pm )}=M_n^2\mathrm{\Psi }_n^{(\pm )},`$ (2.3)
$`P_{n,\mu }^2=M_n^2,P_{n,0}=\sqrt{M_n^2+𝐏_n^2}`$; and the label $`(\pm )`$ represents the positive (negative) frequency part; and $`n`$ does a freedom of excitation. We have the following field theoretical expression for the WF in mind as a guide for developing the present semi-phenomenological approach:
$`\mathrm{\Phi }_A{}_{}{}^{B}(x_1,x_2)`$ $`=`$ $`{\displaystyle \underset{n}{}}[0|\psi _A(x_1)\overline{\psi }^B(x_2)|M_n+M_n^c|\psi _A(x_1)\overline{\psi }^B(x_2)|0],`$ (2.4)
where $`\psi _A(\overline{\psi }^B)`$ denotes the quark field (its Pauli-conjugate) and $`|M_n(M_n^c|)`$ does the composite meson (its charge conjugate) state, and the first (second) term in the RHS corresponds to the positive (negative) frequency part in Eq. (2.2). The internal WF is, concerning the Dirac-spinor-dependence, expanded in terms of a complete set $`\{W_i\}`$ of free bi-Dirac spinors of quarks and anti-quarks; and the internal WF is expressed as
$`\mathrm{\Psi }_A^{(\pm )B}(x,P_n)={\displaystyle \underset{i}{}}W_{i\alpha }^{(\pm )\beta }(P_n)\varphi _a^{(\pm )b}(x,P_n),`$
$`\varphi ^{(\pm )}(x,P_n)=ϵ_i\overline{W}_i^{()}\mathrm{\Psi }^{(\pm )},`$ (2.5)
where $`A`$ means trace of $`A`$. The ortho-normal relations $`\overline{W}_i^{()}W_j^{(\pm )}=ϵ_i\delta _{ij}`$ holds for the Pauli-conjugate of WF, defined by $`\overline{W}_i^{()}\gamma _4W_i^{(\pm )}\gamma _4`$, where $`\overline{W}_i^{()}`$ is related with the $`W_i^{}`$, and the $`ϵ_i`$ and $`\delta _{ij}`$ denote,respectively, the sign and the Kronecker symbols, see the appendix A).
## 3 Complete Set of Spin Wave Function <br>and Composite Mesons with Definite Spin
We set up the conventional “free” Dirac spinors with four-momentum of composite meson itself $`P=P_M`$, $`D_{q,\alpha }(P)(u_{q,\alpha }(P,s_q),v_{q,\alpha }(P,s_q)`$ ($`s_q=\pm `$ representing the spin up-down)) for quarks and $`\overline{D}_{\overline{q}}{}_{}{}^{\beta }(P)(\overline{v}_{\overline{q}}^\beta (P,s_{\overline{q}}),\overline{u}_{\overline{q}}^\beta (P,s_{\overline{q}})`$($`s_{\overline{q}}=\pm `$ representing spin up-down)) for anti-quarks. It is to be noted that all four spinors for both “quarks and anti-quarks” are necessary<sup>\**)</sup><sup>\**)</sup>\**) For understanding this it may be useful to take an analogy of the Bethe-Salpeter amplitude of deuteron. In expanding the amplitude all $`4\times 4`$ members of direct product of both Dirac spinors for constituent proton and neutron are necessary. to describe the spin WF of mesons. Then the complete set of bi-Dirac spinors is given by<sup>\***)</sup><sup>\***)</sup>\***) In the following from §3 to §4.3 we give only the expressions of $`(+)`$-frequency parts, and consider only the ground states of composite system, disregarding the relative coordinates. In Eq.(3.1) considerations on the freedom of Pauli-spin are neglected. We have given detailed considerations on this problem and useful formulas in Appendix A.
$`\{W^{(+)}(P)\}`$ $`:`$
$`U(P)`$ $`=`$ $`u_q(p_1,s_q)\overline{v}_{\overline{q}}(p_2,s_{\overline{q}})|_{p_{i,\mu }=\kappa _iP_\mu }=u_+(𝑷,s_q)\overline{v}_+(𝑷,s_{\overline{q}}),`$
$`C(P)`$ $`=`$ $`u_q(p_1,s_q)\overline{u}_{\overline{q}}(p_2,s_{\overline{q}})|_{p_{i,\mu }=\kappa _iP_\mu }=u_+(𝑷,s_q)\overline{v}_{}(𝑷,s_{\overline{q}}),`$
$`D(P)`$ $`=`$ $`v_q(p_1,s_q)\overline{v}_{\overline{q}}(p_2,s_{\overline{q}})|_{p_{i,\mu }=\kappa _iP_\mu }=u_{}(𝑷,s_q)\overline{v}_+(𝑷,s_{\overline{q}}),`$
$`V(P)`$ $`=`$ $`v_q(p_1,s_q)\overline{u}_{\overline{q}}(p_2,s_{\overline{q}})|_{p_{i,\mu }=\kappa _iP_\mu }=u_{}(𝑷,s_q)\overline{v}_{}(𝑷,s_{\overline{q}}),`$ (3.1)
where $`u_+(𝑷)(\overline{v}_+(𝑷))`$ and $`u_{}(𝑷)(\overline{v}_{}(𝑷))`$ denote the Dirac spinors with positive energy and momentum $`𝑷`$ and with negative energy and momentum $`𝑷`$, respectively, describing quarks (anti-quarks). These energy and momentum concern with the total meson, while in Eq.(3.1) we have defined technically the momenta of “constituent quarks”<sup>\****)</sup><sup>\****)</sup>\****) In so far as concerned with Eqs. (3.1) and (3.2) the quantities $`\kappa _i`$ and accordingly $`m_i`$ are arbitrary and have no physical meaning. However, $`m_i`$ have proved to be the effective masses of constituent quarks through the phenomenological applications of COQM so far made. as
$`p_{i,\mu }`$ $``$ $`\kappa _iP_\mu ,p_{i,\mu }^2=m_i^2;P_\mu ^2=M^2,M=m_1+m_2`$
$`(\kappa _{1,2}`$ $``$ $`m_{1,2}/(m_1+m_2);\kappa _1+\kappa _2=1).`$ (3.2)
The respective members in Eq.(3.1) satisfy a couple of the corresponding free Dirac equations in momentum space (which are equivalent to the (conventional or new-type of) Bargman-Wigner Equations) and are expressed in terms of their irreducible composite meson WF as follows:
$`(\mathrm{Non}\mathrm{Rela}.\mathrm{comp}.)`$
$`(iP\gamma ^{(1)}`$ $`+`$ $`M)U(P)=0,U(P)(iP\gamma ^{(2)}+M)=0;`$
$`U_A{}_{}{}^{B}(P)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}[(i\gamma _5P_{s,a}^{(NR)b}(P)+i\gamma _\mu V_{\mu ,a}^{(NR)b}(P))(1+{\displaystyle \frac{iP\gamma }{M}})]_\alpha {}_{}{}^{\beta },`$
$`(\mathrm{Semi}\mathrm{Rela}.\mathrm{comp}.)`$
$`\overline{q}\mathrm{type}(iP\gamma ^{(1)}`$ $`+`$ $`M)C(P)=0,C(P)(iP\gamma ^{(2)}+M)=0;`$
$`C_A{}_{}{}^{B}(P)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}[(S_a^{(\overline{q})b}(P)+i\gamma _5\gamma _\mu A_{\mu ,a}^{(\overline{q})b}(P))(1{\displaystyle \frac{iP\gamma }{M}})]_\alpha {}_{}{}^{\beta },`$
$`q\mathrm{type}(iP\gamma ^{(1)}`$ $`+`$ $`M)D(P)=0,D(P)(iP\gamma ^{(2)}+M)=0;`$
$`D_A{}_{}{}^{B}(P)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}[(S_a^{(q)b}(P)+i\gamma _5\gamma _\mu A_{\mu ,a}^{(q)b}(P))(1+{\displaystyle \frac{iP\gamma }{M}})]_\alpha {}_{}{}^{\beta },`$
$`(\mathrm{Extrly}.\mathrm{Rela}.\mathrm{comp}.)`$
$`(iP\gamma ^{(1)}`$ $`+`$ $`M)V(P)=0,V(P)(iP\gamma ^{(2)}+M)=0;`$
$`V_A{}_{}{}^{B}(P)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}[(i\gamma _5P_{s,a}^{(ER)b}(P)+i\gamma _\mu V_{\mu ,a}^{(ER)b}(P))(1{\displaystyle \frac{iP\gamma }{M}})]_\alpha {}_{}{}^{\beta },`$ (3.3)
where all vector and axial-vector mesons satisfy the Lorentz conditions, $`P_\mu V_\mu (P)=P_\mu A_\mu (P)=0`$. Here it is to be noted that, in each type of the above members, the number of freedom counted both in the quark representation and in the meson representation is equal, as it should be ($`2\times 2=4`$ and $`1+3=4`$, respectively). Also it may be amusing to note that each constituent quarks in all the above members is in “parton-like motion,” having the same 3-dimentional velocity as that of total mesons. (For example, in $`V(P),𝒗_{1,2}=\frac{𝒑^{(1,2)}}{p_0^{(1,2)}}=\frac{\kappa _{1,2}𝑷_M}{\kappa _{1,2}P_{M,0}}=𝒗_M`$.)
## 4 Transformation Properties of Composite Meson Systems <br>and Chiral Symmetry
For benefit of covariant quark-representation of mesons given in §3 we can deduce automatically the transformation rules of composite mesons for general symmetry operations as follows:
### 4.1 Charge conjugation
Charge conjugation properties of the bi-spinors and, correspondingly, of the composite mesons are derived from those of quarks as follows:
$`\mathrm{quark}\mathrm{field}:`$ $`\psi _\alpha (x)\psi _\alpha ^c(x)=U_C^1\psi _\alpha (x)U_C=(C^1)_{\alpha \alpha ^{}}\overline{\psi }^\alpha ^{}(x)`$
$`\overline{\psi }^\beta (x)\overline{\psi ^c}^\beta (x)=U_C^1\overline{\psi }^\beta (x)U_C=C^{\beta \beta ^{}}\psi _\beta ^{}(x)`$
$`(CC^{}=1,C=\gamma _4\gamma _2,C\gamma _\mu C^1=^t\gamma _\mu ).`$
$`\mathrm{Internal}\mathrm{meson}`$ $`\mathrm{WF}:`$
$`\mathrm{\Psi }_A^{(+)B}(P,x)(`$ $``$ $`0|\psi _A(x_1)\overline{\psi }^B(x_2)|M)`$
$``$ $`\mathrm{\Psi }_A^{c,(+)B}(P,x)(0|\psi _A(x_1)\overline{\psi }^B(x_2)|M^c`$
$`=`$ $`0|U_C^1\psi _A(x_1)U_CU_C^1\overline{\psi }^B(x_2)U_C|M`$
$`=`$ $`0|(C^1)_{AA^{}}\overline{\psi }^A^{}(x_1)C^{BB^{}}\psi _B^{}(x_2)|M`$
$`=`$ $`(C^1)_{AA^{}}(1)0|\psi _B^{}(x_2)\overline{\psi }^A^{}(x_1)|M(^tC)^{B^{}B})`$
$`=`$ $`(C^1)_{AA^{}}{}_{}{}^{t}\mathrm{\Psi }_{}^{(+)}(P,x)^A^{}{}_{B^{}}{}^{}C_{}^{B^{}B},`$
$`\mathrm{Spinor}`$ $`\mathrm{WF}:`$
$`U_{P_s}^{(NR)}U_{P_s}^{(NR)},U_{V_\mu }^{(NR)}U_{V_\mu }^{(NR)},`$
$`C_S^{\overline{q}}D_S^q,C_{A_\mu }^{\overline{q}}D_{A_\mu }^q,`$
$`U_{P_s}^{(ER)}U_{P_s}^{(ER)},U_{V_\mu }^{(ER)}U_{V_\mu }^{(ER)},`$
$`\mathrm{Composite}`$ $`\mathrm{meson}`$ $`\mathrm{WF}:`$
$`P_{s,a}^{(NR)b}`$ $``$ $`P_{s,b}^{(NR)a},V_{\mu ,a}^{(NR)b}V_{\mu ,b}^{(NR)a},`$
$`S_a^{(\overline{q})b}`$ $``$ $`S_b^{(q)a},A_{\mu ,a}^{(\overline{q})b}A_{\mu ,b}^{(q)a},`$
$`P_{s,a}^{(ER)b}`$ $``$ $`P_{s,b}^{(ER)a},V_{\mu ,a}^{(ER)b}V_{\mu ,b}^{(ER)a}.`$ (4.2)
### 4.2 Chiral transformation
Chiral transformation properties of composite mesons are also derived straightforwardly from those of the bi-spinors. For example, for $`SU(3)`$ chiral transformation
$`\mathrm{\Psi }_A{}_{}{}^{B}(P,x)`$ $``$ $`[e^{i\alpha ^i\frac{\lambda ^i}{2}\gamma _5}\mathrm{\Psi }(P,x)e^{i\alpha ^i\frac{\lambda ^i}{2}\gamma _5}]_A{}_{}{}^{B}.`$ (4.3)
leading to the results:
$`[M^{(J)}(P)]_i^{}`$ $`=`$ $`[^tS^{(J)}(P,\alpha )]_{ij}[M^{(J)}(P)]_jJ=(0,1)`$
$`M_i^{(0)}`$ $``$ $`{}_{}{}^{t}[P_s^{(NR)},S^{(R\overline{q})},S^{(Rq)},P_s^{(ER)}]_{i}^{},`$
$`M_i^{(1)}`$ $``$ $`{}_{}{}^{t}[V_\mu ^{(NR)},A_\mu ^{(R\overline{q})},A_\mu ^{(Rq)},V_\mu ^{(ER)}]_{i}^{},`$ (4.4)
where $`[S^{(J)}]`$ denotes a unitary matrix, of which explicit form is omitted here. (For infinitesimal transformation, see appendix B).
### 4.3 Light-quark meson system—“chiral SU(6) multiplet”
The quark representation applying to the light-quark mesons is obtained by the linear transformation of the bi-Dirac spinors given in §3 as follows:
$`U_{P_s,\alpha }^{(N,E)\beta }`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(U_{P_s}\pm V_{P_s})_\alpha {}_{}{}^{\beta }={\displaystyle \frac{1}{2}}[(i\gamma _5,\gamma _5v\gamma )]_\alpha {}_{}{}^{\beta };P_s^{(N,E)};C=(+,+)`$
$`C_{S,\alpha }^{(N,E)\beta }`$ $``$ $`{\displaystyle \frac{(1,i)}{\sqrt{2}}}(D_S\pm C_S)_\alpha {}_{}{}^{\beta }={\displaystyle \frac{1}{2}}[(1,v\gamma )]_\alpha {}_{}{}^{\beta };S^{(N,E)};C=(+,)`$
$`U_{V,\alpha }^{(N,E)\beta }`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(U_V\pm V_V)_\alpha {}_{}{}^{\beta }={\displaystyle \frac{1}{2}}[(i\stackrel{~}{\gamma }_\mu ,i\sigma _{\mu \nu }v_\nu )]_\alpha {}_{}{}^{\beta };V^{(N,E)};C=(,)`$
$`C_{A,\alpha }^{(N,E)\beta }`$ $``$ $`{\displaystyle \frac{(1,i)}{\sqrt{2}}}(D_A\pm C_A)_\alpha {}_{}{}^{\beta }={\displaystyle \frac{1}{2}}[(i\gamma _5\stackrel{~}{\gamma }_\mu ,\gamma _5\sigma _{\mu \nu }v_\nu )]_\alpha {}_{}{}^{\beta };A^{(N,E)};C=(+,)`$
(4.5)
($`v_\mu P_\mu /M,\stackrel{~}{\gamma }_\mu v_\mu 0`$; and $`U_{P_s}`$ denotes the coefficient bi-spinors of $`P_s`$ and so on ), where we have given also the charge-conjugation parity of the corresponding (hidden flavor) composite mesons. The chiral transformation properties of the new bi-spinors are easily seen to be similar as the conventional ones as
$`1i\gamma _5,v\gamma \gamma _5v\gamma ,`$ (4.6)
$`i\gamma _5\stackrel{~}{\gamma }_\mu i\stackrel{~}{\gamma }_\mu ,i\sigma _{\mu \nu }v_\nu \gamma _5\sigma _{\mu \nu }v_\nu .`$
### 4.4 “Local chiral SU(6) field”
Extending our considerations to include the $`()`$-frequency part, we are led to a unified expression of what to be called, Local Chiral SU(6) field, as
$`\mathrm{\Psi }_A{}_{}{}^{B}(X)`$ $`=`$ $`\mathrm{\Psi }_A^{(N)B}(X)+\mathrm{\Psi }_A^{(E)B}(X)`$
$`\mathrm{\Psi }_A^{(N)B}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[i\gamma _{5\alpha }{}_{}{}^{\beta }P_{s,a}^{(N)b}+i\stackrel{~}{\gamma }_{\mu ,\alpha }{}_{}{}^{\beta }V_{\mu ,a}^{(N)b}+1_\alpha {}_{}{}^{\beta }S_{a}^{(N)b}+(i\gamma _5\stackrel{~}{\gamma }_\mu )_\alpha {}_{}{}^{\beta }A_{\mu ,a}^{(N)b}]`$
$`\mathrm{\Psi }_A^{(E)B}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(i\gamma _5\gamma _\mu )_\alpha ^\beta {\displaystyle \frac{_\mu }{\sqrt{^2}}}P_{s,a}^{(E)b}+(\sigma _{\mu \nu })_\alpha ^\beta {\displaystyle \frac{_\mu }{\sqrt{^2}}}V_{\nu ,a}^{(E)b}`$ (4.7)
$`+i\gamma _{\mu ,\alpha }{}_{}{}^{\beta }{\displaystyle \frac{_\mu }{\sqrt{^2}}}S_a^{(E)b}+(i\gamma _5\sigma _{\mu \nu })_\alpha {}_{}{}^{\beta }{\displaystyle \frac{_\mu }{\sqrt{^2}}}A_{\nu ,a}^{(E)b}].`$
## 5 Covariant Classification and Spectroscopy of Mesons and Chiral Symmetry
### 5.1 Level classification of ground states
In the previous sections §2 and §3 we have presented a general covariant kinematical framework for describing the (ground states of) composite meson systems with a definite total quark spin. However, what kinds of mesons do really exist or not, that is, the meson spectroscopy, is a dynamical (still unsolved) problem of QCD.
For this problem, it is useful to apply in our scheme a physical consideration of dynamically broken chiral symmetry of QCD, typically displayed in the NJL model: In the (ground state of) heavy quarkonium ($`Q\overline{Q}`$) system both quarks($`Q`$) and antiquarks($`\overline{Q}`$) are possible to do, since $`m_Q>\mathrm{\Lambda }_{\mathrm{conf}}`$, only non-relativistic motions with positive energy, and the non-relativistic $`LS`$-symmetry is good. Accordingly the bi-spinor $`U`$ is considered to be applied to $`Q\overline{Q}`$ system as a covariant spin WF. In the (ground state of) heavy-light quark meson $`Q\overline{q}`$($`q\overline{Q}`$) system the anti-quarks(quarks) make, since $`m_q\mathrm{\Lambda }_{\mathrm{conf}}`$, relativistic motions both with positive and negative energies, and the relativistic chiral symmetry concerning light antiquarks (quarks) is good. Accordingly both the bi-spinors $`U`$ and $`C`$ ($`U`$ and $`D`$) are to be applied to the $`Q\overline{q}`$($`q\overline{Q}`$) system, and in this system there is a possibility of existence of new composite scalar and axial-vector mesons(see Eq.(3.3)). In the (ground state of) light quark $`q\overline{q}`$-meson system both quarks $`q`$ and anti-quarks $`\overline{q}`$ make, since $`m_q\mathrm{\Lambda }_{\mathrm{conf}}`$, relativistic motions with both positive and negative energies, and chiral symmetry is good. Accordingly the linear combinations (specified in §4.3) of bispinors $`U`$ and $`V`$ are applied to the $`q\overline{q}`$-system, and in this system there is a possibility of existence of an extra(, in addition, to a normal) set of composite pseudo-scalar and vector mesons. Furthermore, the linear combinations of $`C`$ and $`D`$ are also applied, and normal and extra sets of composite scalar and axial-vector mesons possibly exist as relativistic $`S`$-wave bound states.
In the above discussion on the validity of chiral symmetry, we have supposed that $`\mathrm{\Lambda }_{\mathrm{conf}}1`$GeV regardless of quark-flavor. Here it may be useful to note that the positive (negative) energy Dirac spinors with momentum $`𝑷`$ for quarks and antiquarks change, respectively, into negative (positive) energy ones with momentum -$`𝑷`$ under the operation of $`\gamma _5`$ as
$`\gamma _5u_+(𝑷)`$ $`=`$ $`u_{}(𝑷),\overline{v}_+(𝑷)\gamma _5=\overline{v}_{}(𝑷).`$ (5.1)
The above expected level structure of the ground states of general light-through-heavy quark meson systems is summarized in Table I.
### 5.2 Level classification of excited states
In the last subsection we have stated on our level classification scheme, focusing on the ground-state mesons. In classifying the excited-state mesons we can proceed essentially similarly as the case of ground state mesons. In (the present extended version of) COQM in the pure-confining force limit, the masses of N-th excited states are given by the formula $`M_N^2=M_G^2+N\mathrm{\Omega }(M_GM_0,\mathrm{\Omega }`$ being the inverse Regge slope), and their covariant spin wave functions are defined by the same formulas as given in §3 with substitution of constituent exciton-quark mass $`m_i`$ by
$`m_i^{}`$ $`=`$ $`\gamma _Nm_i(\gamma _NM_N/M_G).`$ (5.2)
The value of confinement momentum $`\mathrm{\Lambda }_{\mathrm{conf}}`$ is conventionally considered to be
$`\mathrm{\Lambda }_{\mathrm{conf}}`$ $``$ $`1\mathrm{G}\mathrm{e}\mathrm{V}.`$ (5.3)
The value of $`m_i^{}`$ for respective quark-configuration meson systems obtained by the formula (5.2) are given in Table II. From this table we see that $`m^{}\mathrm{\Lambda }_{\mathrm{conf}}`$ for the lower levels, especially for the ground and first-excited states in the light-quark meson systems, and accordingly we can expect that chiral symmetry for these states may be still good. Similarly, in the light/heavy quark meson systems, the chiral symmetry concerning the light quark is also good for the several lower levels.
The quantum numbers $`P`$ and $`C`$ are given for the respective orbitally $`L`$-th excited light-quark mesons as follows:
$`P_s^{(N,E)}{\displaystyle \{L\}}`$ $`P=(1)^{L+1},C=(1)^L`$
$`V_\mu ^{(N,E)}{\displaystyle \{L\}}`$ $`P=(1)^{L+1},C=(1)^{L+1}`$
$`S^{(N,E)}{\displaystyle \{L\}}`$ $`P=(1)^L,C=\pm (1)^L`$
$`A_\mu ^{(N,E)}{\displaystyle \{L\}}`$ $`P=(1)^L,C=\pm (1)^L.`$ (5.4)
### 5.3 Expected spectroscopy of mesons
In our fundamental equations Eqs.(2.1) to (2.3) it was supposed that the squared-mass spectra $`^2`$ is, as a first step in the pure-confining force limit, Dirac-spinor independent and also quark-flavor independent.<sup>\*****)</sup><sup>\*****)</sup>\*****) Of course we consider, separately, the light-quark system, light/heavy-quark system and heavy-quark system. Actually we must take into account the various effects due to one-gluon-exchange potential<sup>\******)</sup><sup>\******)</sup>\******) We must also take the other non-perturbative QCD effects like quark-condensation and instantons. and also the effects due to quark mass difference.
$`\underset{¯}{lightquark(q\overline{q})mesonsystem}`$: In the pure-confining force limit all the ground state mesons expected in §3.1 , $`P_s^{(N,E)},V_\mu ^{(N,E)},S^{(N,E)}`$ and $`A_\mu ^{(N,E)}`$ are degenerate and have the same mass, $`M_0=m_1+m_2`$. Actually the mass of $`P_s^{(N)}`$, to be assigned $`\pi `$-nonet, should be exceptionally low because of its nature<sup>\*******)</sup><sup>\*******)</sup>\*******) The $`q\overline{q}`$-condensation corresponds to the non-zero vacuum expectation value of $`S^{(N)}`$, $`S^{(N)}_0`$. Thus $`P_s^{(N)}`$ other than $`P_s^{(E)}`$ is a Nambu-Goldstone boson. as a Nambu-Goldstone boson. The masses of the $`V_\mu ^{(N)}`$-nonet, which is assigned to be $`\rho `$-meson nonet, are almost equal to the corresponding masses $`M_0=m_1+m_2`$ in the pure-confining force limit. The masses of all the other ground-state mesons are expected to be almost equal to those of the corresponding normal vector mesons, and to be lower than those of the corresponding first-excited states.
For the first-excited states, the chiral symmetry is expected to be still effective, as is seen from Table II given in the last subsection §5.2 , so we expect the existence of a series of the first excited $`P`$-wave states of the ground state multiplets. They are expected to have the masses, which are almost equal to the first excited states of normal vector mesons, and which are lower than the second excited states of those. Among the multiplets newly predicted in the present scheme, to be called “chiralons,” the especially interesting mesons are the ones with $`J^{PC}`$=$`0^+(S^{(E)}(Swave)),1^+(S^{(E)}(Pwave))`$ and $`1^+(A_\mu ^{(E)}(Pwave))`$, which are “exotic particles” out of the conventional non-relativistic $`q\overline{q}`$-mesons. Their masses, by the above mentioned estimate, is expected to be, respectively,
$`m(0^+)`$ $`\stackrel{<}{}`$ $`1.3\mathrm{GeV},`$
$`1.3\mathrm{GeV}`$ $``$ $`m(1^+,S^{(E)})m(1^+,A_\mu ^{(E)})\stackrel{<}{}1.7\mathrm{GeV}`$ (5.5)
$`\underset{¯}{heavylightquark(Q\overline{q}andq\overline{Q})mesonsystem}`$: As is seen from Table I we are able to expect the existence of new multiplets (at least the ground states of scalar and axial-vector triplet).
$`\underset{¯}{heavyquark(Q\overline{Q})mesonsystem}`$: No new multiplets are expected to exist.
## 6 Experimental Evidences and Concluding Remarks
In this paper we have presented a kinematical framework for describing covariantly the ground states as well as excited states of light-through-heavy quark mesons. For light-quark mesons our scheme gives a theoretical basis to classify the composite meson systems unifying the two contrasting viewpoints based on non-relativistic quark model with $`LS`$ symmetry and on NJL model with chiral symmetry. The essential physical assumption is to set up the Klein-Gordon type of Yukawa equation on the bi-local meson wave function with the squared-mass operator, which is, in the pure-confining force limit, independent of Dirac-spinor suffix, and accordingly is chiral symmetric. As a result is pointed out a possibility of existence of rather an abundant new nonets, chiralons, with masses lower than about 2 GeV; several new ground state meson nonets and some new excited meson nonets.
For heavy/light quark meson systems we have similarly pointed out a possibility of existence of new multiplets(triplets), chiralons.
In Table III we have summarized the expected level structure of general ground and the first-excited quark-antiquark mesons. Presently we can give a few experimental candidates for the predicted members of new multiplets: One of the most important and interesting ones is the scalar $`\sigma `$ nonet; the members<sup>\********)</sup><sup>\********)</sup>\********) This assignment of $`\sigma `$ nonet was first proposed in Ref. ? and afterwards in Refs. ?, ?, and by one of the present authors in Ref. ?. Rather strong experimental evidences for $`\sigma `$(600) have been given recently by many authors; from the $`\pi \pi `$ scattering process Refs. ?, ?, ?, ?, ? . and from the $`\pi \pi `$ production processes Ref. ?, ?, ?, ?. Possible evidences for $`\kappa `$(900) was pointed out by reanalyzing the $`K\pi `$ scattering phase shift in Refs. ?, ? and ?. are $`\sigma (600)`$, $`\kappa `$(900), $`a_0(980)`$ and $`f_0(980)`$, which constitute, with the members of $`\pi `$-nonet, a linear representation of the chiral $`SU(3)`$ symmetry. It is notable that the $`\sigma `$ nonet is the relativistic $`S`$-wave states, which should be discriminated from the non-relativistic $`{}_{}{}^{3}P_{0}^{}`$ states.
Another example suggesting possible validity of the present scheme is existence of the three pseudoscalars with mass between 1 GeV$``$1.5GeV, $`\eta `$(1295),
$`\eta `$(1420) and $`\eta `$(1460). The two out of them may be the members of the radially excited $`\pi `$-nonet, while the one extra may belong to the ground states of the extra pseudoscalar nonet newly predicted.
Also we have the other candidates for chiralons: It was a problem for experimental groups for long time whether a possible resonance observed in the $`\eta \pi `$ system, with an exotic quantum number $`J^{PC}=1^+`$ and with a mass around 1.5 GeV, really exists or does not. Recently it seems that the existence of two such particles $`\pi _1(1400)`$ and $`\pi _1(1600)`$ have been accepted widely. These two particles have the mass in the region estimated in Eq.(5.5) to be assigned as the respective excited $`P`$-wave states of $`S^{(E)}`$ and $`A_\mu ^{(E)}`$.
We have the other longstanding problem in hadron spectroscopy: the mass and width of $`a_1(1260)`$ seem to be variant depending on the production process and/or decay channel.
In connection to this problem we have made recently a preliminary analysis of the data obtained by GAMS group WA102 experiment on process $`\pi ^{}p3\pi ^0n`$. As a result we have obtained an evidence of existence of two $`a_1(J^{PC}=1^{++})`$ particles: $`a_1^c(m=0.9\mathrm{GeV},\mathrm{\Gamma }=200\mathrm{M}\mathrm{e}\mathrm{V})`$, and $`a_1^N(m=1.2\mathrm{GeV},\mathrm{\Gamma }=440\mathrm{M}\mathrm{e}\mathrm{V})`$. The former may be assigned to be $`A_\mu ^{(N)}(^3S_1)`$, while the latter be conventional $`a_1`$ particle (to be $`V_\mu ^{(N)}(^3P_1)`$ in our classification scheme).
All the above candidates for chiralons are concerned with the light-quark mesons. Here, we should like to refer to a preliminary result concerning the heavy/light quark meson systems that a scalar chiralons $`B_0^c`$ with $`M`$=5.52 GeV and $`\mathrm{\Gamma }`$=44 MeV may be observed in the $`B\pi `$ channel produced through the $`Z`$-boson decay.
We have also given the above mentioned experimental candidates for chiralons in Table III, where we have listed, for reference, also our assignment to the conventional $`S`$-wave and $`P`$-wave states.
In concluding we should like to remark that in this paper we have dared to present a very “brave” attempt for unified classification scheme of mesons, which predicts the possible existence of a lot of new meson multiplets. Further serious investigations and search for them will be required to test validity of the present scheme.
## Acknowledgements
The authors would like to express their sincere gratitude to prof. K. Takamatsu, T. Tsuru and T. Sawada for encouragements and useful informations. They would like to thank prof. K. Yamada for useful comments. They are also grateful to Dr. T. Ishida for encouragements and comments.
## A Spin wave functions of composite meson systems <br>and fundamental crossing rules for constituent quarks
In the free local field its annihilation (positive frequency) and creation (negative frequency) parts are related with each other by the conventional crossing rule. In this appendix we shall derive a similar crossing rule for our covariant spin WF of composite meson systems by applying the fundamental crossing rule of our extended Dirac spinors for constituent quarks. In the following we first make the annihilation part of the composite meson WF by decomposition of total spin of the constituent quarks and antiquarks. Next by using the fundamental crossing rules for the constituent spinors we construct the creation part of the composite meson spinor WF.
We use the following conventional “free” Dirac $`u`$ and $`v`$ spinors:
$`u(𝒑,h)`$ $`=`$ $`\left(\begin{array}{c}\sqrt{E+m}\chi ^{(h)}\hfill \\ \sqrt{Em}𝒏𝝈\chi ^{(h)}\hfill \end{array}\right),v(𝒑,h)=\left(\begin{array}{c}\sqrt{Em}𝒏𝝈\chi ^{(h)}\hfill \\ \sqrt{E+m}\chi ^{(h)}\hfill \end{array}\right).`$ (A.5)
In this Appendix A we choose the $`z`$-axis parallel to the momentum of composite mesons as $`𝒑=p𝒏=p\widehat{𝒛}`$, and accordingly
$`\chi ^{(+)}`$ $`=`$ $`\left(\begin{array}{c}1\\ 0\end{array}\right),\chi ^{()}=\left(\begin{array}{c}0\\ 1\end{array}\right),\chi ^{(h)}i\sigma _2\chi ^{(h)},\begin{array}{c}\chi ^{(+)}\\ \chi ^{()}\end{array}\begin{array}{c}=\\ =\end{array}\begin{array}{c}\hfill \chi ^{()}\\ \hfill \chi ^{(+)}\end{array}.`$ (A.16)
### A.1 quark and antiquark spinor inside of mesons
Here we give the annihilation part of the constituent quark and antiquark spinor WF. As was explained in the text the constituent quark inside of mesons has the freedom of positive and negative energy, as well as that of up and down spin, and totally four degrees of freedom.
The positive energy quark spinor $`u_+`$ with the meson momentum $`𝒑`$, which is related with the “free” Dirac $`u`$ spinor, is defined by
$`u_+(𝒑,\pm )`$ $`=`$ $`\left(\begin{array}{c}\sqrt{E+m}\chi ^{(\pm )}\hfill \\ \sqrt{Em}𝒏𝝈\chi ^{(\pm )}\hfill \end{array}\right)=u(𝒑,\pm ).`$ (A.19)
The negative energy quark spinor $`u_{}`$ with the momentum, $`𝒑`$, anti-parallel to the meson momentum $`𝒑`$, is defined by using the same spin operator $`𝒏𝝈`$ as for $`u_+`$. They are related with free “Dirac” $`v`$ spinors.
$`u_{}(𝒑,\pm )`$ $`=`$ $`\left(\begin{array}{c}\sqrt{Em}𝒏𝝈\chi ^{(\pm )}\hfill \\ \sqrt{E+m}\chi ^{(\pm )}\hfill \end{array}\right)=v(𝒑,).`$ (A.22)
The positive energy antiquark spinor $`\overline{v}_+`$ with momentum $`𝒑`$, which is related with the “free” Dirac $`\overline{v}`$ spinor, is defined by
$`\overline{v}_+(𝒑,\pm )`$ $`=`$ $`(\sqrt{Em}\chi ^{(\pm )}𝒏𝝈,\sqrt{E+m}\chi ^{(\pm )})=\overline{v}(𝒑,\pm ).`$ (A.23)
The negative energy antiquark spinor $`\overline{v}_{}`$ with momentum $`𝒑`$, is defined, similarly by using the same spin operator $`𝒏𝝈`$, as in the case of $`\overline{v}_+`$. They are related with free “Dirac” $`\overline{u}`$ spinors.
$`\overline{v}_{}(𝒑,\pm )`$ $`=`$ $`(\sqrt{E+m}\chi ^{(\pm )},\sqrt{Em}\chi ^{(\pm )}𝒏𝝈)=\pm \overline{u}(𝒑,).`$ (A.24)
### A.2 Fundamental crossing rules for constituent spinor inside of mesons
Under the crossing operation for mesons the annhilation (positive frequency) part of meson WF is transformed into the creation (negative frequency) part of that. The corresponding fundamental crossing rule is set up as follows: For example, the positive energy constituent quark spinor in the annihilation part of meson WF is transformed into the positive energy antiquark spinor in the the creation part of meson WF.
$`u_+(𝒑,s)`$ $``$ $`v_+(𝒑,s)\mathrm{that}\mathrm{is}u(𝒑,\pm )v(𝒑,\pm ).`$ (A.25)
The other kinds of constituent quark spinors are transformed similarly as
$`u_{}(𝒑,s)`$ $``$ $`v_{}(𝒑,s)\mathrm{that}\mathrm{is}v(𝒑,)\pm u(𝒑,),`$ (A.26)
$`\overline{v}_+(𝒑,s)`$ $``$ $`\overline{u}_+(𝒑,s)\mathrm{that}\mathrm{is}\overline{v}(𝒑,\pm )\overline{u}(𝒑,\pm ),`$ (A.27)
$`\overline{v}_{}(𝒑,s)`$ $``$ $`\overline{u}_{}(𝒑,s)\mathrm{that}\mathrm{is}\pm \overline{u}(𝒑,)\overline{v}(𝒑,).`$ (A.28)
### A.3 Annihilation part of the composite meson spinor WF
As was explained in the text, depending upon the energy sign of the constituent spinor WF, there are four different types of (the annihilation part of) bi-spinor WF, $`U^{(+)}`$, $`C^{(+)}`$, $`D^{(+)}`$ and $`V^{(+)}`$. Each type of the WF is decomposed into the irreducible components, which are constructed by using the usual spin composition.
The $`U^{(+)}u_+\overline{v}_+(u\overline{v})`$ is decomposed into the pseudoscalar and vector components, which are constructed from the constituent spinors as
$`U_{P_s}^{(+)}`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(u_+(𝒑,+)\overline{v}_+(𝒑,)u_+(𝒑,)\overline{v}_+(𝒑,+))`$
$`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(u(𝒑,+)\overline{v}(𝒑,)u(𝒑,)\overline{v}(𝒑,+))={\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5(1+iv\gamma )`$
$`U_{V_\mu ^{(\pm ,0)}}^{(+)}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}u_+(𝒑,+)\overline{v}_+(𝒑,+)\hfill \\ u_+(𝒑,)\overline{v}_+(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(u_+(𝒑,+)\overline{v}_+(𝒑,)+u_+(𝒑,)\overline{v}_+(𝒑,+))\hfill \end{array}\right]`$ (A.33)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}u(𝒑,+)\overline{v}(𝒑,+)\hfill \\ u(𝒑,)\overline{v}(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(u(𝒑,+)\overline{v}(𝒑,)+u(𝒑,)\overline{v}(𝒑,+))\hfill \end{array}\right]`$
$`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}(1iv\gamma )iϵ^{(\pm ,0)}\gamma (1+iv\gamma )={\displaystyle \frac{1}{2\sqrt{2}}}iϵ^{(\pm ,0)}\gamma (1+iv\gamma ).`$
The $`C^{(+)}u_+\overline{v}_{}(u\overline{u})`$ is decomposed into the scalar and axialvector components, which are constructed from the constituent spinors as
$`C_S^{(+)}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(u_+(𝒑,+)\overline{v}_{}(𝒑,)u_+(𝒑,)\overline{v}_{}(𝒑,+))`$
$`=`$ $`{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(u(𝒑,+)\overline{u}(𝒑,+)u(𝒑,)\overline{u}(𝒑,))={\displaystyle \frac{1}{2\sqrt{2}}}(1iv\gamma )`$
$`C_{A_\mu ^{(\pm ,0)}}^{(+)}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}u_+(𝒑,+)\overline{v}_{}(𝒑,+)\hfill \\ u_+(𝒑,)\overline{v}_{}(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(u_+(𝒑,+)\overline{v}_{}(𝒑,)+u_+(𝒑,)\overline{v}_{}(𝒑,+))\hfill \end{array}\right]`$ (A.42)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}u(𝒑,+)\overline{u}(𝒑,)\hfill \\ u(𝒑,)\overline{u}(𝒑,+)\hfill \\ \frac{1}{\sqrt{2}}(u(𝒑,+)\overline{u}(𝒑,+)+u(𝒑,)\overline{u}(𝒑,))\hfill \end{array}\right]`$
$`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}(1iv\gamma )i\gamma _5ϵ^{(\pm ,0)}\gamma (1iv\gamma )={\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5ϵ^{(\pm ,0)}\gamma (1iv\gamma ).`$
The $`D^{(+)}u_{}\overline{v}_+(v\overline{v})`$ is decomposed into the scalar and axialvector components, which are constructed from the constituent spinors as
$`D_S^{(+)}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(u_{}(𝒑,+)\overline{v}_+(𝒑,)u_{}(𝒑,)\overline{v}_+(𝒑,+))`$
$`=`$ $`{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(v(𝒑,)\overline{v}(𝒑,)v(𝒑,+)\overline{v}(𝒑,+))={\displaystyle \frac{1}{2\sqrt{2}}}(1+iv\gamma )`$
$`D_{A_\mu ^{(\pm ,0)}}^{(+)}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}u_{}(𝒑,+)\overline{v}_+(𝒑,+)\hfill \\ u_{}(𝒑,)\overline{v}_+(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(u_{}(𝒑,+)\overline{v}_+(𝒑,)+u_{}(𝒑,)\overline{v}_+(𝒑,+))\hfill \end{array}\right]`$ (A.51)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}v(𝒑,)\overline{v}(𝒑,+)\hfill \\ v(𝒑,+)\overline{v}(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(v(𝒑,)\overline{v}(𝒑,)+v(𝒑,+)\overline{v}(𝒑,+))\hfill \end{array}\right]`$
$`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}(1+iv\gamma )i\gamma _5ϵ^{(\pm ,0)}\gamma (1+iv\gamma )={\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5ϵ^{(\pm ,0)}\gamma (1+iv\gamma ).`$
The $`V^{(+)}u_{}\overline{v}_{}(v\overline{u})`$ is decomposed into the pseudoscalar and vector components, which are constructed from the constituent spinors as
$`V_{P_s}^{(+)}`$ $`=`$ $`i{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(u_{}(𝒑,+)\overline{v}_{}(𝒑,)u_{}(𝒑,)\overline{v}_{}(𝒑,+))`$
$`=`$ $`i{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(v(𝒑,)\overline{u}(𝒑,+)v(𝒑,+)\overline{u}(𝒑,))={\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5(1iv\gamma )`$
$`V_{V_\mu ^{(\pm ,0)}}^{(+)}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}u_{}(𝒑,+)\overline{v}_{}(𝒑,+)\hfill \\ u_{}(𝒑,)\overline{v}_{}(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(u_{}(𝒑,+)\overline{v}_{}(𝒑,)+u_{}(𝒑,)\overline{v}_{}(𝒑,+))\hfill \end{array}\right]`$ (A.60)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}v(𝒑,)\overline{u}(𝒑,)\hfill \\ v(𝒑,+)\overline{u}(𝒑,+)\hfill \\ \frac{1}{\sqrt{2}}(v(𝒑,)\overline{u}(𝒑,+)+v(𝒑,+)\overline{u}(𝒑,))\hfill \end{array}\right]`$
$`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}(1+iv\gamma )iϵ^{(\pm ,0)}\gamma (1iv\gamma )={\displaystyle \frac{1}{2\sqrt{2}}}iϵ^{(\pm ,0)}\gamma (1iv\gamma ).`$
### A.4 Creation part of the composite meson spinor WF
Creation part of composite meson spinor WF is obtained from the annihilation part, by applying the fundamental crossing rule Eqs.(A.25), (A.26), (A.27) and (A.28), as follows:
For $`U^{()}v_+\overline{u}_+(v\overline{u})`$
$`U_{P_s}^{()}`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(v_+(𝒑,+)\overline{u}_+(𝒑,)v_+(𝒑,)\overline{u}_+(𝒑,+))`$
$`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(v(𝒑,+)\overline{u}(𝒑,)v(𝒑,)\overline{u}(𝒑,+))={\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5(1iv\gamma )`$
$`U_{V_\mu ^{(\pm ,0)}}^{()}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}v_+(𝒑,+)\overline{u}_+(𝒑,+)\hfill \\ v_+(𝒑,)\overline{u}_+(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(v_+(𝒑,+)\overline{u}_+(𝒑,)+v_+(𝒑,)\overline{u}_+(𝒑,+))\hfill \end{array}\right]`$ (A.69)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}v(𝒑,+)\overline{u}(𝒑,+)\hfill \\ v(𝒑,)\overline{u}(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(v(𝒑,+)\overline{u}(𝒑,)+v(𝒑,)\overline{u}(𝒑,+))\hfill \end{array}\right]`$
$`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}(1+iv\gamma )i\stackrel{~}{ϵ}^{(\pm ,0)}\gamma (1iv\gamma )={\displaystyle \frac{1}{2\sqrt{2}}}i\stackrel{~}{ϵ}^{(\pm ,0)}\gamma (1iv\gamma ).`$
For $`C^{()}v_+\overline{u}_{}(v\overline{v})`$
$`C_S^{()}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(v_+(𝒑,+)\overline{u}_{}(𝒑,)v_+(𝒑,)\overline{u}_{}(𝒑,+))`$
$`=`$ $`{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(v(𝒑,+)\overline{v}(𝒑,+)+v(𝒑,)\overline{v}(𝒑,))={\displaystyle \frac{1}{2\sqrt{2}}}(1+iv\gamma )`$
$`C_{A_\mu ^{(\pm ,0)}}^{()}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}v_+(𝒑,+)\overline{u}_{}(𝒑,+)\hfill \\ v_+(𝒑,)\overline{u}_{}(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(v_+(𝒑,+)\overline{u}_{}(𝒑,)+v_+(𝒑,)\overline{u}_{}(𝒑,+))\hfill \end{array}\right]`$ (A.78)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}v(𝒑,+)\overline{v}(𝒑,)\hfill \\ v(𝒑,)\overline{v}(𝒑,+)\hfill \\ \frac{1}{\sqrt{2}}(v(𝒑,+)\overline{v}(𝒑,+)v(𝒑,)\overline{v}(𝒑,))\hfill \end{array}\right]`$
$`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}(1+iv\gamma )i\gamma _5\stackrel{~}{ϵ}^{(\pm ,0)}\gamma (1+iv\gamma )={\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5\stackrel{~}{ϵ}^{(\pm ,0)}\gamma (1+iv\gamma ).`$
For $`D^{()}v_{}\overline{u}_+(u\overline{u})`$
$`D_S^{()}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(v_{}(𝒑,+)\overline{u}_+(𝒑,)v_{}(𝒑,)\overline{u}_+(𝒑,+))`$
$`=`$ $`{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(u(𝒑,)\overline{u}(𝒑,)+u(𝒑,+)\overline{u}(𝒑,+))={\displaystyle \frac{1}{2\sqrt{2}}}(1iv\gamma )`$
$`D_{A_\mu ^{(\pm ,0)}}^{()}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}v_{}(𝒑,+)\overline{u}_+(𝒑,+)\hfill \\ v_{}(𝒑,)\overline{u}_+(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(v_{}(𝒑,+)\overline{u}_+(𝒑,)+v_{}(𝒑,)\overline{u}_+(𝒑,+))\hfill \end{array}\right]`$ (A.87)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}u(𝒑,)\overline{u}(𝒑,+)\hfill \\ u(𝒑,+)\overline{u}(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(u(𝒑,)\overline{u}(𝒑,)u(𝒑,+)\overline{u}(𝒑,+))\hfill \end{array}\right]`$
$`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}(1iv\gamma )i\gamma _5\stackrel{~}{ϵ}^{(\pm ,0)}\gamma (1iv\gamma )={\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5\stackrel{~}{ϵ}^{(\pm ,0)}\gamma (1iv\gamma ).`$
For $`V^{()}v_{}\overline{u}_{}(u\overline{v})`$
$`V_{P_s}^{()}`$ $`=`$ $`i{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(v_{}(𝒑,+)\overline{u}_{}(𝒑,)v_{}(𝒑,)\overline{u}_{}(𝒑,+))`$
$`=`$ $`i{\displaystyle \frac{1}{2m}}{\displaystyle \frac{1}{\sqrt{2}}}(u(𝒑,)\overline{v}(𝒑,+)u(𝒑,+)\overline{v}(𝒑,))={\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5(1+iv\gamma )`$
$`V_{V_\mu ^{(\pm ,0)}}^{()}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}v_{}(𝒑,+)\overline{u}_{}(𝒑,+)\hfill \\ v_{}(𝒑,)\overline{u}_{}(𝒑,)\hfill \\ \frac{1}{\sqrt{2}}(v_{}(𝒑,+)\overline{u}_{}(𝒑,)+v_{}(𝒑,)\overline{u}_{}(𝒑,+))\hfill \end{array}\right]`$ (A.96)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left[\begin{array}{c}u(𝒑,)\overline{v}(𝒑,)\hfill \\ u(𝒑,+)\overline{v}(𝒑,+)\hfill \\ \frac{1}{\sqrt{2}}(u(𝒑,)\overline{v}(𝒑,+)+u(𝒑,+)\overline{v}(𝒑,))\hfill \end{array}\right]`$
$`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}(1iv\gamma )i\stackrel{~}{ϵ}^{(\pm ,0)}\gamma (1+iv\gamma )={\displaystyle \frac{1}{2\sqrt{2}}}i\stackrel{~}{ϵ}^{(\pm ,0)}\gamma (1+iv\gamma ).`$
### A.5 Representation by local meson field
We can simply combine the annihilation and creation parts of the composite spinor WF into the local meson field(, neglecting the freedom of internal space-time $`x`$).
For non-relativistic type spinors
$`{\displaystyle \frac{1}{2\sqrt{2}}}`$ $`i`$ $`\gamma _5(1+\gamma _\nu {\displaystyle \frac{_\nu }{\sqrt{^2}}})P_s^{(NR)}(X),{\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _\mu (1+\gamma _\nu {\displaystyle \frac{_\nu }{\sqrt{^2}}})V_\mu ^{(NR)}(X),`$ (A.101)
where $`P_s^{(NR)}(X)`$ and $`V_\mu ^{(NR)}(X)`$ represents the local pseudoscalar and vector meson field operators of non-relativistic-type, respectively.
For relativistic $`\overline{q}`$-type spinors
$`{\displaystyle \frac{1}{2\sqrt{2}}}`$ $`(`$ $`1\gamma _\nu {\displaystyle \frac{_\nu }{\sqrt{^2}}})S^{\overline{q}}(X),{\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5\gamma _\mu (1\gamma _\nu {\displaystyle \frac{_\nu }{\sqrt{^2}}})A_\mu ^{\overline{q}}(X),`$ (A.102)
where $`S^{\overline{q}}(X)`$ and $`A_\mu ^{\overline{q}}(X)`$ represents the local pseudoscalar and vector meson field operators of $`\overline{q}`$-type, respectively.
For relativistic $`q`$-type spinors
$`{\displaystyle \frac{1}{2\sqrt{2}}}`$ $`(`$ $`1+\gamma _\nu {\displaystyle \frac{_\nu }{\sqrt{^2}}})S^q(X),{\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _5\gamma _\mu (1+\gamma _\nu {\displaystyle \frac{_\nu }{\sqrt{^2}}})A_\mu ^q(X),`$ (A.103)
where $`S^q(X)`$ and $`A_\mu ^q(X)`$ represents the local pseudoscalar and vector meson field operators of $`D`$-type, respectively.
For extremely relativistic-type spinors
$`{\displaystyle \frac{1}{2\sqrt{2}}}`$ $`i`$ $`\gamma _5(1\gamma _\nu {\displaystyle \frac{_\nu }{\sqrt{^2}}})P_s^{(ER)}(X),{\displaystyle \frac{1}{2\sqrt{2}}}i\gamma _\mu (1\gamma _\nu {\displaystyle \frac{_\nu }{\sqrt{^2}}})V_\mu ^{(ER)}(X),`$ (A.104)
where $`P_s^{(ER)}(X)`$ and $`V_\mu ^{(ER)}(X)`$ represents the local pseudoscalar and vector meson field operators of extremely relativistic-type, respectively.
These local meson field representation of the composite quark spinor WF guarantees the crossing symmetry of the interaction among composite mesons, which are induced from the crossing symmetric interaction between constituent quarks and antiquarks.
## B orthonormality relation for spinor WF
### B.1 Bargmann-Wigner (BW) bases
The Pauli-conjugate of the annihilation parts $`W_i^{(+)}`$ of WF, $`\overline{W}_i^{()}(\gamma _4W_i^{(+)}\gamma _4)`$, are related with the creation parts as
$`(\overline{W}_i^{()})`$ $`=`$ $`(\overline{U}_{P_s}^{()},\overline{U}_{V_\mu }^{()},\overline{C}_S^{()},\overline{C}_{A_\mu }^{()},\overline{D}_S^{()},\overline{D}_{A_\mu }^{()},\overline{V}_{P_s}^{()},\overline{V}_{V_\mu }^{()})`$ (B.1)
$`=`$ $`(U_{P_s}^{()},U_{V_\mu }^{()},D_S^{()},D_{A_\mu }^{()},C_S^{()},C_{A_\mu }^{()},V_{P_s}^{()},V_{V_\mu }^{()}).`$
They are orthonormal to the annihilation WF $`W_i^{(+)}`$ as
$`\overline{W}_i^{()}W_j^{(+)}`$ $`=`$ $`ϵ_i\delta _{ij}`$ (B.2)
$`(ϵ_{P_s^{(NR)}},ϵ_{V_\mu ^{(NR)}},ϵ_{S^{(\overline{q})}},`$ $`ϵ_{A_\mu ^{(\overline{q})}},ϵ_{S^{(q)}},ϵ_{A_\mu ^{(q)}},ϵ_{P_s^{(ER)}},ϵ_{V_\mu ^{(ER)}})=(1,1,1,1,1,1,1,1).`$
By using this orthonormality relation we can decompose the general spinor WF $`\mathrm{\Psi }^{(+)}`$ into the meson components, $`(\varphi _i^{(+)})=(P_s^{(+)(NR,ER)}`$, $`S^{(+)(\overline{q},q)}`$, $`V_\mu ^{(+)(NR,ER)}`$, $`A_\mu ^{(+)(\overline{q},q)})`$.
$`\mathrm{\Psi }^{(+)}`$ $`=`$ $`{\displaystyle \underset{i}{}}W_i^{(+)}\varphi _i^{(+)}.\varphi _i^{(+)}=ϵ_i\overline{W}_i^{()}\mathrm{\Psi }^{(+)}`$ (B.3)
### B.2 chiral bases—light quark $`q\overline{q}`$-meson systems
For description of the light quark $`q\overline{q}`$-meson systems the chiral $`(N,E)`$ bases of spinor WF are expected to be more effective. They are obtained by the linear combination of BW bases as explained in the text.
The annihilation WF $`W_i^{(+)}`$ and creation WF $`W_i^{()}`$ are given, respectively, by
$`(W_i^{(+)}(P))`$ $`=`$ $`(U_{P_s}^{(+)(N)},C_S^{(+)(N)},U_{V_\mu }^{(+)(N)},C_{A_\mu }^{(+)(N)};U_{P_s}^{(+)(E)},C_S^{(+)(N)},U_{V_\mu }^{(+)(E)},C_{A_\mu }^{(+)(E)})`$
$`=`$ $`{\displaystyle \frac{1}{2}}(i\gamma _5,1,i\stackrel{~}{\gamma }_\mu ,i\gamma _5\stackrel{~}{\gamma }_\mu ;\gamma _5v\gamma ,v\gamma ,i\sigma _{\mu \nu }v_\nu ,\gamma _5\sigma _{\mu \nu }v_\nu )`$
$`(W_i^{()}(P))`$ $`=`$ $`(U_{P_s}^{()(N)},C_S^{()(N)},U_{V_\mu }^{()(N)},C_{A_\mu }^{()(N)};U_{P_s}^{()(E)},C_S^{()(N)},U_{V_\mu }^{()(E)},C_{A_\mu }^{()(E)})`$ (B.4)
$`=`$ $`{\displaystyle \frac{1}{2}}(i\gamma _5,1,i\stackrel{~}{\gamma }_\mu ,i\gamma _5\stackrel{~}{\gamma }_\mu ;\gamma _5v\gamma ,v\gamma ,i\sigma _{\mu \nu }v_\nu ,\gamma _5\sigma _{\mu \nu }v_\nu ),`$
which are related through crossing rules of constituent spinors with each other.
The Pauli-conjugate of the creation spinor $`\overline{W}_i^{()}(\gamma _4W_i^{(+)}\gamma _4)`$ are equal to the annihilation $`W_i^{()}`$.
$`\overline{W}_i^{()}(P)`$ $`(`$ $`\gamma _4W_i^{(+)}\gamma _4)=W_i^{()}(P)`$ (B.5)
$`\mathrm{for}`$ $`\varphi _i=P_s^{(N)},S^{(N)},V_\mu ^{(N)},A_\mu ^{(N)};P_s^{(E)},S^{(N)},V_\mu ^{(E)},A_\mu ^{(E)}.`$
They satisfy orthonormality relation:
$`\overline{W}_i^{()}W_j^{(+)}`$ $`=`$ $`ϵ_i\delta _{ij}`$
$`(ϵ_i)`$ $`=`$ $`(ϵ_{P_s^{(N)}},ϵ_{S^{(N)}},ϵ_{V_\mu ^{(N)}},ϵ_{A_\mu ^{(N)}};ϵ_{P_s^{(E)}},ϵ_{S^{(N)}},ϵ_{V_\mu ^{(E)}},ϵ_{A_\mu ^{(E)}})`$ (B.6)
$`=`$ $`(1,1,1,1;1,1,1,1)`$
By using this orthonormality relation the general WF $`\mathrm{\Psi }^{(+)}`$ are decomposed into the meson components $`(\varphi _i^{(+)})=(P_s^{(+)(N,E)}`$, $`S^{(+)(N,E)}`$, $`V_\mu ^{(+)(N,E)}`$, $`A_\mu ^{(+)(N,E)})`$.
$`\mathrm{\Psi }^{(+)}`$ $`=`$ $`{\displaystyle \underset{i}{}}W_i^{(+)}\varphi _i^{(+)}.\varphi _i^{(+)}=ϵ_i\overline{W}_i^{()}\mathrm{\Psi }^{(+)}`$ (B.7)
## C Chiral transformation for spinor WF of light quark $`q\overline{q}`$ mesons
$`SU(3)`$ chiral transformation for the annihilation part of spinor WF of light quark $`q\overline{q}`$ mesons is given by
$`\mathrm{\Psi }_A{}_{}{}^{(+)B}(P,x)`$ $``$ $`[e^{i\frac{\alpha ^i\lambda ^i}{2}\gamma _5}\mathrm{\Psi }^{(+)}(P,x)e^{i\frac{\alpha ^i\lambda ^i}{2}\gamma _5}]_A{}_{}{}^{B}.`$ (C.1)
For the infinitesimal transformation $`\{i\frac{\alpha ^i\lambda ^i}{2}\gamma _5,\mathrm{\Psi }^{(+)}(P,x)\}`$ the respective meson spinor WF are transformed as
$`P_s^{(N)}U_{P_s}^{(N)}`$ $``$ $`\{{\displaystyle \frac{\alpha ^i\lambda ^i}{2}},P_s^{(N)}\}{\displaystyle \frac{1}{2}}\{i\gamma _5,U_{P_s}^{(N)}\}=d^{ijk}\alpha ^iP_s^{(N)j}{\displaystyle \frac{\lambda ^k}{2}}C_S^{(N)}`$
$`S^{(N)}C_S^{(N)}`$ $``$ $`\{{\displaystyle \frac{\alpha ^i\lambda ^i}{2}},S^{(N)}\}{\displaystyle \frac{1}{2}}\{i\gamma _5,C_S^{(N)}\}=d^{ijk}\alpha ^iS^{(N)j}{\displaystyle \frac{\lambda ^k}{2}}U_{P_s}^{(N)}`$
$`V_\mu ^{(N)}U_{V_\mu }^{(N)}`$ $``$ $`[{\displaystyle \frac{\alpha ^i\lambda ^i}{2}},V_\mu ^{(N)}]{\displaystyle \frac{1}{2}}[i\gamma _5,U_{V_\mu }^{(N)}]=f^{ijk}\alpha ^iV_\mu ^{(N)j}{\displaystyle \frac{\lambda ^k}{2}}C_{A_\mu }^{(N)}`$
$`A_\mu ^{(N)}C_{A_\mu }^{(N)}`$ $``$ $`[{\displaystyle \frac{\alpha ^i\lambda ^i}{2}},A_\mu ^{(N)}]{\displaystyle \frac{1}{2}}[i\gamma _5,C_{A_\mu }^{(N)}]=f^{ijk}\alpha ^iA_\mu ^{(N)j}{\displaystyle \frac{\lambda ^k}{2}}U_{V_\mu }^{(N)}`$
$`P_s^{(E)}U_{P_s}^{(E)}`$ $``$ $`[{\displaystyle \frac{\alpha ^i\lambda ^i}{2}},P_s^{(E)}]{\displaystyle \frac{1}{2}}[i\gamma _5,U_{P_s}^{(E)}]=f^{ijk}\alpha ^iP_s^{(E)j}{\displaystyle \frac{\lambda ^k}{2}}C_S^{(E)}`$
$`S^{(E)}C_S^{(E)}`$ $``$ $`[{\displaystyle \frac{\alpha ^i\lambda ^i}{2}},S^{(E)}]{\displaystyle \frac{1}{2}}[i\gamma _5,C_S^{(E)}]=f^{ijk}\alpha ^iS^{(E)j}{\displaystyle \frac{\lambda ^k}{2}}U_{P_s}^{(E)}`$
$`V_\mu ^{(E)}U_{V_\mu }^{(E)}`$ $``$ $`\{{\displaystyle \frac{\alpha ^i\lambda ^i}{2}},V_\mu ^{(E)}\}{\displaystyle \frac{1}{2}}\{i\gamma _5,U_{V_\mu }^{(E)}\}=d^{ijk}\alpha ^iV_\mu ^{(E)j}{\displaystyle \frac{\lambda ^k}{2}}C_{A_\mu }^{(E)}`$
$`A_\mu ^{(E)}C_{A_\mu }^{(E)}`$ $``$ $`\{{\displaystyle \frac{\alpha ^i\lambda ^i}{2}},A_\mu ^{(E)}\}{\displaystyle \frac{1}{2}}\{i\gamma _5,C_{A_\mu }^{(E)}\}=d^{ijk}\alpha ^iA_\mu ^{(E)j}{\displaystyle \frac{\lambda ^k}{2}}U_{V_\mu }^{(E)}.`$ (C.2)
For the finite $`U(1)`$ chiral transformation
$`\mathrm{\Psi }_A{}_{}{}^{(+)B}(P,x)`$ $``$ $`[e^{i\alpha \gamma _5}\mathrm{\Psi }^{(+)}(P,x)e^{i\alpha \gamma _5}]_A{}_{}{}^{B},`$ (C.3)
the spinor WF for $`V_\mu ^{(N)},A_\mu ^{(N)},P_s^{(E)}`$ and $`S^{(E)}`$ are invariant, while
$`P_s^{(N)}U_{P_s}^{(N)}`$ $``$ $`P_s^{(N)}(U_{P_s}^{(N)}\mathrm{cos}2\alpha C_s^{(N)}\mathrm{sin}2\alpha ),`$
$`S^{(N)}C_S^{(N)}`$ $``$ $`S^{(N)}(C_s^{(N)}\mathrm{cos}2\alpha +U_{P_s}^{(N)}\mathrm{sin}2\alpha ),`$
$`V_\mu ^{(E)}U_{V_\mu }^{(E)}`$ $``$ $`V_\mu ^{(E)}(U_{V_\mu }^{(E)}\mathrm{cos}2\alpha +C_{A_\mu }^{(E)}\mathrm{sin}2\alpha ),`$
$`A_\mu (E)C_{A_\mu }^{(E)}`$ $``$ $`A_\mu ^{(E)}(C_{A_\mu }^{(E)}\mathrm{cos}2\alpha U_{V_\mu }^{(E)}\mathrm{sin}2\alpha ).`$ (C.4)
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# Cosmic crystallography in a circle
## 1 Introduction
Cosmic crystallography (CC) is a method to unveil the topology of the universe, and initially looked for spikes in a pair separation histogram (PSH) . Since spikes are absent in hyperbolic spaces, it appeared that the method was useless in such spaces. However, it was soon shown that not only a Clifford translation (responsible for a spike) press its fingerprint on a PSH, but also the other isometries of the space .
When spikes are absent, the PSH of a ball containing repeated images – the $`\varphi ^m(l)`$ – is very similar to that of a ball with same radius and same geometry, but without duplication of images – the $`\varphi ^s(l)`$. A suggestion was then made, of studying the difference of the $`m`$ultiply and the $`s`$imply connected histograms, $`\varphi ^m(l)\varphi ^s(l)`$ .
To improve the method, $`exp`$ected functions $`\varphi _{exp}^s(l)`$ were derived to replace the histograms $`\varphi ^s(l)`$ obtained from computer simulations, for all three geometries with constant curvature . Graphs of $`\varphi ^m(l)\varphi _{exp}^s(l)`$ were obtained, clearly evincing the topology of an euclidian, an elliptic, and a hyperbolic three-space . The contribution of each individual isometry $`g`$ to a PSH was examined, and normalized histograms $`\varphi ^g(l)`$ (defined in ref.) were obtained from computer simulations ; these simulations also gave histograms of $`\varphi ^u(l)\varphi _{exp}^s(l)`$, a previously unsuspected quantity .
Recently the exact (noiseless) functions $`\varphi _{exp}^g(l)`$ were given for the euclidian isometries . In the present report we finally have a first acquaintance with functions $`\varphi _{exp}^u(l)`$, the exact (noiseless) counterparts of the ’uncorrelated’ normalized histograms $`\varphi ^u(l)`$ defined in . We examine a one-dimensional system: a universe with topology $`S^1`$, a circle with circumference 1; we assume the horizon at a distance $`L/2`$ on each side of an observer, so the visible universe has total length $`L`$; clearly if $`L>1`$ then there are repeated images in this visible universe. In section 2 we give a detailed description of how to obtain the expected uncorrelated signature $`\phi _{Lexp}^u(l)`$ when $`1<L<2`$. In section 3 we exhibit the generalization for arbitrary horizon $`L/2`$. In the Conclusion we make a few comments, and in four Appendices we derive a few somehow lengthy mathematical results stated in the report.
## 2 When $`1<L<2`$
In a computer simulation, we usually execute the following set of prescriptions to obtain the uncorrelated signature $`\phi _L^u(l)`$:
Figure 1 The distribution of objects in the interval $`(1,L)`$ is an exact copy of the distribution in $`(0,x)`$; here $`p=3`$ and $`m=8`$.
1. in an interval (0, 1) randomly distribute $`m`$ objects; see figure 1;
2. in the side interval $`(1,L)`$ make an exact replica of the $`p`$ objects laying in $`(0,x)`$;
3. measure the $`(m+p)(m+p1)/2`$ separations $`l`$ between the total $`m+p`$ objects, and discard the $`p`$ correlated separations (those which have $`l=1`$ exactly);
4. make a normalized histogram of the
$`𝒟_{mp}={\displaystyle \frac{1}{2}}(m+p)(m+p1)p(1<L<2)`$ (1)
uncorrelated separations;
5. make a large number of new normalized histograms, by repeating the steps 1 to 4 with same $`m`$ (although $`p`$ usually varies);
6. take the mean of these histograms, $`<\varphi _{mL}^u(l)>`$, and construct the quantity
$`<\phi _{mL}^u(l)>=(n1{\displaystyle \underset{g\stackrel{~}{\mathrm{\Gamma }}}{}}\nu _g)\left[<\varphi _{mL}^u(l)>\varphi _L^s(l)\right],`$ (2)
where
$`\varphi _L^s(l)={\displaystyle \frac{2}{L}}(1{\displaystyle \frac{l}{L}}),0<l<L,`$ (3)
and where the factor $`n1\nu _g=(m1)Lx(1x)/L`$ is explained in the appendix 1;
7. the (computer simulated) $`u`$ncorrelated signature $`<\phi _L^u(l)>`$ is the quantity $`<\phi _{mL}^u(l)>`$ when $`m\mathrm{}`$; in practice $`m>50`$ usually suffices. See figure 2.
Figure 2 Computer simulated functions $`<\phi _{mL}^u(l)>`$ for $`\{m=2,L=1.7\}`$ and $`\{m=30,L=1.3\}`$.
We now develop an analytical method to obtain the uncorrelated signature $`\phi _L^u(l)`$. We are dropping the subscript $`exp`$ in all expected (theoretic, analytic, mean) probability distributions. Initially define the lengths $`x`$ and $`y`$ (see figure 1)
$`x=L1,y=1x(1<L<2),`$ (4)
and assume that $`m`$ objects are randomly distributed in (0, 1); the probability that $`p`$ objects be in the interval $`(0,x)`$ and $`mp`$ objects be in the interval $`(x,1)`$ clearly is
$`𝒫_{mpx}=C_m^px^py^{mp},C_m^p={\displaystyle \frac{m!}{p!(mp)!}};`$ (5)
irrespective of the values of $`m`$ and $`x`$ we have
$`{\displaystyle \underset{p=0}{\overset{m}{}}}𝒫_{mpx}=1.`$ (6)
We denote as $`\varphi _{mpL}^u(l)dl`$ the probability of finding in $`(0,L)`$ an uncorrelated pair with separation between $`l`$ and $`l+dl`$, when there are $`m`$ objects in (0, 1) and $`p`$ objects in $`(0,x)`$; clearly it satisfies
$`{\displaystyle _0^L}\varphi _{mpL}^u(l)𝑑l=1.`$ (7)
Recall that a pair $`(P,Q)`$ is said $`g`$-correlated when the isometry $`g`$ brings one of the members to the other; the pair is uncorrelated when no such $`g`$ exists. To investigate $`\varphi _{mpL}^u(l)`$ when $`1<L<2`$ we first call $`A`$ the interval $`(0,x)`$, call $`B=(x,1)`$, and call $`C=(1,L)`$, and note that there are
* $`w_{AA}=p(p1)/2`$ pairs with both members in $`A`$;
* $`w_{AB}=p(mp)`$ pairs with a member in $`A`$ and the other in $`B`$;
* $`w_{AC}=p(p1)`$ uncorrelated pairs, with a member in $`A`$ and the other in $`C`$;
* $`w_{BB}=(mp)(mp1)/2`$ pairs with both members in $`B`$;
* $`w_{BC}(=w_{AB})`$ pairs with a member in $`B`$ and the other in $`C`$;
* $`w_{CC}(=w_{AA})`$ pairs with both members in $`C`$.
In total, there are $`𝒟_{mp}`$ (eq.(1)) pair separations to be considered.
A short reflection gives that the density $`\varphi _{mpL}^u(l)`$ can be decomposed as
$`\varphi _{mpL}^u(l)={\displaystyle \frac{1}{𝒟_{mp}}}[w_{AA}\varphi _{AA}(l)`$ $`+`$ $`w_{AB}\varphi _{AB}(l)+w_{AC}\varphi _{AC}(l)+`$
$`w_{BB}\varphi _{BB}(l)`$ $`+`$ $`w_{BC}\varphi _{BC}(l)+w_{CC}\varphi _{CC}(l)],`$ (8)
where each $`\varphi _{XY}(l)`$ is the probability density of finding an uncorrelated pair of objects separated by $`l`$, one in $`X`$ and the other in $`Y`$; clearly all obey
$`{\displaystyle _0^L}\varphi _{XY}(l)𝑑l=1.`$ (9)
There are two basic types of $`\varphi _{XY}(l)`$, according as $`X=Y`$ or $`XY`$. When $`X=Y`$, suppose a segment of length $`\mu `$, and randomly select two points of it; the probability that their separation lie between $`l`$ and $`l+dl`$ is $`\varphi _\mu ^s(l)dl`$ with (see figure 3)
$`\varphi _\mu ^s(l)={\displaystyle \frac{2}{\mu }}(1{\displaystyle \frac{l}{\mu }}),0<l<\mu .`$ (10)
Figure 3 Pair separation density function for an interval $`\mu `$. The underlying area is 1.
When $`XY`$, consider two intervals with lengths $`\alpha `$ and $`\beta `$, with separation $`\delta `$ (see figure 4); randomly select one point in each $`\alpha `$ and $`\beta `$; the probability that the separation between these points lie between $`l`$ and $`l+dl`$ is $`\varphi _{\delta (\alpha \beta )}(l)dl`$, with the density $`\varphi _{\delta (\alpha \beta )}(l)`$ as depicted in figure 5.
Figure 4 Intervals with lengths $`\alpha `$ and $`\beta `$, with separation $`\delta `$; assume $`\alpha \beta `$.
Figure 5 The probability density $`\varphi _{\delta (\alpha \beta )}(l)`$ for $`\alpha \beta `$ (see figure 4); three particular cases are also displayed; all underlying areas are $`=1`$.
The functions $`\varphi _{XY}(l)`$ appearing in eq.(2) are as displayed in the figure 6, for the case with $`xy`$; for $`xy`$ a similar set has to be constructed, see figure 7.
Figure 6 The normalized functions $`\varphi _{XY}(l)`$ when $`1<L<2`$ and $`x0.5`$ .
Figure 7 The same functions when $`x0.5`$ .
When $`x0.5`$ the density $`\varphi _{mpL}^u(l)`$, eq.(2), is a sequence of four straight segments with endpoints at $`l=0,x,y,1`$, and $`L`$ (in this order), and values
$`\varphi _{mpL}^u(0)`$ $`=`$ $`{\displaystyle \frac{1}{𝒟_{mp}}}[w_{BB}{\displaystyle \frac{2}{y}}+2w_{AA}{\displaystyle \frac{2}{x}}],`$ (11)
$`\varphi _{mpL}^u(x)`$ $`=`$ $`{\displaystyle \frac{1}{𝒟_{mp}}}[w_{BB}{\displaystyle \frac{2(yx)}{y^2}}+2w_{AB}{\displaystyle \frac{1}{x}}](x0.5),`$
$`\varphi _{mpL}^u(y)`$ $`=`$ $`{\displaystyle \frac{1}{𝒟_{mp}}}[2w_{AB}{\displaystyle \frac{1}{y}}](x0.5),`$
$`\varphi _{mpL}^u(1)`$ $`=`$ $`{\displaystyle \frac{1}{𝒟_{mp}}}[w_{AC}{\displaystyle \frac{1}{x}}],\varphi _{mpL}^u(L)=0.`$
When $`x0.5`$ the sequence of endpoints changes to $`l=0,y,x,1`$, and $`L`$, and the values of $`\varphi _{mpL}^u(l)`$ at $`l=y`$ and $`l=x`$ become
$`\varphi _{mpL}^u(y)`$ $`=`$ $`{\displaystyle \frac{1}{𝒟_{mp}}}[2w_{AA}{\displaystyle \frac{2(xy)}{x^2}}+2w_{AB}{\displaystyle \frac{1}{x}}](x0.5)`$ (12)
$`\varphi _{mpL}^u(x)`$ $`=`$ $`{\displaystyle \frac{1}{𝒟_{mp}}}[2w_{AB}{\displaystyle \frac{1}{x}}+w_{AC}{\displaystyle \frac{xy}{x^2}}](x0.5).`$
Two examples of functions $`\varphi _{mpL}^u(l)`$ for $`1<L<2`$ are shown in figure 8.
Figure 8 The probability density $`\varphi _{mpL}^u(l)`$ for $`m=3,p=2`$, and two values of $`L`$: 1.4 and 1.6 . Both underlying areas are 1.
Having the $`m+1`$ functions $`\varphi _{mpL}^u(l)`$, $`p=0,\mathrm{},m`$, we introduce the probability density
$`\varphi _{mL}^u(l)={\displaystyle \underset{p=0}{\overset{m}{}}}𝒫_{mpx}\varphi _{mpL}^u(l),`$ (13)
whose interpretation is obvious: $`\varphi _{mL}^u(l)dl`$ is the probability that two uncorrelated objects randomly selected in $`L`$ have separation between $`l`$ and $`l+dl`$, when $`m`$ objects were randomly distributed in the interval (0, 1). Examples of $`\varphi _{mL}^u(l)`$ are given in figure 9.
Figure 9 Probability densities $`\varphi _{mL}^u(l)`$ for $`1<L<2`$. The graph of $`\varphi _L^s(l)`$ is given in dotted line, for comparison.
Cosmic crystallography is mostly interested in systems with $`m>>1`$. In this limit the function $`\varphi _{mL}^u(l)`$ closely resembles the simple triangular function $`\varphi _L^s(l)`$ (eq.(3), fig. 3), so one is led to define the difference
$`\phi _L^u(l)=\underset{m\mathrm{}}{lim}mL\left[\varphi _{mL}^u(l)\varphi _L^s(l)\right],`$ (14)
the asymptotic uncorrelated signature of $`L`$.
We soon find that the function $`\phi _L^u(l)`$ has a number of symmetries:
$`\phi _L^u(0)=\phi _L^u(L/2)=\phi _L^u(L)=0,\phi _L^u(x)=\phi _L^u(1).`$ (15)
In other words, every $`\phi _L^u(l)`$ with $`1<L<2`$ is composed of three line segments, with the first segment parallel to the third (see figure 10). As expected, the entire graph of $`\phi _L^u(l)`$ is uniquely fixed by the number $`f(L)`$, the value of $`\phi _L^u(l)`$ at $`l=x`$; in the appendix 2 we show that
$`f(L)={\displaystyle \frac{8xy}{L^3}}(1<L<2).`$ (16)
A plot of $`f(L)`$ valid for arbitrary $`L>1`$ is given in figure 11.
Figure 10 Geometro-topological signature $`\phi _L^u(l)`$ for $`L=`$1.1 , 1.5, and 1.9 .
Figure 11 The function $`f(L)`$, the absolute maximum of $`\phi _L^u(l)`$ (which occurs at $`l=x`$); three particular values of $`L`$ are marked, those used in figure 10.
## 3 When $`L>2`$
The generalization of the previous results for arbitrary values of $`L`$ is straightforward but lengthy, so we only state the final results in this section. See the appendix 3 for details. The graph of the uncorrelated signature (14) with
$`L=\lambda +x,\lambda Z_+,0<x<1`$ (17)
has the aspect of a slanted saw; see figure 12, drawn for $`\lambda =5`$ and $`x=0.2`$.
Figure 12 The function $`\phi _L^u(l)`$ for $`L=5.2`$ .
There are $`\lambda `$ maxima, which occur in the positions $`l=x,1+x,\mathrm{},(\lambda 1)+x`$ , and there are $`\lambda `$ minima, which lay in the positions $`l=1,2,\mathrm{},\lambda `$. A straight line connects the maxima, another one connects the minima, both have angular coefficient $`8xy/L^3`$. The $`\lambda +1`$ segments with positive angular coefficient are parallel, as well as the $`\lambda `$ segments with negative slope. As expected, the value of $`L`$ is the sufficient datum to draw $`\phi _L^u(l)`$, since
$`\phi _L^u(x)=\phi _L^u(\lambda )={\displaystyle \frac{8\lambda xy}{L^3}},`$ (18)
as shown in the appendix 4. The graph of $`f(L)=8\lambda xy/L^3`$ is given in Figure 11.
## 4 Conclusion
In our first contact with the cosmic crystallography it appeared plausible that the normalized expected functions $`\varphi _{exp}^u(l)`$ and $`\varphi _{exp}^s(l)`$ were the same, since both are concerned with separations between objects isometrically unrelated . However, in our computer simulations a persistent non-nullity of the difference $`<\varphi ^u(l)>\varphi _{exp}^s(l)`$ made imperative a more close exam. It soon became evident that a difference indeed existed, and that it diminished as the number $`n`$ of objects present in the sample increased.
Further investigation suggested to define the uncorrelated signature
$`\phi _{exp}^u(l)=(n1{\displaystyle \nu _g})\left[\varphi _{exp}^u(l)\varphi _{exp}^s(l)\right],`$ (19)
where $`\nu _g=N_g/n`$ , with $`N_g=`$number of $`g`$-pairs in the observed universe; for the cosmic crystallography we usually have $`n>>1+\nu _g`$.
Earlier attempts to find $`\varphi _{exp}^u(l)`$ for three-dimensional balls failed, and also for $`2D`$ balls; we then focussed our attention on a $`1D`$ ball, this report. When we compare the final theoretical result (14) with the mean of an increasing number of histograms obtained from computer simulations, we note a rapid agreement of the two approaches in the region of large separations $`l>L/2`$, while in the region where $`l<L/2`$ a quite larger number of simulated catalogs is demanded. This can be seen in Figure 2, where we observe that the statistical fluctuations for $`l`$ large are sensibly less pronounced than those for small $`l`$.
When $`L<1`$, then there is no replication of objects; in this case $`\varphi _L^u(l)=\varphi _L^s(l)`$ and clearly $`\phi _L^u(l)=0`$. When $`L>1`$ is an integer, then objects are replicated; nevertheless still $`\varphi _L^u(l)=\varphi _L^s(l)`$ and $`\phi _L^u(l)=0`$. This can be seen in the figure 11, where we note that $`f(L)`$ vanishes for $`L=`$ integer $`>0`$.
Appendix 1
We evaluate the quantity $`n1\nu _g`$ for a universe $`S^1`$ with circumference 1 and observed universe with total amplitude $`L=\lambda +x`$, being $`\lambda `$ a positive integer and $`0<x<1`$.
Assuming $`m`$ objects along the circle $`S^1`$ with radius $`1/(2\pi )`$, then the expected number of objects in $`L`$ is $`n=mL`$. The sum $`\nu _g=\nu _\lambda +\nu _{\lambda +1}+\mathrm{}+\nu _1+\nu _1+\mathrm{}+\nu _{\lambda 1}+\nu _\lambda `$ indeed simplifies to $`2(\nu _1+\nu _2+\mathrm{}+\nu _\lambda )`$, since $`\nu _i=\nu _i`$.
Now remember that for $`i`$ a positive integer $`n\nu _i`$ is the expected number of pairs of objects in the observed universe whose separation is $`i\lambda `$ ; its value is
$`n\nu _i=m(Li).`$ (20)
As a consequence $`\nu _g=\lambda (Ly)/L`$, and finally
$`n1{\displaystyle \nu _g}=(m1)L{\displaystyle \frac{xy}{L}}.`$ (21)
Appendix 2
We show that $`\phi _L^u(x)=8xy/L^3`$ when $`1<L<2`$: from (11) or (12) we have at $`l=1`$
$`\varphi _{mpL}^u(1)={\displaystyle \frac{1}{𝒟_{mp}}}{\displaystyle \frac{p(p1)}{x}},`$ (22)
so we have from (13)
$`\varphi _{mL}^u(1)={\displaystyle \frac{1}{x}}{\displaystyle \underset{p=0}{\overset{m}{}}}{\displaystyle \frac{𝒫_{mpx}}{𝒟_{mp}}}p(p1),`$ (23)
whose value is sought, correct to order $`m^1`$ when $`m>>1`$. In this limit we have
$`{\displaystyle \underset{p=0}{\overset{m}{}}}𝒫_{mpx}(p/m)^k=x^k+{\displaystyle \frac{k(k1)}{2m}}yx^{k1}+O(m^2),`$ (24)
and consequently
$`{\displaystyle \underset{p=0}{\overset{m}{}}}𝒫_{mpx}F(p/m)=F(x)+{\displaystyle \frac{xy}{2m}}{\displaystyle \frac{d^2}{dx^2}}F(x)+O(m^2).`$ (25)
For $`m>>1`$ in eq.(1) we find that
$`{\displaystyle \frac{p(p1)}{𝒟_{mp}}}={\displaystyle \frac{2\xi ^2}{(\xi +1)^2}}+{\displaystyle \frac{2\xi (2\xi +1)(\xi 1)}{m(\xi +1)^4}}+O(m^2),\xi :=p/m,`$ (26)
so from (25) we obtain
$`{\displaystyle \underset{p=0}{\overset{m}{}}}𝒫_{mpx}{\displaystyle \frac{p(p1)}{𝒟_{mp}}}`$ $`=`$ $`{\displaystyle \frac{2x^2}{(1+x)^2}}+{\displaystyle \frac{2x(2x+1)(x1)}{m(1+x)^4}}+{\displaystyle \frac{xy}{2m}}{\displaystyle \frac{d^2}{dx^2}}\left[{\displaystyle \frac{2x^2}{(1+x)^2}}\right]+O(m^2)`$ (27)
$`=`$ $`{\displaystyle \frac{2x^2}{L^2}}{\displaystyle \frac{8x^2y}{mL^4}}+O(m^2).`$
Since $`\varphi _L^s(1)=2x/L^2`$, we finally have from (14), (23), and (27)
$`\phi _L^u(1)={\displaystyle \frac{8xy}{L^3}}(1<L<2).`$ (28)
Appendix 3
We generalize for arbitrary $`L>1`$ the results obtained for $`1<L<2`$, in particular the equations (1) and (16). We first decompose the total interval $`(0,L)`$ into $`2\lambda +1`$ subintervals according to figure 13, drawn for $`\lambda =5`$.
Figure 13 The one-dimensional observed universe with length $`L=\lambda +x`$, partitioned into $`\lambda +1`$ intervals $`A_i`$ with length $`x`$ and $`\lambda `$ intervals $`B_i`$ measuring $`y=1x`$.
For $`m`$ objects randomly distributed in the universe $`(0,1)`$ we expect $`p=mx`$ objects in each interval $`A_i`$ and $`mp=my`$ objects in each $`B_i`$. The number of objects in the observed universe $`(0,L)`$ being $`m\lambda +p`$, the total number of pairs of objects in it is $`(m\lambda +p)(m\lambda +p1)/2`$; if we deduct the $`p\lambda (\lambda +1)/2`$ correlated pairs with members in the $`A`$’s, and the $`(mp)(\lambda 1)\lambda /2`$ correlated pairs with members in the $`B`$’s, then we obtain the expected number of uncorrelated separations (cf eq.(1)):
$`𝒟_{mp}={\displaystyle \frac{1}{2}}(m\lambda +p)(m\lambda +p1){\displaystyle \frac{1}{2}}\lambda (\lambda +1)p{\displaystyle \frac{1}{2}}\lambda (\lambda 1)(mp).`$ (29)
We next note in $`(0,L)`$ the existence of
* $`w_{A_iA_i}=p(p1)/2`$ pairs with both members in $`A_i`$;
* $`w_{B_iB_i}=(mp)(mp1)/2`$ pairs with both members in $`B_i`$;
* $`w_{A_iA_{j>i}}=2w_{A_iA_i}`$ uncorrelated pairs, with a member in $`A_i`$ and the other in $`A_{j>i}`$;
* $`w_{B_iB_{j>i}}=2w_{B_iB_i}`$ uncorrelated pairs, with a member in $`B_i`$ and the other in $`B_{j>i}`$;
* $`w_{A_iB_{ji}}=p(mp)`$ pairs, with a member in $`A_i`$ and the other in $`B_{ji}`$;
* $`w_{B_iA_{j>i}}=w_{A_iB_{ji}}`$ pairs, with a member in $`B_i`$ and the other in $`A_{j>i}`$.
There is a total of $`(2\lambda +1)(\lambda +1)`$ such numbers $`w_{XY}`$, and their sum clearly is the $`𝒟_{mp}`$ given in (29).
With the probability densities $`\varphi _{XY}(l)`$ defined as before, the normalized probability density $`\varphi _{mpL}^u(l)`$ is written similarly to eq. (8),
$`\varphi _{mpL}^u(l)={\displaystyle \frac{1}{𝒟_{mp}}}{\displaystyle \underset{X,Y}{}}w_{XY}\varphi _{XY}(l).`$ (30)
As a matter of fact, there are only three essentially different $`w_{XY}`$, which we dub $`w_{AA},w_{AB}`$, and $`w_{BB}`$, as in sec. 2. Also, there are indeed only $`3\lambda +1`$ different functions $`\varphi _{XY}^{}(l)`$, each appearing with variable multiplicity $`m_{XY}`$. These functions, together with the corresponding $`m_{XY}`$ and weights $`w_{XY}`$, are displayed in figure 14, drawn for $`L=5.2`$.
Figure 14 The $`3\lambda +1`$ different functions $`\varphi _{XY}^{}(l)`$ when $`x0.5`$. On top of each function the corresponding multiplicity $`m_{XY}`$ is written. On the left side the corresponding weight $`w_{XY}`$ is also given. The value $`L=5.2`$ was taken for definiteness.
When $`x>0.5`$ the set of functions $`\varphi _{XY}^{}(l)`$ has a different aspect; see figure 15, drawn for $`L=5.8`$.
Figure 15 The $`3\lambda +1`$ different functions $`\varphi _{XY}^{}(l)`$ when $`x0.5`$. The multiplicities $`m_{XY}`$ and weights $`w_{XY}`$ are indicated as in figure 14. The value $`L=5.8`$ was taken for definiteness.
It is now clear that the functions $`\varphi _{mpL}^u(l)`$ are a sequence of $`3\lambda +1`$ segments, each segment having endpoints either at an integer or separated $`x`$ from an integer; as a consequence, also the functions $`\varphi _{mL}^u(l)`$ (eq.(13)) have that behavior, as well as the functions $`\phi _{mL}^u(l)`$ (eq.(2)). See figure 16.
Figure 16 The function $`\phi _{mL}^u(l)`$ for $`m=2`$ and $`L=5.2`$. A straight line connects the points with abscissa $`l=i+x(i=0,\mathrm{},\lambda )`$; another, parallel, connects those with $`l=i+y(i=0,\mathrm{},\lambda 1)`$; also the points with $`l=`$integer are aligned.
Appendix 4
We generalize eq.(28) for arbitrary $`L>1`$. For $`m>>1`$ and $`𝒟_{mp}`$ as in eq.(29) we have
$`{\displaystyle \frac{p\left(p1\right)}{𝒟_{mp}}}={\displaystyle \frac{2\xi ^2}{\left(\xi +\lambda \right)^2}}+{\displaystyle \frac{2\lambda \xi \left(2\xi +\lambda \right)\left(\xi 1\right)}{m\left(\xi +\lambda \right)^4}}+O\left(m^2\right),`$ (31)
while eq.(27) now reads
$`\varphi _{mL}^u\left(\lambda \right)={\displaystyle \frac{2x^2}{L^2}}{\displaystyle \frac{8\lambda x^2y}{mL^4}}+O\left(m^2\right).`$ (32)
Finally (28) becomes
$`\phi _L^u\left(\lambda \right)={\displaystyle \frac{8\lambda xy}{L^3}},L>1.`$ (33)
The graph of $`f(L)=8\lambda xy/L^3`$ is given in figure 11.
## Acknowledgments
Conversations with Armando Bernui, Germán I. Gomero, and Marcelo J. Rebouças are hearty acknowledged.
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# Mass Hierarchy from 𝑺𝑼(𝟏,𝟏) Horizontal Symmetry
## 1 Introduction
In the low-energy particle physics, Nature shows two remarkable aspects. First, quarks and leptons appear with the repetition of three generations. It is not yet clear enough where the definite number of chiral generations comes from. Second, the generations have the hierarchical mass structure. This structure is translated to the hierarchical Yukawa coupling structure whose origin is not yet well established. The characteristic structure of Yukawa coupling matrix seems to suggest some inter-generation symmetry, that is, horizontal symmetry.
For the possible approach to understanding these problems, the model was proposed based on the noncompacet horizontal gauge symmetry $`G_H=SU(1,1)`$. This model is a vector-like generalization of the minimal supersymmetric standard model(MSSM), and the three chiral generations and the hierarchical Yukawa coupling structure are realized through the mechanism called “spontaneous generation of generations”.
In the model, a MSSM chiral superfield($`f`$) is extended to a chiral superfield($`F`$) which belongs to an infinite dimensional unitary representation of the horizontal gauge symmetry $`SU(1,1)`$. For example, the three generations of quark doublets $`q`$ is extended to a $`SU(1,1)`$ multiplet;
$$Q_\alpha =\{q_\alpha ,q_{\alpha +1},q_{\alpha +2},q_{\alpha +3},\mathrm{}\},$$
(1.1)
where $`\alpha `$ is the lowest weight of $`SU(1,1)`$ and assumed to be real positive. The three chiral $`q`$’s are contained in $`Q_\alpha `$. Precisely speaking, $`q`$’s are realized as linear combinations of infinite number of $`q_{\alpha +i}`$’s. According to the vector-like(left-right) symmetry, all $`F`$’s are accompanied by $`\overline{F}`$’s which are conjugate to $`F`$’s under $`SU(1,1)\times SU(3)\times SU(2)\times U(1)`$. So $`Q_\alpha `$ has a partner ;
$$\overline{Q}_\alpha =\{\overline{q}_\alpha ,\overline{q}_{\alpha 1},\overline{q}_{\alpha 2},\overline{q}_{\alpha 3},\mathrm{}\}.$$
(1.2)
The horizontal symmetry is spontaneously broken by a non-vanishing vacuum expectation value(VEV) of some component $`\psi _i`$ of $`\mathrm{\Psi }_R`$, which is a multiplet belonging to a finite dimensional non-unitary representation of $`SU(1,1)`$ with the integer highest weight $`R`$,
$$\mathrm{\Psi }_R=\{\psi _R,\psi _{R+1},\mathrm{},\psi _{R1},\psi _R\},$$
(1.3)
and assumed to be singlet under the standard model(SM) gauge group.
The coupling of the non-vanishing VEV $`\psi _3`$ to $`F_\alpha `$ and $`\overline{F}_\alpha `$ provides the three chiral generations of $`f`$’s. The relevant part of the superpotential taken in Reference 2) is
$`M_F\overline{F}_\alpha F_\alpha `$ $`+`$ $`v_F\mathrm{\Psi }_R\overline{F}_\alpha F_\alpha `$ (1.4)
$`=`$ $`M_F{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}(1)^i\overline{f}_{\alpha i}f_{\alpha +i}+v_F{\displaystyle \underset{i,j=0}{\overset{\mathrm{}}{}}}A_{i,j}^{\alpha ,R}\psi _{ij}\overline{f}_{\alpha i}f_{\alpha +j},`$
where $`M_F`$ is a $`SU(1,1)`$ invariant mass and the Clebsch-Gordan coefficient $`A_{i,j}^{\alpha ,R}`$ are given, with the normalization factor $`N_A`$, as<sup>)</sup><sup>)</sup>)We correct the type miss of the equation (A8) in Reference 2).
$`A_{i,j}^{\alpha ,R}`$ $`=`$ $`N_A(1)^j\sqrt{{\displaystyle \frac{i!j!(ij+R)!(i+j+R)!}{\mathrm{\Gamma }(2\alpha +i)\mathrm{\Gamma }(2\alpha +j)}}}`$ (1.5)
$`\times {\displaystyle \underset{r=0}{\overset{ij+R}{}}}{\displaystyle \frac{\mathrm{\Gamma }(2\alpha +j+r)}{r!(ji+r)!(Rr)!(ij+Rr)!(r+jR)!}}.`$
When the $`SU(1,1)`$ is spontaneously broken and $`\mathrm{\Psi }_R`$ is replaced with its VEV, the second term of (1.4) does not contain the first three components of $`F_\alpha `$. Under the presence of the first term of (1.4), this absence of three components is retained in the unitary transformed way due to the fact that the second term of (1.4) is always dominant for large $`i`$ and $`j`$. This disappearance means a generation of the three chiral generations. The other components acquire huge Dirac masses with all components of $`\overline{F}_\alpha `$ and decouple from the low-energy phenomena. The SM Yukawa couplings are obtained by extracting these three chiral generations from the $`SU(1,1)`$ invariant Yukawa couplings. The mixing effect due to the unitary transformation mentioned above and group theoretical structure of $`SU(1,1)`$ Clebsch-Gordan coefficients reproduce the hierarchical structure of Yukawa coupling matrices.
At the high-energy scale, where the horizontal symmetry is manifest, it is natural to expect that the SM gauge group is also unified into the higher symmetry(grand unified) group $`GSU(3)\times SU(2)\times U(1)`$. Thus, it would be reasonable to extend the present model to the model based on the gauge group $`SU(1,1)\times G`$. Grand unified theory(GUT) usually requires us to introduce extra particles which acquire super heavy masses through spontaneous break down of $`G`$. The present model is basically vector-like, so that it has a desirable potential for making such extra particles super heavy. However the “spontaneous generation of generations” mechanism based on the superpotential (1.4) encounters serious difficulty under the grand unification. This is because the extra particles, such as colored Higgses, are unified with the MSSM particles into the same grand unified multiplet, and consequently the chiral realization of the MSSM particles is always accompanied with that of the GUT partners. Therefore we must abandon (1.4) as the “spontaneous generation of generations” mechanism and search for an available alternative.
In this paper we propose the solution to this problem by a simple extension. The extension is to replace the $`SU(1,1)`$ invariant mass term $`M_F\overline{F}_\alpha F_\alpha `$ in (1.4) by a Yukawa coupling $`u_F\mathrm{\Phi }_R\overline{F}_\alpha F_\alpha `$ with a non-unitary multiplet $`\mathrm{\Phi }_R`$. Under this replacement, the vacuum has various phase structure. When $`v_F\mathrm{\Psi }`$ is much larger than $`u_F\mathrm{\Phi }`$, the chiral structure is determined by $`\mathrm{\Psi }\psi _g`$ and we have $`g`$ chiral generations. On the contrary, when $`v_F\mathrm{\Psi }u_F\mathrm{\Phi }`$, we have no chiral generations of $`F`$, if $`\mathrm{\Phi }\varphi _0`$. The number of chiral generations is integer which cannot vary continuously. This means that the phase space of the vacuum is divided into two domains whose boundary is characterized by the critical values of the ratio $`u_F\mathrm{\Phi }/v_F\mathrm{\Psi }`$. This novel mechanism enables us to give heavy masses to superfluous particles like colored Higgses in the supersymmetric $`SU(5)`$ model without fine tuning. The purpose of this paper is to present a general framework of the model and to clarify the characteristics which the model reveals.
## 2 The New Mechanism for Generating Chiral Generations
Introducing the horizontal gauge symmetry $`SU(1,1)`$ in a vector-like manner, we extend the MSSM superfields $`(q,\overline{u},\overline{d},l,\overline{e},h,h^{})`$ and their conjugates $`(\overline{q},u,d,\overline{l},e,\overline{h},\overline{h^{}})`$ to
$$\begin{array}{ccccccc}Q_\alpha ,& \overline{U}_\beta ,& \overline{D}_\gamma ,& L_\eta ,& \overline{E}_\lambda ,& H_\rho ,& H_\sigma ^{},\\ \overline{Q}_\alpha ,& U_\beta ,& D_\gamma ,& \overline{L}_\eta ,& E_\lambda ,& \overline{H}_\rho ,& \overline{H^{}}_\sigma ,\end{array}$$
where each of the fields belong to an infinite dimensional representation of $`SU(1,1)`$ and $`\alpha ,\beta ,\gamma ,\eta ,\lambda ,\rho ,\sigma `$ are real positive.
Let us discuss the chiral structure of quark doublets $`Q_\alpha `$ and $`\overline{Q}_\alpha `$. We examine the superpotential
$$u_Q\mathrm{\Phi }_S\overline{Q}_\alpha Q_\alpha +v_Q\mathrm{\Psi }_R\overline{Q}_\alpha Q_\alpha ,$$
(2.1)
where $`v_Q`$ and $`u_Q`$ are real coupling constants and $`\mathrm{\Psi }_R`$ and $`\mathrm{\Phi }_S`$ belong to finite dimensional representations of $`SU(1,1)`$ with highest weight $`R`$ and $`S`$ respectively. We assume both $`\mathrm{\Psi }_R`$ and $`\mathrm{\Phi }_S`$ acquire non-vanishing VEV’s $`\psi _g`$ and $`\varphi _g^{}`$ simultaneously. We must set $`g^{}=0`$ to give both possibilities that all modes become massive and $`g`$ modes become massless. Thus, the $`SU(1,1)`$ broken superpotential(mass terms) takes a form as
$`u_Q\varphi _0{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}A_{i,i}^{\alpha ,S}\overline{q}_{\alpha i}q_{\alpha +i}+v_Q\psi _g{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}A_{i,i+g}^{\alpha ,R}\overline{q}_{\alpha i}q_{\alpha +i+g}.`$ (2.2)
The generation of chiral generations means the appearance of massless modes in (2.2).
Let us assume that we have $`g`$ massless modes $`\{q^{(0)},q^{(1)},q^{(2)},\mathrm{},q^{(g1)}\}`$ and write
$$q_{\alpha +j}=\underset{i=0}{\overset{g1}{}}a_j^{q(i)}q^{(i)}+\text{ massive modes}.$$
(2.3)
The coefficients $`a_j^{q(i)}`$ must be normalizable as $`_{j=0}^{\mathrm{}}|a_j^{q(i)}|^2=1`$. Inserting (2.3) into (2.2), we obtain the massless condition for $`q^{(i)}`$ as
$$u_Q\varphi _0A_{j,j}^{\alpha ,S}a_j^{q(i)}+v_Q\psi _gA_{j,j+g}^{\alpha ,R}a_{j+g}^{q(i)}=0.$$
(2.4)
This equation gives recursion relation for coefficients $`a_j^{q(i)}`$,
$$a_{j+g}^{q(i)}=\frac{u_Q\varphi _0A_{j,j}^{\alpha ,S}}{v_Q\psi _gA_{j,j+g}^{\alpha ,R}}a_j^{q(i)}.$$
(2.5)
This recursion relation means that $`a_j^{q(i)}`$ goes to zero if the magnitude of the coefficient of $`a_j^{q(i)}`$ in the right-hand side is smaller than unity in the limit $`j\mathrm{}`$. In this case $`a_j^{q(i)}`$ are normalizable, so massless modes appear.
For fixed $`ij`$, the asymptotic behavior of the Clebsch-Gordan coefficient (1.5) is
$$A_{i,j}^{\alpha ,R}N_A(1)^j\frac{(2R)!}{(R!)^2\sqrt{(Ri+j)!(R+ji)!}}(j)^R,(i,j\mathrm{}).$$
(2.6)
Therefore, we find, in the limit $`j\mathrm{}`$,
$$\frac{u_Q\varphi _0A_{j,j}^{\alpha ,S}}{v_Q\psi _gA_{j,j+g}^{\alpha ,R}}(1)^g\frac{u_Q\varphi _0}{v_Q\psi _g}\frac{(R!)^2(2S)!}{(S!)^2(2R)!}\sqrt{\frac{(Rg)!(R+g)!}{R!R!}}\frac{j^S}{j^R}.$$
(2.7)
For $`S<R`$, the right-hand side goes to zero. This means that all $`a_j^{q(i)}`$ are well defined, so $`g`$ massless generations appear. By setting $`g=3`$, the three chiral generations of quark doublets are realized. On the other hand, for $`S>R`$, $`a_j^{q(i)}`$ diverges for $`j\mathrm{}`$ and we can not define normalizable massless modes. This means that all components in (2.2) become massive. Obviously, the present model is a natural extension of the previous model based on (1.4), which is recovered by the choice $`S=0`$. This extension, nevertheless, gives the model with the remarkable feature, which is realized in the special case $`S=R`$.
For $`S=R`$, the situation in (2.2) is subtle. Since $`j^S/j^R`$ is unity, the coefficient of $`j^S/j^R`$ in the right-hand side of (2.7) determines whether the chiral generations appear or not. Let us define the ratio $`r_Q`$ by
$$r_Q=\frac{u_Q\varphi _0}{v_Q\psi _g}.$$
(2.8)
The asymptotic behavior of $`a_j^{q(i)}`$ is governed by the value of this ratio. From (2.5) and (2.7) we find the critical value of $`r_Q`$ is given by
$$r_g^c=\sqrt{\frac{R!R!}{(Rg)!(R+g)!}}.$$
(2.9)
For $`|r_Q|<r_g^c`$, $`a_j^{q(i)}`$ is normalizable and the $`g`$ chiral generations appear. The chiral generation structure does not appear for $`|r_Q|r_g^c`$. The generation of generations depends on the magnitude of $`r_Q`$ which is controlled by couplings $`u`$ and $`v`$. That is, two dimensional parameter space of $`u`$ and $`v`$ is divided into two domains, each of which has the quite different chiral structure.
The choice $`S=R`$ not only accomplishes simplification purpose, but also gives us an attractive mechanism. Some of the MSSM particles are unified in a large GUT multiplet with extra particles at high energy. If each particles in the GUT multiplet have different ratio $`r`$’s (2.8), it is possible to realize without fine tuning that the MSSM particles have the chiral structure and the extra particles are super heavy.
For instance, let us consider the Higgs sector of the supersymmetric $`SU(5)`$ grand unified model. A Higgs doublet $`h`$ is embedded in a $`\mathrm{𝟓}`$-dimensional representation together with an extra colored Higgs $`h^C`$. The $`SU(1,1)`$ Higgs multiplet is denoted as
$$K_\rho =\left(\begin{array}{c}H_\rho ^C\\ H_\rho \end{array}\right).$$
(2.10)
We introduce two multiplet $`\mathrm{\Phi }_R^{}^{}`$ and $`\mathrm{\Psi }_{}^{}{}_{R^{}}{}^{}`$ which belong to $`\mathrm{𝟏}`$ and $`\mathrm{𝟐𝟒}`$ of $`SU(5)`$ respectively, and couple them to $`K_\rho `$ and its conjugate $`\overline{K}_\rho `$ by
$$u_K\mathrm{\Phi }_R^{}^{}\overline{K}_\rho K_\rho +v_K\overline{K}_\rho \mathrm{\Psi }_R^{}^{}K_\rho .$$
(2.11)
The gauge symmetry $`SU(1,1)\times SU(5)`$ is spontaneously broken to $`SU(3)\times SU(2)\times U(1)`$ via non-vanishing VEV’s of $`\mathrm{\Phi }_R^{}^{}`$ and $`\mathrm{\Psi }_{}^{}{}_{R^{}}{}^{}`$
$$\varphi _0^{}0,\psi _1^{}\left(\begin{array}{ccccc}1& & & & \\ & 1& & & \\ & & 1& & \\ & & & \frac{3}{2}& \\ & & & & \frac{3}{2}\end{array}\right)0,$$
(2.12)
where we have set $`g=1`$ for $`\mathrm{\Psi }_R^{}`$ to realize one chiral generation for Higgs because of its negative weights of $`SU(1,1)`$. The ratios (2.8) for $`h`$ and $`h^C`$ are
$$r_{H^C}=\frac{u_K\varphi _0^{}}{v_K\psi _1^{}},$$
(2.13)
$$r_H=\frac{2}{3}\frac{u_K\varphi _0^{}}{v_K\psi _1^{}}.$$
(2.14)
The critical value in the present case is given by (2.9) with $`g=1`$$`r_1^c`$. When the ratio $`r`$’s satisfies the condition $`|r_H|<r_1^c|r_{H^C}|`$, $`H`$ has one chiral generation but $`H^C`$ does not. This condition is consistent with (2.13) and (2.14). Therefore, we can remove the colored Higgs in a low-energy theory by choosing such $`r`$’s. There is no fine tuning. Like this example, we can use this mechanism to make superfluous particles heavy.
Following above discussion, we describe the generation generating superpotential as
$`(u_Q\mathrm{\Phi }_R`$ $`+`$ $`v_Q\mathrm{\Psi }_R)\overline{Q}_\alpha Q_\alpha +(u_U\mathrm{\Phi }_R+v_U\mathrm{\Psi }_R)U_\beta \overline{U}_\beta +(u_D\mathrm{\Phi }_R+v_D\mathrm{\Psi }_R)D_\gamma \overline{D}_\gamma `$ (2.15)
$`+`$ $`(u_L\mathrm{\Phi }_R+v_L\mathrm{\Psi }_R)\overline{L}_\eta L_\eta +(u_E\mathrm{\Phi }_R+v_E\mathrm{\Psi }_R)E_\lambda \overline{E}_\lambda `$
$`+`$ $`(u_H\mathrm{\Phi }_R^{}^{}+v_H\mathrm{\Psi }_R^{}^{})\overline{H}_\rho H_\rho +(u_H^{}\mathrm{\Phi }_R^{}^{}+v_H^{}\mathrm{\Psi }_R^{}^{})\overline{H^{}}_\sigma H_\sigma ^{}.`$
For simplicity and economical purpose, we assume the same non-unitary multiplets $`\mathrm{\Phi }_R`$ and $`\mathrm{\Psi }_R`$ couple to quarks and leptons and generate the three chiral generations through non-vanishing VEV’s $`\varphi _0`$ and $`\psi _3`$. Also, the same $`\mathrm{\Phi }_R^{}^{}`$ and $`\mathrm{\Psi }_R^{}^{}`$ with $`\varphi _0^{}0`$ and $`\psi _1^{}0`$ realize one chiral generation for Higgses, The potential (2.15) allows the selective generation of the chiral generations and the extension to grand unified models.
## 3 The Hierarchical Structure of Yukawa Coupling Matrices
Let us now proceed to the Yukawa couplings of Higgses to quarks and leptons. The relevant terms in the superpotential are given as
$$y_U\overline{U}_\beta Q_\alpha H_\rho +y_D\overline{D}_\gamma Q_\alpha H_\sigma ^{}+y_E\overline{E}_\lambda L_\eta H_\sigma ^{}.$$
(3.1)
The $`SU(1,1)`$ invariance requires
$`\rho `$ $`=`$ $`\alpha +\beta +\mathrm{\Delta },`$ (3.2)
$`\sigma `$ $`=`$ $`\alpha +\gamma +\mathrm{\Delta }^{}=\lambda +\eta +\mathrm{\Delta }^{\prime \prime },`$ (3.3)
where $`\mathrm{\Delta }`$, $`\mathrm{\Delta }^{}`$ and $`\mathrm{\Delta }^{\prime \prime }`$ are non-negative integers. The Clebsch-Gordan decomposition of the first term of (3.1) takes the form
$$y_U\underset{i,j=0}{\overset{\mathrm{}}{}}B_{i,j}^U\overline{u}_{\beta +i}q_{\alpha +j}h_{\rho ij+\mathrm{\Delta }}.$$
(3.4)
The Clebsch-Gordan coefficient is
$`B_{i,j}^U=`$ $`N^U(1)^{i+j}\sqrt{{\displaystyle \frac{i!j!\mathrm{\Gamma }(2\beta +i)\mathrm{\Gamma }(2\alpha +j)}{(i+j\mathrm{\Delta })!\mathrm{\Gamma }(2\rho +i+j\mathrm{\Delta })}}}`$ (3.5)
$`\times {\displaystyle \underset{r=0}{\overset{\mathrm{\Delta }}{}}}{\displaystyle \frac{(1)^r(i+j\mathrm{\Delta })!}{(ir)!(j\mathrm{\Delta }+r)!r!(\mathrm{\Delta }r)!\mathrm{\Gamma }(2\beta +r)\mathrm{\Gamma }(2\alpha +\mathrm{\Delta }r)}},`$
where $`N^U`$ is a normalization constant.
To find the structure of Yukawa couplings for massless modes, we must pick up them from original components. As mentioned in section 2, the massless modes of $`Q_\alpha `$ which consist of $`g=3`$ chiral generations are given from (2.3) as
$$q^{(i)}=\underset{j=0}{\overset{\mathrm{}}{}}a_{j}^{q(i)}{}_{}{}^{}q_{\alpha +j}(i=0,1,2).$$
(3.6)
The coefficients $`a_j^{q(i)}`$ are determined by the equation (2.5) as
$$a_{3k+i}^{q(i)}=(r_Q)^ka_i^{q(i)}\underset{n=0}{\overset{k1}{}}\frac{A_{3n+i,3n+i}^{\alpha ,R}}{A_{3n+i,3(n+1)+i}^{\alpha ,R}},$$
(3.7)
where
$$r_Q=\frac{u_Q\varphi _0}{v_Q\psi _3}.$$
(3.8)
To guarantee that the three chiral generations of quark doublets are realized, we assume that the ratio $`r_Q`$ satisfies the relation
$$|r_Q|<r_3^c=\sqrt{\frac{R(R1)(R2)}{(R+1)(R+2)(R+3)}}.$$
(3.9)
The mixing contribution of higher components $`q_{\alpha +j}`$ to massless modes is generally suppressed, for reasonable value of $`|r_Q|<1`$, due to the suppression factor $`(r_Q)^k`$ in (3.7). There exists an additional suppressing mechanism which becomes effective for large value of $`R`$. From the asymptotic form of Clebsch-Gordan coefficient $`A_{i,j}^{\alpha ,R}`$ in the limit $`R\mathrm{}`$,
$$A_{i,j}^{\alpha ,R}N_A(1)^j\frac{1}{\sqrt{i!j!\mathrm{\Gamma }(2\alpha +i)\mathrm{\Gamma }(2\alpha +j)}}R^{2\alpha +i+j1},$$
(3.10)
we find the suppression factor
$$a_{3k+i}^{q(i)}(r_Q)^ka_i^{q(i)}\left(\frac{1}{R^3}\right)^k.$$
(3.11)
For all quarks and leptons which have the three chiral generations, we have similar relations (3.6)-(3.9) with appropriate replacements of the weights and couplings $`u`$, $`v`$.
Higgses have negative weights so that the Yukawa couplings (3.1) are invariant under $`SU(1,1)`$. We can obtain one massless mode by setting the non-vanishing VEV for positive component of $`\mathrm{\Psi }^{}`$ as $`\psi _1^{}0`$. So the massless mode of $`H_\rho `$ is given as
$$h=\underset{j=0}{\overset{\mathrm{}}{}}a_{j}^{h}{}_{}{}^{}h_{\rho j},$$
(3.12)
where
$$a_j^h=(r_H)^ja_0^h\underset{n=0}{\overset{j1}{}}\frac{A_{n,n}^{\rho ,R^{}}}{A_{n+1,n}^{\rho ,R^{}}},$$
(3.13)
$$r_H=\frac{u_H\varphi _0^{}}{v_H\psi _1^{}}.$$
(3.14)
Here we assume
$$|r_H|<r_1^c=\sqrt{\frac{R^{}}{R^{}+1}},$$
(3.15)
to realize one chiral generation. Similar relations apply to $`H_\sigma ^{}`$. The mixing effect in (3.12) is also suppressed for $`|r_H|<1`$ and we have additional suppression in (3.13) for large value of $`R^{}`$,
$$a_j^h(r_H)^ja^h\left(\frac{1}{R^{}}\right)^j.$$
(3.16)
From the $`SU(1,1)`$ invariant couplings (3.1), we can extract Yukawa couplings among massless modes. For example, for up-type quarks, we have
$$\underset{i,j=0}{\overset{2}{}}\mathrm{\Gamma }_{ij}^u\overline{u}^{(i)}q^{(j)}h,$$
(3.17)
where the Yukawa coupling matrix $`\mathrm{\Gamma }^u`$ is given as
$$\mathrm{\Gamma }_{ij}^u=y_U\underset{m,n=0}{\overset{\mathrm{}}{}}B_{3m+i,3n+j}^Ua_{3m+i}^{u(i)}a_{3n+j}^{q(j)}a_{3(m+n)+i+j\mathrm{\Delta }}^h.$$
(3.18)
Other Yukawa coupling matrices $`\mathrm{\Gamma }^d`$ and $`\mathrm{\Gamma }^e`$ are obtained similarly.
This general formula (3.18) is very complex, and it is difficult to explore the general property analytically. However, the suppressing effect of mixing in (3.6) and (3.12) realizes the hierarchy among the generations generally. We can expand $`\mathrm{\Gamma }_{ij}^u`$ in terms of the power series of $`r`$’s and extract the lowest order contribution of $`r`$’s. We find that the phenomenologically attractive structure is realized when we choose $`\mathrm{\Delta }=0`$. In this case the leading contribution to $`\mathrm{\Gamma }^u`$ comes from the term with $`m=n=0`$ in (3.18). Omitting the numerical coefficients, it takes the form
$$\mathrm{\Gamma }^u(\mathrm{\Delta }=0)\left(\begin{array}{ccc}1& r_H& r_H^2\\ r_H& r_H^2& r_H^3\\ r_H^2& r_H^3& r_H^4\end{array}\right).$$
(3.19)
The terms with $`m>0`$ and/or $`n>0`$ in (3.18) give the modification to the lowest order term with $`m=n=0`$. In this simplification, $`r_H`$ governs the matrix structure. The masses of each generation are $`m_0𝒪\left(1\right)`$, $`m_1𝒪\left(r_H^2\right)`$ and $`m_2𝒪\left(r_H^4\right)`$.
## 4 Numerical Analysis Under the Minimal Parameter Set
In this section, we show the result of numerical analysis for the typical parameter region to see whether the observed hierarchical structure is realized or not. There are many free parameters in the model; seven weights of unitary fields ($`\alpha `$, $`\beta `$, $`\gamma `$, $`\eta `$, $`\lambda `$, $`\rho `$, $`\sigma `$), two weighs of non-unitary field ($`R`$, $`R^{}`$), seven ratios of VEV’s and couplings ($`r_Q`$, $`r_U`$, $`r_D`$, $`r_L`$, $`r_E`$, $`r_H`$, $`r_H^{}`$), and three Yukawa coupling constants ($`y_U`$, $`y_D`$, $`y_E`$). In practice, it is impossible to make numerical analysis retaining all of this freedom. Based on the consideration of grand unification, we decrease the number of these parameters.
First of all, we assign the same lowest weight $`\alpha `$ to all quarks and leptons, $`\alpha =\beta =\gamma =\eta =\lambda `$. Then the weights of Higgses are $`\rho =\sigma =2\alpha `$. We further assume $`R=R^{}`$ and make the identification $`\mathrm{\Phi }_R=\mathrm{\Phi }_R^{}^{}`$ in the potential (2.15). To check whether our framework has capacity for describing the observed hierarchical structure, we calculate mass ratios of quarks and leptons among the generations and the CKM matrix under this parameter set. In the calculation of these quantities, we do not need values of Yukawa couplings $`y_U,y_D,y_E`$ and normalization constants of Clebsch-Gordan coefficients.
First, we set $`\alpha =0.5`$ and $`R=3`$ which is minimal allowed value of $`R`$, and investigate the dependence of the masses on the ratios $`r`$’s. In the case of $`R=3`$, to realize one generation for Higgses and three generations for quarks, we must restrict values of $`r^{}s`$ to $`|r_H|,|r_H^{}|<\sqrt{3}/2`$ and $`|r_Q|,|r_U|,|r_D|<1/\sqrt{20}`$.
Fig. 1 shows mass rations $`m_1/m_0`$ and $`m_2/m_0`$ of up-type quarks(down-type quarks) for various values of $`r_H`$($`r_H^{}`$) when $`r_Q=0`$ and $`r_U=0`$($`r_D=0`$). It is surprising that the resulting mass hierarchy is much larger than what is naively expected from (3.19) ( $`m_1/m_0𝒪(r_H^2)`$, $`m_2/m_0𝒪(r_H^4)`$). In fact, the order of the observed mass ratios of up-type quarks, $`m_c/m_t=m_1/m_0𝒪(10^3)`$, $`m_u/m_t=m_2/m_0𝒪(10^6)`$, and down-type quarks, $`m_s/m_b=m_1/m_0𝒪(10^2)`$, $`m_b/m_d=m_2/m_0𝒪(10^4)`$), is realized by setting $`r_H0.20.25`$ and $`r_H^{}0.40.65`$.
In Fig. 2, we show the effect of non-vanishing values of $`r_Q`$ and $`r_U`$ with keeping $`r_Q=r_U`$. Evidently the dependence of masses on the values of $`r_Q`$, $`r_U`$ and $`r_D`$ is small. This situation has been vaguely anticipated in the end of the previous section.
Generally the value of ratio $`r`$ is complex. From (2.15) we observe that $`r_Q`$, $`r_U`$ and $`r_D`$ have the common phase $`\mathrm{arg}[\varphi _0/\psi _3]`$, and also $`r_H`$ and $`r_H^{}`$ have the phase $`\mathrm{arg}[\varphi _0/\psi _1^{}]`$. The physically meaningful phase is $`SU(1,1)`$ invariant phase
$$\theta =\mathrm{arg}\left[\psi _3\right]+3\mathrm{arg}\left[\psi _1^{}\right].$$
(4.1)
The result of Fig. 2 (insensitivity to $`\psi _3`$) implies that the effect of this phase $`\theta `$ on masses is small.
The $`R`$ dependence of mass ratios is given in Fig. 3. We see the mass hierarchy is enhanced slowly with an increase in the value of $`R`$. The $`\alpha `$ dependence is given in Fig. 4. The variation of $`\alpha `$ gives sizable effect, although not so large. We reproduce the observed mass hierarchy for quarks by choosing adequate values of $`r_H`$ and $`r_H^{}`$ under the minimal parameter set.
The hierarchical structure of Yukawa coupling matrices $`\mathrm{\Gamma }^u`$ and $`\mathrm{\Gamma }^d`$ reflects itself in the characteristic form of the CKM matrix. Naively, the CKM matrix is expected to take the form
$$V_{\text{CKM}}\left(\begin{array}{ccc}1& r& r^2\\ r& 1& r\\ r^2& r& 1\end{array}\right).$$
(4.2)
As a typical example, we give the numerical result in the case $`R=3`$, $`\alpha =0.5`$, $`r_H=0.2`$, $`r_H^{}=0.45`$, $`r_Q=0.2`$, $`r_U=0.2`$ and $`r_D=0.2`$;
$$\left(\begin{array}{ccc}|V_{ud}|& |V_{us}|& |V_{ub}|\\ |V_{cd}|& |V_{cs}|& |V_{cb}|\\ |V_{td}|& |V_{ts}|& |V_{tb}|\end{array}\right)\left(\begin{array}{ccc}0.975& 0.221& 0.006\\ 0.221& 0.973& 0.074\\ 0.010& 0.088& 0.997\end{array}\right).$$
(4.3)
The mass ratios in this case are $`m_c/m_t=2.5\times 10^3`$, $`m_u/m_t=7.0\times 10^6`$, $`m_s/m_b=1.4\times 10^2`$ and $`m_d/m_b=1.7\times 10^4`$. Although these mass ratios are not enough satisfactory, it is rather surprising that such a simple choice of parameters reproduces the qualitative feature of $`V_{\text{CKM}}`$.
To see the magnitude of the $`CP`$-violating phase, we calculate the Jarlskog determinant $`J`$ which is the rephaseing-invariant parameter. In Fig. 5, we give a result for $`J`$ in the parameter region which gives above CKM matrix. Here, $`\theta `$ defined in (4.1) is $`\theta =\mathrm{arg}[r_Q]3\mathrm{arg}[r_H]`$. The magnitude of $`J`$, $`|J|𝒪(10^7)`$, looks much smaller than what is expected in the SM, $`|J|𝒪(10^5)`$.
Finally, let us mention on masses of charged leptons. Since charged leptons acquire masses through Yukawa couplings to Higgs $`h^{}`$, their mass ratios are given by those of down-type quarks with the replacement $`r_Qr_L`$ and $`r_Dr_E`$. As show in Fig. 2, mass ratios are almost insensitive to these parameters. Therefore the present minimal model predicts essentially equivalent mass ratios for charged leptons and down-type quarks. This prediction is unfavorable to the observation. The model cannot reproduce both of the mass ratios simultaneously.
## 5 Summary and Discussions
In this paper we have proposed a new mechanism for “spontaneous generation of generations” based on the $`SU(1,1)`$ gauge symmetry. This mechanism works well in the grand unified theories. We can naturally understand the coexistence of chiral massless particles and super heavy GUT partners without fine tuning. At the same time, this mechanism yields the hierarchal Yukawa interactions among chiral generations naturally.
We investigated the characteristic feature of the model through the numerical analysis under the minimal parameter set which is motivated by the grand unification and also by the economical purpose. We found that the observed hierarchies of both up-type and down-type quarks are qualitatively reproduced under the reasonable values of the parameters. This minimal model is also capable to realize the characteristic structure of the CKM matrix.
The most serious difficulty of the minimal model is that the model cannot simultaneously reproduced the mass ratios of down-type quarks and charged leptons. The smallness of the predicted value of $`J`$ may be also the problem if we seek the origin of the $`CP`$-violation at the phase of the CKM matrix.
From this survey, we conclude that the minimal model based on the severely restricted set of parameters does not have enough capacity to reproduce the reality. The model must be modified by some relaxation of the constraints. At this stage, the compatibility of the model with the underling grand unification symmetry will give useful guiding principle. We expect that the further attempts along this line will give us deeper insight on the generations of quarks and leptons.
## Acknowledgements
The work of K. I. was supported in part by the Grant-in-Aid of the Ministry of Education, Science, Sports and Culture, Government of Japan (No. 11640282), and Priority area “Supersymmetry and Unified Theory of Elementary Particles” (No. 707).
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# First results of UVES at VLT: abundances in the Sgr dSph Based on public data released from the UVES commissioning at the VLT Kueyen telescope, European Southern Observatory, Paranal, Chile.
## 1 Introduction
In the recent years our ideas on galaxy formation and evolution have considerably developed, and it is generally acknowledged that it is a complex process which may well take different paths in different galaxies. Much attention is being devoted to dwarf spheroidal galaxies, essentially for two reasons: 1) they seem to be relatively simple systems, typically characterized by a single stellar population; in such a system we hope to be able to isolate some of the key ingredients of the phenomenon; 2) Interaction of these dwarf galaxies with large galaxies (such as our own or the Andromeda galaxy) could, in principle, play an important role in shaping the morphology of the large galaxies. The nearest members of the class, the dwarf spheroidals of the Local Group, are close enough that their stars are amenable to detailed analysis with the same techniques employed to study Galactic stars, with the advent of the new 8m class telescopes. In this paper we report on such an observation: the first detailed chemical analysis of two stars in the Sgr dwarf spheroidal based on high resolution spectra obtained with the UVES spectrograph on the ESO 8.2m Kueyen telescope. Ever since the discovery of Sgr (Ibata et al 1994) photometric studies have shown the red giant branch (RGB) of Sgr to be wider than expected for a population with a single age and metallicity. This has been generally interpreted as evidence that Sgr displays a spread in metallicity which is likely due to different bursts of star formation. Ibata et al (1995) found a mean metallicity of \[Fe/H\]=-1.25 and their metallicity distribution displays a spread of over 1 dex. Sarajedini & Layden (1995) found a main population with \[Fe/H\]= -0.52 and suggested the possible existence of a population of \[Fe/H\]$`1.3`$. Mateo et al (1995) provide a mean metallicity of $`1.1\pm 0.3`$, Ibata et al (1997) estimate metallicities in the range $`1.00.8`$, Marconi et al (1998) $`1.580.7`$ , Bellazzini et al (1999) $`2.1,0.7`$. The age of Sgr may not be disentangled from its metallicity, from Main Sequence fitting, Fahlman et al (1996) found acceptable solutions for an age of 10 Gyr and metallicity $`0.8`$ or an age of 14 Gyr and a metallicity $`1.3`$. Clearly if Sgr may not be described as a single population the concepts of age and metallicity loose some of their significance; Bellazzini et al (1999) proposed an extreme scenario in which star formation began rather early and continued for a period longer than 4 Gyr. Foreseeing the potentiality of UVES to perform detailed abundance analysis of these stars, to confirm or refute the photometrically inferred spread in metallicity, we undertook already in 1995 observations of low resolution spectra of photometrically identified (Marconi et al 1998) Sgr candidates. The main purpose was to obtain confirmed radial velocity members of Sgr for subsequent high resolution follow-up with UVES. From the low resolution spectra we also devised a method to obtain crude metallicity estimates from spectral indices defined in the Mg I b triplet region. The two stars were selected from this low-resolution study of Sgr with photometry and abundance estimates which suggested these stars to differ by at least 0.5 dex in metallicity.
## 2 Observations and data reduction
The data were obtained during the Commissioning of UVES and have been released by ESO for public use. The spectra of the two stars, whose basic data is given in Table 1, were taken on the nights 2,3,4 and 6 October 1999. The slit was $`1^{\prime \prime }`$, which provided a resolution of $`43000`$ at 565 nm. We used only the red arm of UVES with the standard setting with central wavelength at 580 nm, which provides spectral coverage from 480 nm to 680 nm The detector was the mosaic of two CCDs composed of a EEV CCD-44 for the bluemost part of the echellogramme and an MIT CCID-20 for the redmost part. Both CCDs are composed of $`4096\times 2048`$ square pixels of 15 $`\mu m`$ side. We used a $`2\times 2`$ on–chip binning, without any loss of resolution, given the relatively wide slit used. For each of the two stars three one-hour exposures were collected, under median seeing conditions, allowing to reach a signal to noise ratio of $`30`$ at 510 nm, confirming the excellent performance of the instrument (D’Odorico et al, 2000).
The data was reduced using the ECHELLE context of MIDAS; reduction included background subtraction, cosmic ray filtering, flat fielding, extraction, wavelength calibration and order merging. Each of the CCDs of the mosaic was reduced independently. The r.m.s. of the calibration was typically of the order of 0.2 pm for each order and in all cases less than 0.3 pm. Flat-fielding was highly successful and the single echelle orders were rectified to within $`5`$% by this process, except in the vicinity of CCD defects, where the correction was not always satisfactory. The three spectra available for each star were coadded without any shift in wavelength. By cross-correlation we estimated the shift between any pair of spectra to be less than 0.5 pixel, so we decided not to perform any shift. This results in a very slight degradation of the resolution, which is not an issue for our analysis. The differences in barycentric correction were at most of 0.1 kms<sup>-1</sup>, so no appreciable shift was expected from this cause either. The normalized and merged spectra were plotted superimposed on preliminary synthetic spectra for the purpose of line identification. A single velocity shift was adequate for all the spectra, confirming that the internal accuracy of our wavelength scale is better than 0.2 km/s. Portions of the normalized spectra of the two stars are displayed in figures 1 and 2.
## 3 Analysis
Our analysis is standard and essentially based on LTE model atmospheres. For each star we estimated effective temperature from $`(VI)_0`$ through the colour–T<sub>eff</sub> calibration for giants of Alonso et al (1999). We adopted log g = 2.5 for both stars, this value is compatible with the position of the stars in the colour–magnitude diagram, whichever the isochrone set considered. Furthermore we verified a posteriori that this gravity nearly satisfied the ionization equilibrium of both Fe I/Fe II and Ti I/Ti II . The atmospheric parameters are summarized in Table 1. The mean abundances for all the elements are given in Table 2.
With these parameters we computed model atmospheres using the ATLAS9 code (Kurucz 1993) using opacity distribution functions with microturbulence of 2 km/s and suitable metallicity.
We began by deriving the Fe abundance for both stars. We picked a set of lines which our preliminary synthetic spectra predicted to be substantially unblended, which had accurate laboratory or theoretical gf values and which spanned a range of line strengths. We excluded from the analysis lines stronger than $`\mathrm{log}(EW/\lambda )=4.8`$ to avoid an excessive dependence on microturbulence and lines weaker than $`\mathrm{log}(EW/\lambda )=6.1`$ to avoid lines which are too noisy.
For these lines we measured equivalent widths by fitting a gaussian with the iraf task splot and used these and the model atmosphere as input to the WIDTH9 code (Kurucz 1993). Some of the lines were removed from the analysis because they provided highly discrepant abundances. The microturbulence was determined by imposing that strong lines and weak lines give the same abundance. The results are given in Table 3. The excitation equilibrium for Fe I is very nearly satisfied for both stars (slopes are -0.05 dex/eV for star 139 0.13 dex/eV for star 143), we did not adopt an excitation temperature, so that the equilibria are better satisfied, since our lines cover a range of only about 2.7 eV, this means, that in the worst case (star 143) the slope predicts a difference of slightly over 0.3 dex between the highest and lowest excitation lines, this is $`2\times `$r.m.s.
As a check of our metallicity we derived Fe abundances also using MARCS models computed by Plez et al (1992). The difference is not significative: for both stars the MARCS models provide an Fe I abundance which is 0.02 dex larger than that provided by the ATLAS models, while for Fe II it is 0.07 dex larger for star 139 and 0.05 dex larger for star 143.
Having fixed the metallicity and the microturbulence we proceeded to determine all the other abundances with the same method, again we disregarded the too strong and too weak lines, with a few exceptions, such as Sc, Cu and Ba, for which only one or two lines were available, in order to get information on as many elements as possible. The line data and abundances are given in Table 4. In addition for some blended lines we resorted to spectrum synthesis using the same model-atmosphere and the SYNTHE code (Kurucz 1993). We did not take into account hyperfine splitting (HFS) for Sc, V, Mn, Co and Cu, however given that the abundances of these elements are coherent with those of other elements we do not expect corrections due to HFS to be very large. For Eu we determined the abundance from the Eu II 664.5 nm line, taking into account HFS splitting, the relevant data is given in Table 5. The error on the Eu abundances estimated from the quality of the fit is 0.15 dex.
The main result of this analysis confirms the impression gathered by a direct comparison of the spectra of the two stars: the stars are very nearly identical, the few differences in their spectra are quite likely determined by slightly different T<sub>eff</sub> and log g, but not by chemical composition.
From the measure of the line centers of the unblendend lines used for abundances we determined the radial velocity for the two stars. We obtain the following heliocentric radial velocities $`133.8\pm 0.8`$ kms<sup>-1</sup> from 60 lines for star 139 and $`143.8\pm 0.8`$ kms<sup>-1</sup> from 57 lines for star 143, the quoted error is just the r.m.s. The measurement of the position of the atmospheric Na I D emission lines allowed to estimate the zero point shift to be less than 0.1 kms<sup>-1</sup>, this, coupled with the excellent reproducibility of the wavelength scale from night to night induced us to assume a null zero–point shift. These heliocentric radial velocities support membership to Sagittarius, Ibata et al (1995) give a mean heliocentric radial velocity for Sgr of $`140\pm 2`$ kms<sup>-1</sup> and Ibata et al (1997) find the intrinsic velocity dispersion to be $`11.4\pm 0.7`$ kms<sup>-1</sup> and constant across the face of the galaxy. The N–body model of Sgr computed by Helmi & White (2000) displays a similar velocity dispersion, if only the stars with 100 kms$`{}_{}{}^{1}v_{hel}180`$ kms<sup>-1</sup> are included, as done by Ibata et al. However if this condition is relaxed the velocity dispersion turns out to be much larger, due to the contribution of stars in the debris streams. Given the above considerations it is not surprising that we find a difference of 10 kms<sup>-1</sup> between our stars. Our measured radial velocities compare quite well with those measured from our EMMI low resolution spectra (147 kms<sup>-1</sup> for star 139 and 154 kms<sup>-1</sup> for star 143, both are accurate to $`\pm 15`$ kms<sup>-1</sup>).
We cannot rule out the possibility that the stars studied here belong to the Bulge, rather than Sgr. If they were at a distance of 8.5 Kpc, rather than 25 Kpc their log g should be $`0.5`$ dex higher than what we assumed, but such a difference is within the errors of the analysis. However the radial velocity ought to be a very good discriminant. By looking at Figure 1 of Ibata et al (1995) we see that the distribution of radial velocity of Bulge stars in directions which do not intercept the Sgr dSph, show a vanishingly small number of stars at the radial velocity of Sgr. Also the chemical composition suggests that the two stars do not belong to the Bulge: in fact Bulge stars are expected, theoretically, to have \[$`\alpha `$/Fe\]$`>0`$, even at solar metallicities. Observationally the situation is not so clear, however our distinctly solar \[$`\alpha `$/Fe\] is a clue against Bulge membership.
## 4 Discussion.
The metallicity of the two stars examined here is higher than all previous photometric estimates. Although it is possible that we happened to select two members of the high–metallicity tail of Sgr, this position is hardly tenable, the event of finding two such stars in a 9 square arcmin field must be quite rare. It is more likely that Sgr actually possesses a population, perhaps the main population, this metal–rich. For our two stars the Schlegel et al (1998) maps provide $`E(BV)=0.14`$. By comparison Marconi et al (1998) used E(B-V)=0.18. The fact that the actual reddening could be 0.04 less than this could explain why the metallicity we find is 0.3 dex higher than the highest metallicity estimated by Marconi et al.
Quite obviously our results do not rule out the existence of a more metal–poor population. Preliminary results of abundance analysis in Sgr are given by Smecker-Hane & McWilliam (1999), who find two metal–poor Sgr member stars, with \[Fe/H\]$`=1.41`$ and \[Fe/H\]$`=1.14`$. It is interesting that out of 11 stars analyzed by them 7 have metallicities in the range $`0.6<[\mathrm{Fe}/\mathrm{H}]<0.2`$, two are metal-poor and two are metal rich (\[Fe/H\]$`0.`$). The 7 stars of intermediate metallicity, which should be analogous to the two under study here, show solar abundance ratios and no enhancement of $`\alpha `$ elements, in agreement with our findings. Also the Na abundance displays a similar pattern: for all their stars, except the two metal–poor ones Na is over-deficient with respect to iron by 0.3 – 0.5 dex. Unfortunately, these results have not been published in a more detailed form and we lack information on the temperatures and luminosities of the the stars considered by Smecker-Hane & McWilliam so we do not know if we are comparing similar giants.
It is also interesting to compare the present results with the abundances of the two Sgr planetary nebulae He 2-436 and Wray 16-423, studied by Walsh et al (1998). The only element in common in the two analysis is O, for which Walsh et al find \[O/H\]$`=0.64\pm 0.08`$ and \[O/H\]$`=0.62\pm 0.07`$ for He 2-436 and Wray 16-423, respectively. Our result for Sgr 143 is about 0.2 dex higher, but it is also more uncertain, because it is based on a single weak line, which is also very sensitive to gravity. An increase of gravity of 0.5 dex results in an increase of O abundance of 0.25 dex. O should be only marginally affected during AGB evolution, so that the O abundance in the PN ought to be quite close to that in the progenitor star. Walsh et al (1998) argued that their abundances suggested a mild enhancement of O over Fe, because they assumed $`0.8`$ to be the mean \[Fe/H\] of Sgr. Another scenario appears more likely, in view of our results: a solar O/Fe ratio, which suggests that the PNe have \[Fe/H\]$`0.5`$.
Having established that the two stars are quite similar in atmospheric parameters and abundances we must explain why their photometry is different and why the metallicity estimated from the low resolution spectra for star 143 is far lower than the one derived here. We consider 5 possibilities: 1) errors in $`V`$; 2) errors in $`VI`$; 3) different reddening; 4) different age; 5) different distance. Let us examine all of these cases.
That a difference of 0.18 mag in $`V`$ may be due to the photometric error may be discarded since this is a factor of ten larger than the photometric error of Marconi et al (1998). An error in $`VI`$ is more likely; a 0.03 – 0.04 mag error in $`VI`$ would allow to slide sideways one of the two stars in the colour-magnitude diagram in such a way that both stars lie on the same isochrone, since the RGB, in this range of $`VI`$ is very steep. The implied difference in T<sub>eff</sub> is of $`100`$ K, the errors of our analysis.
Differential reddening seems unlikely for three reasons. The dust maps of Schlegel et al (1998) give $`E(BV)=0.14`$ for both stars, suggesting that the reddening of the two stars is the same within 0.01 mag. The absence of detectable amounts of HI in Sgr (Burton & Lockman, 2000) also argues against a differential reddening. If the 0.18 mag difference in $`V`$ were due to reddening it would imply a difference of almost 0.08 mag in $`VI`$, i.e. $`200`$ K in T<sub>eff</sub>. Although such a difference is within the errors of the present analysis and cannot be ruled out, it does seem somewhat unlikely, given the similarity of the two spectra.
An age difference of $`1`$ Gyr could be enough to explain the difference in the photometry of the two stars. A larger age spread would be necessary to explain the width of the RGB, like in the scenario proposed by Bellazzini et al (1999). Although such a possibility is attractive, it appears somewhat contrived and it is not so clear that star formation may continue for several Gyrs without resulting in a spread in metallicity, as well as ages.
A distance difference of about 2Kpc would be enough to explain the difference in $`V`$. This value is not unreasonable, Ibata et al (1997), estimate the half-brightness depth of Sgr to be about 1.2 kpc . It is interesting that recent N–body simulations by Helmi & White (2000) support a considerable depth of Sgr: inspection of their figure 2 shows that the bulk of their model for Sgr has a depth of about 2 Kpc, however considering the debris shed during previous orbits, one has a sizeable population over a depth of 10 Kpc. Further support to the possibility that the two stars have a different distance comes from inspection of the Na I D lines (Fig. 2), three interstellar components belonging to our Galaxy are evident in both the spectra of Sgr 143 and of Sgr 139 at radial velocity +16.5 kms<sup>-1</sup>, $`+28.0`$ kms<sup>-1</sup> and $`+47.3`$ kms<sup>-1</sup>; while the Na I D lines of star 143 appear symmetric and there is no hint of an interstellar component at the radial velocity of Sgr, the lines of star 139 show a weak but definite asymmetry, which we interpret as a weak interstellar line associated with Sgr. Star 139 is in fact the fainter of the two and hence the most distant, according to this interpretation, this would explain why the interstellar Na I D lines appear in its spectrum but not in the spectrum of star 143, which would then be in the side of Sgr nearer to us.
So of the five possibilities considered only the photometric error in $`V`$ and the differential reddening are discarded. We may not decide which is the correct one with the present data, new accurate photometric measurements will allow to settle at least the issue of errors in $`VI`$. However we consider that the distance difference is the most likely explanation, because it is the simplest and is supported by several arguments. This suggests that the non-negligible line of sight depth of Sgr could explain at least a part of the width of the RGB of Sgr. Up to now all investigators have adopted a unique distance modulus for Sgr, in order to compare their photometry to fiducial ridge lines of Galactic clusters or to theoretical isochrones. This assumption may prove to be bit too naive.
A full discussion of the metallicity estimates from low resolution spectra shall be given elsewhere. Suffice to say here that the method of estimating abundances from low resolution needs a relatively high S/N ratio. In the case of star 139 the metallicity derived from high resolution analysis coincides with that estimated from low resolution to within the errors of the latter. We verified that the degraded UVES spectrum is very similar to the low resolution EMMI spectrum. The indices measured on this degraded spectrum yield in fact almost the same abundance provided by those measured on the low resolution spectrum. In the case of star 143 instead the method has been applied to a spectrum of too low signal to noise ratio, in this case the degraded UVES spectrum bears almost no resemblance to the low resolution spectrum, except for the strongest feature of the Mg I b triplet, which was enough to provide the correct radial velocity for this star.
The ratios of all elements are essentially solar, noticeable exceptions are: Na which is overdeficient with respect to iron, and the heaviest elements Ba, La, Ce, Nd, Eu, which appear over-abundant while Y appears underabundant. Such anomalies are not readily interpretable, deep mixing would produce an enhanced Na and low O and Mg, at variance to what is observed. While it would be tempting to interpret the overabundance of heavy elements as due to s-process enrichment, the stars do not appear luminous enough to be on the thermally pulsating asymptotic giant branch, where this mechanism is operative. Moreover, s-process enrichement would produce also a high Y abundance and no Eu (which is thought to be a “pure” r-process element), at variance to the low Y and high Eu abundances observed here. On the other hand, these stars could have been born in r-process enhanced material (sugested by the Eu enhancement). Howevever this seems also quite implausible since the r-process is thought to take place in Type II supernovae which also produce large amounts of O and other $`\alpha `$-elements, which are not observed to be enhanced in our stars. This surprising pattern is reminiscent to what is observed in the young supergiants in both Magellanic Clouds, where the ratios of the moderate-mass s-process elements Y and Zr to iron are essentially solar, whereas the heavier species Ba to Eu are overabundant by ratios \[X/Fe\] of the order of 0.3 and 0.5 dex respectively in the LMC and SMC (Hill et al. 1995, Hill 1997, Luck et al 1998). In the Magellanic Clouds also, we are at loss of an explanation for this behaviour (see discussion in Hill 1997). Note that our two Sgr giants have the same overall metallicity as the LMC young population, and that the heavy elements overabundances are also of the same order as in the LMC.
###### Acknowledgements.
We are extremely grateful to S. D’Odorico, H. Dekker and the whole UVES team for the conception and construction of this wonderful instrument.
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# A THEORETICAL LIGHT-CURVE MODEL FOR THE 1985 OUTBURST OF RS OPHIUCHI
## 1. INTRODUCTION
RS Oph is one of the well observed recurrent novae and characterized by a long orbital period of 460 days (Dobrzycka & Kenyon (1994)), a relatively short recurrence period of $`10`$—20 yrs, and a companion (M-giant) star underfilling the Roche lobe but losing its mass by massive stellar winds (e.g., Dobrzycka et al. (1996)). Historically, RS Oph underwent five outbursts, in 1898, 1933, 1958, 1967, and 1985, with the light curves very similar each other (e.g., Rosino (1987)). The latest 1985 outburst has been observed at all wave lengths from radio to X-rays (e.g., papers in the 1985 RS Oph conference proceedings ed. by Bode (1987)).
Although there had been intense debates on the mechanism of RS Oph outbursts (e.g., Livio et al. 1986b ; Webbink et al. (1987)), various observational aspects favor thermonuclear runaway (TNR) models on a very massive white dwarf (WD) (e.g., Anupama & Mikołajewska (1999) for a recent summary). Rapid decline rates of the light curves indicate a very massive WD close to the Chandrasekhar limit. Kato (1991) has first calculated RS Oph light curves for the WD masses of 1.33, 1.35, 1.36 and 1.37 $`M_{}`$ and found that the light curve of the $`1.36M_{}`$ model is in better agreement with the observational light curve of RS Oph than the other lower mass models.
However, structures of bloated envelopes on white dwarfs have been drastically changed (Kato & Hachisu (1994); Kato (1999)) after the advent of the new opacity (e.g., Iglesias & Rogers (1996)). Thus, we have been recalculating light curves of the recurrent novae and have again determined the mass, composition, and distance (Hachisu & Kato (1999) for T CrB; Hachisu et al. 2000a , 2000b for U Sco; Hachisu & Kato (2000) for V394 CrA). Moreover, we have added new parts to Kato’s (1991) simple model in order to follow later stages of the light curves, i.e., the irradiation effects of the red giant (RG) companion star and the accretion disk (ACDK). As a result, much more reliable quantities are derived from the light curve fittings than Kato’s ones.
Type Ia supernovae (SNe Ia) are one of the most luminous explosive events of stars. Recently, SNe Ia have been used as good distance indicators which provide a promising tool for determining cosmological parameters because of their almost uniform maximum luminosities (Riess et al. (1998); Perlmutter et al. (1999)). These both groups derived the maximum luminosities ($`L_{\mathrm{max}}`$) of SNe Ia completely empirically from the shape of the light curve (LCS) of nearby SNe Ia, and assumed that the same $`L_{\mathrm{max}}`$–LCS relation holds for high red-shift SNe Ia. To be sure of any systematic biases, the physics of SNe Ia must be understood completely. By far, one of the greatest problems facing SN Ia theorists is the lack of a real progenitor (e.g., Livio (1999) for a recent review). Finding a reliable progenitor is urgently required in SN Ia research. Recurrent novae are probably the best candidate for this target (e.g., Starrfield, Sparks, & Truran 1985; Hachisu et al. 1999b; Hachisu, Kato, & Nomoto 1999a).
In §2, we present our methods for calculating theoretical light curves based on the optically thick wind theory developed by Kato & Hachisu (1994). Fitting results of the modeled light curves to the observations are shown in §3. Our results indicate that RS Oph is a strong candidate for Type Ia supernovae. Discussion follows in §4 especially for the distance to RS Oph because we obtain a rather short distance of 0.6 kpc.
## 2. THEORETICAL LIGHT CURVES
In the TNR model, it has been established that the WD photosphere expands to a red giant size and then gradually shrinks to the original size of the WD in quiescent phase (e.g., $`0.0039R_{}`$ for $`M_{\mathrm{WD}}=1.35M_{}`$) with the bolometric luminosity being kept near the Eddington limit. The visual luminosity reaches its maximum at the maximum expansion of the photosphere and then gradually decays to the level in quiescent phase. Accordingly, the main emitting region of the WD photosphere moves from visual into soft X-rays through ultraviolet. Optically thick winds, which are blowing from the WD during the outbursts, play a key role in determining the nova duration because a large part of the envelope mass is blown in the outburst wind. The development of the WD photosphere can be followed by a unique sequence of steady-state, optically thick, wind solutions as shown by Kato & Hachisu (1994).
We have calculated such sequences of $`M_{\mathrm{WD}}=1.3`$, 1.35, 1.36, 1.37, and $`1.377M_{}`$ by assuming the hydrogen content $`X=0.70`$, 0.50, 0.35, the helium content $`Y=0.28`$, 0.48, 0.63, respectively, and the metallicity $`Z=0.02`$ (solar metallicity) and then, obtained the optical light curves. Here, we have used the updated OPAL opacity (Iglesias & Rogers (1996)). The envelope mass is decreasing due to wind mass loss ($`\dot{M}_{\mathrm{wind}}`$) and hydrogen shell burning ($`\dot{M}_{\mathrm{nuc}}`$), i.e.,
$$\frac{d}{dt}\mathrm{\Delta }M=\dot{M}_{\mathrm{acc}}\dot{M}_{\mathrm{wind}}\dot{M}_{\mathrm{nuc}},$$
(1)
where $`\dot{M}_{\mathrm{acc}}`$ is the mass accretion rate of the WD and assumed to be $`\dot{M}_{\mathrm{acc}}=1\times 10^7M_{}`$ yr<sup>-1</sup> during the outburst. The photospheric temperature $`T_{\mathrm{ph}}`$, the photospheric radius $`R_{\mathrm{ph}}`$, the photospheric velocity $`v_{\mathrm{ph}}`$, the wind mass loss rate $`\dot{M}_{\mathrm{wind}}`$, and the nuclear burning rate $`\dot{M}_{\mathrm{nuc}}`$ are unique functions of the envelope mass $`\mathrm{\Delta }M`$. Integrating equation (1), we can follow the development of the envelope mass $`\mathrm{\Delta }M`$ and obtain various physical quantities of the WD envelope. The mass lost by the wind, $`\mathrm{\Delta }M_{\mathrm{wind}}`$, and the mass added to the helium layer of the WD, $`\mathrm{\Delta }M_{\mathrm{He}}`$, are calculated from
$$\mathrm{\Delta }M_{\mathrm{wind}}=\dot{M}_{\mathrm{wind}}𝑑t,\mathrm{\Delta }M_{\mathrm{He}}=\dot{M}_{\mathrm{nuc}}𝑑t.$$
(2)
Assuming a black-body photosphere, we have calculated the $`V`$-magnitude of the WD photosphere with a response function given by Allen (1973). For simplicity, a limb-darkening effect is neglected. By fitting with the observational points in the early 4 days of the superposed $`V`$ light curve (e.g., Rosino (1987)) and the 1985 $`V`$ light curve, we have determined the WD mass of $`M_{\mathrm{WD}}=1.35\pm 0.01M_{}`$ and the apparent distance modulus of $`m_0=11.09`$ as shown in Figure 1. The early 4 days $`V`$-magnitude is determined mainly by the bloated WD photosphere, and depends sensitively on the WD mass, but depends hardly on the hydrogen content, $`X`$, the binary parameters such as the mass ratio, or the accretion disk shape, so that the WD mass determination itself is rather robust (see also Kato (1999)).
To reproduce the $`V`$ light curve in the late stage ($`t4`$—100 days after the optical maximum), we have to include the contributions of the irradiated RG and the irradiated ACDK. As suggested by Dobrzycka et al. (1996), we have assumed that the RG lies well within the inner critical Roche lobe, i.e.,
$$R_{\mathrm{RG}}=\gamma R_2^{}(\gamma <1),$$
(3)
where $`R_2^{}`$ is the effective radius of the inner critical Roche lobe for the RG component (e.g., Eggleton (1983)) and $`\gamma `$ is a numerical factor. Here, we assume 50% efficiency of the irradiation ($`\eta _{\mathrm{RG}}=0.5`$). The nonirradiated photospheric temperature of the RG is a parameter for fitting, and has been determined to be $`T_{\mathrm{ph},\mathrm{RG}}=3100`$ K for $`\gamma =0.4`$.
The orbit of the companion star is assumed to be circular. The ephemeris for the inferior conjunction of the giant is 2,444,999.9$`+460\times E`$ (Dobrzycka & Kenyon (1994)). The light curves are calculated for six cases of the companion mass, i.e., $`M_{\mathrm{RG}}=0.5`$, 0.6, 0.7, 0.8, 1.0, and $`1.2M_{}`$. Since we obtain similar light curves for all of these six masses, we show here only the results for $`M_{\mathrm{RG}}=0.7M_{}`$. In this case, the separation is $`a=318.5R_{}`$, the effective radii of the inner critical Roche lobes for the WD component and the RG component are $`R_1^{}=139.1R_{}`$ and $`R_2^{}=103.1R_{}`$. If $`\gamma =0.4`$, then we have $`R_2=0.4R_2^{}40R_{}`$.
We have also included the luminosity coming from the ACDK irradiated by the WD photosphere when the accretion disk exists during the outburst. Here, we assume 50% efficiency of the ACDK irradiation ($`\eta _{\mathrm{DK}}=0.5`$). The viscous heating is neglected because it is rather smaller than that of the irradiation effects. The temperature of the unheated surface of the ACDK including the rim is assumed to be $`T_{\mathrm{disk}}=2000`$ K. We have checked two other cases of $`T_{\mathrm{disk}}=1000`$ and 0 K and have found no significant differences in the light curves. The shape of the ACDK is assumed to be axisymmetric and approximated by
$$R_{\mathrm{disk}}=\alpha R_1^{},$$
(4)
and
$$h=\beta R_{\mathrm{disk}}\left(\frac{\varpi }{R_{\mathrm{disk}}}\right)^\nu ,$$
(5)
where $`R_{\mathrm{disk}}`$ is the outer edge of the accretion disk, $`h`$ the height of the surface from the equatorial plane, $`\varpi `$ the distance on the equatorial plane from the center of the WD, and $`\alpha `$ and $`\beta `$ are both numerical factors which should be determined by fitting. The power of $`\nu `$ is assumed to be $`\nu =9/8`$ from the standard disk model. We have checked the dependency of the light curves on the parameter $`\nu `$ by changing from $`\nu =9/8`$ to $`\nu =2`$ but cannot found any significant differences when the disk rim is as small as $`\beta =0.01`$—0.05.
## 3. RESULTS
The early 4 days $`V`$ light curve can be reproduced with the WD mass of $`M_{\mathrm{WD}}=1.35\pm 0.01M_{}`$ as already mentioned in the previous section. The late phase $`V`$ light curve ($`t4`$—100 days after maximum) indicates strong irradiations both of the ACDK and of the RG. If we assume a lobe-filling RG companion ($`\gamma =1`$), then the irradiation of the RG is too luminous to be compatible with the observational light curves. On the other hand, a smaller size of the RG such as $`\gamma =0.4`$ ($`40R_{}`$) gives a reasonable fit with the observational one as shown in Figure 2. Here, we assume the inclination angle of $`i=30\mathrm{°}`$ (Dobrzycka & Kenyon (1994)).
To fit the late phase light curve, we have finally adopted the disk parameters of $`\alpha =0.01`$ and $`\beta =0.01`$ as shown in Figure 2. Here, the thick and thin solid lines denote the cases of ($`\alpha =0.01`$, $`\beta =0.01`$) and ($`\alpha =0.1`$, $`\beta =0.01`$), respectively. Something between these two can roughly reproduce the light curve of the 1985 outburst but the case of ($`\alpha =0.01`$, $`\beta =0.01`$) is much better to reproduce the superposed light curve given by Rosino (1987).
We have also calculated UV light curves with a response function of 911Å—3250Å to fit the UV data (Snijders 1987a ). Figure 3 shows three cases of UV light curves for $`X=0.35`$, 0.50, and 0.70, with the parameters of $`\alpha =0.01`$ and $`\beta =0.01`$. The calculated UV light curves can reproduce well the UV observations for $`X=0.70`$ and the distance to RS Oph of $`d=0.57`$ kpc. Thus, we have obtained the absorption of $`A_V=11.095\mathrm{log}(570/10)=2.3`$ for $`d=0.57`$ kpc. This absorption of $`A_V=2.3`$ is consistent with the color excess of $`E(BV)=0.73`$ (Snijders 1987b ) because of $`A_V=3.1E(BV)`$, although this distance of $`0.6`$ kpc is not consistent with another estimation of 1.6 kpc from the hydrogen column density (Hjellming et al. (1986)). The UV light curves depend hardly on the disk parameters of $`\alpha `$, $`\beta `$, or $`\nu `$ when $`\alpha 0.01`$ mainly because the UV light is coming from the innermost part of the ACDK, but depend on the irradiation efficiency of $`\eta _{\mathrm{DK}}`$. We have examined that the distance becomes 0.70 kpc for $`\eta _{\mathrm{DK}}=1.0`$ or becomes 0.50 kpc for $`\eta _{\mathrm{DK}}=0.25`$.
The envelope mass at the optical maximum is estimated to be $`\mathrm{\Delta }M=2.1\times 10^6M_{}`$, which is indicating a mass accretion rate of $`1.2\times 10^7M_{}`$ yr<sup>-1</sup> during the quiescent phase between the 1968 and the 1985 outbursts. About 90% of the envelope mass ($`1.9\times 10^6M_{}`$) has been blown off in the optically thick wind and the residual 10% ($`0.2\times 10^6M_{}`$) has been left and added to the helium layer of the WD. The residual mass itself depends on both the hydrogen content $`X`$ and the WD mass. It is roughly ranging from 7% ($`X=0.70`$) to 14%($`X=0.50`$) including ambiguities of the WD mass ($`M_{\mathrm{WD}}=1.35\pm 0.01M_{}`$). Therefore, the net mass increasing rate of the WD is $`1.2\times 10^8M_{}`$ yr<sup>-1</sup>, which meets the condition for Type Ia supernova explosion if the WD core consists of carbon and oxygen. Thus, we conclude that RS Oph is an immediate progenitor of Type Ia supernova if the donor is massive enough to supply fuel to the WD until the WD reaches $`1.378M_{}`$ for explosion (Nomoto, Thielemann, & Yokoi 1984).
## 4. DISCUSSION
Assuming a tilting accretion disk around the WD, Hachisu & Kato (1999) have reproduced theoretically the second peak of T CrB outbursts. Such a radiation induced, tilting disk instability sets in if the condition
$`{\displaystyle \frac{\dot{M}_{\mathrm{acc}}}{3\times 10^8M_{}\text{ yr}^1}}`$ $``$ $`\left({\displaystyle \frac{R_{\mathrm{disk}}}{R_{}}}\right)^{1/2}\left({\displaystyle \frac{L_{\mathrm{bol}}}{2\times 10^{38}\text{ erg s}^1}}\right)`$ (6)
$`\times `$ $`\left({\displaystyle \frac{R_{\mathrm{WD}}}{0.004R_{}}}\right)^{1/2}\left({\displaystyle \frac{M_{\mathrm{WD}}}{1.35M_{}}}\right)^{1/2}`$ (7)
is satisfied (Southwell et al. (1997)). Our estimated accretion rate of the WD in RS Oph is about $`\dot{M}_{\mathrm{acc}}1.2\times 10^7M_{}`$ yr<sup>-1</sup>, which does not meet the above condition, so that we do not expect the radiation induced instability in RS Oph system. Thus we can explain the reason why a second peak as seen in T CrB system does not appear in RS Oph system.
Very soft X-rays were observed 251 days after the optical maximum in the 1985 outburst of RS Oph (Mason et al. (1987)). Mason et al. estimated a black-body temperature of $`3.5\times 10^5`$ K, a total energy flux of $`L_\mathrm{X}1\times 10^{37}`$ erg s$`{}_{}{}^{1}(d/1.6`$ kpc)<sup>2</sup> for a hydrogen column density of $`N_\mathrm{H}=3\times 10^{21}`$ cm<sup>-1</sup>, and concluded that these soft X-rays come from a WD photosphere with a steady hydrogen shell-burning. However, steady hydrogen shell-burning has stopped 122 days after the optical maximum in our model of $`X=0.70`$ and $`M_{\mathrm{WD}}=1.35M_{}`$ (Fig. 3), thus indicating an accretion luminosity instead of hydrogen shell-burning. Using the black-body temperature of $`T3.5\times 10^5`$ K and the WD radius of 0.0039 $`R_{}`$, we estimate a total luminosity of $`L_\mathrm{X}=4\pi R_{\mathrm{WD}}^2\sigma T^4200L_{}`$. This low luminosity corresponds to the accretion luminosity of $`\dot{M}_{\mathrm{acc}}=2L_\mathrm{X}R_{\mathrm{WD}}/GM_{\mathrm{WD}}0.5\times 10^7M_{}`$ yr<sup>-1</sup>, which is roughly consistent with our estimated value of $`\dot{M}_{\mathrm{acc}}1.2\times 10^7M_{}`$ yr<sup>-1</sup>. Then, a rather short distance of 0.45 kpc to RS Oph can be derived from $`8\times 10^{35}1\times 10^{37}(d/1.6`$ kpc)<sup>2</sup>. Thus, our short distance of 0.6 kpc is consistent with the soft X-ray observation. Moreover, the recent optical and ultraviolet observations also suggest 100—600 $`L_{}`$ for a hot component of RS Oph (e.g., Dobrzycka et al. (1996)).
X-rays were also observed at quiescent phase (Orio (1993)), but the flux is too low to be compatible with the TNR model. One possible explanation is an absorption by the massive cool wind (e.g., Shore et al. (1996)) from the RG component as suggested by Anupama & Mikołajewska (1999).
Our ejected mass ($`1.9\times 10^6M_{}`$) is consistent with the observed ejected mass of (1.2—1.8)$`\times 10^6M_{}`$ estimated by Bohigas et al. (1989) except their distance of 1.6—2.0 kpc to RS Oph. Using their equation (8) and $`V_s=200`$ and $`V_e=2000`$ km s<sup>-1</sup>, we obtain the cool wind mass of $`6\times 10^5M_{}`$, indicating the wind mass loss rate from the RG of $`3\times 10^6M_{}`$ yr<sup>-1</sup>. It is reasonable that about one twentieth of the cool wind accretes to the WD (e.g., Livio et al. 1986a ) to power the TNR explosion.
Hjellming et al. (1986) estimated the distance to RS Oph of $`d=1.6`$ kpc from H I absorption line measurements, using the hydrogen column density of $`N_\mathrm{H}=(2.4\pm 0.6)\times 10^{21}`$ cm<sup>-2</sup> and the relation of $`N_\mathrm{H}/(T_\mathrm{s}d)=1.59\times 10^{19}`$ cm<sup>-2</sup> K<sup>-1</sup> kpc<sup>-1</sup> together with $`T_\mathrm{s}=100`$ K. If a large part of this hydrogen column density stems not from the Galactic one but from the local one belonging to RS Oph, the distance is overestimated. For example, in the direction of the recurrent nova U Sco, Kahabka et al. (1999) reported a hydrogen column density of (3—4)$`\times 10^{21}`$ cm<sup>-2</sup>, which is much higher than the Galactic one of $`1.4\times 10^{21}`$ cm<sup>-2</sup>. Here, U Sco is at least 2 kpc above the Galactic plane ($`b=22\mathrm{°}`$). If RS Oph is the same case as U Sco, the distance of 1.6 kpc is an upper limit and a much shorter distance is possible.
We thank the anonymous referee for many critical comments that helped to improve the content of the manuscript. This research has been supported in part by the Grant-in-Aid for Scientific Research (09640325, 11640226) of the Japanese Ministry of Education, Science, Culture, and Sports.
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# Chiral vacuum alignment in dense QCD
## Abstract
We discuss an interesting possibility of nontrivial, quark-mass induced chiral vacuum alignment in color-flavor locking phase of cold, dense QCD. With the simplifying assumption that the gaps for quarks are identical to those of antiquarks and the light quark masses are given by $`m_u=m_d`$ and $`m_s/m_d=15`$, we find the true chiral vacuum can align only to one of the discrete number of directions in the continuum of chiral vacua. The alignment depends on the size of the diquark condensates, and the vacuum transitions between the discrete vacua caused by the evolution of the diquark condensates can be first order phase transition with vanishing or nonvanishing latent heat, depending on the vacua involved. It is also shown that as $`\mu \mathrm{}`$, where $`\mu `$ is the baryon chemical potential, parity is spontaneously broken through the vacuum alignment.
The chiral symmetry $`\text{SU}_L(3)\times \text{SU}_R(3)`$ of quantum chromodynamics (QCD) with three massless quark flavors is spontaneously broken to $`\text{SU}_{L+R}(3)`$ at low energies. Any point in the continuum of the chiral vacua, the coset space $`[\text{SU}_L(3)\times \text{SU}_R(3)]/\text{SU}_{L+R}(3)`$, can be chosen as the vacuum, since any two vacua are equivalent as far as physics is concerned. However when quarks receive small (current) masses the chiral symmetry is then explicitly broken, and the degeneracy of the vacua is lifted. The quark mass term then picks up the lowest energy state in the continuum of the vacua as the true chiral vacuum. This is called Dashen’s chiral vacuum alignment . To determine the true vacuum one has to find the quark-mass induced potential that lifts the degeneracy of the vacua. Since the Nambu-Goldstone bosons of the chiral symmetry breaking are small fluctuations about a point in the continuum of the vacua, the potential is nothing but the meson mass term in the chiral Lagrangian of the nonlinearly realized Nambu-Goldstone bosons. Therefore, the true vacuum is the one that minimizes the potential
$`V(\mathrm{\Sigma })=_m(\mathrm{\Sigma })\text{Re}[\text{Tr}(m\mathrm{\Sigma })],`$ (1)
where $`\mathrm{\Sigma }\text{SU(3)}`$ denotes the chiral fields for the Nambu-Goldstone bosons, while $`_m(\mathrm{\Sigma })`$ is the quark-mass induced meson mass term in the chiral Lagrangian.
When the true vacuum is nontrivial, that is, $`\mathrm{\Sigma }_0I`$, where $`\mathrm{\Sigma }_0`$ is the vacuum that minimizes the potential, some interesting phenomena could arise. For instance, if the quark mass matrix were given as $`m\text{Diag}(1,1,\delta )`$ with $`\delta >1/2`$, then the vacuum $`\mathrm{\Sigma }_0`$ would have an imaginary part in its matrix elements . Then writing $`\mathrm{\Sigma }=\mathrm{\Sigma }_0\mathrm{exp}(i\mathrm{\Pi }^AT^A)`$, where $`\mathrm{\Pi }^A`$ and $`T^A`$ denote the octet pseudo scalar mesons and the SU(3) generators, respectively, it can be easily seen that the mass term (1) could give rise to a spontaneous CP violation, which enables, for example, $`\eta `$ decay into two pions. Of course, with the quark mass matrix realized in nature, the true vacuum is at $`\mathrm{\Sigma }_0=I`$, and there is no such CP violation.
Recently, cold, dense quark system with large baryon chemical potential $`\mu `$ has received strong interest . As well known, the Fermi surface of such a dense system is unstable against Cooper pairing when an attractive force is present. At large $`\mu `$ a dressed gluon exchange (in hard dense loop approximation) provides such a force, and causes the system to be in color-flavor locking phase in which quarks condensate in a pattern
$`\chi _i^a\chi _j^b`$ $``$ $`k_1(U_0)_{}^{a}{}_{i}{}^{}(U_0)_{}^{b}{}_{j}{}^{}+k_2(U_0)_{}^{a}{}_{j}{}^{}(U_0)_{}^{b}{}_{i}{}^{}`$ (2)
$`\overline{\phi }_i^a\overline{\phi }_j^b`$ $``$ $`[k_1(V_0)_{}^{a}{}_{i}{}^{}(V_0)_{}^{b}{}_{j}{}^{}+k_2(V_0)_{}^{a}{}_{j}{}^{}(V_0)_{}^{b}{}_{i}{}^{}],`$ (3)
where $`\chi _i^a,\phi _i^a`$, with $`a`$ the color index, $`i=1,\mathrm{},3`$, the flavor index, denote the two-component Weyl fermions for the left-handed quarks and the complex conjugate of the right-handed quarks, respectively. The unitary matrices $`U_0,V_0\text{U(3)}`$ can be arbitrary in the absence of axial $`\text{U}(1)`$ anomaly, but the anomaly effect, even though small because of the suppression of instanton effects at high density , chooses a parity even vacuum in which
$$U_0=V_0.$$
(4)
Upon the condensation of the quarks the symmetry of dense, massless QCD, $`\text{SU}_L(3)\times \text{SU}_R(3)\times \text{U}_A(1)\times \text{U}_B(1)`$, where $`\text{U}_A(1)`$ and $`\text{U}_B(1)`$ denote the axial and the baryon number symmetry, respectively, is spontaneously broken to $`\text{SU}_{L+R}(3)`$, generating 10 Nambu-Goldstone bosons (mesons). The $`\text{U}_A(1)`$ is not an exact symmetry, but at large $`\mu `$, due to the suppression of anomaly, can be regarded a good approximate symmetry. As in vacuum QCD a continuum of chiral vacua $`[\text{SU}_L(3)\times \text{SU}_R(3)\times \text{U}_A(1)\times \text{U}_B(1)]/\text{SU}_{L+R}(3)`$ arises upon the spontaneous symmetry breaking. Note, however, that the vacua connected by the $`\text{U}_A(1)`$ rotation are only approximately degenerate, because the $`\text{U}_A(1)`$ is not an exact symmetry. It turns out that the vacuum of lowest energy is parity even, hence the relation (4).
Now, when the current quark masses are turned on, the degeneracy of the chiral vacua is lifted. The true vacuum can be picked up, as in vacuum QCD, by minimizing the quark mass induced potential $`V(\mathrm{\Sigma })=_m(\mathrm{\Sigma })`$ of dense QCD.
The meson mass term $`_m(\mathrm{\Sigma })`$ in color-flavor locking phase has a different form than in vacuum QCD because of the absence of left-right quark condensates. For small quark masses, $`_m(\mathrm{\Sigma })`$ can be expanded in powers of the quark mass matrix. The absence of a left-right quark condensate due to the suppression of instanton effects renders the leading term to be quadratic in quark mass . The most general form for the leading $`_m(\mathrm{\Sigma })`$, consistent with the chiral symmetry and the condensates (3), is given by
$`_m(\mathrm{\Sigma })`$ $`=`$ $`A[\text{Tr}(m^t\mathrm{\Sigma })]^2+B\text{Tr}[(m^t\mathrm{\Sigma })^2]+C\text{Tr}(m^t\mathrm{\Sigma })\text{Tr}(m^{}\mathrm{\Sigma }^{})+\text{H.c.},`$ (5)
where $`\mathrm{\Sigma }\text{U}(3)`$ denotes the chiral fields for the 9 mesons. Note that the Nambu-Goldstone boson associated with the baryon number symmetry remains exactly massless, and thus does not appear in the mass term.
The coefficients $`A,B`$ and $`C`$ can be determined either by matching the vacuum energy of (5) with that computed in the microscopic theory or by integrating out quark fields in the effective Lagrangian of quarks and the mesons . Here we take the latter approach. Using the global color-chiral-axial-baryon number symmetry of dense QCD and the pattern of the diquark condensates (3) one can write an effective chiral Lagrangian for the quarks and the Nambu-Goldstone bosons as
$``$ $`=`$ $`i\overline{\chi }_i^a\overline{\sigma }^\nu _\nu \chi _i^a+\mu \overline{\chi }_i^a\overline{\sigma }^0\chi _i^a+i\overline{\phi }_i^a\overline{\sigma }^\nu _\nu \phi _i^a\mu \overline{\phi }_i^a\overline{\sigma }^0\phi _i^a[m_{ij}\chi _i^a\phi _j^a+\text{H.c.}]`$ (7)
$`+\left[\chi _i^a(\mathrm{\Delta }_\chi ^{})_{ij}^{ab}\chi _j^b+\phi _i^a(\mathrm{\Delta }_\phi )_{ij}^{ab}\phi _j^b+\text{H.c.}\right]+_{\text{NG}}(U,V),`$
where
$`(\mathrm{\Delta }_\chi ^{})_{ij}^{ab}`$ $`=`$ $`k_1U_i^aU_j^b+k_2U_j^aU_i^b,`$ (8)
$`(\mathrm{\Delta }_\phi )_{ij}^{ab}`$ $`=`$ $`[k_1V_i^aV_j^b+k_2V_j^aV_i^b].`$ (9)
The unitary matrices $`U`$ and $`V`$ denote the nonlinearly realized Nambu-Goldstone bosons arising from the symmetry breaking $`SU_c(3)\times SU_L(3)\times U_{A+B}(1)SU_{c+L}(3)`$ through the $`\chi \chi `$ condensation and $`SU_c(3)\times SU_R(3)\times U_{BA}(1)SU_{c+R}(3)`$ through the $`\overline{\phi }\overline{\phi }`$ condensation, respectively. The $`_{\text{NG}}(U,V)`$ is the usual chiral Lagrangian for the Nambu-Goldstone bosons alone. We also note that with (9) it was assumed, for simplicity, the gaps for antiquarks are identical to those for quarks . It can be easily seen that the most important interactions for our calculation, those between the quarks and the Nambu-Goldstone bosons, satisfy the Goldberger-Treiman relation. This Lagrangian can be regarded as an effective Lagrangian before color is gauged for the quarks and the 18 Nambu-Goldstone bosons $`U,V`$. Upon gauging color it can be seen that 8 out of the 18 Nambu-Goldstone bosons are eaten by the gluons via Higgs mechanism and there remain 10 color-singlet mesons as low energy excitations.
Integrating out the quark fields, which corresponds to the evaluation of quark one-loop diagrams in the background of constant $`U,V`$ fields with two Dirac mass insertions (see Fig. 1), we obtain the mass term (5) in which $`\mathrm{\Sigma }`$ is given by
$$\mathrm{\Sigma }=UV^{}$$
(10)
and
$`A`$ $`=`$ $`i{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}\{(k_1^2+k_2^2)I_8(p)I_{8+}(p)+k_1(k_1+k_2/3)[I_8(p)I_{8+}(p)`$ (14)
$`I_{1+}(p)I_8(p)]+k_1(k_1+k_2/3)[I_8(p)I_{8+}(p)I_1(p)I_{8+}(p)]`$
$`+(k_1+k_2/3)^2[I_8(p)I_{8+}(p)I_1(p)I_{8+}(p)I_{1+}(p)I_8(p)`$
$`+I_1(p)I_{1+}(p)]\},`$
$`B`$ $`=`$ $`i{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}\{2k_1k_2I_8(p)I_{8+}(p)+k_2(k_1+k_2/3)[I_8(p)I_{8+}(p)`$ (16)
$`I_{1+}(p)I_8(p)]+k_2(k_1+k_2/3)[I_8(p)I_{8+}(p)I_1(p)I_{8+}(p)]\},`$
$`C`$ $`=`$ $`i{\displaystyle \frac{1}{9}}{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}[(p_0\mu )^2+|\stackrel{}{p}|^2]\{I_8(p)I_{8+}(p)I_1(p)I_{8+}(p)`$ (18)
$`I_{1+}(p)I_8(p)+I_1(p)I_{1+}(p)\},`$
where
$`I_1(p)=1/[p_0^2+(|\stackrel{}{p}|\mu )^2+m_1^2],I_8(p)=1/[p_0^2+(|\stackrel{}{p}|\mu )^2+m_8^2].`$ (19)
Here $`m_1`$ and $`m_8`$ are the Majorana masses for the singlet and octet quarks, respectively, under the unbroken $`\text{SU}_{L+R}(3)`$ and are given as
$`m_1^2=(3k_1+k_2)^2,m_8^2=k_2^2.`$ (20)
Generally the diquark condensates depend on energy, but here we shall ignore this fact and treat $`k_i`$ as constants. Then the integration over the loop momentum can be easily performed by doing contour integration over $`p_0`$ first and then replacing $`d^3p4\pi \mu ^2d|\stackrel{}{p}|`$ for the integration over the spatial components. It can be easily seen that $`A`$ and $`B`$ arise from the first diagram in Fig.1 with $`A,B\mathrm{\Delta }^2\mathrm{ln}(\mu ^2/\mathrm{\Delta }^2)`$, while $`C`$, coming from the second diagram, is given by $`C\mathrm{\Delta }^4/\mu ^2\mathrm{ln}(\mu ^2/\mathrm{\Delta }^2)`$, where $`\mathrm{\Delta }k_i`$. Note that $`C`$ is suppressed by a factor $`\mathrm{\Delta }^2/\mu ^2`$ compared to $`A,B`$.
To determine the true chiral vacuum at a given chemical potential we have to minimize the potential $`V(\mathrm{\Sigma })=_m(\mathrm{\Sigma })`$, with $`_m(\mathrm{\Sigma })`$ given by (5) and (18). To demonstrate a nontrivial vacuum alignment, we shall now take a simplified form for the quark mass matrix, in which the up and down quark masses are identical, as
$$m=m_d\text{Diag}(1,1,\delta ),\delta m_s/m_d.$$
(21)
With this quark mass matrix it is easy to see that the $`\mathrm{\Sigma }`$ that minimizes the potential must be of the form
$$\mathrm{\Sigma }=\text{Diag}(\alpha ,\alpha ,\beta )\text{with}|\alpha |^2=|\beta |^2=1,$$
(22)
where $`\alpha ,\beta `$ are complex variables.
Substituting (21) and (22) into (5) we obtain
$$V=2m_d^2\left\{\text{Re}[A(2\alpha +\beta \delta )^2+B(2\alpha ^2+\beta ^2\delta ^2)]+C\left|2\alpha +\beta \delta \right|^2\right\},$$
(23)
which can also be written as
$$V=8m_d^2\left\{(2A+B)\mathrm{cos}^2\theta +(A+B)\delta ^2/2\mathrm{cos}^2\varphi +A\delta \mathrm{cos}(\theta +\varphi )+C\delta \mathrm{cos}(\theta \varphi )\right\}$$
(24)
by putting $`\alpha =\mathrm{exp}(i\theta )`$ and $`\beta =\mathrm{exp}(i\varphi )`$.
We first notice that the potential is symmetric under the transformation $`(\alpha ,\beta )(\alpha ,\beta )`$, so the minima of the potential must occur in pairs. Secondly we observe that for arbitrary coefficients $`A,B`$, and $`C`$ the potential is stationary when $`\theta ,\varphi `$ satisfy
$`\mathrm{sin}(\theta \pm \varphi )=0,\mathrm{sin}(2\theta )=0,\mathrm{sin}(2\varphi )=0,`$ (25)
which have 8 common solutions corresponding to the following $`(\alpha ,\beta )`$ pairs
$$(i,i),(1,1),(1,1),(i,i)$$
(26)
and their partners of opposite sign. Of course none of these pairs needs necessarily minimize the potential, but it can be shown numerically that the minima occur always on one of these solutions. This shows that the true chiral vacuum can align only to one of these eight directions.
In ideal situation we may know the dependence of the gap parameters $`k_i`$ on $`\mu `$, could choose the true chiral vacuum at a given chemical potential from the above discrete vacua, and investigate vacuum transitions as the chemical potential evolves. However, presently there is no reliable calculation to determine the gaps as functions of the chemical potential except when the chemical potential is extremely large , in which case Shwinger-Dyson equation can be used as an approximate gap equation and solved . Even in this case the absolute magnitudes of the gaps are yet to be determined, but, it may not be so unreasonable to assume that $`k_i/\mu `$ are in the order of 0.1 at large chemical potential .
Taking into account this uncertainty we shall here treat $`k_i`$ as free parameters and study the vacuum transitions as $`k_i`$ vary. For definiteness, we shall put $`\delta =15`$ and scan numerically the minima of the potential in the $`k_i`$ parameter space defined by $`0k_1/\mu 0.5`$ and $`0.5k_2/\mu 0.5`$. Note that $`k_1`$ can always be assumed positive, since its phase can be rotated away by the baryon number symmetry of the Lagrangian (7).
The result of the numerical scanning is shown in Fig. 2. As we see, the parameter space is divided into four domains according to their vacuum directions. In this figure the corresponding vacua of opposite sign to those in the group (26) are not included. The reason for this is that the numerical scanning shows there is always a potential barrier between a vacuum in the group (26) and the partners of those in the group (26), so in infinite volume limit the transitions between the vacua (26) and their partners are negligibly small and can be ignored. Therefore, only the transitions within the group of vacua in (26) need be considered.
Because the vacuum can align only to a discrete number of directions the transitions between vacua will be of first order phase transition. To understand what happens at the transition, we look more carefully at the potential on the boundaries between the domains. For convenience, we shall designate the domains associated with the vacua defined in (26) by (I),(II),(III), and (IV), respectively.
First, consider the transition between the domain (I) and (II). As one approaches the boundary from either side of (I) and (II) it can be shown numerically that a potential valley opens up in $`(\theta ,\varphi )`$ space along the line $`\varphi =\theta \pi `$, with the bottom of the valley connecting the two points $`(\pi /2,\pi /2)`$ and $`(\pi ,0)`$. When exactly on the boundary, we have from $`V(\mathrm{\Sigma }_\text{I})=V(\mathrm{\Sigma }_{\text{II}})`$, with $`\mathrm{\Sigma }_{\text{I,II}}`$ defined by (22) and (26),
$$(2A+B)+(A+B)\delta ^2/22C\delta =0.$$
(27)
Then it is easy to see that on the bottom of the valley (i.e. along the direction $`\varphi =\theta \pi `$) the potential is constant with $`V=8m_d^2(AC)\delta `$. Thus, in this case there will be no latent heat released at the phase transition.
Similarly, for the transition between (II) and (III) it can be shown that a valley opens up along the line $`\varphi =0`$ as the boundary is approached, and that the relation
$$A+C=0,$$
(28)
holds on the boundary. On the bottom of the valley the potential is given by
$$V=8m_d^2\left\{(2A+B)\mathrm{cos}^2\theta +(A+B)\delta ^2/2\right\},$$
(29)
which shows a barrier between the two vacua. This then indicates that there will be latent heat released at the phase transition. This may have some important phenomenological consequence in dense systems.
In a similar fashion we can easily show that for the transition across the domain (III) and (IV) a potential valley opens up along the direction $`\varphi =\theta `$, and on the bottom of the valley the potential is constant with $`V=8m_d^2(AC)\delta `$. In this case there will be no latent heat released as in the transition between (I) and (II). Also for the transition between (I) and (IV) a valley opens up along $`\varphi =\pi /2`$ direction and on the bottom of the valley the potential is given by
$$V=4m_d^2(A+B)\delta ^2\mathrm{cos}^2\theta ,$$
(30)
which shows a barrier between the two vacua, and consequently, nonzero latent heat at the transition.
Could the phase transition of this kind occur when the chemical potential increases from zero to an asymptotic value? Although it is difficult to answer this question conclusively until we have the quark-mass induced meson potential at an arbitrary value of $`\mu `$, which at smaller $`\mu `$ would probably contain the instanton-induced potential that gives mass to the $`\eta ^{}`$ meson and the old $`\text{Tr}(m\mathrm{\Sigma })`$ as well as the quark-mass quadratic $`_m(\mathrm{\Sigma })`$ in (5), there is an interesting observation concerning this question. It is well known that at large chemical potential the sextet components of the condensates (3) is suppressed , thus $`k_1k_2`$, which then suggests the system must be in domain (I) at high density. The vacuum associated with domain (I) is $`\mathrm{\Sigma }_0=\text{Diag}(i,i,i)`$ that has an overall factor $`i`$ compared to the vacuum $`\mathrm{\Sigma }_0=I`$ at zero density. Although this does not imply a first order phase transition of the kind hitherto discussed, at least it does suggest that the vacuum must shift from the unit matrix at zero density to something else as the chemical potential increases.
This also leads to an interesting consequence of the vacuum alignment, namely that parity must be spontaneously broken at high density. The possibility of parity violation in dense QCD was pointed out by Pisarski and Rischke , and noted also in , who observed that the axial U(1) is almost a good symmetry at large $`\mu `$ due to the suppression of instanton effects, so a vacuum with no parity symmetry can be as easily formed as the parity even vacuum. Since the axial U(1) is never exactly restored this possibility, however, cannot be realized with massless quarks. The lowest energy state is always the parity even vacuum. However, when quarks are massive, the explicit breaking of the axial U(1) by quark masses can be more important than the small anomalous breaking, so the true vacuum could break parity.
Indeed, we see that this possibility is realized through the chiral vacuum alignment. Since the diquark condensates satisfy (4) in parity even vacuum it can be seen using (10) that only the vacuum $`\mathrm{\Sigma }_0=I`$ associated with the domain (III) is parity even. Therefore, when the system is in domain (I), (II), and (IV) parity is spontaneously broken. Moreover, since the system should be in domain (I) as $`\mu \mathrm{}`$, parity must be broken at high density.
To conclude, we have studied quark-mass induced chiral vacuum alignment in cold, dense QCD. When quarks are massless any point in the continuum of chiral vacua can be chosen as the vacuum, but when nonzero quark masses are introduced, the true vacuum must be found via Dashen’s procedure. We have shown that in color-flavor locking phase at large chemical potential the vacuum can align only to a discrete number of directions, and the nature of the vacuum transitions is of first order with vanishing or nonvanishing latent heat. It was also shown that at high density parity is spontaneously broken through the vacuum alignment. The consequence of the first order phase transitions and the parity violation may become important in dense systems like quark stars and heavy ion collisions.
Finally, we remark on the calculation of the meson masses. Usually in meson mass calculation it was assumed that the vacuum is at $`\mathrm{\Sigma }_0=I`$, and the meson potential was expanded around this vacuum to pick up the meson spectrum. However, as we have seen, the vacuum is not necessarily at the unit matrix and so this is not always correct. For correct meson masses one should first find the true vacuum and then expand the potential about it.
Useful comments by D.K. Hong are acknowledged. This work was supported in part by the Korea Science and Engineering Foundation (KOSEF).
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# Untitled Document
UTTG-07-00
The Cosmological Constant Problems<sup>*</sup><sup>*</sup>*This research was supported in part by the Robert A. Welch Foundation and NSF Grant PHY-9511632.
(Talk given at Dark Matter 2000, Marina del Rey, CA, February 2000)
Steven Weinberg
Department of Physics, University of Texas
Austin, Texas 78712
## Abstract
The old cosmological constant problem is to understand why the vacuum energy is so small; the new problem is to understand why it is comparable to the present mass density. Several approaches to these problems are reviewed. Quintessence does not help with either; anthropic considerations offer a possibility of solving both. In theories with a scalar field that takes random initial values, the anthropic principle may apply to the cosmological constant, but probably to nothing else.
1. Introduction
There are now two cosmological constant problems. The old cosmological constant problem is to understand in a natural way why the vacuum energy density $`\rho _V`$ is not very much larger. We can reliably calculate some contributions to $`\rho _V`$, like the energy density in fluctuations in the gravitational field at graviton energies nearly up to the Planck scale, which is larger than is observationally allowed by some 120 orders of magnitude. Such terms in $`\rho _V`$ can be cancelled by other contributions that we can’t calculate, but the cancellation then has to be accurate to 120 decimal places. The new cosmological constant problem is to understand why $`\rho _V`$ is not only small, but also, as current Type Ia supernova observations seem to indicate,<sup>2</sup><sup>2</sup>2 A. G. Riess et al.: Astron. J. 116, 1009 (1998): P. M. Garnavich et al.: Astrophys. J. 509, 74 (1998); S. Perlmutter et al.: Astrophys. J. 517, 565 (1999). of the same order of magnitude as the present mass density of the universe.
The efforts to understand these problems can be grouped into four general classes. The first approach is to imagine some scalar field coupled to gravity in such a way that $`\rho _V`$ is automatically cancelled or nearly cancelled when the scalar field reaches its equilibrium value. In a review article over a decade ago<sup>3</sup><sup>3</sup>3S. Weinberg: Rev. Mod. Phys. 61, 1 (1989). I gave a sort of ‘no go’ theorem, showing why such attempts would not work without the need for a fine tuning of parameters that is just as mysterious as the problem we started with. I wouldn’t claim that this is conclusive — other no-go theorems have been evaded in the past — but so far no one has found a way out of this one. The second approach is to imagine some sort of deep symmetry, one that is not apparent in the effective field theory that governs phenomena at accessible energies, but that nevertheless constrains the parameters of this effective theory so that $`\rho _V`$ is zero or very small. I leave this to be covered in the talk by Edward Witten. In this talk I will concentrate on the third and fourth of these approaches, based respectively on the idea of quintessence and on versions of the anthropic principle.
2. Quintessence
The idea of quintessence<sup>4</sup><sup>4</sup>4P. J. E. Peebles and B. Ratra: Astrophys. J. 325, L17 (1988); B. Ratra and P. J. E. Peebles: Phys. Rev. D 37, 3406 (1988); C. Wetterich: Nucl. Phys. B302, 668 (1988). is that the cosmological constant is small because the universe is old. One imagines a uniform scalar field $`\varphi (t)`$ that rolls down a potential $`V(\varphi )`$, at a rate governed by the field equation
$$\ddot{\varphi }+3H\dot{\varphi }+V^{}(\varphi )=0,$$
(1)
where $`H`$ is the expansion rate
$$H=\sqrt{\left(\frac{3}{8\pi G}\right)\left(\rho _\varphi +\rho _M\right)}.$$
(2)
Here $`\rho _\varphi `$ is the energy density of the scalar field
$$\rho _\varphi =\frac{1}{2}\dot{\varphi }^2+V(\varphi ),$$
(3)
while $`\rho _M`$ is the energy density of matter and radiation, which decreases as
$$\dot{\rho }_M=3H\left(\rho _M+p_M\right),$$
(4)
with $`p_M`$ the pressure of matter and radiation.
If there is some value of $`\varphi `$ (typically, $`\varphi `$ infinite) where $`V^{}(\varphi )=0`$, then it is natural that $`\varphi `$ should approach this value, so that it eventually changes only slowly with time. Meanwhile $`\rho _M`$ is steadily decreasing, so that eventually the universe starts an exponential expansion with a slowly varying expansion rate $`H\sqrt{8\pi GV(\varphi )/3}`$. The problem, of course, is to explain why $`V(\varphi )`$ is small or zero at the value of $`\varphi `$ where $`V^{}(\varphi )=0`$.
Recently this approach has been studied in the context of so-called ‘tracker’ solutions.<sup>5</sup><sup>5</sup>5I. Zlatev, L. Wang, and P. J. Steinhardt: Phys. Rev. Lett. 82, 896 (1999); Phys. Rev. D 59, 123504 (1999). The simplest case arises for a potential of the form
$$V(\varphi )=M^{4+\alpha }\varphi ^\alpha ,$$
(5)
where $`\alpha >0`$, and $`M`$ is an adjustable constant. If the scalar field begins at a value much less than the Planck mass and with $`V(\varphi )`$ and $`\dot{\varphi }^2`$ much less than $`\rho _M`$, then the field $`\varphi (t)`$ initially increases as $`t^{2/(2+\alpha )}`$, so that $`\rho _\varphi `$ decreases as $`t^{2\alpha /(2+\alpha )}`$, while $`\rho _M`$ is decreasing faster, as $`t^2`$. (The existence of this phase is important, because the success of cosmic nucleosynthesis calculations would be lost if the cosmic energy density were not dominated by $`\rho _M`$ at temperatures of order $`10^9^{}`$K to $`10^{10}^{}`$K.) Eventually a time is reached when $`\rho _M`$ becomes as small as $`\rho _\varphi `$, after which the character of the solution changes. Now $`\rho _\varphi `$ becomes larger than $`\rho _M`$, and $`\rho _\varphi `$ decreases more slowly, as $`t^{2/(4+\alpha )}`$. The expansion rate $`H`$ now goes as $`H\sqrt{V(\varphi )}t^{\alpha /(4+\alpha )}`$, so the Robertson–Walker scale factor $`R(t)`$ grows almost exponentially, with $`\mathrm{log}R(t)t^{4/(4+\alpha )}`$. In this approach, the transition from $`\rho _M`$-dominance to $`\rho _\varphi `$-dominance is supposed to take place near the present time, so that both $`\rho _M`$ and $`\rho _\varphi `$ are now both contributing appreciably to the cosmic expansion rate.
The nice thing about these tracker solutions is that the existence of a cross-over from an early $`\rho _M`$-dominated expansion to a later $`\rho _\varphi `$-dominated expansion does not depend on any fine-tuning of the initial conditions. But it should not be thought that either of the two cosmological constant problems are solved in this way. Obviously, the decrease of $`\rho _\varphi `$ at late times would be spoiled if we added a constant of order $`m_{\mathrm{Planck}}^4`$ (or $`m_W^4`$, or $`m_e^4`$) to the potential (5). What is perhaps less clear is that, even if we take the potential in the form (5) without any such added constant, we still need a fine-tuning to make the value of $`\rho _\varphi `$ at which $`\rho _\varphi \rho _M`$ close to the present critical density $`\rho _{c0}`$. The value of the field $`\varphi (t)`$ at this crossover can easily be seen to be of the order of the Planck mass, so in order for $`\rho _\varphi `$ to be comparable to $`\rho _M`$ at the present time we need
$$M^{4+\alpha }(8\pi G)^{\alpha /2}\rho _{c0}(8\pi G)^{1\alpha /2}H_0^2.$$
(6)
Theories of quintessence offer no explanation why this should be the case. (An interesting suggestion has been made after Dark Matter 2000.<sup>6</sup><sup>6</sup>6 C. Armendariz-Picon, V. Mukhanov, and P. J. Steinhardt: astro-ph/0004134.)
3. Anthropic Considerations
In several cosmological theories the observed big bang is just one member of an ensemble. The ensemble may consist of different expanding regions at different times and locations in the same spacetime,<sup>7</sup><sup>7</sup>7A. Vilenkin: Phys. Rev. D 27, 2848 (1983); A. D. Linde: Phys. Lett. B175, 395 (1986). or of different terms in the wave function of the universe.<sup>8</sup><sup>8</sup>8E. Baum: Phys. Lett. B133, 185 (1984); S. W. Hawking: in Shelter Island II – Proceedings of the 1983 Shelter Island Conference on Quantum Field Theory and the Fundamental Problems of Physics, ed. by R. Jackiw et al. (MIT Press, Cambridge, 1985); Phys. Lett. B134, 403 (1984); S. Coleman: Nucl. Phys. B 307, 867 (1988). If the vacuum energy density $`\rho _V`$ varies among the different members of this ensemble, then the value observed by any species of astronomers will be conditioned by the necessity that this value of $`\rho _V`$ should be suitable for the evolution of intelligent life.
It would be a disappointment if this were the solution of the cosmological constant problems, because we would like to be able to calculate all the constants of nature from first principles, but it may be a disappointment that we will have to live with. We have learned to live with similar disappointments in the past. For instance, Kepler tried to derive the relative distances of the planets from the sun by a geometrical construction involving Platonic solids nested within each other, and it was somewhat disappointing when Newton’s theory of the solar system failed to constrain the radii of planetary orbits, but by now we have gotten used to the fact that these radii are what they are because of historical accidents. This is a pretty good analogy, because we do have an anthropic explanation why the planet on which we live is in the narrow range of distances from the sun at which the surface temperature allows the existence of liquid water: if the radius of our planet’s orbit was not in this range, then we would not be here. This would not be a satisfying explanation if the earth were the only planet in the universe, for then the fact that it is just the right distance from the sun to allow water to be liquid on its surface would be quite amazing. But with nine planets in our solar system and vast numbers of planets in the rest of the universe, at different distances from their respective stars, this sort of anthropic explanation is just common sense. In the same way, an anthropic explanation of the value of $`\rho _V`$ makes sense if and only if there is a very large number of big bangs, with different values for $`\rho _V`$.
The anthropic bound on a positive vacuum energy density is set by the requirement that $`\rho _V`$ should not be so large as to prevent the formation of galaxies.<sup>9</sup><sup>9</sup>9S. Weinberg: Phys. Rev. Lett. 59, 2607 (1987). Using the simple spherical infall model of Peebles<sup>10</sup><sup>10</sup>10P. J. E. Peebles: Astrophys. J. 147, 859 (1967). to follow the nonlinear growth of inhomogeneities in the matter density, one finds an upper bound
$$\rho _V<\frac{500\rho _R\delta _R^3}{729}$$
(7)
where $`\rho _R`$ is the mass density and $`\delta _R`$ is a typical fractional density perturbation, both taken at the time of recombination. This is roughly the same as requiring that $`\rho _V`$ should be no larger than the cosmic mass density at the earliest time of galaxy formation, which for a maximum galactic redshift of 5 would be about 200 times the present mass density. This is a big improvement over missing by 120 orders of magnitude, but not good enough.
However, we would not expect to live in a big bang in which galaxy formation is just barely possible. Much more reasonable is what Vilenkin calls a principle of mediocrity,<sup>11</sup><sup>11</sup>11A. Vilenkin: Phys. Rev. Lett. 74, 846 (1995); in Cosmological Constant and the Evolution of the Universe, ed. by K. Sato et al. (Universal Academy Press, Tokyo, 1996). which suggests that we should expect to find ourselves in a big bang that is typical of those in which intelligent life is possible. To be specific, if $`𝒫_{\mathrm{a}\mathrm{priori}}(\rho _V)d\rho _V`$ is the a priori probability of a particular big bang having vacuum energy density between $`\rho _V`$ and $`\rho _V+d\rho _V`$, and $`𝒩(\rho _V)`$ is the average number of scientific civilizations in big bangs with energy density $`\rho _V`$, then the actual (unnormalized) probability of a scientific civilization observing an energy density between $`\rho _V`$ and $`\rho _V+d\rho _V`$ is
$$d𝒫(\rho _V)=𝒩(\rho _V)𝒫_{\mathrm{a}\mathrm{priori}}(\rho _V)d\rho _V.$$
(8)
We don’t know how to calculate $`𝒩(\rho _V)`$, but it seems reasonable to take it as proportional to the number of baryons that wind up in galaxies, with an unknown proportionality factor that is independent of $`\rho _V`$. There is a complication, that the total number of baryons in a big bang may be infinite, and may also depend on $`\rho _V`$. In practice, we take $`𝒩(\rho _V)`$ as the fraction of baryons that wind up in galaxies, which we can hope to calculate, and include the total baryon number as a factor in $`𝒫_{\mathrm{a}\mathrm{priori}}(\rho _V)`$.
The one thing that offers some hope of actually calculating $`d𝒫(\rho _V)`$ is that $`𝒩(\rho _V)`$ is non-zero in only a narrow range of values of $`\rho _V`$, values that are much smaller than the energy densities typical of elementary particle physics, so that $`𝒫_{\mathrm{a}\mathrm{priori}}(\rho _V)`$ is likely to be constant within this range. <sup>12</sup><sup>12</sup>12S. Weinberg: in Critical Dialogs in Cosmology, ed. by N. Turok (World Scientific, Singapore, 1997). The value of this constant is fixed by the requirement that the total probability should be one, so
$$d𝒫(\rho _V)=\frac{𝒩(\rho _V)d\rho _V}{𝒩(\rho _V^{})𝑑\rho _V^{}}.$$
(9)
The fraction $`𝒩(\rho _V)`$ of baryons in galaxies has been calculated by Martel, Shapiro and myself,<sup>13</sup><sup>13</sup>13H. Martel, P. Shapiro, and S. Weinberg: Astrophys. J. 492, 29 (1998). using the well-known spherical infall model of Gunn and Gott,<sup>14</sup><sup>14</sup>14J. Gunn and J. Gott: Astrophys. J. 176, 1 (1972). in which one starts with a fractional density perturbation that is positive within a sphere, and compensated by a negative fractional density perturbation in a surrounding spherical shell. The results are quite insensitive to the relative radii of the sphere and shell. Taking the shell thickness to equal the sphere’s radius, the integrated probability distribution function for finding a vacuum energy less than or equal to $`\rho _V`$ is
$`𝒫(\rho _V)`$ $``$ $`{\displaystyle _0^{\rho _V}}𝑑𝒫`$ (10)
$`=`$ $`1+(1+\beta )e^\beta `$
$`+`$ $`{\displaystyle \frac{1}{2\mathrm{ln}21}}{\displaystyle _\beta ^{\mathrm{}}}e^x𝑑x\left\{2\sqrt{\beta x}+\beta +2x\mathrm{ln}\left[\sqrt{\beta /x}+1\right]\right\}`$
where
$$\beta \frac{1}{2\sigma ^2}\left(\frac{729\rho _V}{500\rho _R}\right)^{2/3}$$
(11)
with $`\sigma `$ the rms fractional density perturbation at recombination, and $`\rho _R`$ the average mass density at recombination. The probability of finding ourselves in a big bang with a vacuum energy density large enough to give a present value of $`\mathrm{\Omega }_V`$ of 0.7 or less turns out to be 5% to 12%, depending on the assumptions used to estimate $`\sigma `$. In other words, the vacuum energy in our big bang still seems a little low, but not implausibly so. These anthropic considerations can therefore provide a solution to both the old and the new cosmological constant problems, provided of course that the underlying assumptions are valid. Related anthropic calculations have been carried out by several other authors.<sup>15</sup><sup>15</sup>15G. Efstathiou: Mon. Not. Roy. Astron. Soc. 274, L73 (1995); M. Tegmark and M. J. Rees: Astrophys. J. 499, 526 (1998); J. Garriga, M. Livio, and A. Vilenkin: astro-ph/9906210; S. Bludman: astro-ph/0002204.
I should add that when anthropic considerations were first applied to the cosmological constant, counts of galaxies as a function of redshift<sup>16</sup><sup>16</sup>16E. D. Loh: Phys. Rev. Lett. 57, 2865 (1986). indicated that $`\mathrm{\Omega }_\mathrm{\Lambda }`$ is $`0.1_{0.4}^{+0.2}`$, and this was recognized to be too small to be explained anthropically. The subsequent discovery in studies of type Ia supernova distances and redshifts that $`\mathrm{\Omega }_\mathrm{\Lambda }`$ is quite large does not of course prove that anthropic considerations are relevant, but it is encouraging.
Recently the assumptions underlying these calculations have been challenged by Garriga and Vilenkin.<sup>17</sup><sup>17</sup>17J. Garriga and A. Vilenkin: astro-ph/9908115. They adopt a plausible model for generating an ensemble of big bangs with different values of $`\rho _V`$, by supposing that there is a scalar field $`\varphi `$ that initially can take values anywhere in a broad range in which the potential $`V(\varphi )`$ is very flat. Specifically, in this range
$$\left|\frac{V^{}(\varphi )}{V(\varphi )}\right|\sqrt{8\pi G}\mathrm{and}\left|\frac{V^{\prime \prime }(\varphi )}{V(\varphi )}\right|8\pi G.$$
(12)
It is also assumed that in this range $`V(\varphi )`$ is much less than the initial value of the energy density $`\rho _M`$ of matter and radiation. For initial values of $`\varphi `$ in this range, the vacuum energy density $`\rho _\varphi `$ stays roughly constant while $`\rho _M`$ drops to a value of order $`\rho _\varphi `$. To see this, note that during this period the expansion rate behaved as $`H=\eta /t`$, with $`\eta =2/3`$ or $`\eta =1/2`$ during times of matter or radiation dominance, respectively. If we tentatively assume that $`\varphi `$ is roughly constant, then the field equation (1) gives
$$\dot{\varphi }\frac{tV^{}(\varphi )}{1+3\eta }.$$
(13)
During the time that $`\rho _M\rho _\varphi `$, the ratio of the kinetic to the potential terms in Eq. (3) for $`\rho _\varphi `$ is
$$\frac{\dot{\varphi }^2}{2V(\varphi )}\frac{t^2V^2(\varphi )}{2(1+3\eta )^2V(\varphi )}\frac{8\pi Gt^2V(\varphi )}{2(1+3\eta )^2}\frac{3\eta ^2V(\varphi )}{2(1+3\eta )^2\rho _M}1,$$
(14)
so $`\rho _\varphi `$ is dominated by the potential term. The fractional change in $`\rho _\varphi `$ until the time $`t_c`$ when $`\rho _M`$ becomes equal to $`\rho _\varphi `$ is then
$$\frac{|\mathrm{\Delta }\rho _\varphi |}{\rho _\varphi }=\frac{1}{\rho _\varphi }\left|_0^{t_c}V^{}(\varphi )\dot{\varphi }𝑑t\right|\frac{V^2(\varphi )t_c^2}{2(1+3\eta )\rho _\varphi }\frac{V^2(\varphi )}{8\pi G\rho _\varphi ^2}1.$$
(15)
Following this period, $`\rho _\varphi `$ becomes dominant, and the inequalities (12) ensure that the expansion becomes essentially exponential, just as in theories with the ‘tracker’ solutions discussed in the previous section. Hence in this class of models, $`V(\varphi )`$ plays the role of a constant vacuum energy, whose values are governed by the a priori probability distribution for the initial values of $`\varphi `$. In particular, if one assumes that all initial values of $`\varphi `$ are equally probable, then the a priori distribution of the vacuum energy is
$$𝒫_{\mathrm{a}\mathrm{priori}}\left(V(\varphi )\right)\frac{1}{\left|V^{}(\varphi )\right|}.$$
(16)
The point made by Garriga and Vilenkin was that, because $`V(\varphi )`$ is so flat, the field $`\varphi `$ can vary appreciably even when $`\rho _VV(\varphi )`$ is restricted to the very narrow anthropically allowed range of values in which galaxy formation is possible. They concluded that it would also be possible for the a priori probability (16) to vary appreciably in this range, which if true would require modifications in the calculation of $`𝒫(\rho _V)`$ described above. The potential they used as an example was
$$V(\varphi )=V_1+A(\varphi /M)+B\mathrm{sin}\left(\varphi /M\right),$$
with $`V_1`$ large, of order $`M^4`$, $`A`$ and $`B`$ much smaller, and $`M`$ a large mass, but not larger than the Planck mass. This yields an a priori probability distribution (16) that varies appreciably in the anthropically allowed range of $`\varphi `$.
It turns out<sup>18</sup><sup>18</sup>18S. Weinberg: astro-ph/0002387. that the issue of whether the a priori probability (16) is flat in the anthropically allowed range of $`\varphi `$ depends on the way we impose the slow roll conditions (12). There is a large class of potentials for which the probability is flat in this range. Suppose for instance that, unlike the example chosen by Garriga and Vilenkin, the potential is of the general form
$$V(\varphi )=V_1f(\lambda \varphi )$$
(17)
where $`V_1`$ is a large energy density, in the range $`m_W^4`$ to $`m_{\mathrm{Planck}}^4`$, $`\lambda >0`$ is a very small constant, and $`f(x)`$ is a function involving no very small or very large parameters. Anthropically allowed values of $`\varphi /\lambda `$ must be near a zero of $`f(x)`$, say a simple zero at $`x=a`$. Then $`V^{}(\varphi )\lambda V_1f^{}(a)\lambda V_1`$ and $`V^{\prime \prime }(\varphi )\lambda ^2V_1f^{\prime \prime }(a)\lambda ^2V_1`$, so both inequalities (12) are satisfied if
$$\lambda \sqrt{8\pi G}\left(\frac{\rho _V}{V_1}\right).$$
(18)
Galaxy formation is only possible for $`|V(\varphi )|`$ less than an upper bound $`V_{\mathrm{max}}`$, of the order of the mass density of the universe at the earliest time of galaxy formation, which is very much less than $`V_1`$, so the anthropically allowed range of values of $`\varphi `$ is
$$\left|\varphi a/\lambda \right|_{\mathrm{max}}\frac{V_{\mathrm{max}}}{\lambda V_1|f^{}(a)|}.$$
(19)
The fractional variation in the a priori probability density (16) as $`\varphi `$ varies in the range (19) is then
$$\left|\frac{V^{\prime \prime }(\varphi )}{V^{}(\varphi )}\right|\left|\varphi a/\lambda \right|_{\mathrm{max}}\left|\frac{V_{\mathrm{max}}}{V_1}\right|\left|\frac{f^{\prime \prime }(a)}{f^2(a)}\right|\left|\frac{V_{\mathrm{max}}}{V_1}\right|1$$
(20)
justifying the assumptions made in the calculation of Eq. (10).
I should emphasize that no fine-tuning is needed in potentials of type (16). It is only necessary that $`V_1`$ be sufficiently large, $`\lambda `$ be sufficiently small, and $`f(x)`$ have a simple zero somewhere, with derivatives of order unity at this zero. These properties are not upset if for instance we add a large constant to the potential. But why should each appearance of the field $`\varphi `$ be accompanied with a tiny factor $`\lambda `$? As we have been using it, derivatives of the field $`\varphi `$ appear in the Lagrangian density in the form $`\frac{1}{2}_\mu \varphi ^\mu \varphi `$, as shown by the coefficient unity of the second derivative in the field equation (1). In general, we might expect the Lagrangian density for $`\varphi `$ to take the form
$$=\frac{Z}{2}_\mu \varphi ^\mu \varphi V_1f(\varphi /M)$$
(21)
where $`f(x)`$ is a function of the sort we have been considering, involving no large or small parameters, $`M`$ is a mass perhaps of order $`(8\pi G)^{1/2}`$, and $`V_1`$ is a large constant, of order $`M^4`$. With an arbitrary field-renormalization constant $`Z`$ in the Lagrangian, the field $`\varphi `$ is not canonically normalized, and does not obey Eq. (1). We may define a canonically normalized field as $`\varphi ^{}\sqrt{Z}\varphi `$; writing the Lagrangian in terms of $`\varphi ^{}`$, and dropping the prime, we get a potential of the form (16), with $`\lambda =1/M\sqrt{Z}`$. Thus we can understand a very small $`\lambda `$ if we can explain why the field renormalization constant $`Z`$ is very large. Perhaps this has something to do with the running of $`Z`$ as the length scale at which it is measured grows to astronomical dimensions.
There is a problem with this sort of implementation of the anthropic principle, that may prevent its application to anything other than the cosmological constant. When quantized, a scalar field with a very flat potential leads to very light bosons, that might be expected to have been already observed. If we want to explain the masses and charges of elementary particles anthropically, by supposing that these masses and charges arise from expectation values of a scalar field in a flat potential with random initial values, then the scalar field would have to couple to these elementary particles, and would therefore be created in their collisions and decays. This problem does not arise for a scalar field that couples only to itself and gravitation (and perhaps also to a hidden sector of other fields that couple only to other fields in the hidden sector and to gravitation). It is true that such a scalar would couple to observed particles through multi-graviton exchange, and with a cutoff at the Planck mass the Yukawa couplings of dimensionality four that are generated in this way would in general not be suppressed by factors of $`G`$. But in our case the non-derivative interactions of the scalars with gravitation are suppressed by a factor $`V^{}(\varphi )\lambda `$, which according to Eq. (18) is much less than $`\sqrt{8\pi G}`$, yielding Yukawa couplings that are very much less than unity. Thus it may be that anthropic considerations are relevant for the cosmological constant, but for nothing else.
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# 1 Introduction
## 1 Introduction
The transitions from the 5p levels to 5s levels in neutral Kr give rise to the most prominent lines in the Kr I emission spectrum (see Fig. 1 for a simplified energy level diagram). Lifetimes with sub-percent uncertainties for six of the ten 5p levels have recently been measured with beam-gas laser spectroscopy (BGLS) by Schmoranzer and Volz and Schmitt et al. . We will show below how the six lifetimes from BGLS can be used in an intermediate coupling calculation to predict accurately the lifetimes of those four levels that were not accessible to BGLS. The resulting complete set of lifetimes for the 5p levels inspired us to make new, accurate measurements of branching fractions for all 5p–5s transitions in Kr I to obtain electric dipole transition rates with uncertainties limited only by the accuracy of the radiometric calibration.
To date, the most extensive measurements of transition rates in Kr I were carried out by Chang, Horiguchi and Setser who have measured transition rates for all 5p–5s transitions but with an accuracy of only 30%. Similar measurements were made by Fonseca and Campos who used a low-pressure spectral lamp as an excitation source and lifetimes measured in an electron excitation experiment for absolute measurements of transition rates. A number of other experiments used thermal plasma sources, either wall-stabilized electric arcs in the experiments by Ernst and Schulz–Gulde and Brandt, Helbig and Nick or a shock-tube in the experiment by Kaschek, Ernst and Bötticher . These experiments depended on plasma diagnostics and the transition rates have relative uncertainties that are generally not much better than $`\pm `$10% even for strong transitions.
## 2 Upper level lifetimes
The lifetimes of six of the ten 5p upper levels are known with relative standard uncertainties of the order of 0.2% or better from recent beam-gas-laser spectroscopy (BGLS) measurements by Schmitt et al. . Our primary purpose here is to find reliable estimates for the lifetimes of the remaining four 5p levels. We will do this with a semi-empirical theoretical approach that is based upon intermediate coupling theory. In addition, we have evaluated the published experimental data to have an alternate set of lifetimes. Although these lifetimes are less accurate than the semi-empirical ones they provide bounds for the semi-empirical lifetimes.
### 2.1 Experimental lifetimes
A comparison with the six reference lifetimes from BGLS divides the experimental data sets from literature in two classes of different reliability. The results from experiments employing selective laser excitation and from the only wall-stabilized arc emission experiment carried out so far generally fall (with a few explainable exceptions) into a $`\pm `$8% tolerance band around the six reference lifetimes (see Table 1). The pulsed-laser lifetime measurements encountered some problems for the closely-spaced levels 2p<sub>8</sub> and 2p<sub>9</sub> due to fast collisional mixing that resulted in non-exponential decay curves. Apart from these two levels the results agree within $`\pm `$8% with the BGLS lifetimes. A tendency towards underestimated error bars, however, is obvious for all three experiments. The lifetimes resulting from the arc emission experiment also agree within $`\pm `$8% with the BGLS results despite of an uncertainty of $`\pm `$30% the authors quote for their absolute intensity scale. The exception here is the level 2p<sub>9</sub> for which saturation problems were not adequately treated.
These measurements (summarized in Table 1) are the best measured lifetimes for the remaining four 5p levels. The estimated relative uncertainties are around $`\pm `$8%. Other experiments which employed pulsed electron excitation or the Hanle effect all have produced at least one result far outside of the $`\pm `$8% tolerance range and will therefore not be considered further.
### 2.2 Semi-empirical lifetimes
The 5p levels in krypton decay exclusively (apart from some very weak far-IR channels for the four highest 5p levels) through the transitions of the 5p–5s array. In a typical semi-empirical calculation for this transition array (e.g. Lilly , see column CA in Table 2) the wavefunctions of the initial and final configurations $`\gamma =`$4p<sup>5</sup>5p and $`\gamma ^{}=`$4p<sup>5</sup>5s in intermediate coupling are expressed in terms of LS-coupled wavefunctions $`|\gamma LSJM>`$:
$`|i,JM>`$ $`=`$ $`{\displaystyle \underset{LS}{}}|\gamma LSJM>a(\gamma LSJ,i)`$ (1)
$`|f,J^{}M^{}>`$ $`=`$ $`{\displaystyle \underset{L^{}S^{}}{}}|\gamma L^{}S^{}J^{}M^{}>a(\gamma ^{}L^{}S^{}J^{},f)`$ (2)
($`i`$ and $`f`$ denote the coupled initial and final states). The mixing coëfficients $`a(\gamma LSJ,)`$ can be determined from experimental energies with a semi-empirical fit procedure in the manner described by Lilly . Once the mixing coefficients have been determined, the reduced dipole matrix elements $`<i||D||f>`$, which are proportional to the transition rates $`A_{if}`$, may be expressed in terms of the reduced matrix elements in LS-coupling. The latter can be reduced further, using angular momentum theory, to
$$<\gamma LSJ||D||\gamma ^{}L^{}S^{}J^{}>=\delta _{SS^{}}(2J+1)^{1/2}(2J^{}+1)^{1/2}(1)^{L+1+S+J^{}}\left\{\genfrac{}{}{0pt}{}{SLJ}{1J^{}L^{}}\right\}<\gamma LS||D||\gamma ^{}L^{}S^{}>.$$
(3)
In the single electron (or Coulomb) approximation the reduced dipole matrix element is proportional to the dipole transition moment $`\sigma _{\gamma \gamma ^{}}`$:
$$<\gamma LS||D||\gamma ^{}L^{}S^{}>=(2L+1)^{1/2}\sigma _{\gamma \gamma ^{}}$$
(4)
where
$$\sigma _{\gamma \gamma ^{}}=\sqrt{3}_0^{\mathrm{}}u_{5\mathrm{p}}(r)eru_{5\mathrm{s}}(r)𝑑r$$
(5)
and $`u_{5\mathrm{p}}(r)`$, $`u_{5\mathrm{s}}(r)`$ are the radial wavefunctions of the valence electron. In this simple semi-empirical model the relative transition rates, and thus the lifetime ratios, depend on the intermediate coupling coefficients $`a(\gamma LSJ,)`$ of the 5p and 5s configurations, and the absolute scale is given by one single transition moment $`\sigma _{\gamma \gamma ^{}}`$ for the entire 5p–5s transition array.
When we use the transition moment $`\sigma _{\gamma \gamma ^{}}=3.04`$ a.u. that was obtained in the semi-empirical calculation by Lilly we find that, on average, the six experimental BGLS lifetimes can be reproduced no better than within 7%. The predictions from this semi-empirical model for the remaining lifetimes are presumably not more accurate than the recommended experimental values are (see Table 1). It appears unlikely that the mixing coefficients of the 5s and 5p configurations are responsible for the lesser accuracy of the semi-empirical lifetimes since the reproduction of the experimental energies by the semi-empirical intermediate coupling method is quite good . The problem is the assumption of one single transition moment $`\sigma _{\gamma \gamma ^{}}`$ for the entire transition array. To refine the semi-empirical model, we assumed $`LS`$-dependent transition moments $`\sigma _{\gamma \gamma ^{}}(L,S,L^{},S^{})`$ which correspond to LS-dependent radial functions $`u_{5\mathrm{p}}(LS)`$ and $`u_{5\mathrm{s}}(L^{}S^{})`$ in Eq. 5. These are similar to those used, for example, in Hartree-Fock calculations. We further assumed that spin-orbit interaction only results in a mixture of LS-terms but not in a modification of the radial wavefunctions. Since a calculation of these transition moments from first principles or from experimental energies would not have been accurate enough for our purposes we determined the transition moments from the six reference lifetimes from BGLS.
In total six different non-zero transition moments are needed for the description of the 5p-5s array. They correspond to the six allowed transitions in $`LS`$-coupling (see Table 3). In one case, for the transition 5p <sup>3</sup>$``$ 5s <sup>3</sup>P, the transition moment may be calculated directly from the lifetime of the level 2p<sub>9</sub> since both the initial (<sup>3</sup>D<sub>3</sub>) and the final state (<sup>3</sup>P<sub>2</sub>) of the only decay channel are pure states in $`LS`$-coupling. Generally, the transition moments have to be determined by means of a nonlinear least-squares fit procedure that adjusts the transition moments so as to get best agreement of the calculated lifetimes with the six reference lifetimes from BGLS. The results are summarized in Table 3. The quoted standard uncertainties of the semi-empirical transition moments and lifetimes were obtained by Gaussian propagation of the uncertainties of the reference lifetimes, the uncertainties in the energy parameters (see ), and the uncertainties of the contributing branching ratios.
The energy matrices and the intermediate coupling coefficients of the 5p and 5s configurations (see Table 2) were recalculated from the Slater- and spin-orbit parameters (including the $`\alpha L(L+1)`$ correction) given by Lilly . The six reference states (2p<sub>3,4,6..9</sub>) for which the lifetimes are known very accurately are mostly built from the four $`LS`$-terms <sup>1</sup>P, <sup>1</sup>D, <sup>3</sup>P, and <sup>3</sup>D. The four corresponding transition moments could thus be deduced with high accuracy (see Table 3). The transition moment of the transition 5p <sup>3</sup>$``$ 5s <sup>3</sup>P was determined with a somewhat greater uncertainty from the lifetime of the 2p<sub>3</sub> state. This state is the only one in the set of reference states that contains a relevant contribution (16%) from the <sup>3</sup>S term. The <sup>1</sup>S term only contributes to the states 2p<sub>1</sub> and 2p<sub>5</sub> which are not included in the set of reference states. For the determination of the transition moment of the transition 5p <sup>1</sup>$``$ 5s <sup>1</sup>P we resorted to branching fractions as additional criteria. Particularly, we used branching fractions of the weak decay channels 2p<sub>1</sub> $``$ 1s<sub>4</sub> and 2p<sub>5</sub> $``$ 1s<sub>2</sub> which are sensitive to the transition moment sought after. The dependence of the branching fractions for these transitions and the upper level lifetimes on the transition moment is shown in figure Fig. 3. The attainable accuracy for the transition moment, however, is limited by the uncertainty in the branching fractions.
As a last technical detail we note that for the four highest 5p levels (2p<sub>1</sub> $`\mathrm{}`$ 2p<sub>4</sub>) there are weak far-IR decay channels to states of the 4d configuration which have to be accounted for. For this purpose we used the theoretical transition rates calculated by Aymar and Coulombe . Because of the huge discrepancies between length- and velocity-form results for the 5p–4d transition rates we used the greater velocity-form results with a pessimistic uncertainty estimate of $`\pm `$100% (see Table 2).
The uncertainties of the four semi-empirical lifetimes (see Table 2) for the states 2p<sub>1</sub>, 2p<sub>2</sub>, 2p<sub>5</sub>, and 2p<sub>10</sub> vary between 0.1% and 3.7% depending on the $`LS`$-terms they are built from. The 2p<sub>2</sub> state allows for a very precise lifetime calculation because it relates to the four accurately determined transition moments only. The lifetime predictions for the other three levels are less accurate because they include significant contributions from the two less accurate transition moments. Our semi-empirical predictions agree very well with the best previous experimental values (see column BE in Table 1) but they are of superior accuracy and we used them for the normalization of our transition rates.
## 3 Branching fractions
The measurement of branching fractions for the transition from 5p levels presents a formidable task owing to the metastable nature of the lower 5s and 5s levels (see Fig. 1) which may render the lamp discharge column optically thick for transitions to those levels. We have measured branching fractions for 30 lines arising from 5p levels in the wavelength range from 556.2 nm to 1878.5 nm in two separate experiments. The spectral lines in the visible part of the Kr I spectrum were measured in air with a wall-stabilized arc discharge and a 2m – Czerny-Turner monochromator. The infrared portion of the spectrum was measured with a hollow-cathode lamp and the NIST 2m – Fourier transform spectrometer. The comparison of the results from the two different experiments made it easier for us to notice systematic errors due to optically thick transitions in the light sources. The four 5p – 4d transitions (see Fig. 1) near 10000 nm were outside the range of either experiment.
### 3.1 Wall-stabilized arc measurements
The experimental setup for the measurements is shown schematically in Fig. 2. In our experiment we used the wall-stabilized arc previously described in detail by Musielok et al. . The space near the electrodes was operated in argon while the midsection of the arc channel contained helium with a small admixture of krypton. The fraction of krypton in helium was maintained below 0.3% to avoid self-absorption of krypton lines. The arc was operated at a current of 50 A. To check for optical thickness, the krypton spectra were measured with varying amounts of krypton in the discharge.
When the wall-stabilized arc is operated in helium, spectral lines remain narrow and continuum emission is low because of the low electron density in a helium arc. This facilitates more accurate line intensity measurements because spectral lines are well isolated and the ratio of line to continuum intensity is high. It was not necessary to achieve LTE conditions in the arc plasma, because we were only interested in the measurement of branching ratios of spectral lines. The measurements were performed in a side-on configuration to avoid interloping argon lines and argon plasma continuum radiation that are emitted at the ends of the arc. As indicated in Fig. 2, either the wall–stabilized arc or a tungsten strip standard lamp were imaged onto the entrance slit of a 2 m Czerny–Turner monochromator by a concave mirror with a magnification factor of approximately 1.3. A beam splitter was placed in the beam path to reflect a fraction of the light into the 0.25 m monochromator that was used to monitor the discharge stability. This monochromator was set to the 760.2 nm line of Kr I. The total intensity of this line was measured with a photomultiplier tube and a chart recorder and showed less than 1% fluctuation during our measurements. The krypton spectra were recorded with a CCD camera that was mounted at the exit plane of the monochromator. The measured spectral line profiles were first corrected for the spectral response of the experimental system, as determined with the standard lamp, and the residual continuum was subtracted. The lines were then integrated by fitting a spline function to the data using a program package published by Renka which yields the integral of the spectral line without requiring that the apparatus function be known analytically.
### 3.2 Hollow cathode lamp measurements
The experiment described in the previous section was unsuited for measurements in the infrared because it was set up in air. A second experiment in a purged environment was therefore carried out to measure the intensity of lines in the infrared. This used a high-resolution Fourier transform spectrometer to observe spectra of a hollow cathode lamp.
The high-current hollow cathode lamp we used was developed by Danzmann et al. . For our measurements it was equipped with a cathode made of oxygen-free copper which is easy to operate and has no lines that blend with the krypton lines of interest. The hollow cathode lamp was operated with between 130 Pa and 250 Pa of argon or neon as a carrier gas for the discharge with an admixture of between 0.5 Pa and 10 Pa. The discharge current was varied between 100 mA and 500 mA. The experimental setup was similar to the one used with the wall-stabilized arc. The entire imaging system was enclosed in a purge box that was continuously purged with water vapor and carbon dioxide free air to suppress absorption by these gases in the near infrared.
Many lines were strongly self-absorbed when pure krypton was used as a carrier gas in the hollow cathode lamp discharge. This problem was partly overcome when the partial pressure of krypton in the discharge was reduced by using a neon-krypton mixture in the hollow cathode lamp. We also found that the spectra obtained with high currents where the copper density in the discharge is high show self absorption only in the very strongest lines. We assume that the metastable 5s states were depopulated by charge-transfer collisions with copper atoms in the hollow cathode lamp discharge.
The NIST 2m – Fourier transform spectrometer (described in Nave et al. ) was used to measure the spectra of the hollow cathode lamp and the standard lamp. A resolution of around 0.01 cm<sup>-1</sup> was used for the measurements of the krypton spectra. For the near IR region, a liquid nitrogen cooled indium-antimonide detector was used whereas silicon photodiodes were used to record spectra below 1000 nm. To improve the signal-to-noise ratio, colored glass filters were employed to restrict the bandpass of the spectrometer to the wavenumber range of interest. The spectral sensitivity of the optical and detection systems was calibrated with the standard lamp before and after measurements of the spectra of the light from the hollow cathode discharge. Some residual self-absorption was evident for the strongest lines even at low krypton partial pressures. For those lines we relied on the results from the experiment with the wall-stabilized arc, where these lines remained optically thin.
### 3.3 Data analysis and uncertainties
It is common that the the uncertainty of experimental transition rates is limited by the uncertainty of the measurement of the upper level lifetimes and not by the uncertainty of the branching fraction measurement. In our case, the situation is reversed. The uncertainty of the branching fraction measurement is limited by the uncertainty of the radiometric calibration which is around 2%. The uncertainties of the upper level lifetimes are generally much lower. In this section we will describe in detail how the uncertainties for the branching fractions were calculated.
The transition rate $`A_{ki}`$ of a transition from a particular upper level $`k`$ to lower level $`i`$ can be calculated from a measurement of the upper level lifetime $`\tau _k`$ and a measurement of the branching fraction $`F_{ki}`$ – the fraction that the transition to $`i`$ contributes to the total decay rate:
$$A_{ki}=\frac{1}{\tau _k}F_{ki},\mathrm{where}F_{ki}=\frac{A_{ki}}{_jA_{ki}}.$$
(6)
The branching fractions can in turn be calculated from the relative intensities $`I_{ki}`$ of the lines (in photons/s) by
$$F_{ki}=\frac{I_{ki}}{_jI_{ki}}$$
(7)
where the sum is over all the lower levels to which the upper level can decay.
Several independent measurements of the Kr I spectrum were made with different operating conditions for the hollow cathode lamp. The relative intensity $`\widehat{I}_{ki}^\alpha `$ of each spectral line in each measured spectrum $`\alpha `$ was calculated from the observed intensity $`I_{ki}^\alpha `$ and the relative efficiency of the spectrometer $`ϵ(\sigma )`$ at the wavenumber $`\sigma `$ of the spectral line by:
$$\widehat{I}_{ki}^\alpha =\frac{I_{ki}^\alpha (\sigma )}{ϵ(\sigma )}$$
(8)
where the relative efficiency of the optical system was assumed to be constant over the width of the spectral line. These intensities were then divided by a normalizing factor $`\widehat{I}_{\mathrm{norm}}^\alpha `$ to put all the intensities in all spectra on the same relative intensity scale. This normalizing factor was usually chosen such that the intensity of one strong line common to all spectra was 1, hence making the intensities relative with respect to that strong line. This approach was found to be more reliable than using a weighted mean of the intensities, as lines in some of the spectra may be affected by self-absorption, or be too weak to be measured. The weighted mean relative intensity of the line $`\overline{I}_{ki}`$ was then found using:
$$\overline{I}_{ki}=\frac{1}{_\alpha w_{ki}^\alpha }\underset{\alpha }{}\frac{w_{ki}^\alpha \widehat{I}_{ki}^\alpha }{\widehat{I}_{\mathrm{norm}}^\alpha }$$
(9)
where w$`{}_{ki}{}^{}{}_{}{}^{\alpha }`$ is a weighting factor. The weighting factor chosen for the hollow cathode measurements was the signal-to-noise ratio of the line. The small uncertainty of the normalization line intensity is due to its high signal-to-noise ratio. These branching ratios are converted to branching fractions and transition rates using equations 6 and 7.
Absolute transition rates are then determined from the mean values of between 5 and 9 independent measurements of the relative intensities and experimental lifetime data. They are presented in table 4, along with the lifetimes of the upper levels used to determine the transition rates. The uncertainties given in the table result from the estimated standard deviation of the branching fractions and the uncertainty of the lifetime data. The estimated standard deviation of the branching fractions depends on the uncertainty in the weighted mean relative intensity, which in turn depends on the individual measurements of the intensity through equation 9, and the uncertainty in the radiometric calibration of the spectrometer.
The estimated uncertainty in the individual measurements of the intensity was taken as the intensity divided by the signal-to-noise ratio: $`\widehat{I}_{ki}^\alpha /w_{ki}^\alpha `$. When photon noise is the dominant source of uncertainty, the square of the signal-to-noise ratio must be used as the weighting factor in equation 9. We chose to weight the individual intensity measurements with the signal-to-noise ratio to account for a significant systematic component in the uncertainty which may result from self-absorption or line blends.
The statistical component in the uncertainty of the weighted mean relative intensity $`u_{stat}(\widehat{I}_{ki}`$) can then be derived by applying the law of propagation of uncertainty to equation 9:
$$u_{\mathrm{stat}}(\overline{I}_{ki})=\sqrt{\underset{\alpha }{}\left(\frac{\widehat{I}_{ki}^\alpha }{\widehat{I}_{\mathrm{norm}}^\alpha }\frac{1}{w_{ki}^\alpha }\right)^2}$$
where the sum is again over all the observations of the lines. This must be added in quadrature to the uncertainty in the radiometric calibration of the spectrometer, which was estimated at 3.3% for one standard deviation. This estimate includes the uncertainty in the supplied calibration of the standard lamp (1.5% for one standard deviation) and a contribution of 3% for the measurement of the standard lamp spectrum. The total uncertainty in the measurement of the weighted mean relative intensities is thus:
$$u(\overline{I}_{ki})=\sqrt{u_{\mathrm{stat}}^2(\overline{I}_{ki})+(0.033\overline{I}_{ki})^2}$$
(10)
The uncertainty in the measurement of the branching fractions $`u(F_{ki}`$) is derived by applying the law of propagation of uncertainty to equation 7 to give:
$$u(F_{ki})=\sqrt{\frac{u^2(\overline{I}_{ki})}{(_i\overline{I}_{ki})^2}+\frac{\overline{I}_{ki}^2}{(_i\overline{I}_{ki})^4}\underset{j}{}u^2(\overline{I}_{ki})}$$
(11)
This is combined in quadrature with the uncertainty in the lifetime $`u(\tau _k`$) to give the uncertainty in the transition rate $`u(A_{ki}`$):
$$u(A_{ki})=\sqrt{\frac{1}{\tau _k^2}u^2(F_{ki})+\frac{F_{ki}^2}{\tau _k^4}u^2(\tau _k)}$$
(12)
For the wall-stabilized arc measurements, the transition rates and their uncertainties were calculated similarly.
## 4 Discussion of results
Our new transition rates for 5p – 5s transitions in Kr I are listed in Table 4. Also listed in Table 4 are our experimental branching fractions and the lifetimes that were used to calculate the transition rates.
In Table 5 and Fig. 4 we compare our transition rates with several experimental results. The only other measurement that includes the lines in the IR is that of Chang, Horiguchi and Setser and a comparison of those results with our data is shown in Fig. 5. The results by Fonseca and Campos , presented in Fig. 4, were recalculated using the same lifetime data as in our work. For the set of strong lines around 800 nm our transition rates are in good agreement with most of the results obtained by Fonseca and Campos and differ from those obtained by Ernst et al. by about 10%. The results of Kaschek et al. differ from ours by a constant scaling factor of 1.3, on average. Only in the case of the line at 810.4 nm, all other experimental data exceed our result by 20%–30%.
The most striking difference between previous measurements and our results is that our transition rates for the set of weak lines near 600 nm are much lower than all previous measurements with the exception of the experiment by Brandt, Helbig and Nick . This strongly suggests that many of the earlier experiments had problems with self-absorption of the strong lines around 800 nm which would make the weak lines in a set of transitions from a particular upper level appear stronger.
It is also interesting to compare our results with the theoretical calculations because all calculations were intermediate-coupling calculations in the Coulomb approximations whereas our semi-empirical lifetimes were obtained with a modified intermediate-coupling scheme. Fig. 6 compares our results to the most recent calculations. The earlier calculations by Murphy are not included because they were superseded by those of Lilly . For the strongest lines near 800 nm the best agreement, within 10% on average, was found between our data and calculations made by Aymar and Coulombe with a velocity dipole operator, while there is a constant disagreement (a factor of 1/3) when they used a length dipole operator. A similar discrepancy was found with calculations made by Lilly . We note that the discrepancies for the weak lines near 600 nm and in the IR are considerable but there appear to be no conspicuous systematic trends as we found in the experimental data.
Acknowledgments
K. Dzierżega gratefully acknowledges financial support from the Maria Skłodowska-Curie Foundation through grant number MEN–NIST–96–260. We are grateful to W. L. Wiese, NIST, K. Musioł, Jagiellonian University, Kraków, and H. Schmoranzer, University of Kaiserslautern for helpful discussions and continued support. G. Nave and U. Griesmann were supported by NIST contracts number 43SBNB867005 and 43SBNB960002 to Harvard College Observatory.
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# The effects of Majorana phases in three-generation neutrinos
## 1 Introduction
Recent neutrino oscillation experiments suggest the strong evidences of tiny neutrino masses and lepton flavor mixings. Studies of the lepton flavor mixing matrix, which is so-called Maki-Nakagawa-Sakata(MNS) matrix, will give us important cues of the physics beyond the standard model. One of the most important studies is the analysis of the quantum correction on the MNS matrix.
In order to explain both the solar and the atmospheric neutrino problems, two mass squared differences are needed, which implies
$$\mathrm{\Delta }m_{\mathrm{solar}}^2\left|m_2^2m_1^2\right|,\text{ and }\mathrm{\Delta }m_{\mathrm{ATM}}^2\left|m_3^2m_2^2\right|,$$
(1)
where $`m_i`$ is the $`i`$-th ($`i=13`$) generation neutrino mass ($`m_i0`$). $`\mathrm{\Delta }m_{\mathrm{solar}}^2`$ and $`\mathrm{\Delta }m_{\mathrm{ATM}}^2`$ stand for the mass-squared differences of the solar neutrino and the atmospheric neutrino solutions, respectively. Then there are the following three possible types of neutrino mass hierarchies ;
Type A $`:`$ $`m_1m_2m_3,`$
Type B $`:`$ $`m_1m_2m_3,`$ (2)
Type C $`:`$ $`m_1m_2m_3,`$
where $`m_i`$ is the $`i`$-th generation neutrino absolute mass. In Ref. , it has been studied whether the lepton-flavor mixing angles are stable or not against quantum corrections for all three types of mass hierarchies with all considerable relative sign assignments, which are shown below, in the minimal supersymmetric standard model (MSSM) with an effective dimension-five operator which gives the Majorana masses of neutrinos.
1. Type A:
$`\mathrm{case}(\mathrm{a1}):m_\nu ^{\mathrm{a1}}`$ $`=`$ $`diag.(0,m_2,m_3),`$ (3)
$`\mathrm{case}(\mathrm{a2}):m_\nu ^{\mathrm{a2}}`$ $`=`$ $`diag.(0,m_2,m_3).`$ (4)
$$\left(m_1=0,m_2=\sqrt{\mathrm{\Delta }m_{\mathrm{solar}}^2},m_3=\sqrt{\mathrm{\Delta }m_{\mathrm{solar}}^2+\mathrm{\Delta }m_{\mathrm{ATM}}^2}\right)$$
2. Type B:
$`\mathrm{case}(\mathrm{b1}):m_\nu ^{\mathrm{b1}}`$ $`=`$ $`diag.(m_1,m_2,0),`$ (5)
$`\mathrm{case}(\mathrm{b2}):m_\nu ^{\mathrm{b2}}`$ $`=`$ $`diag.(m_1,m_2,0).`$ (6)
$$\left(m_1=\sqrt{\mathrm{\Delta }m_{\mathrm{ATM}}^2},m_2=\sqrt{\mathrm{\Delta }m_{\mathrm{solar}}^2+\mathrm{\Delta }m_{\mathrm{ATM}}^2},m_3=0\right)$$
3. Type C:
$`\text{case (c1): }m_\nu ^{c1}`$ $`=`$ $`diag.(m_1,m_2,m_3),`$ (7)
$`\text{case (c2): }m_\nu ^{c2}`$ $`=`$ $`diag.(m_1,m_2,m_3),`$ (8)
$`\text{case (c3): }m_\nu ^{c3}`$ $`=`$ $`diag.(m_1,m_2,m_3),`$ (9)
$`\text{case (c4): }m_\nu ^{c4}`$ $`=`$ $`diag.(m_1,m_2,m_3).`$ (10)
$$\left(m_1=m_0,\text{ }m_2=\sqrt{m_0^2+\mathrm{\Delta }m_{\mathrm{solar}}^2},\text{ }m_3=\sqrt{m_0^2+\mathrm{\Delta }m_{\mathrm{solar}}^2+\mathrm{\Delta }m_{\mathrm{ATM}}^2}\right)$$
In Ref., it has been found that the above relative sign assignments of neutrino masses in each type play the crucial roles for the stability of the mixing angles against quantum corrections. Actually, two physical Majorana phases in the lepton flavor mixing matrix connect among the above relative sign assignments of neutrino masses. Therefore, in this paper we analyze the stability of mixing angles against quantum corrections according to three types of neutrino mass hierarchies (Type A, B, C) and two Majorana phases. Two phases play the crucial roles for the stability of the mixing angles against the quantum corrections. In Refs. , it has been already analyzed that the effect of a Majorana phase plays an important role for the stability against the quantum corrections in the two-generation neutrinos.
## 2 Quantum corrections to neutrino mass matrix
In the MSSM with the effective dimension-five operator which gives Majorana masses of neutrinos, the superpotential of the lepton-Higgs interactions is given by
$$𝒲=y_{ij}^\mathrm{e}(H_dL_i)E_j\frac{1}{2}\kappa _{ij}(H_uL_i)(H_uL_j).$$
(11)
Here the indices $`i,j`$ $`(=13)`$ stand for the generation number. $`L_i`$ and $`E_i`$ are chiral super-fields of $`i`$-th generation lepton doublet and right-handed charged-lepton, respectively. $`H_u`$ ($`H_d`$) is the Higgs doublet which gives Dirac masses to the up- (down-) type fermions. The neutrino mass matrix of the three generations, $`\kappa `$ is diagonalized as
$$U^T\kappa U=D_\kappa ,$$
(12)
where $`D_\kappa `$ is given by
$$D_\kappa =\left(\begin{array}{ccc}m_1& 0& 0\\ 0& m_2& 0\\ 0& 0& m_3\end{array}\right),$$
(13)
with $`m_i0`$. The unitary matrix $`U`$ is defined as
$$U=\left(\begin{array}{ccc}U_{e1}& U_{e2}& U_{e3}\\ U_{\mu 1}& U_{\mu 2}& U_{\mu 3}\\ U_{\tau 1}& U_{\tau 2}& U_{\tau 3}\end{array}\right)\left(\begin{array}{ccc}e^{i\varphi _1}& 0& 0\\ 0& e^{i\varphi _2}& 0\\ 0& 0& 1\end{array}\right),$$
(14)
where $`\varphi _{1,2}`$ denote the physical Majorana phases of the lepton sector. In the diagonal base of charged lepton masses, $`U`$ is just the MNS matrix. We can easily show that one Majorana phase connects between cases of (a1) and (a2), (b1) and (b2), and two Majorana phases connect among cases of (c1)$``$(c4). Thus, the stabilities of mixing angles against quantum corrections are completely determined by three types of neutrino mass hierarchies (Type A, B, C) and two Majorana phases $`\varphi _{1,2}`$ in stead of the classifications of Eqs. (3) $``$ (10).
We will analyze whether the lepton flavor mixing angles are changed or not by the quantum corrections by fitting the low energy data. We determine the MNS matrix at $`m_Z`$ scale as
$$U=\left(\begin{array}{ccc}\mathrm{cos}\theta _{12}& \mathrm{sin}\theta _{12}& 0\\ \frac{\mathrm{sin}\theta _{12}}{\sqrt{2}}& \frac{\mathrm{cos}\theta _{12}}{\sqrt{2}}& \frac{1}{\sqrt{2}}\\ \frac{\mathrm{sin}\theta _{12}}{\sqrt{2}}& \frac{\mathrm{cos}\theta _{12}}{\sqrt{2}}& \frac{1}{\sqrt{2}}\end{array}\right)\left(\begin{array}{ccc}e^{i\varphi _1}& 0& 0\text{}\\ 0& e^{i\varphi _2}& 0\text{}\\ 0& 0& 1\end{array}\right),$$
(15)
where we input $`\mathrm{sin}\theta _{23}=1/\sqrt{2}`$ and $`\mathrm{sin}\theta _{13}=0`$ which values are suitable for the atmospheric neutrino experiments and for the CHOOZ experiment , respectively. The mixing angle $`\theta _{12}`$ depends on the solar neutrino solutions of the large angle MSW solution (MSW-L), the small angle MSW solution (MSW-S) and vacuum oscillation solution (VO), which are given by
$`\mathrm{sin}\theta _{12}=\{\begin{array}{ccc}0.0042\hfill & (\theta =0.0042)\hfill & \text{(MSW-S), }\hfill \\ \frac{1}{\sqrt{2}}\hfill & (\theta =\frac{\pi }{4})\hfill & \text{(MSW-L),}\hfill \\ \frac{1}{\sqrt{2}}\hfill & (\theta =\frac{\pi }{4})\hfill & \text{(VO).}\hfill \end{array}`$ (19)
We also use the following values of mass-squared differences in the numerical analyses.
$`\mathrm{\Delta }m_{\mathrm{solar}}^2`$ $``$ $`\{\begin{array}{cc}0.8\times 10^5& \text{eV}\text{2}\text{ (MSW-S),}\hfill \\ 1.8\times 10^5& \text{eV}\text{2}\text{ (MSW-L),}\hfill \\ 0.85\times 10^{10}& \text{eV}\text{2}\text{ (VO), }\hfill \end{array}`$ (23)
$`\mathrm{\Delta }m_{\mathrm{ATM}}^2`$ $``$ $`3.7\times 10^3\text{eV}\text{2}.`$ (24)
The quantum corrections change the neutrino mass matrix, and it is given by<sup>1</sup><sup>1</sup>1Hereafter, we denote the mixing angles and the other physical parameters at the $`m_R`$ scale are written with $`\widehat{}`$ mark.
$`\widehat{\kappa }\left(m_R\right)`$ $`=`$ $`{\displaystyle \frac{\widehat{\kappa }\left(m_R\right)_{33}}{\kappa \left(m_Z\right)_{33}}}\left(\begin{array}{ccc}1ϵ& 0& 0\\ 0& 1ϵ& 0\\ 0& 0& 1\end{array}\right)\kappa \left(m_Z\right)\left(\begin{array}{ccc}1ϵ& 0& 0\\ 0& 1ϵ& 0\\ 0& 0& 1\end{array}\right),`$ (31)
at the high energy scale $`m_R`$, where $`ϵ`$ can be estimated as
$`ϵ`$ $``$ $`1\mathrm{exp}\left({\displaystyle \frac{1}{16\pi ^2}}{\displaystyle _{\mathrm{ln}\left(m_Z\right)}^{\mathrm{ln}\left(m_R\right)}}y_\tau ^2𝑑t\right),`$ (32)
$``$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle \frac{m_\tau ^2}{v^2}}\left(1+\mathrm{tan}^2\beta \right)\mathrm{ln}\left({\displaystyle \frac{m_R}{m_Z}}\right).`$
where $`y_\tau `$ is the Yukawa coupling of $`\tau `$, $`v^2v_u^2+v_d^2`$ and $`\mathrm{tan}\beta v_u/v_d`$ ($`v_u`$ and $`v_d`$ are vacuum expectation values of Higgs bosons, $`H_u`$ and $`H_d`$, respectively). We neglect the Yukawa couplings of $`e`$ and $`\mu `$ in Eqs.(31) and (32), since those contributions to the renormalization group equations are negligibly small comparing to that of $`\tau `$ . The magnitude of $`ϵ`$ can be determined by the value of $`\mathrm{tan}\beta `$ and the scale of $`m_R`$. The unitary matrix $`\widehat{U}`$ which diagonalizes $`\widehat{\kappa }`$ shows us whether the lepton flavor mixing angles are stable against quantum corrections or not.
## 3 Type A $`(m_1m_2m_3)`$
In both (a1) and (a2) cases, all mixing angles are stable against quantum corrections in each sign assignment . This is understood from the analogy of two-generation analysis, which shows the mixing angle of $`2\times 2`$ mass matrix is not changed drastically by the quantum corrections when there is the large mass hierarchy between two neutrinos . This situation is not changed when we consider the Majorana phase contribution as shown in two-generation neutrinos . Cases (a1) and (a2) are connected with each other by the Majorana phase of $`\varphi _2`$. Where $`\varphi _1`$ is rotated out, since $`m_1=0`$. The case of $`\varphi _2=0`$ corresponds to (a1), while the case of $`\varphi _2=\pi /2`$ corresponds to (a2). Since Type A has the large mass hierarchies, all mixing angles are supposed to be stable against quantum corrections independently of the value of the Majorana phase $`\varphi _2`$. This is really confirmed by numerical analyses as shown in Table 1, where we use $`m_R=10^{13}`$ GeV and $`\mathrm{tan}\beta =60`$.
## 4 Type B $`(m_1m_2m_3)`$
In Type B mass hierarchy, all mixing angles except for $`\mathrm{sin}\theta _{12}`$ of (b2) are stable against quantum corrections . The analogy of two-generation neutrinos analysis shows that mixing angles of $`\mathrm{sin}\theta _{13}`$ and $`\mathrm{sin}\theta _{23}`$ are stable against quantum corrections, since there are large mass hierarchies between the first and the third generations, and between the second and third generations. This is the same situation as that of Type A. This situation is not changed by including the Majorana phase contributions of $`\varphi _{1,2}`$ as shown in Table 2, which shows the results of the numerical analyses in the case of $`m_R=10^{13}`$ GeV and $`\mathrm{tan}\beta =60`$.
On the other hand, the mixing angle of $`\theta _{12}`$ can receive significant quantum corrections dependently on the relative sign assignment of $`m_2`$ as shown in Ref . The mixing angle of $`\mathrm{sin}\theta _{12}`$ of (b1) receives the quantum correction while that of (b2) does not. Now we understand that two cases of (b1) and (b2) are connected with each other by the phase of $`\varphi \varphi _1\varphi _2`$, which is the only physical phase, since $`m_3=0`$. The case of $`\varphi =0`$ corresponds to (b1), while the case of $`\varphi =\pi /2`$ corresponds to (b2). The phase $`\varphi `$ is the parameter which determines whether the mixing angle $`\theta _{12}`$ is stable against quantum corrections or not.
Now let us show the analytic estimations for the stabilities of the mixing angles in Type B mass hierarchy. The neutrino mass matrix of Type B which is diagonalized is given by
$$D_\kappa ^{(B)}=m_1\left(\begin{array}{ccc}1& 0& 0\\ 0& 1+\xi _b& 0\\ 0& 0& 0\end{array}\right),$$
(33)
where
$$\xi _b\frac{m_2m_1}{m_1}\frac{1}{2}\frac{\mathrm{\Delta }m_{\mathrm{solar}}^2}{\mathrm{\Delta }m_{\mathrm{ATM}}^2}.$$
(34)
We can determine the mass matrix of $`\kappa ^{(B)}`$ by using Eqns.(12) and (15). Then Eq.(31) gives the mass matrix of $`\widehat{\kappa }^{(B)}`$ at the high energy scale $`m_R`$.
The MNS matrix $`\widehat{U}^{(B)}`$ which diagonalizes $`\widehat{\kappa }^{(B)}`$ is given by
$`\widehat{U}^{(B)}`$ $`=`$ $`\left(\begin{array}{ccc}1& 0& 0\\ 0& (1ϵ)/\sqrt{1+(1ϵ)^2}& 1/\sqrt{1+(1ϵ)^2}\\ 0& 1/\sqrt{1+(1ϵ)^2}& (1ϵ)/\sqrt{1+(1ϵ)^2}\end{array}\right)`$ (45)
$`\times \left(\begin{array}{ccc}\mathrm{cos}\widehat{\theta }_{12}& \mathrm{sin}\widehat{\theta }_{12}& 0\\ \mathrm{sin}\widehat{\theta }_{12}& \mathrm{cos}\widehat{\theta }_{12}& 0\\ 0& 0& 1\end{array}\right)\left(\begin{array}{ccc}e^{i\widehat{\varphi }_1}& 0& 0\\ 0& e^{i\widehat{\varphi }_2}& 0\\ 0& 0& 1\end{array}\right),`$
which means that the mixing angle between the first and the third generations, which is zero, is unchanged by quantum corrections. The mixing angle of $`\widehat{\theta }_{23}`$ is given by
$$\mathrm{sin}^22\widehat{\theta }_{23}=\left(\frac{2(1ϵ)}{1+(1ϵ)^2}\right)^2,$$
(46)
which indicates that the large mixing between the second and the third generations is stable with respect to quantum corrections. By using Eq.(46), we can estimate that $`\mathrm{sin}^22\widehat{\theta }_{23}0.99`$ in the case of $`m_R=10^{13}`$ GeV and $`\mathrm{tan}\beta =60`$, which is consistent with the numerical analysis in Table 2. Therefore the mixing between the first and the third generations and the mixing between the second and the third generations are stable with respect to quantum corrections as shown in Table 2
How about the mixing between the first and the second generations?
For the MSW-L and the VO solutions, where $`\mathrm{sin}\theta _{12}=\mathrm{cos}\theta _{12}=1/\sqrt{2}`$ at $`m_Z`$ scale, the mixing angle of $`\mathrm{tan}\widehat{\theta }_{12}`$ is given by
$`\mathrm{tan}2\widehat{\theta }_{12}`$ $``$ $`(1ϵ)\sqrt{1ϵ}{\displaystyle \frac{\sqrt{4\xi _b^2+ϵ^2\mathrm{sin}^22\varphi }}{ϵ(1+\mathrm{cos}2\varphi )}},`$ (47)
where we use the approximation which neglects the higher order corrections of $`ϵ^2`$, $`ϵ\xi _b`$, and $`\xi _b^2`$. When $`\varphi =\pi /2`$, the mixing angle $`\widehat{\theta }_{12}`$ becomes
$$\mathrm{tan}2\widehat{\theta }_{12}=\mathrm{},$$
(48)
which means the maximal mixing is stable against quantum corrections. On the other hand, when $`\varphi =0`$,
$$\mathrm{tan}2\widehat{\theta }_{12}\frac{\xi _b}{ϵ},$$
(49)
which shows that the mixing angle of $`\widehat{\theta }_{12}`$ strongly depends on the magnitude of $`ϵ`$. The large mixing is spoiled when $`\xi _bϵ`$, which corresponds the region of $`\mathrm{tan}\beta 10`$ for the MSW-L solution, and any value of $`\mathrm{tan}\beta `$ for the VO solution when we take $`m_R=10^{13}`$ GeV. In Fig.1, we show the change of $`\mathrm{sin}^22\widehat{\theta }_{12}`$ due to the continuous change of Majorana phase $`\varphi `$ in the case of $`\mathrm{tan}\beta =60`$ and $`m_R=10^{13}`$ GeV. As the Majorana phase $`\varphi `$ changes from $`0`$ to $`\pi /2`$, the value of $`\mathrm{sin}^22\widehat{\theta }_{12}`$ changes from $`0`$ to $`1`$. The large deviation from $`1`$ of $`\mathrm{sin}^22\widehat{\theta }_{12}`$ means that the mixing angle $`\theta _{12}`$ is unstable with respect to the quantum corrections. Figure 1 shows that the mixing angle $`\theta _{12}`$ changes from being unstable to being stable as the change $`\varphi `$ from $`0`$ to $`\pi `$. The lines of the MSW-L and the VO solutions are almost overlapping in Fig.1, since the discrepancy of $`\xi _b`$ ’s for the two solutions is negligible compared with the quantum correction, $`ϵ=0.1`$, when $`\mathrm{tan}\beta =60`$ and $`m_R=10^{13}`$ GeV.
As for the MSW-S solution, $`\varphi =0`$ induces
$$\mathrm{tan}2\widehat{\theta }_{12}\mathrm{tan}2\theta _{12}\left(1+\frac{1}{\mathrm{cos}2\theta _{12}}\frac{ϵ}{\xi _b}\right)^1,$$
(50)
while $`\varphi =\pi /2`$ induces
$$\mathrm{tan}2\widehat{\theta }_{12}\mathrm{tan}2\theta _{12}.$$
(51)
Equations (50) and (51) show that the mixing angle of $`\theta _{12}`$ is not changed in the region of $`\mathrm{tan}\beta 10`$ when $`\varphi =0`$, while it is not changed independently of $`\mathrm{tan}\beta `$ when $`\varphi =\pi /2`$. Above conclusions are the same as those of Ref.. In Fig. 2, we show the value of $`\mathrm{sin}^22\widehat{\theta }_{12}`$ at $`m_R=10^{13}`$ GeV scale in the case of $`\mathrm{tan}\beta =60`$ according to the continuous change of $`\varphi `$ from $`0`$ to $`\pi `$. Figure 2 shows that the mixing angle $`\theta _{12}`$ changes from being unstable to being stable corresponding to the change of $`\varphi `$ from $`0`$ to $`\pi `$.
## 5 Type C $`(m_1m_2m_3)`$
In Type C mass hierarchy, it has been shown in Ref. that the MNS matrix approaches the definite unitary matrix according to the relative sign assignments of the neutrino mass eigenvalues, as the effects of quantum corrections become large enough to neglect the mass-squared differences of neutrinos. Independent parameters of the MNS matrix at the $`m_R`$ scale approach the following fixed values in the large limit of quantum corrections:
case (c1): $`diag.(m_1,m_2,m_3)`$
$$U_{e2}=\frac{\mathrm{sin}\theta _{12}}{\sqrt{1+\mathrm{cos}^2\theta _{12}}},\text{ }U_{e3}=\frac{1}{2}\frac{\mathrm{sin}2\theta _{12}}{\sqrt{1+\mathrm{cos}^2\theta _{12}}},\text{ }U_{\mu 3}=\frac{1}{\sqrt{2}}\frac{\mathrm{sin}^2\theta _{12}}{\sqrt{1+\mathrm{cos}^2\theta _{12}}}.$$
(52)
case (c2): $`diag.(m_1,m_2,m_3)`$
$$U_{e2}=\mathrm{sin}\theta _{12},\text{ }U_{e3}=\frac{1}{2}\frac{\mathrm{sin}2\theta _{12}}{\sqrt{1+\mathrm{sin}^2\theta _{12}}},\text{ }U_{\mu 3}=\frac{1}{\sqrt{2}}\frac{\mathrm{cos}^2\theta _{12}}{\sqrt{1+\mathrm{sin}^2\theta _{12}}}.$$
(53)
case (c3): $`diag.(m_1,m_2,m_3)`$
$$U_{e2}=0,\text{ }U_{e3}=0,\text{ }U_{\mu 3}=\frac{1}{\sqrt{2}}.$$
(54)
case (c4): $`diag.(m_1,m_2,m_3)`$
$$U_{e2}=0,\text{ }U_{e3}=0,\text{ }U_{\mu 3}=0.$$
(55)
We can easily obtain the values of the mixing angles by using relations of ,
$`\mathrm{sin}^22\theta _{12}`$ $`=`$ $`4{\displaystyle \frac{U_{e2}^2}{1|U_{e3}|^2}}\left(1{\displaystyle \frac{U_{e2}^2}{1|U_{e3}|^2}}\right),`$ (56)
$`\mathrm{sin}^22\theta _{13}`$ $`=`$ $`4|U_{e3}|^2\left(1|U_{e3}|^2\right),`$ (57)
$`\mathrm{sin}^22\theta _{23}`$ $`=`$ $`4{\displaystyle \frac{U_{\mu 3}^2}{1|U_{e3}|^2}}\left(1{\displaystyle \frac{U_{\mu 3}^2}{1|U_{e3}|^2}}\right).`$ (58)
As shown above, the cases of (c1)$``$(c4) are connected by Majorana phases of $`\varphi _1`$ and $`\varphi _2`$.
Figure 3 shows that the values of mixing angles at the high energy scale $`m_R=10^{13}`$ GeV for the MSW-L and the VO solutions according to continuous changes of Majorana phases $`\varphi _1`$ and $`\varphi _2`$ in the case of $`\mathrm{tan}\beta =60`$. Under the conditions that the effects of quantum corrections are large enough to neglect the mass-squared differences of neutrinos, the results of the MSW-L solution are the same as those of the VO solution . Table 3 shows the fixed values of the mixing angles for the MSW-L and the VO solutions in the large limit of quantum corrections which are obtain form Eqns.(52) $``$ (55) by using Eqns.(56) $``$ (58).
The deviations from the values at $`m_Z`$ scale, $`\mathrm{sin}^22\theta _{12}=1`$, $`\mathrm{sin}^22\theta _{13}=0`$ and $`\mathrm{sin}^22\theta _{23}=1`$, indicate that mixing angles receive significant quantum corrections. For $`\mathrm{sin}^22\widehat{\theta }_{12}`$, Table 3 shows that the cases of (c1) and (c2) conserve the maximal mixing, while the cases of (c3) and (c4) do not, in the large limit of quantum corrections. From Eq.(15), we can show that the change of $`\varphi _1`$ form $`0`$ to $`\pi /2`$ with the relation $`\left|\varphi _2\varphi _1\right|=0`$ ($`\left|\varphi _2\varphi _1\right|=\pi /2`$) corresponds to the change of (c4) to (c3) ((c2) to (c1)). Figure 3(a) shows that the unstable region of $`\mathrm{sin}^22\widehat{\theta }_{12}^{}{}_{}{}^{<}0.1`$ exists around the line of $`\left|\varphi _2\varphi _1\right|=0`$, and the stable region of $`\mathrm{sin}^22\widehat{\theta }_{12}1.0`$ exists around the line of $`\left|\varphi _2\varphi _1\right|=\pi /2`$. Since the cases of (c1) and (c2) have masses with opposite signs between the first and second generations, the mixing angle is stable from the analogy of two-generation neutrinos. Therefore the maximal mixing between the first and second generations is conserved in the continuous region preserving the relation of $`\left|\varphi _2\varphi _1\right|=\pi /2`$. As for stability of $`\mathrm{sin}^22\widehat{\theta }_{13}`$, Table 3 shows that the cases of (c3) and (c4) conserve the zero mixing, while the cases of (c1) and (c2) do not. Figure 3(b) shows that the stable region exists around the line of $`\left|\varphi _2\varphi _1\right|=0`$ which connects (c3) and (c4), and the unstable region exists around the line of $`\left|\varphi _2\varphi _1\right|=\pi /2`$ which connects (c1) and (c2). For the stability of $`\mathrm{sin}^22\widehat{\theta }_{23}`$, Table 3 shows that the case of (c3) only conserves the maximal mixing, and the case of (c4) induces the zero mixing. Both cases of (c1) and (c2) induce $`\mathrm{sin}^22\widehat{\theta }_{23}0.36`$. These situations are connected continuously by two Majorana phases $`\varphi _1`$ and $`\varphi _2`$ as shown Fig.3(c).
Figures 4 shows the values of mixing angles at high energy scale $`m_R=10^{13}`$ GeV for the continuous change of the Majorana phases for the MSW-S solution in the case of $`\mathrm{tan}\beta =60`$. Table 4 shows the fixed values of the mixing angles for the MSW-S solution in the large limit of quantum corrections, which are obtained from Eqns.(52) $``$ (55) by using Eqns.(56) $``$ (58).
The deviations from the values at $`m_Z`$ scale, $`\mathrm{sin}^22\theta _{12}=7.1\times 10^5`$, $`\mathrm{sin}^22\theta _{13}=0`$ and $`\mathrm{sin}^22\theta _{23}=1`$, indicate that the mixing angles receive significant quantum corrections. For $`\mathrm{sin}^22\widehat{\theta }_{12}`$, Table 4 shows that all the cases of (c1) $``$ (c4) make it zero in the large limit of quantum corrections. Figure 4(a) shows that the unstable region of $`\mathrm{sin}^22\widehat{\theta }_{12}>0.2`$ exists around the points of $`(\varphi _1,\varphi _2)(\pi /2,\pi /30)`$ and $`(\varphi _1,\varphi _2)(\pi /2,29\pi /30)`$. For the stability of $`\mathrm{sin}^22\widehat{\theta }_{13}`$, Table 4 shows that all the cases of (c1) $``$ (c4) conserve the zero mixing. Figure 4(b) shows that $`\mathrm{sin}^22\widehat{\theta }_{13}`$ is stable with respect to the quantum corrections for any values of two Majorana phases. For the stability of $`\mathrm{sin}^22\widehat{\theta }_{23}`$ Table 4 shows that the cases of (c2) and (c3) conserve the maximal mixing, while the cases of (c1) and (c4) do not, in the large limit of quantum corrections. Figure 4(c) shows that the stable region exists around lines of $`\varphi _2=\pi /2`$ which connect (c2) and (c3) by changing $`\varphi _1`$ from $`0`$ to $`\pi /2`$, and the unstable region exists around the lines $`\varphi _2=0`$ and $`\varphi _2=\pi `$ which connect (c1) and (c4) by changing $`\varphi _1`$ from $`0`$ to $`\pi /2`$.
## 6 Summary
Neutrino-oscillation solutions for the atmospheric neutrino anomaly and the solar neutrino deficit can determine the texture of the neutrino mass matrix according to three types of neutrino mass hierarchies as Type A: $`m_1m_2m_3`$, Type B: $`m_1m_2m_3`$ , and Type C: $`m_1m_2m_3`$. We found that the relative sign assignments of neutrino masses in each type of mass hierarchies play the crucial roles for the stability against quantum corrections. Actually, two physical Majorana phases in the lepton flavor mixing matrix connect among the relative sign assignments of neutrino masses. Therefore, in this paper we analyze the stability of mixing angles against quantum corrections according to three types of neutrino mass hierarchies (Type A, B, C) and two Majorana phases. Two phases play the crucial roles for the stability of the mixing angles against the quantum corrections. The results in Ref., where the stabilities of the mixing angles in (a1) and (a2), (b1) and (b2), (c1) $``$ (c4) with respect to quantum corrections are argued, are reproduced by taking the definite values of two Majorana phases.
## Acknowledgment
The work of NO is supported by the JSPS Research Fellowship for Young Scientists, No.2996.
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# Untitled Document
GOLDBACH‘S RULE
by Metin Aktay with thanks to Clifford H. Taubes for considering the preposterous
Content: A. Overview
B. List of Definitions and Variables
C. The Proof of Goldbach‘s Conjecture
D. Evaluation of Observations and Further Thought
Appendix A
A. OVERVIEW
Goldbach‘s Conjecture, ”every even number greater than $`2`$ can be expressed as the sum of two primes” is renamed Goldbach‘s Rule for it can not be otherwise.
The conjecture is proven by showing that the existence of prime pairs adding to any even number greater than $`2`$ is a natural by-product of the existence of the prime sequence less than that even number. First it is shown that the remainder of cancellations process which identifies primes less than an even number also remainders prime pairs adding to that even number as a natural part of the process. Then a minimum limit for the number of remaindered prime pairs adding to an even number is expressed in terms of that even number and shown to exist for every even number greater than $`2`$. Furthermore, the reasonings and formulations used in the proof are demonstrated to hold against observations.
B. LIST OF DEFINITIONS AND VARIABLES
Let $`E`$ be any even number $`>2`$.
Let $`N`$ be any positive integer $`<E.`$
Let $`i,n`$ be counters, each of integers in natural order beginning with $`1`$.
Let a number couple be a couple where the order of two integers matters.
Let a number pair be a pair where the order of two integers does not matter.
Let $`N`$ symmetric integers be two integers having identical absolute difference with $`N`$.
Let $`P(i)`$ be any prime $`\sqrt{(E1)}`$ in natural order where $`P(1)=2`$.
Let $`P(m)`$ be the largest prime $`\sqrt{(E1)}`$.
Let $`P(i)`$-prime be indivisibility by $`P(i)`$ where $`P(i)`$ is not indivisible and $`1`$ is indivisible.
Let $`P(i)`$-composite be divisibility by $`P(i)`$ where $`P(i)`$ is divisible.
Let $`G1`$ be the number $`E/2`$ symmetric $`N`$ which are $`P(i)`$-prime for all $`P(i)`$.
Let $`G2`$ be the number of $`E/2`$ symmetric primes adding to $`E.`$
Let $`GP`$ be the number of prime pairs adding to $`E`$.
Let $`r.f.P(i)`$ be the $`E/2`$ symmetric $`P(i)`$-prime remaindering frequency for $`E`$ of any $`P(i)`$ divisor.
Let $`stsp.m()`$ be the step truncated series product with steps from $`i=1`$ to $`i=m.`$
Let $`PE`$ be the largest prime $`<E`$.
Let $`NPE`$ be the number of primes $`PE`$.
Let $`GR`$ be Goldbach Ratio,$`(GP/NPE).`$
C. THE PROOF OF GOLDBACH‘S CONJECTURE
C.1 The ”Remainder” Nature Of The Prime Sequence
Primeness of an integer is divisibility by no other than unity and itself. The prime sequence is identified by cancelling divisibilities and retaining indivisibilities. The integers not cancelled as divisible constitute the prime sequence. It is crucial to note that this identification process is an indirect process rather than direct, that primes are remainders, not direct creations but remnants after cancellations.
C.2 Method Of Identification Of The Prime Sequence Less Than $`E`$
The Sieve of Eratosthenes identifies the prime sequence up to any integer by cancelling divisibilities by prime divisors less than the square root of that integer and thus ”remainder”ing indivisibilities. This suffices because a prime larger than that square root is multiplied by a prime less than that square root to produce any composite less than that integer. Therefore the sequence of primes less than $`E`$ are those $`N`$ which are not divisible by the divisors $`P(i)`$ dividing with frequency $`P(i)`$, where $`P(i)`$ were defined to be primes $`\sqrt{(E1)}`$.
C.3 Concepts Of $`E/2`$ Symmetricity And Asymmetricity of Primes and Composites
On the integer line of $`1`$ to $`(E1)`$ of $`N`$, the sum of every $`E/2`$ symmetric integer couple is $`E`$. Both members of any such couple may be prime, or both may be composite, or the larger member prime and the smaller composite, or the other way around.
A prime is defined an $`E/2`$ symmetric prime if it‘s $`E/2`$ symmetric $`N`$ is also a prime, and an $`E/2`$ asymmetric prime if it‘s $`E/2`$ symmetric $`N`$ is a composite.
A composite is defined an $`E/2`$ symmetric composite if it‘s $`E/2`$ symmetric $`N`$ is also a composite, and an $`E/2`$ asymmetric composite if it‘s $`E/2`$ symmetric $`N`$ is a prime.
The above may be visualised as matched integers on two integer lines matched head to tail, $`(E1)`$ to $`1`$ and $`1`$ to $`(E1)`$. These matched integers constitute couples adding to $`E`$. If both members of such couples are prime then they are $`E/2`$ symmetric primes, if both are composites then they are $`E/2`$ symmetric composites. Couples with one prime and one composite members contain $`E/2`$ asymmetric primes and $`E/2`$ asymmetric composites.
C.4 Concepts of $`E/2`$ Symmetricity And Asymmetricity Of $`P(i)`$ Divisors
Any $`P(i)`$ divisor dividing $`N`$ divides either symmetrically or asymmetrically with respect to $`E/2`$. $`E/2`$ symmetric $`P(i)`$ divisors divide $`E/2`$ and $`(E/2)+nP(i)`$ and $`(E/2)nP(i)`$. $`E/2`$ asymmetric $`P(i)`$ divisors never divide $`E/2`$ symmetric $`N`$.
$`E/2`$ symmetricity of any $`P(i)`$ divisor depends on divisibility of $`E/2`$ by that $`P(i)`$. If $`E/2`$ is divisible by a $`P(i)`$ then that $`P(i)`$ divides and ”remainder”s symmetrically with respect to $`E/2`$. If $`E/2`$ is not divisible by a $`P(i)`$ then that $`P(i)`$ divides asymmetrically with respect to $`E/2`$. This dependence applies to all $`P(i)`$ except $`P(1)=2`$ which is $`E/2`$ symmetric independent of $`E/2`$ divisibility, with either two divisibilities or two indivisibilities bracketing $`E/2`$.
If all $`P(i)`$ divisors were $`E/2`$ symmetric, then both composite and prime $`N`$ would be $`E/2`$ symmetric, and thus all primes $`<E`$ would be members of prime couples adding to $`E`$, except for $`P(i)`$ themselves which are symmetric with composites divided by themselves. If all $`P(i)`$ divisors were to be $`E/2`$ asymmetric, except for $`P(1)=2`$ which can not be so, then there would be minimal $`E/2`$ symmetric composites and primes.
C.5 Concept Of $`P(i)`$-Primes And The Primes As Intersecting Sets Of $`P(i)`$-primes
If being $`P(i)`$-prime is indivisibility by $`P(i)`$, where $`1`$ is defined $`P(i)`$-prime and $`P(i)`$ itself is defined divisible, there would be $`1`$ $`P(i)`$-composite and $`(P(i)1)`$ $`P(i)`$-primes every $`P(i)`$ consecutive $`N`$ for any $`P(i)`$ divisor. The $`N`$ which are $`P(i)`$-prime for all $`P(i)`$ can not be but prime except for $`1`$. Therefore, the primes $`<E`$ are the $`N`$ which are $`P(i)`$-prime for all $`P(i)`$, with the addition of the $`P(i)`$ themselves the deduction of $`1`$.
C.6 Existence And Frequency Of $`E/2`$ Symmetric $`P(i)`$-primes For Any $`P(i)`$
Any $`E/2`$ symmetric $`P(i)`$ divisor will divide $`1`$ $`E/2`$ symmetric $`P(i)`$-composite and remainder $`(P(i)1)`$ $`E/2`$ symmetric $`P(i)`$-primes every $`P(i)`$ consecutive $`N.`$
Any $`E/2`$ asymmetric $`P(i)`$ divisor will divide $`1`$ $`E/2`$ asymmetric $`P(i)`$-composite, and will remainder $`1`$ $`E/2`$ asymmetric $`P(i)`$-prime and $`(P(i)2)`$ $`E/2`$ symmetric $`P(i)`$-primes for every $`P(i)`$ consecutive $`N`$.
The above may be visualised with the two lines of integers matched head-to-tail. Any $`P(i)`$ divisor cancelling $`P(i)`$-composites and remaindering $`P(i)`$-primes on each line begins dividing from opposite ends, thus may or may not meet at the midpoint since that $`P(i)`$ may or may not be $`E/2`$ symmetric. Let the two lines be counted by $`n`$ counting from one end only, beginning with $`n=1`$, which counts one line forward and the other line backwards. If the cancellations by any $`P(i)`$ divisor on each line is accounted for with respect to $`n`$, then the frequency with respect to $`n`$ of divisibility by that $`P(i)`$ of both lines is $`1/P(i)`$ since both lines move by $`1`$ for every unit change in $`n`$. However, for an $`E/2`$ asymmetric $`P(i)`$ divisor, the divisibility of each line by that $`P(i)`$ may lead or lag the other line with respect to $`n`$. If a $`P(i)`$ divisor is $`E/2`$ symmetric it will divide one matched integer couple and remainder $`((P(i)1)`$ matched integer couples as $`P(i)`$-primes every $`P(i)`$ $`n`$. If a $`P(i)`$ divisor is $`E/2`$ asymmetric it will divide $`1`$ $`P(i)`$-composite on the first line matched with a $`P(i)`$-prime on the second, and $`1`$ $`P(i)`$-composite on the second line matched with a $`P(i)`$-prime on the first, and will remainder $`(P(i)2)`$ couples with $`P(i)`$-primes on each line every $`P(i)`$ $`n`$.
$`P(1)=2`$ divisor can not be $`E/2`$ asymmetric. $`P(1)=2`$ would render all $`2`$-Primes $`E/2`$ asymmetric if it could be $`E/2`$ asymmetric. $`P(1)=2`$ divisor will always remainder $`1`$ $`E/2`$ symmetric $`2`$-Prime for every $`2`$ consecutive $`N`$.
$`E/2`$ asymmetric $`P(2)=3`$ divisor will remainder $`1`$ $`E/2`$ symmetric $`3`$-prime for every $`3`$ consecutive $`N`$ and $`E/2`$ symmetric $`P(2)=3`$ will remainder $`2`$. $`P(i)`$ divisors larger than $`3`$ will remainder more than $`1`$ $`E/2`$ symmetric $`P(i)`$-primes for every $`P(i)`$ consecutive $`N`$, i.e. $`E/2`$ asymmetric $`P(3)=5`$ divisor will remainder $`3`$ $`E/2`$ symmetric $`5`$-primes every $`5`$ consecutive $`N`$ and $`E/2`$ symmetric $`P(3)=5`$ will remainder $`4`$.
Therefore $`E/2`$ symmetric $`P(i)`$-primes exist for each $`P(i)`$ divisor of $`E`$.
C.7 Lower Limit For The Number of $`E/2`$ symmetric $`P(i)`$-Primes for all $`P(i)`$
To find the number of $`E/2`$ symmetric $`P(i)`$-primes for a $`P(i)`$ divisor, the number of $`N`$, which is $`E1`$, is multiplied with the $`E/2`$ symmetric $`P(i)`$-prime remaindering frequency of that $`P(i)`$. Therefore
$`numberofE/2symmetricP(i)primes=(E1)r.f.P(i)`$
This may result in a fractional result which may need truncation or rounding up. Therefore
$`numberofE/2symmetricP(i)primestruncate((E1)r.f.P(i))`$
To find a lower limit for $`G1`$, the number of $`N`$ which are $`E/2`$ symmetric $`P(i)`$-prime for all $`P(i)`$, $`E1`$ is multiplied consecutively by $`r.f.P(i)`$ with truncations at each step. This reduction with consecutive remaindering frequencies is valid because $`P(i)`$ are indivisible by each other.
Therefore
$`G1stsp.m((E1)r.f.P(i))`$
Expectation of $`G1`$ will be minimum where all $`r.f.P(i)`$ are minimum.
Therefore
$`min.G1stsp.m((E1)sp(min.r.f.P(i)))`$.
All $`r.f.P(i)`$ will be minimum where all $`P(i)`$ divisors are $`E/2`$ asymmetric.
Therefore
$`min.r.f.P(i)=((P(i)2)/P(i))`$, except for $`r.f.P(1)=1/2`$.
If any product $`((a)(b/c))`$, where $`a,b,c`$ are positive integers, and where $`ca`$, is truncated, then the result can not be less than $`b`$. With this logic, lower limits can be found for the successive truncations of the step truncated series product by comparing the denominator of each multiplier with the numerator of the multiplied.
Let $`stsp.m((E1)sp(min.r.f.P(i)))`$ begin from $`min.r.f.P(m)`$ and work backwards, and let the first truncation result be $`T(1)`$ and the last truncation result be $`T(m)`$, where $`T(m)min.G1.`$
Therefore
$`T(1)=truncate((E1)min.r.f.(P(m)))`$.
Given that except for $`P(1)=2`$
$`min.r.f.P(i)=((P(i)2)/P(i))`$
Then
$`T(1)=truncate((E1)((P(m)2)/P(m)))`$
Given that
$`(E1)(P(m))^2`$
Then
$`T(1)(P(m)(P(m)2)).`$
Given that
$`T(2)=truncate(T(1)min.r.f.(P(m1)))`$
and
$`min.r.f.P(m1)=((P(m1)2)/P(m1))`$
Then
$`T(2)=truncate(T(1)(P(m1)2)/P(m1))`$
Then
$`T(2)truncate(P(m)(P(m)2)(P(m1)2)/P(m1))`$
Given that
$`P(i)(P(i+1)2)`$ except for $`P(1)=2`$
Then
$`P(m1)(P(m)2)`$
Therefore
$`T(2)(P(m)(P(m1)2))`$ .
If lower limits for $`T(n)`$ are found consecutively as above
Then
$`T(n)(P(m)(P(mn+1)2))`$.
Therefore
$`T(m1)((P(m))(P(2)2))`$.
Since $`P(2)=3`$ then
$`T(m1)((P(m))`$
Given that $`r.f.P(1)=1/2`$, then
$`T(m)truncate((P(m))1/2)`$
Therefore
$`min.G1truncate(P(m)/2)`$.
C.8 Lower Limit For The Number Of Prime Pairs Adding To $`E`$
Since $`1`$ is defined $`P(i)`$-prime for all $`P(i)`$, $`G1`$, which is the number of $`N`$ $`E/2`$ symmetric $`P(i)`$-prime for all $`P(i)`$, may count $`1`$ if the $`E/2`$ symmetric counterpart of $`1`$ is also $`P(i)`$ prime for all $`P(i)`$. Given that $`G2`$ is the number of $`E/2`$ symmetric primes adding to $`E`$, and given that $`G1`$ may count $`1`$ and $`(E1)`$, and given that $`G1`$ excludes $`P(i)`$, and given that $`min.G1`$ is a truncated result, then
$`G2min.G1(2or0)+2(numberofP(i)addingtoEwithotherprimes)`$
Therefore
$`min.G2(min.G12)`$.
Therefore
$`min.G2(truncate(P(m)/2)2)`$
$`GP`$, the number of prime pairs adding to $`E`$, is half of $`G2`$ since every $`E/2`$ symmetric prime is a member of a pair of primes adding to $`E`$. If the halving of $`G2`$ is fractional, then it is rounded up since an odd numbered $`G2`$ indicates that $`E/2`$ itself is prime and is counted once. $`G2`$ is even numbered where $`E/2`$ itself is not prime.
Therefore
$`min.GProundedup(1/2min.G2)`$.
Then
$`min.GProundedup(1/2(truncate(P(m)/2)2))`$.
Therefore
$`min.GP1`$ for $`E50`$ where $`P(m)7`$.
The logic of the above proof also proves that there is at least one prime pair adding to $`E`$ with primes $`>P(m)`$ for every $`E50`$.
Given the existence of at least one prime pair adding to $`E`$ for every $`E<50`$, then
$`min.GP1`$ for every $`E`$.
Therefore
$`GP1`$ for every $`E`$.
Magnitude of $`min.GP\sqrt{E}/4`$ for large $`E`$ since magnitude of $`P(m)\sqrt{E}`$ for large $`E`$.
C.9 Conclusion Of The Proof of Goldbach‘s Conjecture
Thus it is proven that the remaindering of the prime sequence $`<E`$ by divisors $`P(i)\sqrt{(E1)}`$ can not avoid remaindering $`E/2`$ symmetric primes adding to $`E`$ for every $`E`$ even in minimal expectation conditions, since the number of prime pairs adding to $`E`$ is $`1`$ for every $`E`$, and since a hypothetical minimum limit for this number is $`\sqrt{E}/4`$ for large $`E`$.
Therefore there will always be $`E/2`$ symmetric prime pairs adding to $`E`$ for any $`E`$ since the remaindering process for the prime sequence $`<E`$ can not avoid remaindering $`E/2`$ symmetric primes. Therefore there will be at least one pair of $`E/2`$ symmetric primes adding to $`E`$ for any $`E`$.
Therefore every even number greater than $`2`$ can be expressed as the sum of two primes.
D. EVALUATION OF OBSERVATIONS AND FURTHER THOUGHT
D.1 Table Of Observations And Calculations
Behaviour with respect to $`E`$ of the number of prime pairs adding to $`E`$ may be deduced from Appendix A, Observations and Calculations, where the formulations utilised above and actual counts are tabulated for sample $`E`$ evenly and saliently spread up to $`10,000`$.
D.2 Validation Of Assumed Relationships
$`GP`$ may be approximated as $`(E/2)sp(r.f.P(i))`$, which would undercount since it fails to check the $`P(i)`$ excluded by $`r.f.P(i)`$, and would undercount since it omits rounding up where needed, and would overcount since it omits truncation where needed, and would overcount since it may count $`1`$. In Appendix A, $`GP(E/2)sp(r.f.P(i))`$ is calculated as is without the reduction rationalisations used in the proof. It is observed that the proportion of error in this raw term decreases with larger $`E`$. More important than decreasing error is that it is observed to track flawlessly the volatility with respect to $`E`$ of actual $`GP`$.
D.3 Observations and Evaluations Of Goldbach Ratio
Appendix A shows that $`GR`$ varies as expected. $`GR`$ is low where $`E/2`$ is indivisible by smaller $`P(i)`$ and high where $`E/2`$ is divisible by smaller $`P(i)`$, since divisibility of $`E/2`$ by smaller $`P(i)`$ increases $`E/2`$ symmetricity of primes substantially.
$`E`$ with $`E/2`$ prime or divisible only by $`P(1)=2`$ have low $`GR`$. All $`E`$ with $`E/2`$ divisible by $`3`$ have high $`GR`$ because avoidance of an $`E/2`$ asymmetric $`3`$ divisor doubles symmetricity of primes. For example, $`E=210`$ has a high $`GR`$ , $`41\%`$, since $`E/2=105`$ is divisible by $`3,5,7`$, which thus remainder primes symmetrically with respect to $`E/2`$. A $`GR`$ value of $`41\%`$ is considered high since $`GR`$ is defined by utilising pairs of primes adding to $`E`$, which means that the maximum possible $`GR`$ is $`50\%`$ , where every prime is a member of a pair of primes adding to $`E`$.
D.4 Deductions In Relation To The Goldbach Comet
$`E/2`$ symmetricity of $`P(i)`$ divisors, which is divisibility of $`E/2`$ by $`P(i)`$, explains the dense cluster bands which form when $`GP`$ is plotted against $`E`$, a plot called the Goldbach Comet on account of these cluster bands. The densest asymptotic cluster band is formed by $`E`$ with $`E/2`$ asymmetric smaller prime divisors. $`E/2`$ symmetricity of smaller prime divisors explain other bands. The next dense band are the $`E`$ with $`E/2`$ symmetric $`P(2)=3`$ divisor, at $`(2/3)/(1/3)=2`$ times the heights of the asymptotic lowest band. Then the next dense band are the $`E`$ with $`E/2`$ symmetric $`P(3)=5`$ divisor, at $`(4/5)/(3/5)=4/3`$ times the heights of the asymptotic lowest band. Joint $`E/2`$ symmetricity of the smaller primes also form distinct dense bands, such as the $`E`$ with $`E/2`$ symmetric $`P(2)=3`$ and $`P(3)=5`$ divisors, at $`(2/34/5)/(1/33/5)=8/3`$ times the heights of the asymptotic lowest band.
D.5 Further Thought
The number of prime pairs adding to $`E`$ increases with $`E`$ but with high volatility. It could be that this volatility is fluctuation around a fundamental relationship. Assuming the prime sequence to be a discrete wave function, and assuming the estimate $`N/lnN`$ of J.S.Hadamard to be it‘s frequency for $`N`$, this fundamental relationship is likely to be $`GP=E/(2(lnE)^2)`$.
Metin Aktay
Ihsan Aksoy sok. EVA apt. No:7/2, Camlik, Etiler, Besiktas, Istanbul 80630, Turkey
Phone: +90 212 2651016 Fax: +90 212 2577374 Mobile: +90 532 2741771
E-mail: maktay@superonline.com, or maktay@mba1979.hbs.edu
APPENDIX A: OBSERVATIONS AND CALCULATIONS
E PE NPE P(m) E/2 E/2 Factors observed GR% calcul. Error % of $`P(m)`$ GP GP calcul. GP 128 127 31 11 64 2 3 10 4 33 210 199 46 13 105 2,3,5,7 19 41 17 -11 222 211 47 13 111 P 11 23 5 -55 502 499 95 19 251 P 15 16 10 -33 512 509 97 19 256 2 11 11 10 -9 678 677 123 23 339 3 28 23 24 -14 1,006 997 168 31 503 P 18 11 16 -11 1,024 1,021 172 31 512 2 22 13 16 -27 1,510 1,499 239 37 755 5 33 14 30 -9 2,018 2,017 306 43 1,009 P 28 9 27 -4 2,048 2,039 309 43 1,024 2 25 8 27 8 2,490 2,477 367 47 1,245 3,5 94 26 85 -10 3,022 3,019 433 53 1,511 P 42 10 37 -12 3,514 3,511 490 59 1,757 7 51 10 50 -2 4,006 4,003 552 61 2,003 P 52 9 46 -12 4,096 4,093 564 61 2,048 2 53 9 47 -11 4,690 4,679 633 67 2,345 5,7 95 15 83 -13 5,006 5,003 670 67 2,503 P 63 9 56 -11 5,610 5,591 738 73 2,805 2,3,5,11,17 198 27 186 -6 6,002 5,987 783 73 3,001 P 62 8 63 2 6,578 6,577 851 79 3,289 2,11,13,23 89 10 86 -3 7,022 7,019 903 83 3,511 P 72 8 70 -3 7,314 7,309 932 83 3,657 2,3,23,53 172 18 156 -9 8,002 7,993 1,007 89 4,001 P 80 8 78 -3 8,192 8,191 1,028 89 4,096 2 76 7 80 5 8,610 8,609 1,072 89 4,305 2,3,5,7,41 282 26 276 -2 9,014 9,013 1,021 89 4,507 P 96 9 88 -8 9,510 9,497 1,177 97 4,755 3,5 253 21 243 -4 9,998 9,973 1,229 97 4,999 P 99 8 96 -3
Sample calculations for $`GP`$ calculated as $`((E1)/2)sp(r.f.P(i))`$
$`GP`$ for $`E=2490:(2,489/2)(1/22/34/55/79/1111/1315/1717/1921/2327/2929/3135/3739/4141/4345/47)`$
$`GP`$ for $`E=3022:(3021/2)(1/21/33/55/79/1111/1315/1717/1921/2327/2929/3135/3739/4141/4345/4751/53)`$
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# ┴{SU-4240-709} The Geometrical Structure of 2D Bond-Orientational Order
## 1 Introduction
According to the KTNHY theory , two-dimensional melting in the plane is a two stage defect-driven process involving continuous crystalline-to-hexatic and hexatic-to-fluid transitions. The intermediate hexatic phase is characterized by quasi-long-range bond orientational order and positional disorder. In a wide variety of settings one may expect to encounter geometries more general than the plane. In the theory of membranes the geometry itself fluctuates. External forces may act to bend the geometry to some fixed curved surface. This raises the challenging problem of generalizing the established KTNHY theory to substrates with some arbitrary geometry. In this paper we discuss a complete geometrical formulation of bond-orientational order (hexatic and crystalline) that provides a framework in which the interaction of defects, geometry and topology are easily formulated.
To begin, consider a hypothetical flat monolayer displaying bond-orientational order at zero temperature. Bonds may be represented as six vectors at each point of the plane, with a relative angle of $`\frac{\pi }{3}`$. The six vectors at any given point of the plane are parallel to those at any other point. In other words, if we know the vectors at some point we can reconstruct the vectors at any other point by parallelism. Imagine now adiabatically deforming the plane of the monolayer to an arbitrary curved surface. Intuitively we would expect the bond-orientational order to be stable to this deformation. The curvature of the monolayer, however, now implies an associated Gaussian curvature. As well known from differential geometry there is no intrinsic notion of parallelism on curved surfaces. To define parallel transport we require a rule specified via a connection. In general parallel transport of a vector between two points will depend on the path. For our problem this means that the bond angle at an arbitrary point on the surface has no path-independent meaning. For the standard choice of connection (the Levi-Civita connection) the bond angle at some point reached by two distinct paths from a reference point will depend on the total Gaussian curvature enclosed by the paths. This sensitive dependence on the path suggests that bond-orientational order is incompatible with any curved geometry.
This mathematical argument, however, contradicts our intuition that bond-orientational order should be stable to deformation of the surface. This suggests that a more physical connection exists for defining bond-orientational order on curved surfaces. To see this note that bond-orientational order very generally implies the existence of a frame of six bonds forming a $`\frac{\pi }{3}`$ angle at each point of the sample. Given this data we can define a path-independent parallel transport in the following way. Take a vector $`𝐕`$ at point $`P`$. This forms a certain angle with the frame. Now we can transport $`𝐕`$ to the point $`P^{}`$ by demanding that $`𝐕`$ forms the same angle to the given frame at $`P^{}`$. By construction this physical parallel transport does not depend on the path and is therefore unambiguous <sup>1</sup><sup>1</sup>1The fact that this parallel transport is unambiguous can be made more explicitly if, before starting the experiment, we mark all the bonds pointing, say, positively in the $`x`$-direction, which can be done unambiguously in the absence of free disclinations..
Clearly the connection that corresponds to the parallel transport defined above must be different from the standard Levi-Civita connection. We will argue below that the mathematical realization of this physical connection appropriate for defining bond-orientational order involves a new degree of freedom, a geometrical quantity called torsion. The difference angle between two parallel transported vectors will depend on the enclosed curvature of the new connection with torsion. This connection may be tuned to be curvature-free so that parallel transport is path-independent, as desired.
Although torsion may appear at this stage as merely a mathematical trick, a number of authors have identified torsion as necessary for describing a system with a high density of dislocations, such as in the hexatic phase. This is certainly the case if the density of dislocations is well approximated by a continuous density, corresponding to the Debye-Huckel approximation. We will see more generally, however, that torsion accounts for the intrinsic frustration of a 2d configuration on a curved surface.
This curvature of the connection with torsion also has a clear physical interpretation as being proportional to the density of disclinations. Raising the temperature slightly will excite disclinations. Although the zero curvature condition is violated, it is such that parallel transport is ambiguous only up to an angle $`\pi /3`$, and the hexatic order is preserved. As the temperature is raised further the growing density of free disclinations will eventually melt the hexatic to a fluid. The high density of disclinations in this phase may be represented by a continuous disclination density. In this case the parallel transported bond angle is completely ambiguous, which is equivalent to saying that bond angles are meaningless in the isotropic fluid phase.
The organization of the paper is as follows: in the next section we introduce the necessary background in differential geometry, which will also serve to fix our notation. In section 3 we provide the physical interpretation of torsion and derive some general relations involving the geometry, the defect density and the torsion. In section 4 we reformulate the case of a simple flat monolayer in terms of the new formalism: this is easily generalized to an arbitrary geometry. We end the paper with some conclusions. In the appendix we treat the example of a cone displaying hexatic order as an explicit application of our formalism.
## 2 Differential geometry background
In this section we briefly emphasize those concepts in differential geometry that play a major role in the subsequent analysis. For further details we refer the reader to the literature .
### 2.1 Basic differential geometry
A $`D`$-dimensional surface embedded in $`d`$-dimensional Euclidean space is described by $`d`$-dimensional functions $`R^\alpha (𝐱),\alpha =1,\mathrm{},d`$ of $`D`$ coordinates $`x^\mu ,\mu =1,\mathrm{}D`$, with $`D=2`$ being the physically interesting case.
The fact that the surface is embedded in $`𝐑^d`$ provides a natural metric,
$$ds^2=g_{\mu \nu }dx^\mu dx^\nu ,g_{\mu \nu }=_\mu \stackrel{}{R}_\nu \stackrel{}{R}.$$
(1)
Introducing the vielbeins, $`e_{}^{a}{}_{\mu }{}^{}`$, it is possible to rewrite the previous equation as
$$ds^2=\delta _{ab}\vartheta ^a\vartheta ^b,$$
(2)
where $`\vartheta ^a=e_{}^{a}{}_{\mu }{}^{}dx^\mu `$. It is easy to check that the tangent vectors
$$e_a=e_{}^{\mu }{}_{a}{}^{}\frac{}{x^\mu },a=1,\mathrm{},D,$$
(3)
where $`e_{}^{\mu }{}_{a}{}^{}e_{}^{a}{}_{\nu }{}^{}=g_{\mu \nu }`$, define a basis of orthonormal vectors. The dot product of any two vectors $`𝐮=u^ae_a`$ and $`𝐯=v^ae_a`$ in tangent space is $`𝐮𝐯=\delta _{ab}u^av^b`$, as if the metric were flat. This formalism is called the non-coordinate basis. Hereafter we use the first letters of the Latin alphabet when we work in a non-coordinate basis, and Greek indices for the coordinate basis.
The space of all $`r`$-dimensional forms of a manifold $`M`$ $`\mathrm{\Omega }^r(M)`$ plays an important role. The metric allows one to define a natural isomorphism, the Hodge-star operator
$$\mathrm{\Omega }^r(M)\stackrel{}{}\mathrm{\Omega }^{Dr}(M),$$
(4)
defined as the linear map which in a non-coordinate acts as
$$(\vartheta ^{a_1}\vartheta ^{a_2}\mathrm{}\vartheta ^{a_r})=\frac{1}{(Dr)!}ϵ_{}^{a_1\mathrm{}a_r}{}_{a_{r+1}\mathrm{}a_D}{}^{}(\vartheta ^{a_{r+1}}\mathrm{}\vartheta ^{a_D}),$$
(5)
where $`ϵ^{12\mathrm{}D}=1`$. We may define an inner product in $`\mathrm{\Omega }^r(M)`$. If $`\omega =\frac{1}{r!}\omega _{a_1\mathrm{}a_r}\vartheta ^{a_1}\mathrm{}\vartheta ^{a_r}`$ and $`\nu =\frac{1}{r!}\nu _{a_1\mathrm{}a_r}\vartheta ^{a_1}\mathrm{}\vartheta ^{a_r}`$
$$(\omega |\nu )=_M\omega \nu =_M\omega _{a_1\mathrm{}a_r}\nu ^{a_1\mathrm{}a_r}\sqrt{|g|}dx^1\mathrm{}dx^D,$$
(6)
the integral being over the whole manifold $`M`$.
The adjoint $`d^{}`$ of the exterior derivative is defined as the linear operator satisfying
$$(d\alpha |\beta )=(\alpha |d^{}\beta ).$$
(7)
It is an easy exercise to check that $`d^{}=(1)^{Dr+D+1}d`$. The Laplacian is the operator
$$\mathrm{\Delta }=dd^{}+d^{}d.$$
(8)
Let us recall that the Laplacian is independent of the connection.
### 2.2 Connections
A connection on a manifold $`M`$ specifies how tensors are transported along a curve. It allows one to define a covariant derivative $`_\mu `$ on a vector field $`U`$,
$$_\mu U^\nu =_\mu U^\nu +\mathrm{\Gamma }_{\mu \sigma }^\nu U^\sigma ,$$
(9)
where $`\mathrm{\Gamma }_{\mu \sigma }^\nu `$ are the connection coefficients. The parallel transport of a vector $`V^\mu `$ along a curve $`c^\mu `$ is fixed by demanding that is be covariantly constant:
$$_cV^\mu =\frac{dc^\nu (t)}{dt}_\nu V^\mu =0.$$
(10)
It is convenient to define the connection coefficients in a non-coordinate basis:
$$_ae_b=\mathrm{\Gamma }_{}^{c}{}_{ab}{}^{}e_c,$$
(11)
and a connection one-form via
$$\omega _{}^{a}{}_{b}{}^{}=\mathrm{\Gamma }_{}^{a}{}_{cb}{}^{}\vartheta ^c.$$
(12)
We consider only metric connections that preserve the norm of a vector under parallel transport, a property called metric compatibility. This property constrains the connection one-form to satisfy
$$\omega _{}^{a}{}_{b}{}^{}=\omega _{}^{b}{}_{a}{}^{},$$
(13)
which in $`D=2`$ further simplifies to $`\omega _{}^{a}{}_{b}{}^{}=ϵ_{}^{a}{}_{b}{}^{}\mathrm{\Omega }`$, $`\mathrm{\Omega }=\mathrm{\Omega }_a\vartheta ^a`$. Connection coefficients do not transform as a tensor but it is possible to construct two tensorial objects out of them, the torsion and the curvature. For the sake of completeness, we write them acting on vector fields $`X,Y`$,
$`T(X,Y)`$ $`=`$ $`_XY_YX[X,Y]`$
$`R(X,Y)`$ $`=`$ $`_X_Y_Y_X_{[X,Y]}`$ (14)
Expressions turn out simpler in a non-coordinate basis. Defining torsion and curvature 2-forms by
$`T^a`$ $`=`$ $`{\displaystyle \frac{1}{2}}T_{}^{a}{}_{bc}{}^{}\vartheta ^b\vartheta ^c`$
$`R_{}^{a}{}_{b}{}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}R_{}^{a}{}_{bcd}{}^{}\vartheta ^c\vartheta ^d,`$ (15)
they are related to the connection one-form by Cartan’s structure equations, which read for the case $`D=2`$,
$`T^a`$ $`=`$ $`d\vartheta ^a\mathrm{\Omega }\vartheta ^bϵ_{}^{a}{}_{b}{}^{}`$ (16)
$`{\displaystyle \frac{1}{2}}ϵ_{a}^{}{}_{}{}^{b}R_{}^{a}{}_{b}{}^{}`$ $`=`$ $`d\mathrm{\Omega }.`$
Recall that the curvature 2-form is locally an exact form. The expressions in the coordinate basis follow easily from those in the non-coordinate basis from the tensorial nature of curvature and torsion, namely $`T_{}^{\mu }{}_{\nu \rho }{}^{}=e_{}^{\mu }{}_{a}{}^{}e_{}^{b}{}_{\nu }{}^{}e_{}^{c}{}_{\rho }{}^{}T_{}^{a}{}_{bc}{}^{}`$ and $`R_{}^{\mu }{}_{\nu \rho \sigma }{}^{}=e_{}^{\mu }{}_{a}{}^{}e_{}^{b}{}_{\nu }{}^{}e_{}^{c}{}_{\rho }{}^{}e_{}^{d}{}_{\sigma }{}^{}R_{}^{a}{}_{bcd}{}^{}`$.
### 2.3 Some global considerations
The global structure of the manifold places important constraints on the differential geometry of the surface. From Eq. 16 the curvature 2-form is locally an exact form, but that is not the case globally. For a manifold without boundaries, integrating over the whole manifold we have
$$_M\frac{1}{2}ϵ_{a}^{}{}_{}{}^{b}R_{}^{a}{}_{b}{}^{}=2\pi \underset{\alpha }{}Ind_\alpha =2\pi \chi ,$$
(17)
where $`Ind`$ are the indices of the vector field associated with $`\mathrm{\Omega }`$ and $`\chi `$ is the Euler characteristic of the manifold. If $`R`$ is the Riemannian curvature this is the usual Gauss-Bonnet theorem.
## 3 The geometry of bond-orientational order
We now consider a two-dimensional sample exhibiting bond-orientational order and forming an arbitrary shape in d-dimensional space, specified by the embedding $`R^\alpha (𝐱)`$. Within the sample, distances are measured according to the metric Eq. 1.
The main feature of connections with non-vanishing torsion (non-symmetric connections) is that infinitesimal reference vectors fail to close when parallel transported along each other, as illustrated in Fig. 2. If a crystal is forced to lie on an arbitrary geometry, the atoms in the crystal cannot be all equally separated by a distance $`a`$. This implies that parallel transport along the bonds joining the nearest-neighbors will fail to close, as shown in Fig. 2. Thus we see that torsion is related to the intrinsic frustration experienced by a 2D crystal on an arbitrary surface.
In the hexatic phase there is a similar interpretation of torsion, except that there will be a large number of dislocations. This is specified by a Burgers vector density
$$𝐛(𝐱)=\frac{1}{\sqrt{g}(𝐱)}\underset{i}{}𝐛_i^L\delta (𝐱_i𝐱)$$
(18)
where $`𝐛_i^L`$ is the microscopic Burgers vectors with origin at point $`𝐱_i`$. In a highly dislocated medium, the discreteness of the Burgers vector density $`𝐛(x)`$ may be approximated by a continuous vector density, the so-called Debye-Huckel approximation.
The density $`𝐛(x)`$ is, in fact, directly related to the torsion invariant appearing in Eq. 2.2, as noted by several authors . The torsion two-form is
$$T^a=v^a\vartheta ^1\vartheta ^2=v^a\mathrm{\Omega }_M,$$
(19)
where $`\mathrm{\Omega }_M=\sqrt{|g|(𝐱)}dx^1dx^2`$ is the volume form. Since $`v^a`$ is a vector it can be written in a coordinate basis $`v^\mu =e_{}^{\mu }{}_{a}{}^{}v^a`$. Its geometric interpretation is illustrated in Fig. 2. Let us assume non-zero torsion at point $`A`$. Consider the vector $`\overline{AC}`$ and parallel transport it along vector $`\overline{AB}`$, the resulting vector is $`\overline{BD^{}}`$. If we parallel transport vector $`\overline{AB}`$ along $`\overline{AC}`$ we get $`\overline{CD^{\prime \prime }}`$. Using the parallel transport Eq. 10 and Eq. 9, we may write to lowest order
$$\overline{D^{\prime \prime }D^{}}^\mu =\left\{\mathrm{\Gamma }_{\sigma \nu }^\mu \mathrm{\Gamma }_{\nu \sigma }^\mu \right\}\overline{AB}^\sigma \overline{AC}^\nu =v^\mu \times \text{Area of the region (}ABD^{}D^{\prime \prime }C).$$
(20)
If $`𝐯=0`$ then $`D^{}=D^{\prime \prime }=D`$ and the parallelogram closes.
It is clear from the previous geometrical argument that torsion is related to the continuous density of Burgers vector by
$$𝐛(𝐱)=𝐯(𝐱),$$
(21)
or $`T^a=b^a\mathrm{\Omega }_M`$, in a non-coordinate basis. A highly dislocated media physically is represented geometrically by a surface endowed with a connection with torsion .
Clearly then torsion is a necessary ingredient in a geometrical description of bond-orientational order, as it allows a connection form from which the bond angle is constructed by parallel transport. This interpretation of torsion is completely general. In the Debye-Huckel limit we have the additional identification of the torsion with the dislocation density.
In the absence of free disclinations we pointed out that this result should be independent of the path chosen. This in turn translates into the mathematical requirement of vanishing curvature 2-form
$$R_{}^{a}{}_{b}{}^{}=0.$$
(22)
This is a local condition, not always possible to fulfill globally, as for example in the case of spherical topology. Anyway, Eq. 22 is fulfilled everywhere for the case of a sample having the topology of a plane.
In this regime, the Cartan structure equations Eq. 16 read
$`T^a`$ $`=`$ $`d\vartheta ^a\mathrm{\Omega }\vartheta ^bϵ_{}^{a}{}_{b}{}^{},`$ (23)
$`0`$ $`=`$ $`d\mathrm{\Omega },`$
The last equation implies $`\mathrm{\Omega }=d\theta `$, where $`\theta (𝐱)`$ is a 0-form, a function. Let us interpret its physical meaning. The equation of parallel transport Eq. 10 reads
$$\frac{dV^a}{dt}\mathrm{\Omega }_\mu \frac{dc^\mu }{dt}ϵ_{}^{a}{}_{b}{}^{}V^b=0,$$
(24)
where $`c^\mu (t)`$ parametrizes a curve in the surface joining the points $`P`$ and $`P^{}`$.
$$V_\pm (P^{})=V_\pm (P)e^{i_c\mathrm{\Omega }_\mu 𝑑x^\mu },$$
(25)
In the hexatic phase with no free disclinations, Eq. 23 implies
$$V_\pm (P^{})=V_\pm (P)e^{i\theta (P^{})},$$
(26)
where $`V_\pm =V^1\pm iV^2`$ and $`\theta (P)=0`$. This parallel transport id manifestly path independent and $`\theta `$ clearly is the bond angle at point $`P^{}`$. Dragging $`P^{}`$ along the entire surface and performing the parallel transport Eq. 26 we can unambiguously construct the bond angle at any point from the knowledge of the bond angle at $`P`$. We have succeeded in giving a precise mathematical formulation of the heuristic zero temperature experiment performed in the introduction.
### 3.1 Introduction of free disclinations
The sample may also contain some free disclinations, either as a result of thermal fluctuations or because topological constraints force them to appear. In either case, the parallel transport can only be defined modulo an angle of $`\frac{\pi }{3}`$. The system responds to the raising of the temperature by generating curvature, but only via defects that preserve overall hexatic order, the disclinations. In this way, the curvature is directly related to the density of disclinations .
From these arguments, it is clear that a finite density of $`N`$ free disclinations is represented as
$`{\displaystyle \frac{1}{2}}ϵ_{a}^{}{}_{}{}^{b}R_{}^{a}{}_{b}{}^{}`$ $`=`$ $`s(𝐱)\mathrm{\Omega }_M`$
$`s(𝐱)`$ $`=`$ $`{\displaystyle \frac{\pi }{3}}{\displaystyle \frac{1}{\sqrt{|g|}}}{\displaystyle \underset{j=1}{\overset{N}{}}}q_j\delta (𝐱𝐱_j),`$ (27)
where $`q_i`$ may be $`+1`$ (a five-fold disclination) or $`1`$ (a seven-fold disclination). Higher integer charges are also possible although rarely need to be considered. If the surface is closed, we can integrate the previous equation over the whole manifold. With the aid of the Gauss-Bonnet theorem Eq. 17 we find
$$\underset{j=1}{\overset{N}{}}q_j=6\chi .$$
(28)
This shows that, apart from the torus with $`\chi =0`$, closed manifolds always contain a certain number of free disclinations.
With free disclinations, the Cartan structure equations (Eq. 16) become
$`T^a`$ $`=`$ $`d\vartheta ^a\mathrm{\Omega }\vartheta ^bϵ_{}^{a}{}_{b}{}^{},`$
$`d\mathrm{\Omega }`$ $`=`$ $`s(𝐱)\mathrm{\Omega }_M.`$ (29)
Away from the actual location of the disclinations $`\mathrm{\Omega }=d\theta `$, but now it is generally impossible for $`\theta `$ to be defined as a continuous function. It is easy to check that a vectors parallel transported from point $`P`$ to $`P^{}`$ following two different paths $`𝐜(t)`$ and $`𝐟(t)`$ paths (see Fig. 1) differ by an angle $`\psi `$
$$\psi =\frac{\pi }{3}\underset{j}{}q_j,$$
(30)
where $`j`$ runs over all charges within the region enclosed by the two paths. The ambiguity is a multiple of $`\frac{\pi }{3}`$ and as expected hexatic order is preserved. Mathematically we have a connected manifold with a $`Z_6`$ holonomy. Disclinations may be regarded as orbifolds.
Generally, we will write
$$\mathrm{\Omega }=d\theta +\mathrm{\Omega }_{sing},$$
(31)
explicitly displaying the bond angle $`\theta `$ and a a singular part, obviously corresponding to a vortex contribution. This is nothing but the procedure of separating a regular part and a singular part in the standard treatment of the XY model .
There is more information encoded in the Cartan structure equations (Eq. 16). First of all, the Levi-Civita connection is obtained imposing the vanishing of torsion,
$$d\vartheta ^a\mathrm{\Omega }^L\vartheta ^bϵ_{}^{a}{}_{b}{}^{}=0,$$
(32)
where $`\mathrm{\Omega }^L`$ is the connection form of the Levi-Civita connection. We have the important relation
$$d\mathrm{\Omega }^L=\frac{1}{2}ϵ_{}^{a}{}_{b}{}^{}R_{a}^{G}{}_{}{}^{b}=K(𝐱)\mathrm{\Omega }_M,$$
(33)
where $`R_{a}^{G}{}_{}{}^{b}`$ is the standard Riemann tensor two-form and $`K(𝐱)`$ is the Gaussian curvature of the surface. Using this relation in Eq. 3.1
$$\mathrm{\Omega }^L\vartheta ^bϵ_{}^{a}{}_{b}{}^{}\mathrm{\Omega }\vartheta ^bϵ_{}^{a}{}_{b}{}^{}=T^a.$$
(34)
Introducing the Burgers form $`b=b_a(𝐱)\vartheta ^a=b_\mu (𝐱)dx^\mu `$, together with Eq. 21, this may be rewritten as
$$b=\mathrm{\Omega }+\mathrm{\Omega }^L.$$
(35)
This is a very important equation. It relates the bond angle, computed from the actual connection $`\mathrm{\Omega }`$ as indicated in Eq. 25, to the distribution of dislocations and to the Gaussian curvature of Eq. 33. We can make this statement more apparent by applying the exterior derivative to Eq. 35, and using Eqs. 3.1 and 33:
$`K(𝐱)`$ $`=`$ $`s(𝐱)+{\displaystyle \frac{1}{\sqrt{|g|}}}ϵ^{\mu \nu }_\mu b_\nu (𝐱)`$ (36)
$`{\displaystyle \frac{3}{\pi }}K(𝐱)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{|g|}}}{\displaystyle \underset{j=1}{\overset{N}{}}}q_j\delta (𝐱𝐱_j)+{\displaystyle \frac{3}{\pi }}{\displaystyle \frac{1}{\sqrt{|g|}}}ϵ^{\mu \nu }_\mu b_\nu (𝐱).`$ (37)
Note that the coordination number of an atom $`n_i`$, the number of nearest-neighbors, is given by
$$6n_i=q_i=_\mathrm{\Sigma }𝑑u\sqrt{g}\frac{1}{\sqrt{|g|}}\underset{j=1}{\overset{N}{}}q_j\delta (𝐱𝐱_j)=\frac{3}{\pi }_\mathrm{\Sigma }𝑑u\sqrt{g}K\frac{3}{\pi }_\mathrm{\Sigma }b_\mu 𝑑x^\mu ,$$
(38)
where $`\mathrm{\Sigma }_i`$ is a small region around the atom. This result contains an additional surface term when compared with the existing literature . This surface term will only vanish if defects successfully fully screen out the Gaussian curvature. As we will see in Eq. 59, this is the condition that minimizes the energy of the model. So, at very low temperatures, the surface term will be small.
Comparing Eq. 37 with , there is a somewhat disturbing factor $`\frac{3}{\pi }`$ in the last term. In a fine grained description of the model, one regards a dislocation as a composite object, a tightly bound pair of opposite disclinations (see Fig. 3). Defining the dipole moment $`𝐩`$ as
$$p^\mu =\frac{3}{\pi }ϵ_{}^{\mu }{}_{\nu }{}^{}b^\nu ,$$
(39)
where $`b^\nu `$ is the microscopic Burgers vector, we can Taylor expand the density of two free disclinations of opposite charge as
$$\frac{\pi }{3}\left\{\delta (𝐱+𝐩)\delta (𝐱)\right\}=ϵ^{\mu \nu }b_\mu _\nu \delta (𝐱)+𝒪(a^2),$$
(40)
where we neglect higher order terms in the lattice spacing $`a`$. The right hand side of this last equation corresponds exactly to an isolated dislocation. The factor $`\frac{3}{\pi }`$ is the slight displacement of the center of the disclinations necessary to render the right dislocation density Eq. 40. As a byproduct, we derive the result that the dipole moment $`𝐩`$ is perpendicular to the microscopic Burgers vector, a result which is also apparent from Fig. 3.
For our purposes, however, it will be more convenient to use Eq. 35 in combination with Eq. 31, which yields
$$b=d\theta \mathrm{\Omega }_{sing}+\mathrm{\Omega }^L.$$
(41)
This completes our geometrical development of the structure of bond-orientational order.
## 4 The properties of the hexatic phase
The previous section was somewhat formal. We derived some general results valid in the presence of bond-orientational order. Further progress requires making additional physical assumptions about the interactions of dislocations and disclinations with the underlying geometry and topology.
For the simplest case of a flat monolayer a free energy was proposed in (see for a review). We now examine this proposal within the framework of the geometric formalism developed here. We use a completely covariant language with the aid of the mathematical concepts sketched in Sec. 2. Although this is not necessary for the case of flat space, it is directly applicable to general geometries.
### 4.1 Debye-Huckel approximation on a flat monolayer
On a flat monolayer, the metric of Eq. 1 is the standard Euclidean metric
$$ds^2=dx^2+dy^2.$$
(42)
In this case $`\vartheta ^1=dx`$ and $`\vartheta ^2=dy`$.
Let us first treat the simpler case of vanishing free disclination density. The bond angle as a function of the Burgers vector density is found by inverting Eq. 41:
$$\theta (𝐱)=_\mu (\frac{1}{\mathrm{\Delta }}b_\mu )=\frac{1}{2\pi }𝑑𝐱^{}\frac{𝐛(𝐱^{})(𝐱𝐱^{})}{|𝐱𝐱^{}|^2},$$
(43)
where $`\frac{1}{\mathrm{\Delta }}=1/(2\pi )\mathrm{log}(|xx^{}|)`$. This result is in agreement with the results in . We emphasize here that this this result follows from the Cartan structure equations alone, with no further physical assumptions. It provides a non-trivial consistency check for our geometrical interpretation of the hexatic phase.
Additional constraints follow from Eq. 36. With no free disclinations it implies
$$db=0ϵ^{\mu \nu }_\mu b_\nu =0.$$
(44)
Geometry thus constrains the Burgers vector distribution $`𝐛(𝐱)`$ to be irrotational.
The free energy of the crystal, within a Debye-Huckel approximation, was derived in :
$$/T=\frac{K_B}{2}(db|\frac{1}{\mathrm{\Delta }^2}db)+\frac{E_ca^2}{T}(b|b),$$
(45)
where the inner product between forms was defined in Eq. 6 and the Laplacian in Eq. 8. $`K_B`$ is the 2-dimensional Young modulus, $`a`$ is the original lattice spacing and $`E_c`$ the core energy of a single dislocation. The RG recursion relations of the 2d KTNHY melting transition imply that in the hexatic phase
$$\mathrm{lim}_{𝐩0}K_B(𝐩)=0.$$
(46)
If now we consider the bond angle as the fundamental variable, Eq. 41 and Eq. 44 imply a simple $`XY`$ model for the energy of the bond angle variable
$$/T=\frac{K_A}{2}_\mu \theta ^\mu \theta ,$$
(47)
with hexatic stiffness $`K_A=\frac{2E_ca^2}{T}`$ .
The irrotationality of $`𝐛`$ is a general constraint to be satisfied in any hexatic regime with vanishing free disclination density. It has additional consequences. We already noted that a dislocation may be regarded as a tightly bound disclination pair. The dipole vector $`𝐩`$, defined in Eq. 39, is perpendicular to the Burgers vector $`𝐛`$. If we consider a closed loop $`𝒞`$, we can compute the circulation for $`𝐩`$
$$_𝒞𝐩𝑑𝐬=_𝒞ϵ^{\mu \nu }b_\nu 𝑑x_\mu =_\mathrm{\Sigma }𝑑x𝑑y\mathrm{\Delta }\theta ,$$
(48)
where $`\mathrm{\Sigma }`$ is the two dimensional region enclosed by the closed curve $`𝒞`$. The resultant integral is generally not zero. This means that dipoles form grain boundaries of consecutive fives and sevens. If there are no free disclinations these strings cannot self-intersect.
In the presence of free disclinations, Eq. 44 gets modified to
$$db=sdxdyϵ^{\mu \nu }_\mu b_\nu (𝐱)=s(𝐱).$$
(49)
The free energy for the bond angle variable now reads
$$/T=\frac{K_B}{2}(s|\frac{1}{\mathrm{\Delta }}s)+\frac{K_A}{2}(d\theta +\mathrm{\Omega }_{sing}|d\theta +\mathrm{\Omega }_{sing}).$$
(50)
The first term leads to a strong binding of free disclinations, since it corresponds to a $`|𝐱|^2\mathrm{ln}|𝐱|`$ disclination-disclination interaction . Recalling Eq. 46, however, we see that this term is absent. The long-wavelength description of the hexatic phase is given by the standard XY model .
The hexatic phase can be characterized by the correlator
$$C(𝐫)=\sigma (𝐫)\sigma (\mathrm{𝟎}),$$
(51)
with
$$\sigma (𝐫)=\underset{l}{}\mathrm{Ind}_l\delta ^2(𝐫𝐑_l)=\frac{1}{2\pi }s(𝐱),$$
(52)
where $`\mathrm{Ind}_l`$ is the index of the order parameter. The last identity follows from the definition of the index of a vector field. The order parameter Eq. 51 is thus the correlation function of the free disclination density, the hexatic curvature. To lowest order in the fugacity, we have in the hexatic phase
$$s(𝐱)s(\mathrm{𝟎})\frac{y^2}{|𝐱|^\eta },$$
(53)
where $`y`$ is the fugacity of the vortices. In the isotropic liquid phase
$$s(𝐱)s(\mathrm{𝟎})e^{\frac{|𝐱|}{\xi }},$$
(54)
where $`\xi e^{\frac{a}{TT_c}}`$ is the correlation length of the model.
### 4.2 Hexatic order in an arbitrary geometry
In the previous section we recast standard results for a flat monolayer in the geometric framework developed in this paper. These results may now be extended directly to the case of an arbitrary geometry. The metric
$$ds^2=g_{\mu \nu }dx^\mu dx^\nu ,$$
(55)
takes a general form with metric coefficients given by Eq. 1.
The first step is to find the free energy. This is easily obtained from Eq. 45 using the covariant definition of the operators involved, as described in sect. 2 We start by defining the form
$$\rho (𝐱)=\left\{K(𝐱)s(𝐱)\right\}\mathrm{\Omega }_M,$$
(56)
expressing the difference between Gaussian curvature and free disclination density. In Eq. 45 we deliberately wrote the free energy of a simple flat monolayer in terms of the torsion degrees of freedom. It is now extremely simple to generalize it to the case of a general geometry by making use the general relation Eq. 41 that holds for an arbitrary geometry to obtain
$$/T=\frac{K_B}{8}(\rho |\frac{1}{\mathrm{\Delta }^2}\rho )+\frac{K_A}{2}(d\theta +\mathrm{\Omega }_{sing}\mathrm{\Omega }^L|d\theta +\mathrm{\Omega }_{sing}\mathrm{\Omega }^L),$$
(57)
where $`\mathrm{\Omega }^L`$ is the connection form of the Levi-Civita connection. This corresponds to the free energy of a 2D crystal on an arbitrary geometry. This result is identical to that obtained by integrating out the in-plane phonons in the absence of free disclinations, and generalizing to include these additional degrees of freedom. The free energy Eq. 57 corresponds interacting dislocations and disclinations in the crystalline phase. At zero temperature, one can substitute the bond orientational order parameter $`\theta `$ by its minimum value $`\theta =\frac{1}{\mathrm{\Delta }}_\mu \mathrm{\Omega }^\mu `$, and we get
$$/T=\frac{K_B}{8}(\rho |\frac{1}{\mathrm{\Delta }^2}\rho )+\frac{K_A}{2}(\rho |\frac{1}{\mathrm{\Delta }}\rho ).$$
(58)
Free positive disclinations are attracted to positive curvature regions and negative disclinations to negative curvature regions. The ground state is defined by the equation
$$\rho =0s(𝐱)=K(𝐱).$$
(59)
Defects will arrange themselves to optimally screen out the Gaussian curvature.
As for the case of a flat monolayer (Eq. 46), the hexatic phase requires that in the infrared limit
$$K_B=0.$$
(60)
The free energy then becomes
$`/T`$ $`=`$ $`{\displaystyle \frac{K_A}{2}}(d\theta +\mathrm{\Omega }_{sing}\mathrm{\Omega }^L|d\theta +\mathrm{\Omega }_{sing}\mathrm{\Omega }^L)`$ (61)
$`=`$ $`{\displaystyle \frac{K_A}{2}}{\displaystyle 𝑑𝐱\sqrt{g}g^{\mu \nu }(_\mu \theta +\mathrm{\Omega }_{sing}\mathrm{\Omega }_\mu ^L)(_\nu \theta +\mathrm{\Omega }_{sing}\mathrm{\Omega }_\nu ^L)},`$
which is the hexatic free energy first considered in .
It is beyond the scope of this paper to investigate the precise mechanism by which Eq. 60 is implemented in an arbitrary geometry, but we show instead that it is a necessary condition for the hexatic phase to be realized.
In analogy with the flat monolayer, one may compute the order parameter Eq. 51 to lowest order in the fugacity,
$`s(𝐱)s(\mathrm{𝟎})`$ $`=`$ $`2\left({\displaystyle \frac{\pi }{3}}\right)^2y^2e^{\frac{K_A}{2}_{𝐱_i𝐱_j}\left(\frac{1}{\mathrm{\Delta }}\right)_{𝐱_𝐢𝐱_𝐣}}e^{K_A_{𝐢=1}^2{\scriptscriptstyle d^2𝐱^{}{\scriptscriptstyle \frac{q_i}{\mathrm{\Delta }}}_{\mathrm{𝐱𝐱}^{}}\sqrt{g}K(𝐱)}},`$
which reduces to Eq. 53 in the case of a flat monolayer.
It is worthwhile to recall at this point that if fluctuations in the geometry were allowed, we should include an additional bending rigidity or extrinsic curvature term
$$\frac{_{ex}}{T}=\frac{\kappa }{2}𝑑𝐱\sqrt{g}\stackrel{}{H}^2,$$
(63)
where $`\stackrel{}{H}`$ is the Mean curvature of the surface. In this case we average over all possible fluctuations in geometry. It seems that Dislocations have a finite energy and therefore they proliferate at any temperature, see also .
## 5 Conclusions and Outlook
A full understanding of defects in curved geometries is a challenging subject with many novel features not encountered in the plane. In this paper, we presented a geometrical interpretation of bond-orientational order leading to a variety of relations between the underlying geometry, the topology and the nature of the defects. The formalism developed in this paper has already been applied very successfully to the case of defect arrays in spherical crystals . Further applications are currently being explored. The most challenging open problem may be the generalization of the KTNHY renormalization group flows to fixed or dynamical curved geometries. The spherical crystal itself gives rise to a very rich structure of defects in the ground state and other geometries are very likely to lead to even more novel results.
Acknowledgements The research of M.B. and A.T. was supported by the U.S. Department of Energy under contract No. DE-FG02-85ER40237.
## Appendix A The hexatic cone
A sample in the hexatic phase on a template having the geometry of a cone is the simplest case where the present formalism may be applied. It provides an interesting benchmark as the cone has been the subject of much interest . For some range of parameters, an isolated disclination forces the flat monolayer to buckle; a positive disclination is optimally screened by the Gaussian curvature located at the cone. It is more convenient to work in polar coordinates $`(r,\psi )`$, the metric being
$$ds^2=(1+m^2)dr^2+r^2d\psi ^2,$$
(64)
where $`m`$ is related to the aperture angle of the cone $`\zeta `$ by $`\mathrm{tan}\zeta =\frac{1}{m}`$.
One readily finds $`\vartheta ^1=\sqrt{1+m^2}dr`$, $`\vartheta ^2=rd\psi `$. The Levi-Civita connection is $`\mathrm{\Omega }^L=\frac{1}{\sqrt{1+m^2}}d\psi `$. The actual hexatic connection is flat, corresponding to a plane
$$\mathrm{\Omega }=d\psi +d\varphi \theta (r,\psi )=\psi +\varphi (r,\psi ),$$
(65)
where $`\varphi `$ is regular $`\varphi (r,\psi +2\pi )=\varphi (r,\psi )`$ and corresponds to the different Burgers vectors distributions one may have in a cone. The lowest energy solution is $`\varphi =0`$, and this is the only case we consider. The torsion vector distribution may be computed easily,
$$𝐛(𝐱)=(1\frac{1}{\sqrt{1+m^2}})\frac{1}{r}e_\psi ,$$
(66)
Torsion vectors are tangent to circles centered at the tip of the cone, as depicted in Fig. 4. It is also illuminating to plot the order parameter, as illustrated in Fig. 4. The parallel transported bond angle is completely insensitive to the Gaussian curvature located at the tip of the cone. As a further cross check, we can compute the hexatic energy for this configuration, using the free energy Eq. 61. Assuming a lattice spacing $`a`$ as an ultraviolet cut-off and the radius $`R`$ of the cone as an infrared cut-off, we get
$`E_{hex}`$ $`=`$ $`{\displaystyle \frac{K_A}{2}}{\displaystyle _a^R}{\displaystyle \frac{dr}{r}}(1+m^2)^{\frac{1}{2}}{\displaystyle _0^{2\pi }}𝑑\psi (1+{\displaystyle \frac{1}{\sqrt{1+m^2}}})^2`$ (67)
$`=`$ $`\pi K_A(1{\displaystyle \frac{1}{\sqrt{1+m^2}}})^2(1+m^2)^{1/2}\mathrm{ln}{\displaystyle \frac{R}{a}}.`$
This result has already been obtained by other authors by solving the equations of motion for the order parameter and computing the Green’s function. Instead we have obtained it by constructing the order parameter by parallel transport. As a byproduct, we are able to compute the torsion vector distribution in the cone.
If there is an isolated disclination located at the tip of the cone, the hexatic connection is no longer flat, as there is hexatic curvature located at the tip of the cone. Parallel transport is now ambiguous if we encircle the tip of the cone. We get,
$$\mathrm{\Omega }=(\frac{q_i}{6}1)d\psi \theta (r,\psi )=(\frac{q_i}{6}1)\psi .$$
(68)
Following the same steps as before we obtain a free energy of the form given in Eq. 67.
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# Mechanism of confinement in low-dimensional organic conductors
## 1 Introduction
In low-dimensional organic conductors, repulsive interactions play an important role for electronic states with a gap or a pseudo gap. The anisotropy in electric conductivity is enhanced by interactions since the induced pseudo gap around the Fermi surface of a single chain precludes electrons from hopping between chains \[?\]. There are several arguments as to whether or not the electrons are confined to a chain by the repulsive interaction. Away from half-filling, the confinement needs a large magnitude of the interaction even for the small limit of the interchain hopping \[?\] since the effect of the interchain hopping is much larger than that of the intrachain interaction. However, in the case of half-filling, the electrons can be confined by the interaction with a moderate strength due to umklapp scattering which induces the charge gap \[?, ?, ?\].
Bechgaard salts of organic conductors, TMTSF and TMTTF, can be regarded as effectively half-filling due to dimerization \[?, ?\]. The optical experiments have shown the finite Drude weight for the TMTSF salts but not for the TMTTF salts although the correlation gap exists in both salts \[?\]. This indicates the transition from an insulating state with the electrons confined to chains to a metallic state with deconfined electrons when the correlation gap becomes larger than the interchain hopping \[?, ?\]. In the present study, such a transition is elucidated by applying the renormalization group (RG) method to a model of quarter-filled two-coupled chains with dimerization.
## 2 Formulation
We consider quarter-filled two-coupled chains given by
$``$ $`=`$ $`{\displaystyle \underset{j,\sigma ,l}{}}\left[t+(1)^jt_\mathrm{d}\right]\left(c_{j\sigma l}^{}c_{j+1\sigma l}+\text{h.c.}\right)`$ (1)
$`2t_{}{\displaystyle \underset{j,\sigma }{}}\left(c_{j\sigma 1}^{}c_{j\sigma 2}+\text{h.c.}\right)+U{\displaystyle \underset{j,l}{}}n_{jl}n_{jl},`$
where $`t`$, $`t_{}`$, $`t_\mathrm{d}`$ and $`U`$ denote energies for the intrachain hopping, the interchain hopping, dimerization and on-site repulsion, respectively. $`n_{j\sigma l}=c_{j\sigma l}^{}c_{j\sigma l}`$. The quantity $`c_{j\sigma l}`$ denotes the annihilation operator of the electron at the $`j`$-th site of the $`l`$-th chain ($`l=`$1, 2) with spin $`\sigma `$($`=,`$). We use the Fourier transform, $`c_{j\sigma l}=N^{1/2}\mathrm{\Sigma }_kc_{k\sigma l}\mathrm{exp}[ikja]`$ with the total number of sites $`N`$ and the lattice constant $`a`$. First, the $`t_\mathrm{d}`$-term is diagonalized to obtain two bands in the reduced zone, $`\pi /2a<k<\pi /2a`$, and the lower band becomes effectively half-filled, which band is described with fermion operators, $`d_{k\sigma l}`$, and is examined in the present study. Next, diagonalizing the $`t_{}`$-term by $`a_{k\sigma \pm }=(d_{k\sigma 1}+d_{k\sigma 2})/\sqrt{2}`$, the kinetic term is written as $`_K^d=_{k,\sigma ,\zeta }`$ $`\epsilon (k,\zeta )`$ $`a_{k\sigma \zeta }^{}a_{k\sigma \zeta }`$ $`(\zeta =\pm )`$ with $`\epsilon (k,\pm )=2[t^2\mathrm{cos}^2ka+t_\mathrm{d}^2\mathrm{sin}^2ka]^{1/2}\pm 2t_{}`$. Thus we have the following effective Hamiltonian \[?\]. The kinetic energy with the linearized dispersion around the Fermi surfaces, $`k_{\mathrm{F}\pm }=k_\mathrm{F}t_{}/v_\mathrm{F}`$, is expressed as $`_K^d=_{k,p,\sigma ,\zeta }`$ $`v_\mathrm{F}(pkk_{\mathrm{F}\zeta })a_{kp\sigma \zeta }^{}a_{kp\sigma \zeta }`$ with $`p(=+,)`$ denoting right moving (left moving) electrons and $`v_\mathrm{F}=\sqrt{2}ta`$ $`[1(t_\mathrm{d}/t)^2]`$ $`/[1+(t_\mathrm{d}/t)^2]^{1/2}`$, in which the $`t_{}`$-dependence of the velocity is discarded. Coupling constants of interactions corresponding to forward scattering with the same and opposite directions ($`g_4`$ and $`g_2`$), backward scattering ($`g_1`$) and Umklapp scattering ($`g_3`$) are given by $`g_1^{}=g_2^{}=g_4^{}=Ua`$, $`g_3Ua(t_\mathrm{d}/t)`$ and $`g_1^{}=g_2^{}=g_4^{}=0`$ where $``$ and $``$ denote interactions for the same spin and opposite spin.
Applying the bosonization method to electrons around the new Fermi points, we introduce Bose fields of phase variables, $`\theta _{\rho +}`$ and $`\theta _{\sigma +}`$ ($`\theta _{\mathrm{C}+}`$ and $`\theta _{\mathrm{S}+}`$) \[?, ?\], which express fluctuations for the total (transverse) charge density and spin density, respectively \[?\]. The commutation relation with conjugate phase variables is given by $`[\theta _{\nu +}(x),\theta _\nu ^{}(x^{})]_=i\pi \delta _{\nu ,\nu ^{}}\mathrm{sgn}(xx^{})`$ . In terms of these phase variables, our Hamiltonian is given by
$``$ $`=`$ $`{\displaystyle \underset{\nu =\rho ,\sigma ,\mathrm{C},\mathrm{S}}{}}{\displaystyle \frac{v_\nu }{4\pi }}{\displaystyle dx\left[\frac{1}{K_\nu }\left(\theta _{\nu +}\right)^2+K_\nu \left(\theta _\nu \right)^2\right]}`$ (2)
$`+{\displaystyle \frac{g_\rho }{4\pi ^2\alpha ^2}}{\displaystyle }\mathrm{d}x[\mathrm{cos}(\sqrt{2}\theta _{\mathrm{C}+}8t_{}x/v_\mathrm{F})`$
$`+\mathrm{cos}\sqrt{2}\theta _\mathrm{C}\left]\right[\mathrm{cos}\sqrt{2}\theta _\mathrm{S}\mathrm{cos}\sqrt{2}\theta _\mathrm{S}]`$
$`+{\displaystyle \frac{g_\sigma }{4\pi ^2\alpha ^2}}{\displaystyle }\mathrm{d}x[\mathrm{cos}(\sqrt{2}\theta _{\mathrm{C}+}8t_{}x/v_\mathrm{F})`$
$`\mathrm{cos}\sqrt{2}\theta _\mathrm{C}\left]\right[\mathrm{cos}\sqrt{2}\theta _{\mathrm{S}+}+\mathrm{cos}\sqrt{2}\theta _\mathrm{S}]`$
$`+{\displaystyle \frac{g_1}{2\pi ^2\alpha ^2}}{\displaystyle }\mathrm{d}x\mathrm{cos}\sqrt{2}\theta _{\sigma +}[\mathrm{cos}(\sqrt{2}\theta _{\mathrm{C}+}8t_{}x/v_\mathrm{F})`$
$`\mathrm{cos}\sqrt{2}\theta _\mathrm{C}\mathrm{cos}\sqrt{2}\theta _{\mathrm{S}+}\mathrm{cos}\sqrt{2}\theta _\mathrm{S}]`$
$`{\displaystyle \frac{g_3}{2\pi ^2\alpha ^2}}{\displaystyle }\mathrm{d}x\mathrm{sin}\sqrt{2}\theta _{\rho +}[\mathrm{cos}(\sqrt{2}\theta _{\mathrm{C}+}8t_{}x/v_\mathrm{F})`$
$`+\mathrm{cos}\sqrt{2}\theta _\mathrm{C}\mathrm{cos}\sqrt{2}\theta _{\mathrm{S}+}+\mathrm{cos}\sqrt{2}\theta _\mathrm{S}],`$
where $`v_{\rho (\sigma )}=v_\mathrm{F}[1+()U/\pi v_\mathrm{F}]^{1/2}`$, $`K_{\rho (\sigma )}=[1+()Ua/\pi v_\mathrm{F}]^{1/2}`$, $`v_\mathrm{C}=v_\mathrm{S}=v_\mathrm{F}`$, $`K_\mathrm{C}=K_\mathrm{S}=1`$, $`g_{\rho (\sigma )}=+()Ua`$ and $`g_3=Ua(2t_\mathrm{d}/t)/[1+(t_\mathrm{d}/t)^2]`$. The quantity $`\alpha `$ is a cutoff of the order of lattice constant. In eq. (2), there are twelve nonlinear terms rewritten as
$`{\displaystyle \frac{g_{\nu p,\nu ^{}p^{}}}{2\pi ^2\alpha ^2}}{\displaystyle 𝑑x\mathrm{cos}\sqrt{2}\psi _{\nu p}\mathrm{cos}\sqrt{2}\psi _{\nu ^{}p^{}}},`$ (3)
where $`\psi _{\nu \pm }=\theta _{\nu \pm }`$ except for $`\psi _{\mathrm{C}+}=\theta _{\mathrm{C}+}(8t_{}x/v_\mathrm{F})/\sqrt{2}`$ and $`\psi _{\rho +}=\theta _{\rho +}\pi /(2\sqrt{2})`$. The RG equations for $`K_\nu =K_\nu (l),t_{}=t_{}(l)`$ and $`G_{\nu p,\nu ^{}p^{}}=G_{\nu p,\nu ^{}p^{}}(l)`$ are given, up to the second order, as \[?, ?\]
$`{\displaystyle \frac{d}{dl}}\stackrel{~}{t}_{}`$ $`=`$ $`\stackrel{~}{t}_{}{\displaystyle \frac{1}{8}}K_\mathrm{C}(G_{\rho +,\mathrm{C}+}^2`$ (4)
$`+G_{\sigma +,\mathrm{C}+}^2+G_{\mathrm{C}+,\mathrm{S}+}^2+G_{\mathrm{C}+,\mathrm{S}}^2)J_1(8\stackrel{~}{t}_{}),`$
$`{\displaystyle \frac{d}{dl}}K_\nu `$ $`=`$ $`{\displaystyle \frac{1}{2\stackrel{~}{v}_\nu ^2}}K_\nu ^2[G_{\nu +,\mathrm{C}+}^2J_0(8\stackrel{~}{t}_{})`$ (5)
$`+G_{\nu +,\mathrm{C}}^2+G_{\nu +,\mathrm{S}+}^2+G_{\nu +,\mathrm{S}}^2],`$
$`{\displaystyle \frac{d}{dl}}K_\mathrm{C}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{p=\pm }{}}[(K_\mathrm{C}^2J_0(8\stackrel{~}{t}_{})\delta _{p,+}\delta _{p,})(G_{\rho +,\mathrm{C}p}^2`$ (6)
$`+G_{\sigma +,\mathrm{C}p}^2+G_{\mathrm{C}p,\mathrm{S}+}^2+G_{\mathrm{C}p,\mathrm{S}}^2)],`$
$`{\displaystyle \frac{d}{dl}}K_\mathrm{S}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{p=\pm }{}}[(K_\mathrm{S}^2\delta _{p,+}\delta _{p,})(G_{\rho +,\mathrm{S}p}^2`$ (7)
$`+G_{\sigma +,\mathrm{S}p}^2+G_{\mathrm{C}+,\mathrm{S}p}^2J_0(8\stackrel{~}{t}_{})+G_{\mathrm{C},\mathrm{S}p}^2)],`$
$`{\displaystyle \frac{d}{dl}}G_{\nu +,\mathrm{C}p}`$ $`=`$ $`\left[2K_\nu K_\mathrm{C}^p\right]G_{\nu +,\mathrm{C}p}`$ (8)
$`G_{\nu +,\mathrm{S}+}G_{\mathrm{C}p,\mathrm{S}+}G_{\nu +,\mathrm{S}}G_{\mathrm{C}p,\mathrm{S}},`$
$`{\displaystyle \frac{d}{dl}}G_{\nu +,\mathrm{S}p}`$ $`=`$ $`\left[2K_\nu K_\mathrm{S}^p\right]G_{\nu +,\mathrm{S}p}`$ (9)
$`G_{\nu +,\mathrm{C}+}G_{\mathrm{C}+,\mathrm{S}p}J_0(8\stackrel{~}{t}_{})`$
$`G_{\nu +,\mathrm{C}}G_{\mathrm{C},\mathrm{S}p},`$
$`{\displaystyle \frac{d}{dl}}G_{\mathrm{C}p,\mathrm{S}p^{}}`$ $`=`$ $`\left[2K_\mathrm{C}^pK_\mathrm{S}^p^{}\right]G_{\mathrm{C}p,\mathrm{S}p^{}}`$ (10)
$`{\displaystyle \frac{1}{\stackrel{~}{v}_\rho }}G_{\rho +,\mathrm{C}p}G_{\rho +,\mathrm{S}p^{}}{\displaystyle \frac{1}{\stackrel{~}{v}_\sigma }}G_{\sigma +,\mathrm{C}p}G_{\sigma +,\mathrm{S}p^{}},`$
where $`\stackrel{~}{v}_\nu =v_\nu /v_\mathrm{F}`$ and $`\stackrel{~}{t}_{}=t_{}(l)/v_\mathrm{F}\alpha ^1`$. $`J_n`$ is the $`n`$-th Bessel function. The quantity $`l`$ is related to energy scale $`\omega `$ or temperature $`T`$ by $`l=\mathrm{ln}(W/\omega )`$ or $`\mathrm{ln}(W/T)`$ with $`W(v_\mathrm{F}\alpha ^1)`$ being of the order of band width. The initial condition for the RG equations are given by $`K_\nu (0)=K_\nu `$, $`G_{\nu p,\nu ^{}p^{}}(0)=g_{\nu p,\nu ^{}p^{}}/2\pi v_\mathrm{F}`$ and $`\stackrel{~}{t}_{}(0)=t_{}/(v_\mathrm{F}\alpha ^1)`$ where $`g_{\mathrm{C}+,\mathrm{S}+}=g_{\mathrm{C},\mathrm{S}}=0`$, $`g_{\mathrm{C}+,\mathrm{S}}=g_{\mathrm{C},\mathrm{S}+}=Ua`$, $`g_{\sigma +,\mathrm{C}+}=g_{\sigma +,\mathrm{C}}=g_{\sigma +,\mathrm{S}+}=g_{\sigma +,\mathrm{S}}=Ua`$, and $`g_{\rho +,\mathrm{C}+}=g_{\rho +,\mathrm{C}}=g_{\rho +,\mathrm{S}+}=g_{\rho +,\mathrm{S}}=g_3`$. We take $`\alpha =2a/\pi `$ \[?\] and discard the RG equations for the velocity $`v_\nu `$.
## 3 Confinement-deconfinement transition
We calculate eqs. (4)-(10) numerically by choosing $`U`$, $`t_\mathrm{d}`$, and $`t_{}`$ as parameters.
Figure 1 shows the $`l`$-dependence of $`K_\rho (l)`$, $`K_\sigma (l)`$, $`K_\mathrm{C}(l)`$, $`K_\mathrm{S}(l)`$ and $`t_{}(l)`$ which exhibit four gaps. With increasing $`l`$, $`K_\rho (l)`$ decreases to zero forming a gap in the total charge fluctuation while $`K_\mathrm{C}(l)`$ increases infinity to induce a gap in the transverse charge fluctuation. Quantities $`K_\sigma (l)`$ and $`K_\mathrm{S}(l)`$ decrease also to zero and lead to spin gaps for both the total and transverse spin fluctuations. The rapid increase of $`t_{}(l)`$ comes from a fact that the term with $`K_\mathrm{C}(l)`$ in the r.h.s. of eq. (4) reduces to zero due to a factor $`J_1(8\stackrel{~}{t}_{})`$. Figure 2 displays the corresponding $`l`$-dependence of coupling constants. The main figure shows coupling constants for forward and backward scatterings with $`G_{\mathrm{C}+,\mathrm{S}+}`$ (curve (1)), $`G_{\mathrm{C}+,\mathrm{S}}`$ (curve (2)), $`G_{\mathrm{C},\mathrm{S}+}`$ (curve (3)), $`G_{\mathrm{C},\mathrm{S}}`$ (curve (4)), $`G_{\sigma +,\mathrm{C}+}`$ (curve (5)), $`G_{\sigma +,\mathrm{C}}`$ (curve (6)), $`G_{\sigma +,\mathrm{S}+}`$ (curve (7)) and $`G_{\sigma +,\mathrm{S}}`$ (curve (8)) while the inset shows those for the umklapp scattering with $`G_{\rho +,\mathrm{C}+}`$ (curve (9)), $`G_{\rho +,\mathrm{C}}`$ (curve (10)), $`G_{\rho +,\mathrm{S}+}`$ (curve (11)) and $`G_{\rho +,\mathrm{S}}`$ (curve (12)). Coupling constants $`G_{\rho +,\mathrm{C}}`$ and $`G_{\rho +,\mathrm{S}+}`$ increase rapidly and give rise to the trigger of relevance of the coupling constant $`G_{\mathrm{C},\mathrm{S}+}`$ as seen also in the one-dimensional chain. Note that the relevance of coupling constants $`G_{\mathrm{C},\mathrm{S}+}`$, $`G_{\sigma +,\mathrm{C}}`$ and $`G_{\sigma +,\mathrm{S}+}`$ is also obtained in the absence of umklapp scattering \[?\]. The relevant behaviors found from the zero limit of $`K_\rho `$, $`K_\sigma `$, $`1/K_\mathrm{C}`$ and $`K_\mathrm{S}`$ exhibit the phase locking of $`\theta _{\rho +}`$, $`\theta _{\sigma +}`$, $`\theta _\mathrm{C}`$ and $`\theta _{\mathrm{S}+}`$, which are given by $`\sqrt{2}\theta _{\rho +}=\pi /2`$, $`\sqrt{2}\theta _{\sigma +}=0`$, $`\sqrt{2}\theta _\mathrm{C}=0`$ and $`\sqrt{2}\theta _{\mathrm{S}+}=\pi `$ from relevant behaviors of curves (3), (6), (7), (10) and (11).
Other coupling constants, which are expected to decrease \[?\], are still large due to the second order perturbation. The change of the sign of $`G_{\sigma +,\mathrm{S}+}`$ in the renormalization process comes from the relevance of $`\theta _{\sigma +}`$ and $`\theta _{\mathrm{S}+}`$. The effect of $`t_{}`$ on coupling constants becomes large at low energies where the splitting of magnitudes becomes noticeable for the forward scattering (between curves (1) and (4) and between curves (2) and (3)), the backward scattering (curves (5), (6), (7) and (8)), and the umklapp scattering (curves (9), (10), (11) and (12)).
In Fig. 3, the $`l`$-dependence of $`t_{}`$ is shown for the fixed $`t_\mathrm{d}/t=`$ 0.05, $`t_{\mathrm{d},\mathrm{c}}/t(0.082)`$ and 0.1. The increase of $`t_\mathrm{d}`$ leads to the suppression of $`t_{}(l)`$. The case of $`t_\mathrm{d}/t=0.05`$ shows the relevant behavior, which corresponds to deconfinement. The quantity $`t_{}(l)`$ for $`t_\mathrm{d}/t=0.1`$ does not increase monotonically but decreases to zero after taking a maximum indicating confinement. The quantity $`t_{}(l)`$ with a critical magnitude of $`t_\mathrm{d}=t_{\mathrm{d},\mathrm{c}}`$ denotes the behavior between the confinement and the deconfinement. For comparison, we show the dotted curve (the dashed) curve which is calculated for $`U=0`$ ($`U/t`$ = 5, $`t_\mathrm{d}/t`$ = 0.05 but $`g_3`$ = 0 as a special choice of parameter) where the analytical expression for the dotted curve is given by $`t_{}(l)=t_{}\mathrm{e}^l`$. Note that the dashed curve is different from the case of $`U/t=5`$ and $`t_\mathrm{d}=0`$. The dashed curve is evaluated to obtain $`t_{}^{\mathrm{eff},0}`$, which denotes the interchain hopping renormalized only by the intrachain interaction, i.e., without the umklapp scattering. The effective interchain hopping $`t_{}^{\mathrm{eff},0}`$ is defined by $`t_{}^{\mathrm{eff},0}=t\mathrm{exp}[l_{\mathrm{eff},0}]`$ where $`t_{}(l_{\mathrm{eff},0})/t=1`$ for $`g_3=0`$. The quantity $`t_{}^{\mathrm{eff},0}/t`$ becomes unity for $`U`$ = 0 or $`t_{}/t`$ = 1.
The inset denotes the $`t_{}`$-dependence of $`t_{\mathrm{d},\mathrm{c}}`$ for the fixed $`U/t`$ = 3, 4, and 5. The boundary between confinement ($`t_\mathrm{d}>t_{\mathrm{d},\mathrm{c}}`$) and deconfinement ($`t_\mathrm{d}<t_{\mathrm{d},\mathrm{c}}`$) depends appreciably on $`U`$, which gives rise to the enhancement of the confined region on the plane of $`t_{}`$ and $`t_{\mathrm{d},\mathrm{c}}`$. The limiting form for small $`t_{}`$ is given by $`t_{\mathrm{d},\mathrm{c}}/t(t_{}/t)^F`$ with the $`U`$-dependent $`F`$.
The $`l`$-dependence of $`K_\rho (l)`$ is shown by solid curve for the fixed $`t_\mathrm{d}/t`$ = 0.05 and 0.1 in Fig. 4. The charge gap is defined by $`\mathrm{\Delta }_\rho =v_\mathrm{F}\alpha ^1\mathrm{exp}[l_\mathrm{\Delta }]`$ with $`K_\rho (l_\mathrm{\Delta })=K_\rho /2`$ and $`v_\mathrm{F}\alpha ^1=t(\pi /\sqrt{2})[1(t_\mathrm{d}/t)^2]/[1+(t_\mathrm{d}/t)^2]^{1/2}`$. The dashed curve, which denotes $`K_\rho (l)`$ for $`t_{}=0`$, leads to the charge gap, $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}`$, for one-dimensional (1D) case in the presence of dimerization, $`t_\mathrm{d}`$. The charge gap is suppressed slightly by the interchain hopping since $`K_\rho (l)`$ for the solid curve decreases slowly compared with that for the dashed curve, i.e., $`l_\mathrm{\Delta }`$ is increased by $`t_{}`$. Such a behavior is understood from eq. (5) with $`\nu =\rho `$, in which the first term of the r.h.s. becomes small due to the Bessel function suppressed by the large $`t_{}`$. The inset exhibits the charge gap as a function of $`t_\mathrm{d}`$ with the fixed $`U/t`$ = 3, 4 and 5 where the solid (dashed) curve corresponds to the case of $`t_{}/t`$ = 0.1 (0). Note that the dashed curve is given by $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}W(g_3/W)^{1/(22K_\rho )}`$ \[?\].
In Fig. 5, the $`t_{}`$-dependence of $`t_{}^{\mathrm{eff}}`$ with $`U/t`$ = 5 ($`U/t`$ = 3) is shown by the solid (dotted) curve for $`t_\mathrm{d}/t`$ = 0 (1), 0.05 (2) and 0.1 (3) ($`t_\mathrm{d}/t`$ = 0 (4), 0.05 (5) and 0.1 (6)), where the effective interchain hopping $`t_{}^{\mathrm{eff}}`$, including the effect of the dimerization, is estimated by $`t_{}^{\mathrm{eff}}=t\mathrm{exp}[l_{\mathrm{eff}}]`$ with $`t_{}(l_{\mathrm{eff}})/t=1`$ \[?\]. The curves (1) and (4) for small $`t_{}/t`$ well reproduce the analytical result given by $`t_{}^{\mathrm{eff}}/t_{}(t_{}/t)^{\alpha _0/(1\alpha _0)}`$ with the $`U`$-dependent quantity $`\alpha _0=(K_\rho +K_\rho ^1+K_\sigma +K_\sigma ^14)/4`$ \[?\]. Comparing the slope of curve (1) with that of curve (4), one finds that the renormalization of $`t_{}`$ increases by the intrachain interaction $`U`$. With decreasing $`t_{}/t`$, the ratio $`t_{}^{\mathrm{eff}}/t_{}`$ decreases rapidly and becomes zero for $`t_{}`$ less than a critical value of $`t_{}`$ indicated by the arrow.
From the inset of Fig. 3 and that of Fig. 4, we obtain the phase diagram of confinement (I) and deconfinement (II) on the plane of $`t_{}`$ and $`\mathrm{\Delta }_\rho `$. In Fig. 6, the boundary between these two states is shown for $`U/t`$ = 3 (dashed curve), 4 (dotted curve) and 5 (solid curve). The boundary is rather straight and the $`U`$-dependence becomes small compared with that of the inset of Fig. 3. It has been claimed previously that such a behavior indicates the competition between the charge gap and the interchain hopping $`t_{}`$ \[?, ?, ?\]. However such a statement is not clear enough since the boundaries depend on the choice of $`U`$. This problem can be remedied by the following treatment. We take $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}`$ ($`t_{}^{\mathrm{eff},0}`$) in stead of $`\mathrm{\Delta }_\rho `$ ($`t_{}`$) as the vertical (horizontal) axis on the phase diagram where $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}`$ is obtained from the dashed curve of the inset of Fig. 4. The resultant boundaries are shown in the inset of Fig. 6, where the good coincidence is obtained among these three curves. Thus it is concluded that the boundary between confinement and deconfinement is determined by the competition between the one-dimensional charge gap (i.e., in the absence of the interchain hopping) and the interchain hopping renormalized only by the intrachain interaction (i.e., without the umklapp scattering).
## 4 Discussion
We have examined the mechanism of confinement in terms of quarter-filled two-coupled chains with dimerization as a model of low dimensional systems. The confinement occurs when the charge gap induced by the umklapp scattering becomes larger than the interchain hopping renormalized by the intrachain hopping. It has been found that the ratio for $`U/t=5`$ ($`U/t=3`$) and $`0.05<t_{}/t<0.2`$ is given by $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}/t_{}^{\mathrm{eff},0}1.8`$ ($`\mathrm{\Delta }_\rho ^{1\mathrm{D}}/t_{}^{\mathrm{eff},0}1.9`$) and $`1.1\stackrel{<}{}\mathrm{\Delta }_\rho /t_{}\stackrel{<}{}\mathrm{\hspace{0.17em}1.3}`$ ($`1.3\stackrel{<}{}\mathrm{\Delta }_\rho /t_{}\stackrel{<}{}\mathrm{\hspace{0.17em}1.5}`$).
Here we discuss the effect of nearest neighbor interaction ($`V`$), for the small $`V`$, in which the commensurability energy may be negligibly small \[?, ?\]. In this case, $`V`$ does not contribute to both the backward scattering and the umklapp scattering even in the presence of the dimerization ($`t_\mathrm{d}`$). The effect of $`V`$ appears in the forward scattering where the coupling constant is expressed as $`g_2^{}=g_4^{}=2Va`$ and $`g_2^{}=g_4^{}=(U+2V)a`$. In Fig. 7, the boundary of confinement-deconfinement transition with $`U/t=5`$ is shown for $`V/t`$ = 0, 0.5 and 1 where $`V`$ enhances the confined region.
Finally we comment on the case of two-coupled chain with the conventional half-filled Hubbard model where the phase diagram is shown in the inset of Fig. 7. Compared with those of Fig. 6, the ratio $`\mathrm{\Delta }_\rho /t_{}`$ ( $`1.7`$ for $`t_{}/t=0.1`$) is slightly large and the boundary (solid curve) is rather straight. Such a result is compared with the boundaries for the fixed $`t_\mathrm{d}/t`$ = 0.05 and 0.1, which are shown by the dotted curve and dashed curve, respectively. With increasing the dimerization, the slope of the boundary becomes steep and moves toward the solid curve since the large dimerization may lead the system to half-filling.
## Acknowledgements
This work was supported by a Grant-in-Aid for Scientific Research from the Ministry of Education, Science, Sports and Culture (Grant No.09640429), Japan.
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# A Finite Element Algorithm for High-Lying Eigenvalues and Eigenfunctions with Homogeneous Neumann and Dirichlet Boundary Conditions
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## Abstract
We present a finite element algorithm that computes eigenvalues and eigenfunctions of the Laplace operator for two-dimensional problems with homogeneous Neumann or Dirichlet boundary conditions or combinations of either for different parts of the boundary. In order to solve the generalized eigenvalue problem, we use an inverse power plus Gauss-Seidel algorithm. For Neumann boundary conditions the method is much more efficient than the equivalent finite difference algorithm. We have cheked the algorithm comparing the cumulative level density of the espectrum obtained numerically, with the theoretical prediction given by the Weyl formula. A systematic deviation was found. This deviation is due to the discretisation and not to the algorithm. As an application we calculate the statistical properties of the eigenvalues of the acoustic Bunimovich stadium and compare them with the theoretical results given by random matrix theory.
Presentamos un algoritmo de elementos finitos que calcula eigenvalores y eigenfunciones del operador de Laplace para problemas en dos dimensiones con condiciones a la frontera de Neumann o Dirichlet o combinaciones de ambas en distintas partes de la frontera. Para resolver el problema de eigenvalores generalizado, usamos un algoritmo de potencias inverso más otro de Gauss-Seidel. Para condiciones a la frontera de Neumann, el método es mucho más eficiente que el algoritmo equivalente de diferencias finitas. Hemos probado el algoritmo comparando la densidad acumulada de niveles del espectro obtenido numéricamente, con la predicción teórica dada por la fórmula de Weyl. Se encontró una desviación sistemática. Esta desviación es debida a la discretización y no al algoritmo. Como una aplicación, calculamos las propiedades estadísticas de los eigenvalores del estadio de Bunimovich acústico y las comparamos con los resultados teóricos dados por la teoría de matrices aleatorias.
Subject Classification: 65P25, 81C06, 81C07
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I Introduction
Several years ago Neuberger and Noid presented an algorithm and a FORTRAN program for the successive computation of the high-lying eigenvalues and eigenfunctions of a time independent Schrödinger or Helmholtz equation. They used an inverse power method with a Gauss-Seidel procedure for the inversion, and solved the problem with finite differences on successively finer grids. This program was often used and a two-dimensional version thereof was adapted for the case of Laplace operators with homogeneous boundary conditions . The case of Dirichlet conditions works very well and requires minimal adjustments. This is not the case for Neumann conditions. Computation times increase by orders of magnitude compared to the Dirichlet case; this is not surprising as the treatment of an irregular boundary, and particularly of corners, is very cumbersome.
The need for such programs arises both in acoustic and earthquake research, as well as for other wave phenomena. For example, if we want to make statistics of eigenvalues, say for acoustic systems, we need large numbers of states and therefore efficient algorithms. In particular, geometries whose high-frequency limit (ray dynamics) is chaotic, are of great interest. The starting point in this new field called acoustic chaos is that the time-independent wave equation (Helmholtz equation) is the same for different systems such as: quantum billiards, membranes and flat microwave cavities . Thus the statistical fluctuation measures developed in nuclear physics have been applied to those systems and to a wide variety of more complicated systems such as: Chladni’s plates, quartz crystals, aluminum blocks, quantum dots , quantum corrals, waves in a ripple tank, elastic media with ray splitting,microwave cavities with ray splitting amog others. As rather fine details of the boundary of those systems are believed to be important, a good representation of the boundary conditions is essential. Similar arguments will hold if we wish to study the effect of obstacles inside a cavity or in the old Tenochtitlan lake bed.
We shall use the finite element method (FEM). It is based on the minimization of the functional:
$$\left[\mathrm{\Psi }\right]=_R\left(\mathrm{\Psi }\right)^2𝑑sk^2_R\mathrm{\Psi }^2𝑑s,$$
(1)
where $`R`$ is a two dimensional region and $`\psi `$ is the wave function.
The possibility of solving mixed boundary conditions will be important in systems with mirror symmetries, in which we may solve each non-symmetric part using mixed boundary conditions. An irregular boundary as well as corners can also easily be implemented. The FEM formalism is based on minimizing the functional (1) not in the complete Hilbert space but only in a sub-space spanned by a finite set of piece-wise linear functions, typically defined as pyramids over hexagonal cells. We call an element the triangles that form this cell.
The finite difference method becomes very cumbersome for boundary conditions involving normal derivatives at irregular boundaries. This is particularly troublesome for Neumann conditions for which it leads to poor convergence. Novaro et al. () found in a particular case that computations would be two orders of magnitude slower for Neumann conditions than for Dirichlet conditions.
The minimization of the functional (1) on the subspace of Hilbert space and in the non-orthogonal basis mentioned above yields the generalized eigenvalue equation
$$A𝐱=\lambda B𝐱,$$
(2)
with $`A_{ij}=_R\psi _i\psi _jds`$ and $`B_{ij}=_R\psi _i\psi _j𝑑s`$, where $`\psi _i`$ and $`\psi _j`$ denote the functions defined around the node $`i`$ and $`j`$ of the hexagonal grid respectively. These functions for simplicity, are taken to be linear
$$\psi _i^k=a_i^kx+b_i^ky+c_i^k$$
(3)
with $`a_i^k,b_i^k`$ and $`c_i^k`$ constants to be determined for each triangle $`\mathrm{\Delta }_k`$ of the i-th hexagonal element. The trial functions are then defined as pyramids of unit height over each hexagonal cell, i.e. one function corresponding to each node of the grid. The simplicity of a piecewise linear basis gives as result a non-orthogonal basis. The Neumann boundary conditions are obtained by allowing variations of the trial functions on the boundary. The Dirichlet boundary conditions can be obtained by putting the trial functions to zero in the desired part of the boundary. We refer to the literature for a general and deeper discussion of the finite element formalism.
If the dimension $`N`$ of the matrices $`A`$ and $`B`$ is small enough that they can be diagonalised, standard techniques for non-orthogonal bases may be used. But in a typical application, the dimensions are of several thousand to tenthousands. Yet we are only interested in a fairly small number of eigenvalues and eigenfunctions near the low end of the spectrum. We thus have to use some method that makes explicit use of the sparseness of the matrices $`A`$ and $`B`$. We shall see in the next section, that a combination of the inverse power and Gauss-Seidel methods proposed by Neuberger and Noid can be generalised to solve Eq. (2).
Thus in the following section we show how the inverse power method can be applied when a non-orthogonal basis is used. Next we discuss how to implement this for finite differences as well as a number of tricks that can be used to accelerate the numerical procedure and comment on the performance of the program. In section IV we apply the program to the acoustic stadium, analyze the resulting spectra and the corresponding states. We give a correction to the spectral density based on an analysis of the equations infinite elements for the exactly solvable rectangle discussed in the appendix, Finally we present some conclusions.
II The inverse power method in a non-orthogonal basis
We shall transform the Eq. (2) by left multiplication with $`B^1`$ to the form
$$B^1A𝐱=\lambda 𝐱$$
(4)
As usual a power $`(B^1A)^M(A^1B)^M`$ applied to an arbitrary initial vector $`𝐱^0`$ will successively select the lowest eigenvector corresponding to eigenvalue $`\lambda _1`$ as it will appear with a power $`(1/\lambda _1)^M`$.
The problem now seems to be that $`A^1`$ is no longer a sparse matrix, but this can be averted in the procedure of applying $`A^1B`$. Thus we need to know the product
$$𝐲A^1B𝐱^0.$$
(5)
In order to obtain this product we define $`\stackrel{~}{𝐱}^0B𝐱^0`$ giving for Eq. (5)
$$𝐲=A^1\stackrel{~}{𝐱}^0$$
(6)
Multiplying by A by the left we obtain
$$A𝐲=\stackrel{~}{𝐱}^0$$
(7)
that can be solved alternatively by Jacobi or Gauss-Seidel procedures which will work well as all matrices involved continue to be sparse. Summarising: the Gauss-Seidel Method can be utilized because we do not need the inverse matrix $`A^1`$ but the product $`A^1\stackrel{~}{𝐱}^0`$.
Up to here we have only specified how to obtain the lowest state; for excited states the usual procedure is to orthogonalize the space in which we carry out the variation on all states that have already been calculated. Here again the non-orthogonality of our basis has to be taken into account; indeed, the eigenstates of the Laplace operator are orthogonal and we have to derive the consequences this has for eigenvectors in our non-orthogonal basis. If $`\mathrm{\Psi }_i`$ and $`\mathrm{\Psi }_j`$ are eigenstates of the Laplace operator and are expanded as with coefficients $`\alpha _m^l;l=i,j`$ that form vectors $`𝐱_l`$ we can readily check that
$$_R\mathrm{\Psi }_iB\mathrm{\Psi }_j𝑑s=\delta _{ij}$$
(8)
implies
$$𝐱_i^tB𝐱_j=\delta _{ij}$$
(9)
and vice-versa. Thus we replace the usual orthogonalization by what we may call a $`B`$-orthogonalization, i.e. we require the new vectors to be orthogonal on $`B𝐱_j`$ (where $`j`$ indicates the calculated eigenstates) thus guaranteeing the validity of Eq. (9) and by consequence Eq. (8).
III Implementation of finite elements and convergence
Once we know the grid and the matrix elements of $`A`$ and $`B`$ both of which will be evaluated in the end of this section, we are in principle ready to write our program. As usual a program consists in part of an efficient algorithm which we have presented, and in part of a bag of semi-empirical tricks that tend to repeat themselves in different guises again and again. The efficiency of the latter is quite problem-dependent and the corresponding parameters must be adjusted to optimize operation of the program in every case. We shall give some recommendations that ought to work reasonably, but we urge the user to fiddle around with these parameters.
When using Neumann boundary conditions, the lowest eigenvalue is zero and its corresponding eigenfunction is a constant. Yet the Gauss-Seidel inversion requires positive definite eigenvalues. This is obtained by adding an arbitrary constant $`C`$ multiplied by $`B`$ to the Laplace operator. If we choose this constant large we will need few Gauss-Seidel iterations as the operator is near diagonal. On the other hand, we will loose precision and convergence speed in the inverse power process as neighbouring eigenvalues will have their inverse very close to each other. A good balance seems to be to choose the constant smaller than, but of the order of, the largest eigenvalue which we want to obtain. In the case of Dirichlet conditions this parameter can also be introduced as a means of improving convergence exclusively. Adequate choices of this parameter may improve computation time by a factor of $`2`$.
Superconvergence, or over-relaxation, is another standard tool to improve convergence. We introduce a factor $`1<\alpha <2`$ and at each step of any iteration from $`𝐱_n𝐱_{n+1}`$ we replace $`𝐱_{n+1}`$ by $`𝐱_n+\alpha (𝐱_{n+1}𝐱_n)`$ thus correcting a little more than the iteration warrants. For the Gauss-Seidel iterations values of $`\alpha 1.5`$ yield improvements in computing time of the order of $`2`$, which is consistent with what we found for finite differences. But a careful analysis for the stadium shows that at least in this case for a quite precise value of $`\alpha =1.75`$ we find an acceleration by a factor of $`3`$. On the other hand, in the power iteration, improvements are not very significant and only values of $`\alpha `$ near 1 seem acceptable. We do not use superconvergence in this context, but again we must warn that it might be very significant in other cases. All these parameters were carefully explored by Méndez.
Another option is the stepwise introduction of finer grids, with interpolated results from the rougher grid results as starting point. This idea was extensively exploited in the finite difference programs of Neuberger and Noid and gave excellent results both shortening the computation time by giving a good trial function on the finest grid and allowing a considerable improvement of the eigenvalue upon extrapolation. Unfortunately these advantages diminish even in their case as we go to higher states, for which only the finest grid is acceptable (finer ones would yield too large matrices). For this reason we have not implemented these procedures for finite elements at this point.
Now we return with the problem of establishing the grid that defines our finite elements and of calculating the matrix elements of $`A`$ and $`B`$. For this we use a standard method summarised in the following steps:
1. We immerse the region $`R`$ in a quadratic grid. The triangles of the elements are defined using the sides of the squares and in addition one of the diagonals.
2. We redefine the grid and triangles along the boundary by considering all points that lie just outside our contour. We then move the exterior points along the edges of the squares or the selected diagonal, to the boundary, so as not to change the topology of the grid and its connections.
The corresponding integrals are evaluated using a change of variable with a linear transformation. The transformation can be found solving a linear system of 3 equations. Evaluating the constants of Eq. (3) and the Jacobian of the transformation (in this case the one half area of the triangle, because the transformation is linear) we obtain for the integrals:
$$A_{ij}=\underset{k}{}S(\mathrm{\Delta }_k)(a_i^ka_j^k+b_i^kb_j^k)$$
(10)
and
$$B_{ij}=\{\begin{array}{c}_k2S(\mathrm{\Delta }_k)\frac{1}{24},ij,\\ \\ _k2S(\mathrm{\Delta }_k)\frac{1}{12},i=j\end{array}$$
(11)
Here $`S(\mathrm{\Delta }_k)`$ is the area of triangle $`\mathrm{\Delta }_k`$ and the sum is made over the number of triangles common to both elements $`i`$ and $`j`$.
The routines consist of two main parts. The first computes all elements different from zero. The elements are put in two matrices of $`N\times 7`$ to use the sparseness of matrices. In order to retain the original positions of each element we construct an index matrix. The second part solves the generalized system of equations using the inverse power iteration or a standard diagonalization freeware routine , if the dimensions of the matrices are small (less than $`5000\times 5000`$). We obtain $`500`$ eigenvalues with error less than $`5\%`$ of the average spacing between levels.
The input for the first routine is a set of points localised on the border of the region of interest and a list of intervals of this enumeration in which we want Dirichlet boundary conditions. We assume Neumann boundary conditions in the rest of the border.
For the second routine the inputs are the number of eigenvectors, the tolerance for the Gauss Seidel and inverse power iterations, the constant $`C`$, and the superconvergence factor $`\alpha `$. The output consists of two files with the corresponding eigenvalues and eigenvectors.
IV Application to the acoustic stadium
By way of example we apply the program to some particular two-dimensional region $`R`$; such as the Bunimovich stadium which is completely chaotic in classical mechanics. This geometry consists of two semicircles of radius $`r`$, joined by two straight lines with length $`2l`$ . We will apply pure Neumann boundary conditions. On a grid of $`3000`$ points we obtain, by the inverse power iteration method described above, $`200`$ eigenvectors and eigenvalues with a CPU time of approximately $`200`$ $`\mathrm{sec}`$. on an ALPHA Work-Station.
To analised a discrete sequence of eigenvalues we first define the cumulative level density or staircase function
$$N(E)=\underset{i=1}{\overset{N}{}}\mathrm{\Theta }(EE_i)$$
(12)
and its derivative $`\rho (E)`$ (the level density),
$$\rho (E)=\underset{i=1}{\overset{N}{}}\delta (EE_i).$$
(13)
Here $`\mathrm{\Theta }`$ and $`\delta `$ are the Heaviside and Dirac delta functions respectively. The staircase function is usually divided in a smooth part plus a fluctuating part:
$$N(E)=\overline{N}(E)+N_{fluct}(E).$$
(14)
The cumulative level density obtained by the finite element method for the stadium with Neumann boundary conditions is plotted in Fig. 1a. For comparison the average level density
$$\overline{N}_{Weyl}(E)=(AE+P\sqrt{E}+\kappa )$$
(15)
obtained from the Weyl formula is depicted. Here $`A`$ is the area of the billiard, $`P`$ is the length of its perimeter and $`\kappa `$ is a constant that contains information on the topological nature of the billiard and the curvature of its boundary. From this figure we can see that the finite element method gives a systematic deviation of the eigenvalues. This systematic deviation in the cumulative level density is due to the discretisation of the finite element. To show this we calculate in the appendix the eigenvalues $`\lambda _{n,m}`$ for the equations in finite elements for an $`a\times a`$ square with periodic boundary conditions. The final equation is
$$\lambda _{n,m}=\frac{42\left(\mathrm{cos}\left(k_x\right)+\mathrm{cos}\left(k_y\right)\right)}{\frac{1}{2}+\frac{1}{6}\left(\mathrm{cos}\left(k_x\right)+\mathrm{cos}\left(k_y\right)+\mathrm{cos}\left(k_x+k_y\right)\right)}$$
(16)
where $`k_x=\frac{n\pi }{a}`$ and $`k_y=\frac{m\pi }{a}`$ are the wave number on the $`x`$ and $`y`$ directions, $`a`$ is the size of the side and $`n,m=0,\pm 1,\pm 2,\mathrm{}`$
In Fig. 1b we show the cumulative level density obtained from Eq. (16) for a $`50\times 50`$ square. We also plotted the cumulative level density for the exact square ($`k_x^2+k_y^2`$). The one coming from the diagonalisation is indistinguishable from the obtained by the Eq. (16). The systematic deviation observed in the square with periodic boundary conditions will be used to correct the Weyl formula for arbirary-shaped billiards and arbitray boundary conditions. To do this we calculate the difference between the fits for both cumulative level densities –Eq. (16)– and the exact square. The polynomial for the difference is $`A(4.6005513.75335E10.44418E^2+0.45601E^3)/2500`$, where $`A`$ is the area of the billiard. If this polynomial is added to the Weyl prediction, the resulting curve yields good agreement up to $`400`$ eigenvalues in the stadium (See Fig. 1a).
We shall study the fluctuation properties of the spectrum. In order to do this for a typical sequence of eigenvalues, it is necessary to suppress the secular variation. This “unfolding” of the levels is done through the mapping
$$E_iE_i^{}=\overline{N}(E_i).$$
(17)
The effect of such mapping on the original sequence is that the new sequence has on the average a constant spacing equal to one. Although we should use the Weyl formula corrected by the polynomial, we can use a polynomial fit $`\overline{N}(E)`$ for $`N(E)`$ that takes into account the systematic deviation due to the discretisation. We can then calculate different statistical measures of the fluctuating part of the spectrum. The first statistic we shall use is the nearest-neighbour spacing distribution $`p(s)`$ where $`s=E_{i+1}^{}E_i^{},`$ which gives information on the short range correlations. The $`p(s)`$ for the stadium is given in Fig. 2. Notice that it agrees with the values predicted by the gaussian orthogonal ensemble (GOE) typical for chaotic systems. For completeness, the Poisson case, typical for integrable systems, is also depicted. The spectrum of the acoustic Bunimovich stadium shows $`p(s0)0`$. This behaviour is called level repulsion. In fact the nearest-neighbour spacing distribution agrees with
$$p_{Wigner}(s)=\frac{\pi }{2}s\mathrm{exp}(\frac{\pi }{4}s^2),s0$$
(18)
know in spectral statistics as the Wigner surmise, which is very close to the GOE prediction.
We can also define the $`k`$th-neighbour spacing distribution $`p(k;s_k).`$ Now $`s_k=E_{i+k+1}^{}E_i^{}`$, so that $`p(s)=p(0;s_0s)`$. It is well known that these distributions tend to a normal distribution as $`k`$ grows. Since $`s_k=k+1`$, the only relevant parameter left is the width $`\sigma (k)`$ of the distribution. In Fig. 3 we show $`\sigma (k)`$ for the stadium, GOE and Poisson cases.
The correlation coefficient between adjacent spacings is another short range fluctuation measure. For the acoustic stadium we obtained $`C=0.26`$ near to the GOE value ($`C_{GOE}=0.27`$) and far from the Poisson value ($`C_{Poisson}=0`$).
Another commonly used statistic is the number variance $`\mathrm{\Sigma }^2(L)`$, defined as the second moment of the number of levels $`\nu (L)`$ within an interval of length $`L`$, and given for the stadium, GOE and the Poisson cases in Fig. 4. For the stadium and GOE cases and for large $`L`$, $`\mathrm{\Sigma }^2(L)\mathrm{ln}(L)`$ indicating a very rigid sequence.
The number variance depends exclusively on the two-point function. We consider further moments
$$\mathrm{\Sigma }^k(L)=\left[\nu (L)\nu (L)\right]^k$$
(19)
that depend on higher correlations. To emphasise the three- or four-point properties, it is useful to consider the skewness
$$\gamma _1(L)=\mathrm{\Sigma }^3(L)\times \left[\mathrm{\Sigma }^2(L)\right]^{3/2}$$
(20)
and the excess
$$\gamma _2(rL)=\mathrm{\Sigma }^4(L)\times \left[\mathrm{\Sigma }^2(L)\right]^23.$$
(21)
The numerical values obtained for the stadium are shown in Figs. 5 and 6.
In many instances the Fourier transform of the spectra has also proven useful. It contains the same information as the spectrum itself. On the other hand the power spectrum:
$$\left|C(t)\right|^2=\frac{1}{2\pi }\left|\underset{\mathrm{}}{\overset{\mathrm{}}{}}e^{2i\pi Et}\rho (E)𝑑E\right|^2,$$
(22)
depends exclusively on the two-point function. The short-range part of $`P(t)`$ gives specific information concerning the long-range stiffness. Numerical values are given in Fig. 7.
The eigenfunctions for the rectangle were also successfully compared with the exact solution. In Fig. 8 we show two eigenfunctions of the stadium: one of them with Dirichlet boundary conditions which shows typical feature of the whispering gallery states, and the last one, with Neumann boundary conditions shows a near bouncing ball state. There are other kinds of features in different eigenfunctions, like scars , and others that resemble noise . Some which have been reproduced by the autors with this algorithm can be seen in ref. .
Finally, we performed all the calculation on a quarter stadium and used Neumann (N) or Dirichlet (D) boundary conditions on the symmetry lines. This implies that we studied the symmetric or antisymmetric solutions with respect to both reflection symmetries of the stadium. The full solution, shown in the figures, is recovered making the corresponding reflections respect each symmetry axis.
V Conclusions
We have presented an algorithm based on the finite element method that computes eigenfunctions and eigenvalues of the two-dimensional Helmholtz equation with mixed Neumann and Dirichlet boundary conditions. The algorithm is divided in two parts: one that computes the matrix elements and another that diagonalizes the generalised eigenvalue equation using the inverse power and Gauss-Seidel Methods. The Gauss-Seidel method runs efficiently if we use over-relaxation where the gain in computing time peaks at a factor of $`3`$ around the value 1.75 for the superconvergence (over-relaxation). A systematic error in frequencies was found. This error is due to the discretisation and can be estimated by using the eigenvalues of the equations in finite elements. The programs were used to calculate the eigenvalues of the acoustic stadium. The fluctuation measures of the eigenvalues were compared with the random matrix predictions.
The algorithm is very useful in diverse branches of wave physics. The program can obtain eigenfunctions and eigenvalues of two-dimensional acoustic cavities, two-dimensional microwave cavities, membranes, quantum billiards, elastic valleys in certain approximations, etc., and can be readily generalised to other problems.
Acknowledgements
This work was supported by the UNAM—CRAY Research Inc. project SC101094, by UNAM-DGAPA project IN106894. G. Báez and R. A. Méndez received fellowships by UNAM-DGAPA. We want to thank to the IF-UNAM in which, the main part of this work was developped.
Appendix: Eigenvalues for the equations in finite elements.
In this appendix we deduce the equations in finite elements for a square with periodic boundary conditions. If we denote by
$$\text{x}=e^{i(k_xn+k_ym)}$$
(23)
the amplitude in the grid point $`(n,m)`$, the Eq. (2) for the bulk element is given by
$`(4e^{ik_x}e^{ik_y}e^{ik_x}e^{ik_y})\text{x}`$ (24)
$`=\lambda _{n,m}({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{12}}e^{ik_x}+{\displaystyle \frac{1}{12}}e^{ik_y}+{\displaystyle \frac{1}{12}}e^{ik_x}+{\displaystyle \frac{1}{12}}e^{ik_y}`$ (25)
$`+{\displaystyle \frac{1}{12}}e^{i(k_x+k_y)}+{\displaystyle \frac{1}{12}}e^{i(k_x+k_y)})\text{x.}`$ (26)
Here $`k_x=\frac{n\pi }{a}`$ and $`k_y=\frac{m\pi }{a}`$. We also assumed that the size of the grid is one and that for the bulk $`A_{i,i}=4,`$ $`A_{i,j\pm 1}=A_{i\pm 1,j}=1`$, $`B_{i,j\pm 1}=B_{i\pm 1,j}=B_{i\pm 1,j\pm 1}=\frac{1}{12}`$ and $`B_{i,i}=\frac{1}{2}`$. The final form for the eigenvalue equation in finite elements is given by
$`42\left(\mathrm{cos}(k_x)+\mathrm{cos}(k_y)\right)`$ (27)
$`=\lambda _{n,m}\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{6}}\left(\mathrm{cos}(k_x)+\mathrm{cos}(k_y)+\mathrm{cos}(k_x+k_y)\right)\right).`$ (28)
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# DTP–MSU/00-05 hep-th/0005099 Solitons and black holes in non-Abelian Einstein-Born-Infeld theory
(May 6, 2000)
## Abstract
Recently it was shown that the Born–Infeld–type modification of the quadratic Yang–Mills action gives rise to classical particle-like solutions in the flat space which have a striking similarity with the Bartnik-McKinnon solutions known in the gravity coupled Yang-Mills theory. We show that both families are continuously related within the framework of the Einstein-Born-Infeld theory through interpolating sequences of parameters. We also investigate an internal structure of the associated black holes. It is found that the Born–Infeld non–linearity leads to a drastic modification of the black hole interior typical for the usual Yang-Mills theory. In the latter case a generic solution exhibits violent metric oscillations near the singularity. In the Born-Infeld case a generic interior solution is smooth, the metric has the standard Schwarzschild type singularity, and we did not observe internal horizons. Such smoothing of the ’violent’ EYM singularity may be interpreted as a result of quantum effects.
PACS numbers: 04.20.Jb, 04.50.+h, 46.70.Hg
Physical significance of the particle-like solutions to the Einstein-Yang-Mills (EYM) field equations found by Bartnik and McKinnon (BK) (for more details see a review paper ) as well as their possible role in the string-inspired models remains rather obscure. These ’particles’ where found to have a sphaleronic nature and they could be responsible for fermion number violating effects. However, it is not clear whether these purely classical solutions can survive the embedding into some quantum framework. Another related puzzle is the singularity structure of the associated black holes . It was shown that the metric inside the classical EYM black holes exhibits violent oscillations near the singularity, which go far beyond the classical bounds . Presumably these oscillations must be regularized in the quantum theory, but no relevant explanation was suggested so far.
To probe the effect of quantum corrections to the Yang-Mills lagrangian within the string/M theory framework it is tempting to utilize the Born-Infeld modification of the Yang-Mills action describing the low energy dynamics of D-branes . Classical solutions to both the Abelian and non-Abelian Born-Infeld theories received recently much attention. Although there are certain complications in the definition of the trace over the gauge group generators in the Born-Infeld action , it is believed that the simplified ’square root’ form of the lagrangian gives correct (at least qualitatively) predictions concerning solitons. Magnetic monopoles were shown to persist in such a theory . A new type of a particle-like solution in the non-Abelian Born-Infeld model was obtained by Gal’tsov and Kerner (GK) . It was shown that this theory gives rise to flat space classical glueballs which have a striking similarity with the BK solutions.
More close relationship between these two particle-like configurations becomes clear within the Einstein-Born-Infeld (EBI) theory. As was shown recently by Wirschins, Sood and Kunz , the GK solutions survive when gravity is added, in which case the corresponding black holes also come into play. There is a substantial difference between the above two families which has to be clarified, however. Both the BK and GK particles form the same kind of discrete sequences resulting from the ’quantization’ of the parameter entering the boundary conditions near the origin. But contrary to the BK case for which this parameter sequence is convergent to a limiting value, in the GK case the corresponding sequence is divergent. So it is necessary to check more carefully whether the both families are continuously related indeed.
Here we investigate the non-Abelian Einstein-Born-Infeld solitons and black holes in more detail. We show that within this framework one has a unique family of the regular particle-like solutions smoothly interpolating between the GK and BK solutions while the gravitational coupling constant varies from zero to a large value in units of the Born-Infeld ’critical field’ parameter. We also investigate the interior structure of the non-Abelian EBIYM black holes and find that the problem of violent metric oscillations here is resolved indeed.
Recall that the static spherically symmetric EYM field equations admit three smooth branches of local solutions near the singularity exhibiting the Schwarzschild, Reissner-Nordström and the imaginary-charge Reissner-Nordström type behavior. Neither of these three, however, has a sufficient number of free parameters to be generic. Therefore when one is moving from the event horizon into the black hole interior one can meet smooth local solutions near the singularity at best for discrete values of the black hole mass . The generic EYM black hole interior looks very differently from other explicitly known cases. When the singularity is approached the metric exhibits oscillations with an infinitely growing amplitude and an infinitely decreasing period. Clearly, the classical bounds are exceeded after a few oscillation cycles. In the non-Abelian EBI theory we can also find several local solution branches near the singularity, but the situation is drastically different. The local solution which has a sufficient number of parameters to be generic has a perfectly smooth Schwarzschild type behavior. Therefore the problem of oscillations is resolved in a natural way.
As in we assume the ’square root’ form of the non–Abelian Born–Infeld action which, as we believe, describes well enough particle-like solutions being at the same time much simpler than the favored by strings symmetrized trace action. Thus the EBIYM action is chosen as
$$S=\frac{1}{16\pi G}\left(R+4G\beta ^2(1)\right)d^4x,$$
(1)
where
$$=\sqrt{1+\frac{1}{2\beta ^2}F_{\mu \nu }^aF_a^{\mu \nu }\frac{1}{16\beta ^4}(F_{\mu \nu }^a\stackrel{~}{F}_a^{\mu \nu })^2}.$$
(2)
Without loss of generality the dimensionless gauge coupling constant (in the units $`\mathrm{}=c=1`$) will be set to unity, so we are left with two (dimensionfull) parameters: the BI ’critical field’ $`\beta `$ of dimension $`L^2`$, and the Newton constant $`G`$ of dimension $`L^2`$ (Planck’s length). From these one can form a dimensionless constant
$$g=G\beta ,$$
(3)
which is the only substantial parameter of the theory . The decoupling of gravity corresponds to the limit $`g0`$, while the $`g\mathrm{}`$ limit (after rescaling) gives the EYM theory.
We consider the case of the $`SU(2)`$ gauge group assuming for the YM field a usual spherically symmetric static purely magnetic ansatz
$$A=(1w(r))\left(T_\theta \mathrm{sin}\theta d\phi T_\phi d\theta \right),$$
(4)
where a rotated basis for the gauge group generators is used. The spacetime metric is parameterized as follows:
$$ds^2=N\sigma ^2dt^2\frac{dr^2}{N}r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right).$$
(5)
The equations of motion after a rescaling of the radial coordinate $`(\beta )^{1/2}rr`$ (making $`r`$ dimensionless) take the form of two coupled equations for $`N`$ and $`w`$:
$$\left(\frac{Nw^{}}{}\right)^{}=\frac{w(w^21)}{r^2}\frac{2gw^3N}{^2},$$
(6)
$$(Nr)^{}=1+2gr^2(1),$$
(7)
where now
$$=\sqrt{1+2\frac{Nw^2}{r^2}+\frac{(1w^2)^2}{r^4}},$$
(8)
and primes denotes the derivatives with respect to r. The third is a decoupled equation for $`\sigma `$
$$(\mathrm{ln}\sigma )^{}=\frac{2gw^2}{r}$$
(9)
which can be easily solved once the YM function $`w`$ is found.
We are interested in the asymptotically flat configurations such that the local mass function $`m(r)`$, defined through
$$N=1\frac{2m(r)}{r},$$
(10)
has a finite limit $`mM`$ as $`r\mathrm{}`$, while $`\sigma 1`$. Like in the EYM case, it can be easily derived from the Eq. (7) that finiteness of $`M`$ implies the following asymptotic behavior of the YM function: $`w=\pm 1+O(r^1)`$. In this limit $`1`$ so the BI non-linearity is negligible.
Let us first discuss the globally regular solutions which start at the origin with the following series expansion (its convergence in a non-zero domain may be proved by the standard methods ):
$$w=1br^2+\frac{br^4(3b(44b^2+3)+4g[28b^2+3(4b^2(48b^2+13)+3)_0^1]))}{30(4b^2+1)}+O(r^6),$$
(11)
$$N=1+\frac{2r^2}{3}g\left(1_0\right)\frac{16gb^2r^4}{15(4b^2+1)}\left(g(_01)^2+3b_0\right)+O(r^6),$$
(12)
where $`_0`$ is a limiting value of the square root at the origin
$$_0=\sqrt{1+12b^2}.$$
(13)
This local solution has a unique free parameter $`b`$. In the limiting case $`g=0`$ it coincides with that found in the flat space case , while to make contact with the corresponding expansion of the EYM theory one has to take the limit $`g\mathrm{}`$ with a simultaneous rescaling $`bb/g`$. Near the origin the spacetime is flat.
Like in the EYM case, matching of the local solution departing from the origin as (11,12) and meeting the conditions of asymptotic flatness can be achieved for a discrete sequence of the free parameter values $`b=b_n(g)`$ labeled by the number $`n`$ of zeroes of $`w(r)`$. The proof of existence may be given along the lines of , here we do not enter into mathematical details concentrating rather on the qualitative physical picture.
We have investigated numerical solutions in a large range of $`g`$. Typical behavior for small $`g`$ is shown of Figs. (1,2). The YM curves for $`g=.001`$ are practically undistinguishable from those found previously by Kerner and one of the authors in the flat space BIYM model. The region of oscillations corresponds to an unscreened Coulomb charge inside the particles. One can see that this region expands both in the direction of small and large $`r`$ with growing $`n`$. The metric deviate from the flat one rather weakly for small $`g`$ (weak gravity). One observes a stabilization of the metric in the region of the $`w`$-oscillations.
The amplitude of oscillations of $`w`$ in the intermediate region decreases, so in this region the solution can be regarded as approaching the embedded Abelian solution $`w0`$. As was observed by Wirschins, Sood and Kunz , the metric for small $`g`$ also approaches the corresponding Abelian Bion metric (with zero ’seed’ mass in the singularity). To avoid confusion it is worth noting that for odd $`n`$ the Yang-Mills topology of the EBIYM regular solutions (kink) is essentially different form that of the embedded Abelian one $`w0`$ (trivial). Thus it would be misleading to say that the sequence of non-Abelian EBIYM solutions converges to an Abelian one in the global sense.
For large $`g`$ the regular EBIYM solutions after the coordinate rescaling $`r\sqrt{g}r`$ tend to the BK solutions of the EYM system, see Fig. 3 for $`g=10`$. Comparing this with the weak gravity behavior (Fig. 1) we observe that gravity reduces the region of the unscreened charge. The metric deviates substantially form the flat one. The function $`N(r)`$ for large $`n`$ exhibits a deep well (an ’almost’ horizon), but remains always strictly positive.
The sequences $`b_n`$ of discrete values of the parameter in the expansion (11) behaves very differently for small and large $`g`$. As was found in , $`b_n`$ in the flat space are quite big with respect to the corresponding BK values, and the sequence does not converge with growing $`n`$. Contrary to this, the BK sequence $`b_n`$ rapidly converges to a limiting value $`b_{\mathrm{}}`$. In this case there is an additional limiting solution with different space-time structure .
We have obtained numerically the interpolating functions $`b_n(g)`$ for several $`n`$ clearly demonstrating that these two extrema are continuously related indeed (Fig. 5). For small values of $`g`$ these functions tend to the GK values
$$b_n(0)=b_n^{GK},$$
(14)
while in the opposite limit $`g\mathrm{}`$ one recovers the convergent sequence of rescaled BK values
$$b_n^{BK}=\underset{g\mathrm{}}{lim}b_n(g)/g.$$
(15)
The corresponding solutions of the EYM system with the standard YM lagrangian are recovered in terms of the rescaled coordinate $`\stackrel{~}{r}=r/\sqrt{g}`$. In the intermediate region all parameter functions $`b_n(g)`$ were found to be monotonously varying between the above extrema.
Masses of the regular solutions as functions of the effective gravitational coupling are shown on Fig. 6. For vanishing $`g`$ one recovers the GK masses after a rescaling
$$\underset{g0}{lim}M_n(g)M_n^{GK}g.$$
(16)
We recall that the sequence $`M_n^{GK}`$ converges to the mass of an embedded Abelian solution. The reason is that the main contribution to the mass comes from the region of oscillations where $`w`$ approaches an Abelian value $`w0`$ with growing $`n`$. But for any $`n`$ the limiting values $`w(0),w(\mathrm{})`$ are equal to $`\pm 1`$, so, as we have already noted, the global topology of non-Abelian solutions is entirely different. In the limit of strong gravity one recovers the BK masses after a rescaling
$$\underset{g\mathrm{}}{lim}M_n(g)M_n^{BK}\sqrt{g}.$$
(17)
Now discuss the black holes. These are parameterized by the horizon radius $`r_h`$ and the value of the YM function $`w_h=w(r_h)`$ at the horizon. The series expansions near the regular horizon reads:
$`w`$ $`=`$ $`w_h+{\displaystyle \frac{w_h(w_h^21)}{r_hN_h^{}}}(rr_h)+O\left((rr_h)^2\right)`$
$`N`$ $`=`$ $`N_h^{}(rr_h)+O\left((rr_h)^2\right),`$ (18)
$`N_h^{}`$ $`=`$ $`{\displaystyle \frac{1}{r_h}}\left[1+2g^2r_h^2\left(1+\sqrt{1+(w_h^21)^2/r_h^4}\right)\right].`$
Asymptotically flat solutions with such boundary condition are likely to exist for any horizon radii $`r_h`$. Exterior black hole solutions are very similar to regular solutions, especially for small $`r_h`$. So our main interest is in the interior solutions. We start by listing various series solutions that can be obtained near the singularity.
Let us first explore the series expansion for an embedded Abelian solution . For the unit magnetic charge (what corresponds to $`w0`$) the local mass function satisfies the equation
$$m^{}=g\left(\sqrt{r^4+1}r^2\right).$$
(19)
Expanding the square root at small $`r`$ and integrating one obtains:
$$N=\frac{2m_0}{r}+12g+\frac{2}{3}gr^2\frac{1}{5}gr^4+O(r^8).$$
(20)
The ’seed’ mass parameter $`m_0`$ may be positive, negative or zero. Positive $`m_0`$ corresponds to a timelike singularity of the Schwarzschild type. Negative $`m_0`$ corresponds to a spacelike singularity, in which case an internal horizon also exist (though contrary to a more common example of the Reissner-Nordström metric, the ’local charge’ term $`Q^2/r^2`$ now is absent). Vanishing $`m_0`$ is particularly interesting. The metric near the origin is then locally flat unless $`g=1/2`$ in which case the limiting value $`N(0)`$ shrinks to zero. For other values of $`g`$ there is a conical singularity , the critical value $`g=1/2`$ corresponding to an extremal deficit angle. Within the present framework the embedded Abelian black holes correspond to an identically vanishing function $`w`$. It can be shown that the local series solution starting at $`r=0`$ with zero initial value $`w_0=0`$ and arbitrary $`m_0`$ generates this global embedded Abelian solution.
For non-Abelian solutions the value of the YM function at the singularity $`w(0)=w_0`$ should therefore be non-zero. We have found the generalization of the series expansion (20) with $`w_00`$. It is valid for the non-zero seed mass parameter $`m_0`$:
$$w=w_0\left(1+\frac{er}{2m_0}+\frac{er^2}{16m_0^2}\left[3(2e1)4ge\right]\right)+O(r^3),$$
(21)
$$eN=\frac{2m_0}{r}+12ge+\frac{3gew_0^2r}{2m_0}+gr^2\left(\frac{2}{3}+\frac{ew_0^2}{12m_0^2}\left[3(5e2)8ge\right]\right)+O(r^3),$$
where $`e=1w_0^2`$. This local solution has two free parameters $`w_0,m_0`$. For $`w_0=0`$ it coincides with (20). If $`w_00`$ becomes singular in the limit $`m_0=0`$. The search for local solutions with $`m_0=0`$ shows that in this case either $`w_0=0`$ in which case we come back to the Abelian embedded solution $`w0`$, or $`w_0=\pm 1`$, then we recover the regular solution (11,12). It is also worth noting that, contrary to the EYM case, in the EBIYM theory there are no local solutions with $`Nr^2`$ behavior at the singularity (Reissner-Nordström type).
Since the asymptotic solution (21) fails to contain a sufficient number of free parameters to be a generic solution, (this number is equal to three for the system of equations (6,7)), the question is how a generic solution looks like near the singularity, in particular, whether it admits any series expansion. In the case of the ordinary Yang-Mills lagrangian such an expandable solution does not exist at all, one finds that a generic solution has a non-analytic oscillating behavior . Here the situation is different, although the generic local solution around the singularity still exhibits non-analyticity in terms of the variable $`r`$. It turns out to be series expandable but in terms of the $`\sqrt{r}`$:
$`w`$ $`=`$ $`w_0+a\sqrt{r}{\displaystyle \frac{a^2w_0r}{e}}+{\displaystyle \frac{a\left(24g^2a^232g^2a^2e3c^2\right)r^{3/2}}{32g^2e^2}}`$
$`{\displaystyle \frac{a^2\left(\left(16g^2a^232g^2a^2e15c^2\right)w_08g^2cae\right)r^2}{32g^2e^3}}+O(r^{5/2}),`$
$`N`$ $`=`$ $`12{\displaystyle \frac{e^2}{a^2r}}+c\sqrt{r}{\displaystyle \frac{a\left(3w_0c+ag^2e\right)r}{e}}+O(r^{3/2}).`$ (22)
This local solution contains three free parameters $`w_0,a,c`$. The singularity is of the Schwarzschild type, the seed mass $`m_0`$ is strictly positive. Two leading terms in expansion of ’kinetic’ (negative for $`N<0`$) and ’potential’ (positive) terms in $``$ cancel, so the leading behavior in singularity is
$$\frac{3\mathrm{c}}{4gr^{3/2}}.$$
(23)
We have tested for various $`r_h,g,n`$ that a continuation of the exterior black hole solutions under the horizon meets this generic asymptotic solution indeed. In all numerical experiments the function $`N`$ remained negative under the horizon and no internal Cauchy horizons were met. Typical global black hole solutions are shown on the Figs. 7,8,9 for $`r_h=1,g=1,n=1,2`$. The YM function $`w`$ outside the horizon has qualitatively the same behavior as in the regular case. Inside the horizon it remains perfectly smooth and tends to a finite limit $`w_00`$ at the singularity. Recall that for the ordinary quadratic YM lagrangian the function $`w`$ inside the horizon of the EYM black holes has a rather sophysticated behavior: while $`w`$ itself tends to a finite limit $`w_0`$ as well, its derivative exhibits a sequence of infinitely increasing absolute values at tiny intervals, whose length tends to zero . Such a behavior causes oscillations of the local mass $`m(r)`$ with an infinitely increasing amplitude. At the beginning of each oscillation cycle the metric function $`N(r)`$ takes values very close to zero (’almost’ Cauchy horizons), then an exponential growth of $`m(r)`$ starts. After a few oscillation cycles the maximal values of $`m`$ attained in subsequent cycles become of the googolplexus order, obviously lying beyond any classical bounds. Contrary to this, in the EBIYM case we observe a smooth $`m(r)`$ inside the horizon up to the singularity where $`m(r)`$ has a finite positive value (Fig. 8). Similarly, the second metric function $`\sigma `$ tends smoothly to a finite value at the singularity (Fig. 9)
We conclude with the following remarks. The BK solutions were found for the gravity coupled Yang-Mills theory with the usual quadratic lagrangian. Now they were shown to be a strong gravity limit of the gravity coupled BIYM theory which can be interpreted as an effective YM theory including string quantum corrections. Within this theory there exists a limit of decoupled gravity, in which qualitatively similar solutions continue to exist. This shows the way how the BK solutions can be incorporated into the (quantum) string theory.
A particularly interesting implication of this reasoning is the resolution of the problem of ‘violent’ oscillating singularities typical for hairy EYM black holes. Born-Infeld corrections perfectly regularize the behavior of the metric near the singularity. It is worth noting that the generic singularity is timelike, in conformity with the strong cosmic censorship hypothesis. In numerical experiments we did not observe internal horizons. In principle, such horizons could emerge in pairs, this question is worth to be investigated in more detail.
This work was supported in part by the RFBR grant 00-02-16306.
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# Bose-Einstein condensation of rubidium atoms in a triaxial TOP-trap
## 1 Introduction
Since the first observations of Bose-Einstein condensation (BEC) in dilute alkali gases , experimental as well as theoretical studies of degenerate quantum gases have been published at an astonishing rate . Far beyond the mere realization and detection of BEC, experimenters have investigated the static and dynamic properties of Bose-Einstein condensates and have gained considerable control over these macroscopic quantum objects, up to the point of creating coherent beams of matter waves - atom lasers, in other words. In spite of these early successes, experimental BEC is still a growing and thriving field, and much research needs to be done in order to test the vast number of theoretical predictions made in the last few years.
In this paper, we present the experimental apparatus used to create BECs of rubidium atoms in a triaxial time-orbiting-potential (TOP) trap . To the best of our knowledge, while the triaxial TOP trap has been used in BEC experiments on sodium , no previous application to rubidium has been reported. We describe in some detail the experimental parameters of our system and compare the performance of our apparatus with those of other groups using similar setups. Section 2 presents the experimental set-up, with emphasis on the original parts for the rubidium cooling and transfer between the two magneto-optical traps. Section 3 reports the parameters for the loading and evaporative cooling phases required to produce the condensate. Moreover, the gain in phase-space density achieved during the evaporation phases has been measured. In the following sections the results of various measurements on the condensate are reported. The final phase-space density, number of atoms and temperatures associated to the different condensates are presented. Furthermore, the expansion of the condensate following a switch-off of the magnetic trap has been studied and compared to different theoretical models. Finally, we describe different methods used for precise measurements of the magnetic fields. In this way, we obtained an accurate calibration which was needed as an input parameter for a theoretical model simulating the motion of the atomic cloud .
## 2 Experimental setup
Our experimental apparatus is based on a double-MOT system with a TOP-trap. The design of the vacuum system and the positioning of the coils are shown in figure 1. Owing to the arrangement of the quadrupole coils and the TOP-coils, our trap is triaxial without cylindrical symmetry. In the following, we give a brief overview of the specifications of our system.
Vacuum system: Our vacuum system is composed of two quartz cells connected by a glass tube of inner diameter $`12\mathrm{mm}`$ and length $`20\mathrm{cm}`$ (see figure 1). At the upper end of the glass tube, a graphite tube of length $`6\mathrm{cm}`$ and inner diameter $`5\mathrm{mm}`$ is inserted in order to enhance differential pumping. The upper cell is connected to a $`20\mathrm{ls}^1`$ ion pump, whereas the lower cell is pumped on by a $`40\mathrm{ls}^1`$ ion pump in conjunction with a Ti-sublimation pump. In this way, a pressure gradient is created between the two cells with the pressure in the upper cell being of the order of $`10^8\mathrm{Torr}`$ and that of the lower cell below $`10^{10}\mathrm{Torr}`$. The upper cell also contains two Rb dispensers (SAES getters) which we operate at $`3.0\mathrm{A}`$.
Lasers: The laser light for the upper and lower MOTs is derived from a MOPA (tapered amplifier) injected in turn by a $`50\mathrm{mW}`$ diode laser. Under typical conditions we extract up to $`320\mathrm{mW}`$ of useful output from this system, which is then frequency-shifted by acousto-optic modulators (AOMs) and mode-cleaned by optical fibres. In this way, we create up to $`60\mathrm{mW}`$ of laser power for the upper MOT and $`15\mathrm{mW}`$ for the lower MOT. The repumping light for both the upper and the lower MOT is derived from a $`75\mathrm{mW}`$ diode laser, yielding about $`9\mathrm{mW}`$ of total power after passage through all the optical elements. The injecting laser for the MOPA and the repumping laser are both injected by $`50\mathrm{mW}`$ grating stabilized diode lasers locked to Rb absorption lines.
Magnetic trap: Our TOP-trap consists of a pair of quadrupole coils capable of producing field gradients $`2b^{}`$ (along the symmetry axis) in excess of $`1000\mathrm{Gcm}^1`$ for maximum currents of about $`230\mathrm{A}`$, and two pairs of TOP-coils. The quadrupole coils are water-cooled and are oriented horizontally (along the $`x`$-axis, see fig. 1) about the lower glass cell of our apparatus. A combination of IGBTs and varistors is used for fast switching of the current provided by a programmable current source (HP6882) whilst protecting the circuits from damage due to high voltages induced during switch-off. In this way we are able to switch off the quadrupole field within less than $`50\mu \mathrm{s}`$ even for the largest field gradients. The rotating bias field $`B_0`$ is created by two pairs of coils: One (circular) pair is incorporated into the quadrupole coils, whilst the other (rectangular) pair is mounted along the $`y`$-axis. Within the adiabatic and harmonic approximations, for an atom with mass $`m`$ and magnetic moment $`\mu `$ this results in a triaxial time-orbiting potential $`V_{TOP}`$ given by
$$V_{TOP}=\frac{4\pi ^2m}{2}\left(\nu _x^2x^2+\nu _y^2y^2+\nu _z^2z^2\right)$$
(1)
with the following frequencies along the three axes of the trap in the ratio $`2:1:\sqrt{2}`$, as introduced in :
$`\nu _x`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\sqrt{{\displaystyle \frac{2\mu }{mB_0}}}b^{}`$ (2)
$`\nu _y`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\sqrt{{\displaystyle \frac{\mu }{2mB_0}}}b^{}`$ (3)
$`\nu _z`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\sqrt{{\displaystyle \frac{\mu }{mB_0}}}b^{}.`$ (4)
The anharmonic and gravitational effects neglected in this approximation will be discussed in section 5. The TOP-coils in our experiment can produce a bias field $`B_0`$ of up to $`30\mathrm{G}`$ and are operated at a frequency of $`10\mathrm{kHz}`$.
Imaging: Detection of the condensates is done by shadow imaging using a near-resonant probe beam. The absorptive shadow cast by the atoms is imaged onto a CCD-camera. With a camera pixel size of $`9\mu \mathrm{m}`$ and a magnification of about $`1.2`$, we achieve a resolution of just over $`7\mu \mathrm{m}`$. Most of our measurements are made after a few milliseconds of free fall of the released condensate, when typical dimensions are of the order of $`1030\mu \mathrm{m}`$.
## 3 Evaporative cooling and creation of the condensate
A typical experimental cycle from the initial collection of atoms in the upper MOT to the creation of a BEC is as follows. First, we load about $`5\times 10^7`$ $`\mathrm{Rb}`$ atoms into the lower MOT by repeatedly (up to 200 times) loading the upper MOT for $`160\mathrm{ms}`$ and then flashing on a near-resonant push beam that accelerates the atoms down the connecting tube. Once the lower MOT has been filled, a $`30\mathrm{ms}`$ compressed-MOT phase increases the density of the cloud, which is then cooled further to about $`15\mu \mathrm{K}`$ by a molasses phase of a few milliseconds. At this point, the molasses beams are switched off and an optical pumping beam is flashed on five times for $`20\mu \mathrm{s}`$, synchronized with the rotating bias field of $`1\mathrm{G}`$ to define a quantization axis, in order to transfer the atoms into the $`|F=2,m_F=2`$ Zeeman substate desired for magnetic trapping. Transfer into the TOP-trap is then effected by simultaneously switching on the rotating bias field (at its maximum value of about $`25\mathrm{G}`$) and the quadrupole field (at a value for the gradient chosen such as to achieve mode-matching between the initial cloud of atoms and the resulting magnetic trap frequencies). The subsequent evaporative cooling ramps for the quadrupole and the bias fields are shown schematically in figure 2. After an adiabatic compression phase, during which the quadrupole gradient is increased to its maximum value, the bias field amplitude is ramped down linearly. In this way, we perform circle-of-death evaporative cooling down to a bias field of around $`4\mathrm{G}`$. Next, at a constant bias field, we switch on a radio-frequency field, scanning its frequency exponentially from $`6.5\mathrm{MHz}`$ down to around $`3.2\mathrm{MHz}`$, which we find to be the threshold for condensation for our system. At threshold, we have up to $`3\times 10^4`$ atoms in the condensate/thermal cloud-conglomerate. Continuing rf-evaporation still further yields pure condensates of up to $`12\times 10^4`$ atoms with no discernible thermal fraction. The value for the bias field at which we switch from circle-of-death to rf-evaporation was chosen by maximizing the final condensate number. The approach to BEC is illustrated graphically in figure 3, in which the phase-space density is plotted as a function of the number of atoms.
Before imaging the condensate, we adiabatically change the trap frequency by ramping the bias field and the quadrupole gradient in $`200\mathrm{ms}`$. In this way, we can choose the frequency of the trap in which we wish to study the condensate. Thereafter, both fields are switched off on a timescale of $`2050\mu \mathrm{s}`$ for the quadrupole field and $`100200\mu \mathrm{s}`$ for the bias field. Owing to these short timescales, the change in trap frequency during the switching can essentially be neglected as typical oscillation periods in the trap are larger than $`10\mathrm{ms}`$. In fact, we were able to observe non-adiabatic motion of the trapped condensates at the frequency of the rotating bias field .
## 4 Experimental results
In the following, we briefly summarize some initial measurements made on the condensates obtained with our apparatus.
### 4.1 Evidence for condensation and condensate fraction
In order to find the threshold for condensation, the RF-frequency in the final evaporation step is lowered whilst monitoring the properties of the atom cloud (through shadow imaging after $`3\mathrm{ms}`$ of free expansion). At the threshold, the tell-tale signs of condensation, namely a sudden increase in peak density and the onset of a bimodal distribution, begin to appear. Figure 4 shows plots of the peak density normalized with respect to the number of atoms (which removes the considerable experimental jitter especially in the condensed regime) and the condensate fraction as a function of the final RF-frequency. The condensate fraction is determined from a bimodal fit to single pixel rows of the absorption picture, and it is evident in the two plots that condensation sets in at a final frequency of about $`3.2\mathrm{MHz}`$, corresponding to a temperature of $`365\mathrm{nK}`$ as calculated from the ballistic expansion of the cloud, and a peak density of $`5\times 10^{11}\mathrm{cm}^3`$. From this, we calculate a phase-space density of $`2.5`$ at the threshold, in agreement with theoretical predictions. Using the expression $`k_BT_0=\mathrm{}\overline{\omega }(N/\zeta (3))^{1/3}`$ (valid in the non-interacting approximation and with $`\overline{\omega }`$ equal to the geometric mean of the three trap frequencies) with $`N=10^4`$ atoms at the threshold , we find $`T_0400\mathrm{nK}`$ in good agreement with our observed threshold temperature.
We note here that, unlike in the case of a static trap, for a TOP-trap there is no strict proportionality between $`\nu _{cut}\nu _0`$ and $`k_BT_{cut}`$, where $`\nu _{cut}`$ is the frequency of the RF-field, $`\nu _0`$ is the resonance frequency at the bottom of the trap, and $`T_{cut}`$ is the equivalent cut temperature. A simple calculation considering the maximum instantaneous field at the resonance shell shows that, for low temperatures,
$$\nu _{cut}\nu _0=\frac{g_F}{h}\left(2k_BT_{cut}\mu B_0\right)^{1/2}.$$
(5)
This geometric average between the thermal cut energy $`k_BT_{cut}`$ and the magnetic energy in the bias field $`\mu B_0`$ of the TOP-trap leads to a considerably more accurate control of the cut energy in a TOP-trap. For instance, at a bias field of $`B_0=4\mathrm{G}`$, a frequency difference $`\nu _{cut}\nu _0`$ of $`350\mathrm{kHz}`$ corresponds to a cut energy $`T_{cut}`$ of only $`1.2\mu \mathrm{K}`$, whereas the same frequency difference in a static trap leads to $`T_{cut}=34\mu \mathrm{K}`$.
### 4.2 Free expansion of the condensate
One way of obtaining information on the properties of a Bose-Einstein condensate is to investigate its behaviour after it is released from the trap. Its subsequent evolution is then monitored by taking absorption images after a variable time-of-flight. The results of such measurements on a condensate released from a trap with $`\nu _z=363\mathrm{Hz}`$ are shown in figure 5. Theoretically, the expansion of a condensate has been investigated by several authors, and analytical expressions for the condensate width and its aspect ratio as a function of time can be found in special cases. Figure 5 shows the predictions of a model based on the Thomas-Fermi approximation , in which the energy of the condensate is dominated by the mean-field interaction between the atoms, as well as the theoretical expansion of a ground-state harmonic oscillator wavefunction, for which interactions are neglected entirely. Clearly, our experimental data agree with neither of these two extremes. This is to be expected, as the sizes of our condensates, with typically a few thousand atoms in a pure condensate, are rather small and therefore do not fully satisfy the conditions for a Thomas-Fermi treatment. It is, therefore, necessary to compare our data with a numerical integration of the full Gross-Pitaevskii equation. The results of such an integration are also plotted in figure 5. As expected, they lie between the two extreme models and fit our data reasonably well. It is clear, however, that our condensate number is so low that the interaction term in the Gross-Pitaevskii equation is almost negligible and the numerical results are close to the pure harmonic oscillator case.
## 5 Calibration of the magnetic fields
In many applications of magnetic traps, it is sufficient to describe the trap by its characteristic frequencies for dipolar oscillations of atomic clouds. In such a measurement, one applies a magnetic field along a chosen axis for a short time, thus giving a kick to the (initially stationary) atomic cloud, and monitors the subsequent oscillations of the atoms. With a judicious choice of the points in time at which the position of the cloud is sampled, one can achieve frequency measurements with uncertainties well below the percent level. Deducing absolute values of the magnetic field gradient and the bias field from these measurements with similar accuracy, however, is not so straightforward. The main incentive for us to accurately measure these absolute values was that we needed them as input parameters for numerical simulations of non-adiabatic motion in the TOP-trap . In the following, we shall briefly describe several methods we used to measure absolute values for both the quadrupole gradient and the bias field and indicate the uncertainties associated with these measurements. For the most part, the measurements were carried out with condensates, which facilitated the determination of the position of the atomic cloud.
In the first method, we measure the vibrational frequencies $`\stackrel{~}{\nu _x}`$ and $`\stackrel{~}{\nu _z}`$ along the $`x`$\- and $`z`$-axes, respectively, exciting the dipolar modes along these two directions simultaneously. In order to be able to use theoretical formulas derived in the harmonic approximation taking into account the effect of gravity, we have calculated the anharmonic corrections up to fourth order, including cross-terms, following the scheme presented by Ensher . The results reported in Appendix A allow us to deduce from our measured frequencies the corresponding values in the harmonic limit (equivalent to infinitesimal oscillation amplitudes; typical amplitudes in our experiment are between $`20\mu \mathrm{m}`$ and $`60\mu \mathrm{m}`$.). Those anharmonic corrections can be up to $`1\%`$ of the measured values and are, therefore, essential if an accuracy in the magnetic field below the percent level is desired. The quadrupole gradient can be calculated directly from the ratio $`\stackrel{~}{\nu _x}/\stackrel{~}{\nu }_z`$ given by in the harmonic approximation with the gravitational corrections by
$$\frac{\stackrel{~}{\nu _x}}{\stackrel{~}{\nu _z}}=\sqrt{2}\sqrt{\frac{1+\eta ^2}{1\eta ^2}}$$
(6)
Here, $`\eta `$, defined by
$$\eta =\frac{\mu b^{}}{mg}$$
(7)
measures the ratio of magnetic and gravitational forces along the $`z`$-axis. It is interesting to note that in the triaxial TOP the gravity corrections are equal to those derived for a cylindrically symmetric TOP-trap . Re-substitution of the value for $`b^{}`$ thus retrieved along with either of the two frequencies into the expression for $`\stackrel{~}{\nu _x}`$ or $`\stackrel{~}{\nu _z}`$ then yields a value for $`B_0`$. For instance, $`\stackrel{~}{\nu _z}`$ is given by
$$\stackrel{~}{\nu _z}=\frac{1}{2\pi }\sqrt{\frac{\mu }{mB_0}}b^{}\left(1\eta ^2\right)^{3/4}.$$
(8)
In order to check the obtained values for $`B_0`$ and $`b^{}`$ by independent methods not relying on the calculated frequencies for a TOP-trap, we use two separate strategies. In one method, the quadrupole gradient is measured by first trapping and evaporatively cooling atoms in the presence of both the quadrupole and the bias fields. Then, the bias field is switched off, which shifts the centre of the quadrupole potential with respect to the TOP-potential. The quadrupole gradient is subsequently determined by measuring the acceleration of the atoms and subtracting the acceleration due to gravity. In this way, $`b^{}`$ can be determined with a relative error of less than $`1\%`$. An independent measurement of the bias field $`B_0`$ is made by switching off the quadrupole field after the atoms have been cooled in the TOP whilst leaving the bias field on. A short ($`100500\mu \mathrm{s}`$) RF-pulse is then applied to the atoms at a given frequency, and the number of atoms remaining in the original trapped state is measured after turning the quadrupole field back on (about $`1\mathrm{ms}`$ after switching it off). When the frequency of the RF-pulse matches the Zeeman-splitting due to the bias field, atoms are transferred into untrapped Zeeman-substates and hence lost from the trap. Using this method, we found two different values of the RF-pulses for which atoms were lost from the trap, indicating that there is a slight asymmetry between the magnetic fields produced by the two pairs of TOP-coils. Measuring $`B_0`$ with this method proved to be less reliable than with the method described above, but yielded the same value for the bias field to within $`5\%`$.
## 6 Condensate numbers in TOP-traps
In our experimental apparatus, we obtain condensates containing up to a few $`10^4`$ atoms, starting from typical MOT numbers of about $`5\times 10^7`$. Extrapolating this linearly, one would expect to achieve condensate numbers of up to $`10^6`$ for an initial number of $`5\times 10^9`$ atoms in the MOT. In the literature, however, one typically finds reports of some $`10^5`$ atoms in the condensate under such circumstances. In figure 6 we have plotted typical figures for the MOT and the condensate numbers for a few groups using rubidium TOP-traps. Evidently, the reported condensate numbers do not scale linearly with the MOT numbers. Instead, they can be fitted roughly by a square-root law. Varying the MOT numbers in our own experiment, we find a similar behaviour on a smaller scale. We discovered this when trying to increase the size of our condensates and found that the main limiting factor comes from the compression phase after loading the magnetic trap. Above a certain number of atoms loaded into the MOT, we saw next to no increase in the atom number after compression (or, for that matter, in the condensate) when increasing the initial number of atoms. As in , we attribute this to an unfavourable ratio of the size of the initial cloud and the circle-of-death radius. When the cloud becomes too big, the circle-of-death cuts into it during compression and thus any increase in the atom number is eaten up by this cutting. This may be a limiting mechanism for most groups and could explain the law of diminishing returns that is evident in figure 6. In this context it is interesting to note that, for instance, the JILA group uses a much higher bias field ($`50\mathrm{G}`$) than most other groups and achieves a much better transfer efficiency from the MOT to the condensate , obtaining condensates of $`10^6`$ atoms for initial numbers of the order of $`2\times 10^8`$. Although this may suggest that a larger bias field is the answer, it is not clear whether there are other effects that limit the transfer efficiencies achievable in TOP-traps.
## 7 Conclusion
We have presented the results of preliminary measurements on Bose-Einstein condensates of rubidium atoms obtained in a triaxial TOP-trap. Our experimental data for the condensation threshold and the free expansion of the condensate agree well with theoretical predictions. Increasing the number of atoms in our condensates will allow us to further improve on the quality of our data and investigate the properties of the condensates in more detail.
## Acknowledgments
O.M. gratefully acknowledges financial support from the European Union (TMR Contract-Nr. ERBFMRXCT960002). This work was supported by the INFM ’Progetto di Ricerca Avanzata’ and by the CNR ’Progetto Integrato’. The participation of G. Memoli and D. Wilkowski in the early stages of the experiment is gratefully acknowledged. The authors are grateful to R. Mannella for the numerical integration of the Gross-Pitaveskii equation and to M. Anderlini for help in the calculation of the anharmonic corrections.
## Appendix A Anharmonic corrections in the TOP-trap
For our calibration measurements, we deduced the frequencies in the harmonic limit from the anharmonic corrections in the triaxial TOP-trap. Terms containing the amplitude of the oscillations along the $`y`$-axis have not been calculated as we do not excite oscillations along that direction, but can be obtained in the same manner. The expressions for the frequencies along the axis $`i`$ ($`i=x,z`$) are then
$$\nu _i^{\mathrm{anh}}=\nu _i+\mathrm{\Delta }\nu _i;\mathrm{\Delta }\nu _i=\left(\frac{b^{}}{B_0}\right)^2\underset{j}{}\alpha _{ij}a_j^2$$
(9)
where $`\nu _i`$ is the frequency in the harmonic approximation, as given by Eqs. 4 and $`a_j`$ is the amplitude of the oscillation in the $`j`$-direction. The elements $`\alpha _{ij}`$ of the anharmonic correction matrix are given by
$`\alpha _{xx}={\displaystyle \frac{\nu _x}{4}}\left[6{\displaystyle \frac{1\eta ^2}{1+\eta ^2}}\left(23\eta ^2{\displaystyle \frac{15}{8}}(1\eta ^2)^2\right){\displaystyle \frac{\eta ^2(13\eta ^2)^2}{18(1+\eta ^2)}}\right]`$ (10)
$`\alpha _{xz}={\displaystyle \frac{\nu _x}{4}}\left[{\displaystyle \frac{7\eta ^2(1815\eta ^2)}{12}}{\displaystyle \frac{9\eta ^2(3\eta ^21)(14+8\eta ^2)8\eta ^2}{36(7+9\eta ^2)}}\right]`$ (11)
$`\alpha _{zx}={\displaystyle \frac{\nu _z}{2}}\left[{\displaystyle \frac{13\eta ^2}{2}}{\displaystyle \frac{2\eta ^2(13\eta ^2)}{9(3+5\eta ^2)}}+{\displaystyle \frac{73\eta ^215\eta ^2(1\eta ^2)}{12}}\right]`$ (12)
$`\alpha _{zz}={\displaystyle \frac{\nu _z}{2}}\left[{\displaystyle \frac{3}{8}}(1\eta ^2)(15\eta ^2){\displaystyle \frac{15}{8}}\eta ^2(1\eta ^2)\right].`$ (13)
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# EMBEDDED, SELF-GRAVITATING EQUILIBRIA IN SHEETLIKE AND FILAMENTARY MOLECULAR CLOUDS
## 1. INTRODUCTION
The present observational picture of star-forming regions conveys cloud structures considerably more varied than any one theoretical scenario can explain. Nevertheless, there are certain points of correspondence, when one restricts attention to the molecular, self-gravitating component of the interstellar medium. On scales of tens to hundreds of parsecs, both atomic and molecular gas clouds appear shell-like and filamentary (Scalo 1985; Kulkarni & Heiles 1988). Shells may be formed by several distinct dynamical processes: cloud-cloud collisions (Smith 1980), compressional shock waves from supernovae or OB stars (McCray & Kafatos 1987), and large-scale shocks associated with spiral density waves (Roberts 1969) are but a few possibilities. Filamentary structure persists down to scales of several parsecs, to the regime of individual molecular clouds \[Loren 1989, Nozawa et al 1991 ($`\rho `$ Oph); Heyer et al 1987, Onishi et al 1996 (Taurus); Bally et al 1987, Tatematsu et al 1993 (Orion)\]. Each filament typically contains several distinct subcondensations in close proximity to each other, some of which harbor infrared continuum sources (Onishi et al 1998). In some cases, the embedded clumps are spaced quite regularly along the filament (Schneider & Elmegreen 1979; Dutrey et al 1991). Thus, several authors have speculated that the formation of star-forming clumps proceeds via a hierarchical fragmentation process, in which filaments are formed out of larger structures, and then clumps out of the filaments (Schneider & Elmegreen 1979; Gaida, Ungerechts, & Winnewisser 1984; Hanawa et al 1993; Fiege & Pudritz 2000$`a`$).
Focusing now on the cores themselves, maps in dense tracers such as NH<sub>3</sub> and CS display roughly elliptical intensity contours, with a mean apparent major-to-minor axis ratio of around 2 (Jijina, Myers, & Adams 1999). Statistical arguments applied to the distribution of measured axial ratios for several surveys have prompted some authors to conclude that the cores are more likely to be intrinsically prolate than oblate (David & Verschueren 1987; Myers et al 1991; Ryden 1996). Other statistical arguments imply that the observed elongation is unlikely to be a result of star formation or outflows (Myers et al 1991). Coupled with the additional result that the majority of cores are near virial equilibrium (Jijina et al 1999), this argues against a wholly dynamical origin for prolateness, since the implied lifetimes are then so short that observations of cores without embedded stars would be exceedingly rare (the latter represent about one-half of the ammonia cores detected in the Taurus region, and a somewhat smaller fraction in Ophiuchus and Orion). Rather, it appears more likely that cores are at least quasi-equilibrium structures, and that their shapes therefore offer some clue to the forces responsible for their formation.
If this reasoning is correct, then it leads to a formidable crisis in our current theoretical picture of cloud equilibria, which typically envisions star forming cores as self-gravitating clumps bounded by a zero density, constant pressure—and so high temperature—medium (e.g., McKee et al 1993). That is, it is difficult to conceive of prolate, quasi-equilibrium, gaseous cores as purely isolated structures. Non-gravitational forces, such as rotation or magnetic fields, are not likely to aid in maintaining prolate equilibria, although their role in oblate structures has been made abundantly clear.<sup>1</sup><sup>1</sup>1It is in fact possible to construct isolated, prolate magnetic clouds, but these possess magnetic field structures which are highly unusual and, in any case, lack direct observational justification. See Fiege & Pudritz (2000$`b`$) and Curry & Stahler (2000) for examples. The absence in the literature of even a single, physically acceptable, gaseous, prolate equilibrium solution suggests that a fresh theoretical approach is necessary. We examine a scenario in this paper whereby core morphology is directly attributable to a fragmentation process. As we show in §2 and 3, the existence of prolate structures within cylindrical filaments may be understood in exactly the same way as the existence of the filaments themselves: as a result of the fragmentation of the parent cloud and nonlinear growth of the fragments. What has apparently escaped attention until now, and what we demonstrate explicitly, is that accessible, long-lived states exist wherein cores and their extensive filamentary envelopes occupy the same hydrostatic structure. We now proceed to outline a particular fragmentation hierarchy that might be responsible for such structures.
## 2. NONLINEAR FRAGMENTATION OF AN ISOTHERMAL LAYER
To motivate the new solutions presented in the following section, we begin by outlining a possible formation mechanism for filamentary structures in the molecular interstellar medium. By drawing attention to selected existing results in the literature, we hope to provide a more compelling argument for a particular fragmentation hierarchy alluded to previously by several authors, notably Schneider & Elmegreen (1979) and Larson (1985).
### 2.1. Analytic Solutions
While not realistic in detail, a large, shell-like structure in the interstellar medium may be approximated on sufficiently small scales by a planar, self-gravitating layer, or “sheet” (Elmegreen & Elmegreen 1978). We restrict consideration throughout to an isothermal, self-gravitating gas, whose equilibrium is governed by the Lane-Emden equation,
$$^2\psi =\frac{^2\psi }{x^2}+\frac{^2\psi }{y^2}+\frac{^2\psi }{z^2}=4\pi G\rho _{0,s}e^{\psi /a^2},$$
(1)
where $`\psi `$ is the gravitational potential, $`a`$ is the constant sound speed, and $`\rho _{0,s}`$ is a constant. The equation has been written in cartesian coordinates $`(x,y,z)`$. When $`\psi `$ varies solely in the $`y`$-direction, i.e. $`\psi =\psi _{1\mathrm{D}}(y)`$, and under the usual boundary conditions $`\psi _{1\mathrm{D}}(y=0)=0,(d\psi _{1\mathrm{D}}/dy)_{y=0}=0`$, an exact solution is known (Spitzer 1942):
$$\rho _{1\mathrm{D}}(y)=\rho _{0,s}\mathrm{exp}[\psi _{1\mathrm{D}}(y)/a^2]=\rho _{0,s}\mathrm{sech}^2(y/\mathrm{}_0),$$
(2)
where the scale height $`\mathrm{}_0`$ is defined by $`\mathrm{}_0a/(2\pi G\rho _{0,s})^{1/2}`$. The layer is unbounded in the $`y`$-direction; solutions truncated by an external pressure at constant $`y`$ may also be constructed (e.g., Elmegreen & Elmegreen 1978). However, as these introduce an additional parameter—namely, the external pressure—into the problem, we do not consider them here. The stability of the unbounded solution was examined by Ledoux (1951), who found gravitational instability for infinitesimal sinusoidal perturbations of wavelength exceeding $`\lambda _{x,\mathrm{cr}}=2\pi \mathrm{}_0`$. The maximum growth rate occurs at $`\lambda _{x,\mathrm{MGR}}=2.24\lambda _{x,\mathrm{cr}}=14.1\mathrm{}_0`$. The choice of the $`x`$-direction for the perturbation wavevector here is arbitrary; the same results hold for perturbations in the $`z`$-direction only.
Subsequently, Schmid-Burgk (1967) (hereafter S-B) presented a remarkable two-dimensional (2D) solution of equation (1) (independently discovered, in a hydrodynamic context, by Stuart 1967):
$$\rho _{2\mathrm{D}}(x,y)=\rho _{0,s}\mathrm{exp}\left[\frac{\psi _{2\mathrm{D}}(x,y)}{a^2}\right]=\frac{\rho _{0,s}(1A^2)}{[\mathrm{cosh}(y/\mathrm{}_0)A\mathrm{cos}(x/\mathrm{}_0)]^2}.$$
(3)
The constant $`A,0<A<1`$ describes the amplitude of spatially periodic variations along one direction parallel to the layer (here chosen as $`x`$), with the same critical wavelength, $`\lambda _{x,\mathrm{cr}}`$, as found by Ledoux. As noted by S-B, this result shows that the latter mode is not restricted to infinitesimal amplitudes. Indeed, solutions of all amplitudes $`A`$ are contained in equation (3).<sup>2</sup><sup>2</sup>2Note also that, like other solutions of the Lane-Emden equation, S-B’s solution is homology invariant (Chandrasekhar 1958); i.e. the potential $`\psi _{2\mathrm{D}}^{}\psi _{2\mathrm{D}}(Cx,Cy)2\mathrm{log}C`$, $`C=`$ constant, is also a solution of equation (1). It is noteworthy that the periodicity of this exact, nonlinear solution does not correspond to that of the fastest-growing mode in the linearly unstable layer, but rather to the critically unstable mode. This feature is discussed further in §6.
As $`A0`$, equation (3) reduces to Spitzer’s solution. In the limit $`A1`$, the equilibrium<sup>3</sup><sup>3</sup>3The term “equilibrium” is used here in its strictly limited sense: i.e., a particular equilibrium may be either stable or unstable to perturbations in a direction in which the solution is translation-invariant (in this case, the $`z`$-direction). corresponds to a series of parallel cylindrical “fragments,” infinitely extended in the $`z`$-direction, with density maxima separated by a distance $`\lambda _{x,\mathrm{cr}}`$. For general $`A`$, the fragments are embedded, nearly elliptic cylinders with equidensity contours of eccentricity $`(1A)^{1/2}`$ (see Figure 1 of S-B). The density maxima and minima along the $`x`$-axis (each spaced at an interval of $`\mathrm{\Delta }x=2\pi \mathrm{}_0`$), are given by
$$\rho _c=\rho _{0,s}\left(\frac{1+A}{1A}\right)\mathrm{and}\rho _s=\rho _{0,s}\left(\frac{1A}{1+A}\right),$$
(4)
respectively.
### 2.2. Numerical Solution
There is reason to suspect that analogues of the S-B solution exist in other geometries, in particular in cylindrical symmetry, a case we shall focus on in the following section. Unfortunately, to our knowledge, no other embedded solutions are known in analytic form. However, numerical techniques can be used. As a test of one such scheme used later in this paper, we attempted to find the S-B solution using the self-consistent field method (e.g. Tassoul 1978). That is, equation (1) was solved iteratively on a two-dimensional grid, subject to the boundary conditions
$`{\displaystyle \frac{\psi ^{}}{x^{}}}|_{x^{}=0}`$ $`=`$ $`{\displaystyle \frac{\psi ^{}}{y^{}}}|_{y^{}=0}=0,`$ (5)
$`\psi ^{}(0,0)=\mathrm{ln}\rho ^{}(0,0),`$ $`\mathrm{and}`$ $`{\displaystyle \frac{\psi ^{}}{x^{}}}|_{x^{}=X^{}}=0.`$ (6)
Here a prime indicates a dimensionless quantity—e.g., $`\psi ^{}\psi /a^2`$—and $`0x^{}X^{},0y^{}Y^{}`$ is the extent of the computational region. The nondimensionalization of length will be discussed below; for now it may be assumed arbitrary.
The first of conditions (5) ensures that the $`x`$-component of the gravitational field vanishes on the $`y`$-axis, while the second imposes reflection symmetry about the $`x=0`$ plane. Conditions (6) are consistent with the Spitzer solution, although neither constrains the numerical solution to be identical to the former in any particular limit. In addition, these conditions allow for solutions with a periodic structure in $`x`$. In an attempt to find these, we took $`\psi _{1\mathrm{D},0}^{}(y^{})=\mathrm{ln}\rho _{1\mathrm{D},0}^{}(y^{})`$ (see eq. 2) as an initial guess, and added a perturbation of the form
$$\delta \psi ^{}(x^{},y^{})=ϵ\mathrm{cos}(px^{}/X^{})\mathrm{sin}(py^{}/Y^{}),$$
where $`ϵ`$ is a small, constant amplitude and $`0<p<2\pi `$. For convergence of the numerical code we required $`1\psi ^{(n)}/\psi ^{(n+1)}<\delta =0.005`$, where $`\psi ^{(n)}`$ is the value of $`\psi ^{}`$ at the $`n^{\mathrm{th}}`$ iteration. For $`ϵ<0.005`$, the method converged immediately to the Spitzer solution (2). For any larger $`ϵ`$, however, a unique, 2D structure resulted. The equilibria so obtained constitute a family of solutions in the single parameter $`X^{}`$, such that the density contrast $`\rho _c/\rho _s`$ increases monotonically with increasing $`X^{}`$. Here $`\rho _s\rho (X,0)`$. The vertical extent $`Y^{}`$ was chosen sufficiently large that the effect of different $`Y^{}`$ on 2D solutions of the same $`X^{}`$ was negligible.
In order to compare these solutions to the analytic one of S-B, we need to specify the nondimensionalization. This is, in fact, already implicit in the solution method. At each iteration, the Lane-Emden equation is solved for $`\psi ^{}(x^{},y^{})`$ with the source term
$$\rho ^{(n)}(x^{},y^{})=\mathrm{exp}[\psi ^{(n)}(0,0)\psi ^{(n)}(x^{},y^{})].$$
Thus, $`\rho ^{}(x^{},y^{})`$ is normalized with respect to its central value at each iteration; i.e. $`\rho ^{(n)}(0,0)=1`$ (this is made explicit in the first of equations 6). Consequently, once an exact, 2D equilibrium is found, all “memory” of the original density scale $`\rho _{0,s}`$ is lost. The characteristic density of the solution is $`\rho _c=\rho (0,0)`$, and thus the corresponding unit length is $`\mathrm{}_ca/(2\pi G\rho _c)^{1/2}`$, not $`\mathrm{}_0\rho _{0,s}^{1/2}`$ as in the analytic solution. Hence, the appropriate nondimensionalization is specified a posteriori as
$$(x^{},y^{},z^{})(x/\mathrm{}_c,y/\mathrm{}_c,z/\mathrm{}_c).$$
(7)
An exact relation between the parameter $`X^{}`$ in the numerical problem, and the amplitude $`A`$ in the analytic solution, may now be derived.<sup>4</sup><sup>4</sup>4I thank Dean McLaughlin for pointing out a simplified derivation. Solving equation (3) for $`\psi _{2\mathrm{D}}(x,y)`$ and substituting the result into the second of equations (6) gives
$$X^{}=\frac{X}{\mathrm{}_c}=\pi \left(\frac{\rho _c}{\rho _s}\right)^{1/4}=\pi \left(\frac{1+A}{1A}\right)^{1/2}.$$
(8)
The parameter $`X^{}`$ represents one-half of the spacing between the embedded fragments; i.e., $`\lambda _x^{}=2X^{}`$. In the analytic solution, the spacing of the fragments is constant, i.e. $`\mathrm{\Delta }x/\mathrm{}_0=\lambda _{x,\mathrm{cr}}/\mathrm{}_0=2\pi `$. There the fragment spacing stays the same while the amplitude $`A`$ varies. Just as analytic solutions exist for all $`A`$, numerical solutions exist for all values of $`X^{}>\pi `$.
Figure 1 of S-B (1967) shows a solution with $`A=0.17`$, whose corresponding density contrast is $`\rho _c/\rho _s=1.987`$. Inserting this $`A`$ into equation (8), one finds $`X^{}=3.730`$. Using this $`X^{}`$ as input for the numerical method then yields a solution with $`\rho _c/\rho _s=1.981`$, a difference of 0.3 percent from S-B’s analytic result. This solution is displayed in Figure 1. The discrepancy between the analytic and numerical solutions remains small as $`X^{}`$ increases above 3.73, and is still less than 1 percent for $`X^{}=22,\rho _c/\rho _s2400`$. At smaller $`X^{}`$, the numerical and analytic results slowly diverge, with the relative error in $`\rho _c/\rho _s`$ reaching 5 percent at $`X^{}=3.32`$, and nearly 20 percent at $`X^{}=\pi `$. At the latter value, equation (8) predicts $`A=0`$ and $`\rho _c/\rho _s=1`$. Presumably, the method is not sufficiently sensitive to the slight density contrasts present in these small amplitude ($`A<0.05`$) solutions to render an accurate result. On the other hand, for $`X^{}<3.125`$, the numerical method converges to the Spitzer solution, with a slightly non-uniform density along the midplane (e.g., $`\rho _c/\rho _s0.99`$ at $`X^{}=3.12`$). Overall, these results confirm the validity of the numerical method, with the nondimensional length and density chosen as in equation (7). We may expect the same to be true for other solutions found by the same technique, irrespective of whether an analytic solution is known.
We conclude our discussion of the S-B solution by noting that related states have in fact appeared in the literature, in the context of either linear perturbation theory (Miyama, Narita, & Hayashi 1987$`a`$) or time-dependent nonlinear calculations (Miyama, Narita, & Hayashi 1987$`b`$). The basic character of the solution persists even in the presence of a ($`z`$-independent) magnetic field (Fleischer 1998; Nagai, Inutsuka, & Miyama 1998), or when the layer is truncated by an external medium. The pressure-bounded, magnetized layer has additional modes of fragmentation available to it, as shown by Nagai et al (1998). However, the significance of the basic result—namely, that the embedded filaments are in fact exact, equilibrium solutions—is rarely emphasized.
## 3. NONLINEAR FRAGMENTATION OF AN ISOTHERMAL FILAMENT
### 3.1. Equilibrium and Stability
In cylindrical coordinates $`(r,\varphi ,z)`$, the isothermal Lane-Emden equation reads
$$^2\psi =\frac{1}{r}\frac{}{r}\left(r\frac{\psi }{r}\right)+\frac{^2\psi }{z^2}=4\pi G\rho _{0,f}e^{\psi /a^2},$$
(9)
where $`\rho _{0,f}`$ is a constant. Here the $`z`$-axis is taken to coincide with the principal axis of the cylinder, and we have taken $`/\varphi =0`$, restricting consideration to axisymmetric solutions. When $`\psi `$ depends only on $`r`$, i.e. $`\psi =\psi _{1\mathrm{D}}(r)`$, and under the boundary conditions $`\psi _{1\mathrm{D}}(r=0)=0,(d\psi _{1\mathrm{D}}/dr)_{r=0}=0`$, Stodolkiewicz (1963) and Ostriker (1964) derived the following exact solution of equation (9) (hereafter the S-O solution):
$$\rho _{1\mathrm{D}}(r)=\rho _{0,f}\mathrm{exp}\left[\frac{\psi _{1\mathrm{D}}(r)}{a^2}\right]=\frac{\rho _{0,f}}{(1+r^2/l_0^2)^2},$$
(10)
where $`l_0(2a^2/\pi G\rho _{0,f})^{1/2}`$ is the cylindrical scale radius. Interestingly, solution (10) reduces to that of S-B as $`A1`$ in the latter; details may be found in the Appendix. This suggests that equilibria quite similar to the S-O solution may result directly from the fragmentation of an isothermal layer.
The solution (10) decreases as $`r^4`$ at large $`r`$; as in the planar case, solutions truncated by an external pressure at constant $`r`$ have also been considered (e.g., Inutsuka & Miyama 1997; Fiege & Pudritz 2000$`a`$). The maximum mass per unit length of the filament described by equation (10) is<sup>5</sup><sup>5</sup>5In the S-B solution, the total mass per unit $`z`$-length of each “cell”, $`\pi x/l_0\pi `$, is also $`2a^2/G`$.
$$\mu _{\mathrm{max}}_0^{\mathrm{}}2\pi \rho _{1\mathrm{D}}(r)r𝑑r=\frac{2a^2}{G}.$$
(11)
Isothermal cylinders with density distributions having $`\mu >\mu _{\mathrm{max}}`$ are unstable to radial collapse, while those with $`\mu <\mu _{\mathrm{max}}`$ expand radially outward, unless confined by an external pressure (e.g., Inutsuka & Miyama 1992). The stability of solution (10) to axisymmetric, linear perturbations was examined by Stodolkiewicz (1963), who found instability for perturbations of wavelength exceeding $`\lambda _{z,\mathrm{cr}}=3.94l_0`$. The maximum growth rate occurs at $`\lambda _{z,\mathrm{MGR}}=1.98\lambda _{z,\mathrm{cr}}=7.82l_0`$ (Nagasawa 1987). Further analysis has been carried out by Inutsuka & Miyama (1997) in the nonlinear regime of perturbation growth. However, the nonlinear resolution of the instability remains an open question. That is, what is the final outcome of the fragmentation instability in an isothermal cylinder?
### 3.2. Numerical Solution
We undertake a numerical investigation of the question posed above, using an identical technique to that described in §2.2. That is, rather than following the time evolution of a particular unstable state, we instead search for an exact, static solution having a 2D structure. The boundary conditions on the potential are now
$`{\displaystyle \frac{\psi ^{}}{r^{}}}|_{r^{}=0}`$ $`=`$ $`{\displaystyle \frac{\psi ^{}}{z^{}}}|_{z^{}=0}=0,`$ (12)
$`\psi ^{}(0,0)=\mathrm{ln}\rho ^{}(0,0),`$ $`\mathrm{and}`$ $`{\displaystyle \frac{\psi ^{}}{z^{}}}|_{z^{}=Z^{}}=0,`$ (13)
where $`0r^{}R^{},0z^{}Z^{}`$ is the size of the computational region. The nondimensionalization is the same as that used in §2.2, except with $`\mathrm{}_c`$ replaced by $`l_c(2a^2/\pi G\rho _c)^{1/2}`$, where $`\rho _c\rho (0,0)`$. Again we took the equilibrium potential, $`\psi _{1\mathrm{D}}^{}(r^{})=\mathrm{ln}\rho _{1\mathrm{D}}^{}(r^{})`$ as an initial guess, and added a perturbation of the form
$$\delta \psi ^{}(r^{},z^{})=ϵ\mathrm{sin}(pr^{}/R^{})\mathrm{cos}(pz^{}/Z^{}),$$
with $`0<p<2\pi `$. For any $`ϵ>0.005`$, a family of unique 2D structures was again found, now parameterized by $`Z^{}`$. The corresponding $`R^{}`$ was chosen sufficiently large that its effect on the solutions was negligible. In practice, choosing $`R^{}`$ to be approximately twice the radius of the “tidal lobe”—i.e., the last closed isodensity contour, of density $`\rho _s\rho (0,Z^{})`$—was sufficient.
The density contrast between the center and the tidal lobe, $`\rho _c/\rho _s`$, is a monotonically increasing function of $`Z^{}`$. Models with $`1.2<\rho _c/\rho _s<10^3`$ were generated, corresponding to $`2.1<Z^{}30`$. A few representative equilibria are displayed in Figure 2. At small $`Z^{}`$, the equidensity contours are highly prolate (Fig. 2$`a`$); as $`Z^{}`$ increases, the interior contours become nearly spherical (Fig. 2$`d`$). Note that we were able to find 2D equilibria only down to a minimum $`Z^{}2.1`$; below this value, convergence was not obtained until $`Z^{}=1.8`$. For all $`Z^{}1.8`$, the method converged to the S-O solution, with a slightly non-uniform density along the axis of symmetry (e.g., $`\rho _c/\rho _s=0.997`$ at $`Z^{}=1.8`$). As we discuss below in §6, these results strongly suggest a fragmentation scale $`\lambda _{\mathrm{cr}}`$, as in the planar layer.
Physically, the quantity $`2Z^{}`$ may be interpreted as the separation between any two clumps in a linear chain of identical condensations (e.g. Lizano & Shu 1989; Fiedler & Mouschovias 1992). This feature thus mimics an observed property of star-forming environments: namely, that pre-stellar cores are rarely found in isolation. In §4.2, we examine how the properties of the condensations depend upon $`Z^{}`$, and how the latter may be constrained by observations.
## 4. PHYSICAL PROPERTIES OF THE FRAGMENTS
### 4.1. Virial Theorem Analysis
Insight into the global properties of equilibria may be obtained from the scalar virial theorem. In the present context of embedded cores, however, care must be taken regarding the region of application. For an axisymmetric cloud in equilibrium, we have (McKee et al 1993)
$$2U+\mathrm{\Pi }+W=0,$$
(14)
where
$`U`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle _V}P𝑑V={\displaystyle \frac{3}{2}}a^2M,`$
$`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle _S}P𝐫\mathrm{𝐝𝐒},`$
and
$$W=_V\rho 𝐫\psi dV$$
are the relevant thermal, compressive, and gravitational energies. These integral quantities are usually summed over the volume $`V`$ enclosed within an arbitrary closed surface $`S`$. In the present case, one might choose for $`S`$ any closed isobar. Then the thermal and compressive terms are easily calculated. However, for any choice of $`S`$ the gravitational energy $`W`$ is incomplete because part of the gravitational field arises from outside the region considered. This is in keeping with the nature of embedded equilibria. Indeed, each equilibrium found by the numerical method of the previous section comprises an entire “cell,” i.e., the region $`0rR,\mathrm{\hspace{0.33em}0}\varphi 2\pi ,ZzZ`$, so that the above choice of $`S`$ is inappropriate in any case. The virial theorem may only be consistently applied on the same region. Since the gravitational force vanishes on the plane circular section $`z=\pm Z`$ (by the second of the boundary conditions eq. 13), the corresponding energy $`W`$ summed over the entire cell between $`z=Z`$ and $`z=+Z`$ is therefore complete.
Figure 3 shows the behavior of individual terms in the virial equation (14), expressed in dimensionless units, along with their sum, $`\mathrm{\Sigma }2U+\mathrm{\Pi }+W`$, as one proceeds along a sequence of increasing central density $`\rho _c/\rho _s`$. Here, $`\rho _s`$ is used as a convenient reference density only; the virial terms have been calculated for the entire cell, not just for the region inside the tidal lobe. The figure shows that the compressive term $`\mathrm{\Pi }`$ has a magnitude of only 40 to 50 percent of the gravitational term $`W`$ over most of the sequence. This may be compared with the case of isolated clouds, where external compression exceeds the effect of self-gravity at low density contrast, the latter becoming dominant only for more centrally-condensed clouds. There, the boundary pressure is all that prevents the cloud from dispersing. That there are no such equilibria in the present sequence can also be seen by considering the total energy, $`EW+U`$, also plotted in Fig. 3. Since $`E<0`$ for all $`\rho _c/\rho _s`$, all of the equilibria are gravitationally dominated. Indeed, this feature is consistent with the notion that such structures are the result of nonlinear fragmentation within the self-gravitating parent filament. Note that this also means that the prolate shape of the cores is not primarily due to tidal stretching by adjacent cores; an effect that, in the isolated case, may be mimicked by use of the boundary condition eq. (13) (see, e.g., Curry & Stahler 2000). Finally, note that all three of $`U,\mathrm{\Pi }`$ and $`W`$ approach nonzero limiting values as $`\rho _c/\rho _s1`$. This is to be expected, since these quantities are nonzero in the S-O cylinder. Their limiting values are readily calculated as: $`U_0=3\pi Z^{},\mathrm{\Pi }_0=2\pi Z^{}`$, and $`W_0=4\pi Z^{}`$.
### 4.2. A Sequence of Prolate Equilibria
As indicated by the dashed curves in Fig. 2, for a given $`Z^{}`$ there exists a unique tidal lobe, within which the equidensity countours are closed, and that corresponds to a density minimum along the symmetry axis, $`r=0`$. The intersection of the tidal lobe and the $`z=0`$ plane is a circle whose radius we denote by $`R_t`$. The tidal lobe is thus one possible definition of the “surface” of the dense core, since it marks where the latter may be distinguished from the background filament. However, these solutions make clear the danger of taking such a term too literally. On this interpretation, we may fix the density on the tidal lobe at a value appropriate to the intercore medium. (Defining the extent of a core using a single equidensity contour is in fact in accord with the usual observational definition, which assigns a core’s size on the basis of its associated half-power intensity contour).
Fixing the intercore density $`\rho _s=\overline{m}n_s`$ and temperature $`T`$ allows one to characterize the size and mass of the fragments via the following reference quantities:
$$l_s\left(\frac{a^2}{G\rho _s}\right)^{1/2}=\left(\frac{\pi }{2}\frac{\rho _c}{\rho _s}\right)^{1/2}l_c,m_s\rho _sl_s^3=\frac{a^3}{(G^3\rho _s)^{1/2}}.$$
(15)
In terms of these reference values, we now define the dimensionless lengths and mass:
$$\stackrel{~}{Z}\frac{Z}{l_s}=Z^{}\left(\frac{\pi }{2}\frac{\rho _c}{\rho _s}\right)^{1/2},\stackrel{~}{R}_t\frac{R_t}{l_s},\mathrm{and}\stackrel{~}{M}_t\frac{M_t}{m_s},$$
(16)
where $`M_t`$ is the dimensional mass contained within the tidal lobe.
Figure 4 shows the behavior of $`\stackrel{~}{Z},\stackrel{~}{R}_t`$, and $`\stackrel{~}{M}_t`$ as a function of density contrast for fixed $`a`$ and $`n_s`$. The results are reminiscent of the Bonnor-Ebert (hereafter B-E) sequence, shown by dashed lines, but with a maximum in mass at $`\rho _c/\rho _s10.0`$ instead of at $`(\rho _c/\rho _s)_{\mathrm{BE}}=14.04`$. The maximum mass of the prolate sequence, $`\stackrel{~}{M}_t=1.32`$, exceeds that of the B-E sequence, $`M_{\mathrm{BE}}=1.182`$ by 12 percent. The polar radius $`\stackrel{~}{Z}`$ decreases from a value of 1.53 at $`\rho _c/\rho _s1.2`$ to a minimum of 0.60 at $`\rho _c/\rho _s600`$. The tidal lobe radius $`\stackrel{~}{R}_t`$ has a maximum of 0.45 at $`\rho _c/\rho _s6.4`$, and a minimum of 0.34 at $`\rho _c/\rho _s10^3`$.
Finally, it is instructive to compare the density structure of the prolate fragments with that of known equilibria. The equatorial and polar density profiles for the models displayed in Figs. 2$`c`$ and $`d`$ are shown in Figures 5$`a`$ and $`b`$. These figures also show the radial density profiles of a S-O cylinder and of a marginally stable B-E sphere. At the relatively low $`\rho _c/\rho _s`$ of Fig. 5$`a`$, the polar density profile is shallower than $`z^2`$, while the equatorial density resembles that of a B-E sphere out to about three times its “core radius,” after which it steepens to the $`r^4`$ of the S-O profile. At high $`\rho _c/\rho _s`$ (Fig. 5$`b`$), the similarity of both density profiles to that of the B-E sphere is evident. Thus, except at very low central concentrations or at distances far outside the tidal lobe, the core’s density structure bears little resemblance to that of its parent S-O cylinder. This is consistent with observations of globular filaments, which rarely imply density profiles steeper than $`r^2`$ (Alves et al 1998; Johnstone & Bally 1999; Lada, Alves, & Lada 1999).
### 4.3. Remarks on Stability
In the case of isolated, pressure-bounded clouds, the presence of a maximum in cloud mass as one proceeds from smaller to larger $`\rho _c/\rho _s`$ signifies a transition from stable to unstable equilibria. However, this technique (often referred to as the “static method”; Tassoul 1978), cannot be carried over to the embedded cores characterized by the mass $`\stackrel{~}{M}_t`$, since they are not complete equilibria (§4.1). Thus, the maximum in the $`\stackrel{~}{M}_t`$ vs. $`\rho _c/\rho _s`$ relation seen in Fig. 4, while suggestive, has no direct bearing on the issue of stability.
Another method used to investigate the stability of isolated clouds is that of the Gibbs free energy (Stahler 1983; Tomisaka et al 1988). Unfortunately, this method is also unlikely to succeed here, for the following reason. Along a sequence of equilibria of fixed temperature and surface pressure, the Gibbs free energy is a minimum at the critical stability point. However, in the present context, an equilibrium consists of an entire cell, which lacks a single, isobaric bounding surface. It is therefore impossible to construct a sequence of cells of fixed $`a`$ and $`n_s`$, for which quantities such as mass have definite maxima. Consequently, one finds that the Gibbs free energy is a monotonically decreasing function of $`\rho _c/\rho _s`$; i.e., no minimum value is attained along the sequence. Thus, the free energy approach cannot be applied in the usual manner. Explicit, time-dependent calculations of the evolution of small perturbations within the equilibria may be required to illuminate this important issue.
## 5. COMPARISON WITH OBSERVATIONS
To compare with observations, it is useful to write $`l_s`$ and $`m_s`$ as
$`l_s`$ $`=`$ $`0.382\mathrm{pc}\left({\displaystyle \frac{a}{0.19\mathrm{km}\mathrm{s}^1}}\right)\left({\displaystyle \frac{n_s}{10^3\mathrm{cm}^3}}\right)^{1/2},`$ (17)
$`m_s`$ $`=`$ $`3.19M_{\mathrm{}}\left({\displaystyle \frac{a}{0.19\mathrm{km}\mathrm{s}^1}}\right)^3\left({\displaystyle \frac{n_s}{10^3\mathrm{cm}^3}}\right)^{1/2}.`$ (18)
The scaling for $`n_s`$ is the minimum total number density, $`n=1.2n_{H_2}`$, estimated from <sup>13</sup>CO measurements of filamentary clouds (Nercessian et al 1988), while that for $`a`$ is the thermal sound speed for $`T=10`$ K and a mean molecular weight of $`\overline{m}=2.33m_H`$. In a high mass star-forming region such as Orion, $`n_s`$ can be as high as $`10^4`$ cm<sup>-3</sup> (Dutrey et al 1993), while $`aT^{1/2}`$ can exceed 0.40 km s<sup>-1</sup>. Nonthermal motions, which dominate thermal motions in clouds above scales $``$ few $`\times 0.1`$ pc, may be included in a schematic manner by replacing $`a`$ with a constant velocity dispersion $`\sigma `$. Representative values of $`\sigma `$ are given below.
With the above scalings, the results of Fig. 4 correspond to the following ranges:
$`0<M/M_{\mathrm{}}4.2,`$
$`0.23\mathrm{pc}Z0.58\mathrm{pc},`$
$`0R_t0.17\mathrm{pc}.`$
These ranges are in reasonable agreement with deduced values for cores observed in dense tracers, such as ammonia (Jijina et al 1999). This is true even though NH<sub>3</sub> has a critical excitation density, $`n_{\mathrm{ex}}10^4`$ cm<sup>-3</sup>, an order of magnitude larger than our assumed $`n_s`$. However, this effect is offset in equations (17) and (18) by the fact that some regions have $`T>10`$ K and that all exhibit nonthermal motions.
We now compare the intercore separations found in this simple model with those observed in star-forming regions. Two values of $`\stackrel{~}{Z}`$ may be singled out from Fig. 4 as being of particular interest. The first, denoted by $`\stackrel{~}{Z}_{\mathrm{peak}}`$, corresponds to the peak in $`\stackrel{~}{M}`$, and has the value $`\stackrel{~}{Z}_{\mathrm{peak}}=0.94`$. The second notable value is the minimum in $`\stackrel{~}{Z}`$ occuring near the right-hand side of Fig. 4, $`\stackrel{~}{Z}_{\mathrm{min}}=0.60`$. The dimensional value of each of these quantities depends upon both the intercore density and velocity dispersion through equations (16) and (17). Figure 6 shows both characteristic separations as a function of $`n_s`$ for various $`\sigma `$, where the latter is given in terms of $`a_{10}0.19`$ km s<sup>-1</sup>, the thermal sound speed in $`T=10`$ K gas. The range of observed intercore separations, $`2Z_{\mathrm{obs}}`$, in each of three star-forming regions—Taurus, Ophiuchus, and Orion— is indicated by vertical bars on the graph. These ranges were deduced from molecular line maps of cores embedded within filaments in each region. The characteristic intercore density was estimated as the observed density at a scale corresponding to the mean value of $`2Z_{\mathrm{obs}}`$ for that region. Further details on the data used may be found in the caption to Table 1.
Figure 6 shows that the observed separations agree with $`2Z_{\mathrm{peak}}`$ for $`\sigma /a_{10}12.5`$ in Ophiuchus and Orion, and for $`\sigma /a_{10}1.53.5`$ in Taurus. The corresponding ranges computed for $`2Z_{\mathrm{min}}`$ are $`\sigma /a_{10}1.54`$ in Ophiuchus and Orion, and $`\sigma /a_{10}25`$ in Taurus. These ranges may now be compared with the observed intercore velocity dispersion $`\sigma _{\mathrm{obs}}`$ in each region. In Taurus, the range computed according to $`2Z_{\mathrm{peak}}`$ includes $`\sigma _{\mathrm{obs}}`$, while in Ophiuchus and Orion, $`\sigma _{\mathrm{obs}}`$ lies above the range of theoretical values. On the other hand, the ranges computed for $`2Z_{\mathrm{min}}`$ include $`\sigma _{\mathrm{obs}}`$ for all three regions. However, since the physical relevance of $`Z_{\mathrm{min}}`$ is highly uncertain due to unresolved stability issues (§4.3), henceforth we focus on the results for $`Z_{\mathrm{peak}}`$.
The discrepancy obtained using $`Z_{\mathrm{peak}}`$ can be stated alternatively as follows: given the observed $`\sigma `$ and $`n_s`$ in both Ophiuchus and Orion, the model appears to overestimate $`2Z_{\mathrm{obs}}`$. These estimates, $`2Z_{\mathrm{pred}}`$, are given in Table 1 and plotted as open circles in Fig. 6. However, this disagreement may have a rather simple explanation. The model-derived values assume that the parent filament is aligned exactly perpendicular to the line of sight—an unlikely circumstance. Generally, the filament will be inclined at some angle $`90^{}i`$ to the line of sight, so that $`2Z_{\mathrm{pred}}`$ is reduced to $`2Z_{\mathrm{pred}}`$cos$`i`$. To bring this into agreement with the mean $`2Z_{\mathrm{obs}}`$ then requires $`i=50^{}`$ in Ophiuchus and $`i=53^{}`$ in Orion.
An independent measure of the intercore velocity dispersion may be obtained from the observed properties of filaments. According to equation (11) with $`\sigma `$ in place of $`a`$, the maximum line mass of the parent filament from which the cores could have condensed is:
$$\mu _{\mathrm{max}}=16.4\left(\frac{\sigma _{\mathrm{max}}}{a_{10}}\right)^2\frac{M_{\mathrm{}}}{\mathrm{pc}}.$$
(19)
Using available observations, we computed an approximate line mass, $`\mu _{\mathrm{obs}}`$, for filaments in Taurus, Ophiuchus, and Orion, with the results given in Table 1. The corresponding $`\sigma _{\mathrm{max}}`$, obtained from equation (19) assuming $`\mu _{\mathrm{obs}}=\mu _{\mathrm{max}}`$, is also tabulated. In Taurus and Ophiuchus, $`\sigma _{\mathrm{max}}`$ and $`\sigma _{\mathrm{obs}}`$ are in close agreement, indicating that the filaments themselves are near virial equilibrium. In Orion, $`\sigma _{\mathrm{max}}`$ should not be compared with $`\sigma _{\mathrm{obs}}`$, since the latter is obtained from C<sup>18</sup>O measurements, whereas $`\mu _{\mathrm{obs}}`$ is derived from <sup>13</sup>CO observations (Bally et al 1987). Instead, we note that the corresponding linewidth, $`\mathrm{\Delta }V(^{13}\mathrm{CO})=2.2`$ kms<sup>-1</sup>, gives a dispersion of $`\sigma _{\mathrm{obs}}(^{13}\mathrm{CO})=0.97`$ kms<sup>-1</sup>, again in very good agreement with $`\sigma _{\mathrm{max}}`$.
## 6. DISCUSSION
If the effect of inclination alone does not account for the above discrepancy between $`2Z_{\mathrm{peak}}`$ and $`2Z_{\mathrm{pred}}`$, then this may indicate a more fundamental limitation of the model in its present form. The smaller observed intercore separation (for a given $`n_s`$ and $`\sigma `$) might be the result of further fragmentation of stable cores lying at the low-density end of the equilibrium sequence, Fig. 4. However, the resulting cores would have correspondingly smaller masses than those obtained above in §5, which are already in only marginal agreement with observations.
Although this paper has not included magnetic effects, it is worth mentioning that, in the presence of a longitudinal magnetic field $`B`$ that decreases in strength outward from the parent cylinder axis, the critical wavelength for instability is significantly reduced (Nakamura, Hanawa, & Nakano 1993; Hanawa et al 1993). This occurs because the effective sound speed, $`c_{\mathrm{eff}}=(a^2v_A^2)^{1/2}`$, where $`v_AB/(4\pi G\rho )^{1/2}`$ is the Alfvén speed, is decreased from its non-magnetic value, thereby reducing thermal support against gravity in the longitudinal direction. If a suitable generalization of the solutions found here exists for magnetized, isothermal cylinders, then the constituent fragments should have a smaller spacing for a given density contrast, perhaps offering more satisfactory agreement with observations.
What is the relation between the intercore spacing and the most unstable and critical wavelengths in the S-O cylinder? We saw in §2 that, in the case of the isothermal layer, it is the critical wavelength that determines the inter-filament spacing in S-B’s 2D solution. This feature merits further comment. Let us imagine the temporal growth of an unstable fragment, either filamentary or prolate, to proceed as follows. Fragmentation begins at the scale of $`\lambda _{\mathrm{MGR}}`$, the wavelength of the fastest-growing linearly unstable mode of the parental gas distribution. As the fragment (defined as the material residing within the tidal lobe) grows into the nonlinear regime, it begins to attract, and be attracted by, surrounding fragments with the same properties. Thus, the entire chain of masses becomes more tightly bound as a result of the increased gravitational potential energy due to condensation, and the intercore spacing decreases. There is a limit to this process, however, at the scale of $`\lambda _{\mathrm{cr}}`$, below which the parent cloud is stable even in the linear regime.
The above conjecture is supported in the planar case by the fact that the gravitational potential energy per unit length per cell increases in absolute value with the amplitude $`A`$ (S-B 1967). In the cylinder, the same trend is observed in $`W`$ vs. $`\rho _c/\rho _s`$ (Fig. 3). Unfortunately, the lack of a 2D analytic solution—and therefore an amplitude—in the cylindrical case precludes the derivation of an algebraic relation between the intercore separation and $`\lambda _{z,\mathrm{cr}}`$. However, a close inspection of the numerical results suggests that it is again $`\lambda _{z,\mathrm{cr}}`$ that is of primary importance in filament fragmentation.
Recall from the results of §3.2 that we were able to find 2D equilibria only down to a minimum $`Z^{}2.1`$, where $`\rho _c/\rho _s=1.2`$. Below this value, the method converged to the S-O solution. Now, if the minimum length scale for fragmentation were set by the critical wavelength as in the planar case, then we would expect to find 2D solutions down to a minimum value of $`Z^{}=Z/l_0=\lambda _{z,\mathrm{cr}}/(2l_0)1.97`$3.1), with a corresponding $`\rho _c/\rho _s1`$. Our minimum converged value of $`Z^{}=2.1`$ is consistent with this expectation, given the intrinsic inaccuracy of the numerical method. Moreover, the fact that equilibria with $`\rho _c/\rho _s>1.2`$ were found for $`2.1Z^{}3.91=\lambda _{z,\mathrm{MGR}}/(2l_0)`$ certainly argues against any possible significance of $`\lambda _{z,\mathrm{MGR}}`$ in the final equilibria.
Should a subsequent stability analysis of these solutions reveal that the mass peak in Fig. 4 indeed signifies a stability transition, then an interesting, although highly approximate, criterion for whether a given filament will form stars can be derived. The argument proceeds as follows. The mass peak of Fig. 4 corresponds to $`\stackrel{~}{Z}_{\mathrm{peak}}=0.94`$, whence $`Z_{\mathrm{peak}}=0.94l_s=0.36(\sigma /a_{10})`$ pc. Thus, $`2Z_{\mathrm{peak}}=0.72(\sigma /a_{10})`$ pc represents the scale of a single whole fragment formed out of the parent cloud. But what is the affected region of the original filament? It is likely to be somewhat larger than $`2Z_{\mathrm{peak}}`$, by the argument summarized in the paragraph above. Roughly speaking, the affected region should be larger by a factor $`\lambda _{z,\mathrm{MGR}}/\lambda _{z,\mathrm{cr}}2`$, since $`\lambda _{z,\mathrm{cr}}`$ gives the scale at which the filament first becomes unstable to fragmentation (§3.1). Thus, assuming that the original filament was not radially collapsing ($`\mu <\mu _{\mathrm{max}}`$), we estimate that any filament longer than $`22Z_{\mathrm{peak}}1.44(\sigma /a_{10})`$ pc may be expected to form stars via fragmentation and subsequent dynamical collapse of the prolate fragments. Shorter filaments, which could only harbor cores with $`\rho _c/\rho _s<10`$, would require another agent— e.g., an increase in external pressure, ambipolar diffusion, etc.—to initiate collapse.
## 7. CONCLUSION
We have presented a new family of 2D, numerical solutions of the isothermal Lane-Emden equation in cylindrical symmetry. The equilibria have an embedded, periodic structure, and in this sense are direct counterparts to the 2D solution of Schmid-Burgk (1967) in planar symmetry. Together, the existence of the two equilibria suggest a fragmentation scheme which seems consistent with the observed hierarchical structure in several well known star-forming regions. Moreover, our results constitute a remarkably simple yet robust explanation for the origin and maintenance of the prolate, gaseous cores that represent the lower rung of this hierarchy. Although we have ignored many physical effects that should be included in subsequent studies (particularly nonthermal motions and magnetic fields), we expect that the shift in emphasis from isolated to embedded structures will prove fruitful in future theoretical work on the origin and evolution of dense cores.
It is a pleasure to thank Steve Stahler for discussions that prompted a search for these solutions. I am grateful to Richard Larson and Dean McLaughlin for useful comments on the manuscript, and to an anonymous referee for suggestions which helped to clarify certain properties of the equilibria.
## Appendix A RELATION BETWEEN THE SOLUTIONS OF S-B AND S-O
It is worthwhile noting an intriguing correspondence between the solutions of S-B and S-O. Let us examine the structure of one of the cylindrical fragments of §2 in its own right. That is, consider the limit $`A1`$, so that $`\rho _{0,s}\mathrm{}_0^10`$, giving an isolated cylinder (see equation ). Let $`\mathrm{}_0^1=2\pi /\lambda _{x,\mathrm{cr}}k_{x,\mathrm{cr}}`$. Then in the limit $`k_{x,\mathrm{cr}}0`$ the solution (3) becomes
$$\rho (x,y)=\rho _c\left[1\frac{Ax^2+y^2}{1A}k_{x,\mathrm{cr}}^2+O(k_{x,\mathrm{cr}}^4)\right].$$
(A1)
In terms of $`k_{x,\mathrm{cr}}`$, equation (8) reads
$$k_{x,\mathrm{cr}}^2\mathrm{}_c^2=\frac{1A}{1+A},$$
the substitution of which into equation (A1) gives
$`\rho (x,y)`$ $`=`$ $`\rho _c\left[1{\displaystyle \frac{Ax^2+y^2}{\mathrm{}_c^2(1+A)}}+O(k_{x,\mathrm{cr}}^4)\right]`$ (A2)
$`=`$ $`\rho _c\left[1{\displaystyle \frac{r^2}{2\mathrm{}_c^2}}+O(k_{x,\mathrm{cr}}^4)\right],`$
in the limit as $`A1`$, where we have noted that $`r^2=x^2+y^2`$, appropriate to the geometry of the isolated cylindrical fragment.
The S-O solution, on the other hand, reduces to the following in the limit of small $`r`$ (equation ):
$$\rho (r,z)=\rho _{0,f}\left[1\frac{2r^2}{l_0^2}+O(r^4)\right].$$
(A3)
Equations (A2) and (A3) are identical, at order $`r^2`$, provided that: (i) $`l_0^2=4\mathrm{}_c^2`$; and (ii) $`\rho _{0,f}=\rho _c`$. Recalling the definitions of $`\mathrm{}_c`$ and $`l_0`$ given in §2 and 3, if the latter is true, then so is the former. But equality (ii) does indeed hold, since the planar fragment has a topology indistinguishable from that of a cylinder in the prescribed limit. Thus, the S-O solution appears as the unique $`A1`$ limit of the S-B solution.
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# Evolutionary models for very-low-mass stars and brown dwarfs with dusty atmospheres
## 1 Introduction
Since the unambiguous identification of the first cool brown dwarf (BD) Gl 229B (Oppenheimer et al. 1995), the discovery of objects with temperatures cooler than the latest known M-dwarfs has been steadily increasing. These discoveries revealed two classes of objects requiring new spectral classifications. On one hand, a family of objects shows very red infrared ($`J`$-$`K`$, $`H`$-$`K`$) colors while the signature of metal oxides (TiO, VO), whose band strength index is used to classify M-dwarf spectral types, and hydrides (FeH, CaH bands) disappear gradually from the spectral distribution, as observed e.g. in GD 165B (Kirkpatrick et al. 1999a). Several of these so-called “L” dwarfs (Martín et al. 1999; Kirkpatrick et al. 1999b), have been discovered by the near-IR surveys DENIS (Delfosse et al. 1999) and 2MASS (Kirkpatrick et al. 1999b). On the other hand, a dozen of objects with properties similar to Gl229B have been recently discovered with 2MASS (Burgasser et al. 1999), SDSS (Strauss et al. 1999) and the VLT (Cuby et al. 1999). The IR spectrum of these objects is characterized by the unambiguous signature of methane absorption in the H, K and L bands, the predicted dominant equilibrium form of carbon below a local temperature $`T13001500`$ K in the $`P510`$ bar range (Allard & Hauschildt 1995; Tsuji et al. 1996; Fegley & Lodders 1996). This yields blue near-infrared colors, with $`JK<\mathrm{\hspace{0.17em}0}`$ but $`IJ>\mathrm{\hspace{0.17em}5}`$ (see e.g. Allard 1999). These cool dwarfs are identified as “methane” dwarfs, sometimes called also ”T-dwarfs” (Kirkpatrick et al. 1999b).
The spectroscopic and photometric properties of L-dwarfs and, to less extent, of methane-dwarfs cannot be properly reproduced by dust-free atmosphere models, indicating that below $`T_{\mathrm{eff}}<\mathrm{\hspace{0.17em}2800}K`$, theoretical models must include the formation and the opacity of dust grains (Tsuji et al. 1996, 1999; Ruiz et al. 1997; Jones and Tsuji 1997; Allard 1998; Leggett et al. 1998; Tinney et al. 1998; Basri et al. 1999; Kirkpatrick et al. 1999b).
In previous papers, we have developed a consistent theory of the structure and the evolution of low-mass, dense objects (Chabrier and Baraffe 1997, CB97) based on non-grey dust-free atmosphere models (Hauschildt et al. 1999, HAB99). Evolutionary models derived from this general theory describe successfully various observed properties of M-dwarfs down to the bottom of the main sequence: mass-magnitude relationships, color-magnitude diagrams, mass-spectral type relationships (Chabrier et al. 1996; Baraffe & Chabrier, 1996; Baraffe et al. 1997, 1998, BCAH98). The aim of the present paper is to extend these calculations into the brown dwarf domain by including the effect of grain formation both in the atmospheric equation of state (EOS) and in the opacity for objects with $`T_{\mathrm{eff}}<\mathrm{\hspace{0.17em}2800}K`$. The first comprehensive description of the effect of grain formation on the evolution of BDs was done by Lunine et al. (1989) and has been updated and incorporated in Burrows et al. (1993, 1997) calculations. However only a few grains were included in these calculations and the grain opacity is treated either as a frequency-independent Rosseland mean, implying grey atmosphere conditions, above 1300 K (see Figure 1 of Burrows et al. 1997), or is simply ignored (assuming grains settle immediately below the photosphere) below 1300 K, where non-grey atmosphere models are used. The present calculations extend significantly the number of grain species and include consistently their frequency-dependent opacity in the transfer equation, yielding consistent non-grey atmosphere structure, spectral colors and evolutionary calculations. The models are described in §2, whereas comparison between theory and observation in various color-magnitude diagrams (CMD) is presented in §3. The remaining uncertainties and shortcomings in the theory are discussed in the conclusion.
## 2 Model description
The main input physics involved in the present calculations has been presented in details in CB97 and BCAH98. The outer boundary conditions between the interior and the non-grey atmosphere profiles are described in CB97. Non-grey atmospheres are a necessary condition when absorption coefficients strongly depend upon frequency, which is the case for all objects below $`T_{\mathrm{eff}}5000`$ K, the onset of molecular formation. The formation of grains in the present model atmospheres is described briefly in Leggett et al. (1998) and Allard (1999) and in details in Allard et al. (2000a). The equilibrium abundances of each grain species are determined from the Gibbs energies of formation, following the method developed by Grossman (1972). At each temperature and for each condensed phase under consideration, the variable $`K_c(T)`$ is calculated from the partial pressures of the species $`i`$ forming the condensate (e.g. $`Al`$ and $`O`$ for corundum $`Al_2O_3`$), determined by the vapor phase equilibria, $`P_i=N_iT`$, where $`N_i`$ is the number of moles of species $`i`$ and $``$ is the gas constant. This variable $`K_c`$ is compared to the equilibrium constant $`K_{eq}`$, calculated from the Gibbs energy of formation of the condensate. The abundance of a condensed species is obtained by the condition that this species be in equilibrium with the surrounding gas phase, $`K_cK_{eq}`$ (Grossman, 1972). The opacities of the grains are calculated from Mie theory (see also Alexander & Ferguson, 1994).
The formation of condensed species depletes the gas phase of a number of molecular species (e.g. VO, TiO, FeH, CaH, MgH) and of refractory elements such as Al, Ca, Ti, Fe, V (Fegley & Lodders, 1996; Allard 1998; Lodders 1999; Burrows & Sharp, 1999), modifying significantly the atmospheric structure and the emergent spectrum. For late M-dwarfs and for massive and/or young BDs, the main cloud formation is predicted to occur near the photosphere (see e.g. Allard 1998; Lodders 1999). The present calculations assume a grain-size distribution in the submicron range (Allard et al 2000a). This general grain treatment is consistent with the observations of the objects near the bottom of and below the MS, namely the DENIS and 2MASS objects (Delfosse et al. 1999; Kirkpatrick et al. 1999b), GD 165B (Kirkpatrick et al. 1999a) and Kelu-1 (Ruiz et al. 1997), which all exhibit strong thermal heating and very red colors. Asides from this complex thermochemistry, grain formation is in fact a dynamical process and involves a balance between various timescales, such as the condensation, evaporation, coagulation, coalescence and convection timescales (see e.g. Rossow 1978), not mentioning the fact that non-equilibrium species might be present, as for example in the atmosphere of Jupiter (Fegley & Lodders 1994). Since such dynamical processes are not incorporated in the theory at the present stage, we elected to conduct calculations under various extreme assumptions.
The first set of models are based on the dust-free “NextGen” model atmospheres (Allard et al. 1996; HAB99) and are the same as in BCAH98. In the second set of models, all condensed equilibrium species are included both in the EOS and in the opacity, taking into account dust scattering and absorption in the radiative transfer equation. These models are referred to as “DUSTY” models. The third set of models includes the formation of grains in the atmosphere EOS, thus taking into account the photospheric depletion of dust-forming elements in the gas phase, but not in the transfer equation, therefore ignoring the opacity of these condensates. This case, referred to as “COND” models, mimics a rapid gravitational settling of all grains below the photosphere. This is similar to the Burrows et al. (1997) calculations below $`T_{\mathrm{eff}}1300`$ K (see above).
As mentioned above, the DUSTY and COND models represent extreme situations which bracket the more likely intermediate case resulting from complex, and presently not understood, thermochemical and dynamical processes.
### 2.1 Atmosphere profiles
Figure 1 displays $`P`$-$`T`$ atmosphere profiles for $`T_{\mathrm{eff}}`$ = 1800 K, log $`g`$=5 and solar composition, under the afore-mentioned different assumptions for the dust treatment. The atmospheric heating due to the large grain opacity (the so-called greenhouse or backwarming effect) in the DUSTY model is clearly visible, and yields a significantly hotter structure. The effect is dominant in the outer layers, where the spectrum forms, and decreases in the inner layers, where the connection with the interior profile occurs (see also Tsuji et al. 1996). When dust opacity is not taken into account as in the COND models, the inner atmospheric structure is barely affected by the formation of dust and is similar to the dust-free case (NextGen) for this temperature. Only for significantly cooler $`T_{\mathrm{eff}}`$ does dust formation start to affect substantially the inner structure compared to dust-free models.
As indicated in Figure 1, convection occurs only in optically-thick regions ($`\tau >1`$), in contrast to hotter objects (see BCAH98). However, this optical depth corresponds to only about one pressure scale height $`H_\mathrm{P}`$, or even less, from the photosphere ($`\tau 1`$) near which occurs the condensation temperature of most grains for objects near the bottom of the MS. Even for the cooler models ($`T_{\mathrm{eff}}1000`$ K) where convection retreats to deeper layers, the top of the convective zone lies only at a few $`H_\mathrm{P}`$ from the photosphere. This can have important consequences on the formation and settling of atmospheric grains. Indeed, turbulent diffusion, as produced for example by overshooting from the convection zone or by advection due to rotation, is a much more efficient mechanism than microscopic diffusion and sedimentation, at least for submicron-size particles. Although the temperature at the top of the convective zone is found to be generally above the condensation temperature of all grains, this turbulent diffusion will efficiently bring material upward to the region of condensation and maintain small-grain layers, which otherwise would have settled gravitationally, in this region. Note that for submicron-size species, the condensation time $`\tau _{cond}`$ is much smaller than other dynamical timescales (see e.g. Rossow, 1978; Lunine et al. 1989). As temperature decreases, the photospheric density increases, more species condense so that the particle density $`\rho _p`$ increases. Coalescence will become effective in growing large particles ($`\tau _{coal}/\tau _{cond}T/\rho _p^2`$, Rossow 1978) and the afore-mentioned turbulent mixing could no longer prevent these large particles to fall out rapidly. Moreover the condensation line of some species (e.g. silicates) will drop below the photosphere (see e.g. Lodders 1999). Within this general scheme, the DUSTY models correspond to a situation with very efficient turbulent mixing, whereas the COND models correspond to a situation where mixing is inefficient compared with sedimentation. Whether this general mechanism, and the related decrease of opacity, can explain the color saturation - and thus smaller backwarming - of the coolest L-dwarfs remains to be quantified correctly. Work in this direction is under progress.
As anticipated from Figure 1, spectral and photometric properties are generally more affected by dust formation than the inner atmospheric structure and thus the evolution, as shown in the next section. This reflects the weak dependence of $`L`$ and $`T_{\mathrm{eff}}`$ upon opacity for BDs, with $`L\kappa ^{0.3}`$ and $`T_{\mathrm{eff}}\kappa ^{0.1}`$ (Burrows & Liebert, 1993).
\*** FIGURE 1 \***
### 2.2 Evolutionary calculations
Evolutionary calculations have been performed for $`900T_{\mathrm{eff}}2800`$ K, which corresponds to masses between 0.1 $`M_{}`$ and 0.01 $`M_{}`$ for ages $`10^8<t<\mathrm{\hspace{0.17em}10}^{10}`$ yr, following the method described in CB97. Figure 2 shows the evolution of $`T_{\mathrm{eff}}`$ as a function of time for various masses and dust treatments. This illustrates the uncertainty expected from dust treatment on the cooling of BDs. For a given age, the differences between the two extreme DUSTY and COND models reach up to $``$10% in $`T_{\mathrm{eff}}`$ and $``$25% in luminosity $`L`$. The hydrogen-burning minimum mass (HBMM) is only moderately affected by dust formation. The opacity due to the formation of different grains (e.g. silicates, Al<sub>2</sub>O<sub>3</sub>, CaTiO<sub>3</sub>) produces a blanketing effect which lowers the effective temperature and luminosity for a given mass at the bottom of the main sequence, as noted initially by Lunine et al (1989). Models omitting grain opacity (NextGen and COND) yield $`m_{\mathrm{HBMM}}`$ 0.072 $`M_{}`$, $`T_{\mathrm{eff}}`$=1700K, $`\mathrm{log}L/L_{}=5\times 10^5`$, whereas DUSTY models yield a slightly lower limit, $`m_{\mathrm{HBMM}}`$ 0.07 $`M_{}`$, $`T_{\mathrm{eff}}`$=1550 K, $`\mathrm{log}L/L_{}=4\times 10^5`$, for solar composition. A lower luminosity requires a smaller H-burning nuclear energy in the core to reach thermal equilibrium, thus a lower central temperature and therefore a lower mass, which yields a lower HBMM.
\*** FIGURE 2 \***
The present calculations incorporate new conductive opacities by Potekhin et al. (1999) in the interior. These opacities, initially developed for neutron star envelopes and white dwarf cores, improve previous calculations (Hubbard and Lampe, 1969; Itoh et al. 1983 and Mitake et al., 1984; Brassard and Fontaine, 1994) and cover the range of temperatures and densities characteristic of BD interiors. An interesting new property revealed by the present calculations is the growth of a conductive core for the most massive BDs. As these objects cool and contract, their interior becomes more and more degenerate ($`\psi =kT/kT_F3.3\times 10^6T(\mu _e/\rho )^{2/3}`$, where $`T_F`$ is the electron Fermi temperature and $`\mu _e`$ is the electron mean molecular weight). Electron conductivity becomes more and more important and becomes eventually the main energy transport mechanism in the central regions, instead of convection. This occurs in the core of old ($`t`$ 2 Gyr) BDs in the mass range 0.03-0.07 $`M_{}`$, with a maximum extension of the conductive core $`R_{cond}/R_{}0.60.7`$ for the more massive ones, since the maximum central density occurs near the HBMM (Burrows et al. 1997; Chabrier and Baraffe 2000). In this mass range, the conductive core appears at $`t`$ 2-3 Gyr, for $`\psi <\mathrm{\hspace{0.17em}0.1}`$, and increases rapidly, reaching up to 50%-70% of the total mass at 10 Gyr for masses $`m0.050.06`$ $`M_{}`$. Figure 3 shows the evolution of the conductive core $`M_{cond}/M_{}`$ as a function of time for a 0.06 $`M_{}`$ BD. We verified that these results still hold, although quantitatively different, when using previously published conductive opacities. These properties are essentially unaffected by the dust treatment in the atmosphere, since evolution depends weakly on it.
\*** FIGURE 3 \***
The onset and growth of a conductive core does affect the evolution. Once conduction sets in at the center, it increases dramatically the efficiency of energy transport compared to convection ($`v_F>>v_{conv}`$, where $`v_F=(2kT_F/m_e)^{1/2}`$ is the electron Fermi velocity), decreasing the internal temperature gradient and yielding a cooler central temperature for a given energy flux $``$ ($`T^4T_c^4/R.\kappa _{cond}`$, where $`\kappa _{cond}`$ is the conductive opacity). This increases electron degeneracy at the center ($`\psi kT`$ decreases). In first approximation, retaining only the zero-temperature electron contribution and the kinetic ionic contribution, the radius of a BD can be written $`RR_0(1+\psi )`$ where $`R_0=2.8\times 10^9(\frac{m}{M_{}})^{1/3}\mu _e^{5/3}`$ cm is the zero-temperature (fully degenerate, $`\psi =0`$) radius (Stevenson, 1991). Therefore the more degenerate conductive object contracts towards a slightly smaller radius, yielding a larger release of gravitational energy $`\delta \mathrm{\Omega }=_MP\delta (1/\rho )𝑑m`$ within a Kelvin-Helmoltz time.
This larger energy source increases the total binding energy of the star:
$$\delta B(t)=[\delta \mathrm{\Omega }(t)+\delta U(t)]=_MGm\delta (\frac{1}{r})𝑑m_M\delta \stackrel{~}{u}𝑑m$$
(1)
where $`\stackrel{~}{u}`$ is the specific internal energy, yielding a larger luminosity $`L`$ for a given age, i.e. slowing down the cooling time $`\tau `$ of the BD, since:
$$L(t)=\frac{dB(t)}{dt}\tau (L)=_{t_0}^t\frac{dB(t^{})}{L}$$
(2)
where we have neglected the nuclear source of energy arising from some residual hydrogen-burning (see e.g. Figure 2 of Chabrier and Baraffe 2000).
As an example, for a 10 Gyr old 0.06 $`M_{}`$ BD, $`T_{\mathrm{eff}}`$ is 15% larger and $`L`$ is 30% brighter when conduction is taken into account. An effect difficult to verify observationally, however, given the faintness of such old BDs. Note that the onset of a conductive core occurs only for $`t>\mathrm{\hspace{0.17em}1}`$ Gyr and is thus inconsequential for the application of the lithium-test, which concerns only younger objects (see e.g. CB97).
## 3 Color - Magnitude diagrams
In this section, we compare the present models with available observations of late-M, L- and methane-dwarfs in various CMDs. The comparison is limited to objects with known distances. Although trigonometric parallaxes have been determined for several 2MASS and DENIS L-dwarfs (Dahn et al. 2000), the only methane dwarf with known distance is Gl229B.
### 3.1 Optical CMD
\*** FIGURE 4 \***
Figure 4 compares models and various observations in a $`M_\mathrm{V}`$-$`(V`$-$`I)`$ CMD. Three L-dwarfs of the Dahn et al. (2000) sample are known to be very close binary systems (Martín, Brandner & Basri 1999; Koerner et al. 1999) with nearly equal mass components and their magnitude is corrected in the CMDs to take into account this property. Models with different dust treatments are displayed for ages of 0.1 Gyr (only for the DUSTY models; short-dash line) and 1 Gyr (solid line). For $`t>`$ 1 Gyr, DUSTY isochrones are hardly distinguishable from one another in any of the following optical ($`VRI`$) and near-IR ($`JHK`$) CMDs, although the same masses correspond obviously to different magnitudes, and are not shown for sake of clarity. The properties of the DUSTY models for different ages are displayed in Tables 1-5. Note that we do not give the H-magnitude. Indeed, the present models overestimate the flux in this band and work remains to be done to understand this shortcoming. We also give the lithium-abundance depletion for these models for applications of the lithium-test. The COND models displayed in Figure 4 (long-dash line) cover the mass range between 0.08 $`M_{}`$ ($`T_{\mathrm{eff}}=2350K`$) and 0.03 $`M_{}`$ ($`T_{\mathrm{eff}}=1050K`$) at 1 Gyr.
For $`M_\mathrm{V}`$ brighter than 20 ($`T_{\mathrm{eff}}>\mathrm{\hspace{0.17em}2200}`$ K), the NextGen models (dash-dot) are significantly bluer by 0.3-0.4 mag in $`(VI)`$ than the DUSTY models. This is due to the recent improved TiO line list (Schwenke, 1998) included in the DUSTY atmospheres (Allard et al. 2000a,b), while the NextGen models used the previous line list by Jorgensen (1994). As emphasized in BCAH98, the NextGen models predict too blue (up to $`0.5`$ mag) optical colors compared to observations for $`10<M_\mathrm{V}<\mathrm{\hspace{0.17em}19}`$ (see Figure 6 of BCAH98). As shown in Figure 4, the new TiO linelist improves significantly this situation (see also Allard et al. 2000b) and confirms the fact that the previous discrepancy stemmed primarily from a lacking source of opacity shortward of 1 $`\mu `$m (cf. §4.2 of BCAH98), and not from a major caveat in the treatment of convection. However, as shown in Figure 4 and Figure 5 below, the situation is not completely satisfactory yet, with a remaining $`0.2`$-$`0.3`$ offset in colors at the bottom of the MS.
\*** FIGURE 5 \***
The dust effect is clearly seen for $`M_\mathrm{V}>\mathrm{\hspace{0.17em}19}`$, $`(VI)>\mathrm{\hspace{0.17em}4.8}`$, with a much more modest reddening as $`M_\mathrm{V}`$ increases ($`T_{\mathrm{eff}}`$ decreases) for the DUSTY and COND models than for the dust-free models. This reflects metal and in particular TiO and VO depletion due to grain formation in the COND and DUSTY atmospheres. These species are two main sources of absorption in the optical in the temperature-range of interest. Note that, although the $`(VI)`$ colors are similar, the DUSTY models are significantly fainter than the COND ones for a given mass, revealing the grain opacity which is not included in COND models, and the depressed continuum in the optical due to the opacity-induced backwarming and H<sub>2</sub>O dissociation (Allard et al. 2000a). Interestingly enough, the same trend of change of slope in $`(VI)`$ colors at $`M_\mathrm{V}19`$ appears in the available observations, but at a $`0.5`$ mag bluer limit, i.e. $`(VI)<\mathrm{\hspace{0.17em}4.5}`$, than the theoretical predictions. This builds confidence in our basic treatment of grain formation but suggests that the gas depletion of the main absorbing species in the optical (TiO, VO, FeH) is underestimated in the present theory. Note in passing the limited interest of the V-band for BD search, with most massive BDs older than 1 Gyr being at $`M_\mathrm{V}>\mathrm{\hspace{0.17em}24}`$.
The saturation in optical colors is even more drastic in $`(RI)`$ (Figure 5), where the DUSTY models yield almost constant $`(RI)`$ below $`T_{\mathrm{eff}}2300`$ K ($`M_I>`$ 14) at $`(RI)`$2.4-2.5. As seen in the figure, the distribution of various objects observed in these colors at the bottom of and below the main sequence, either in the field or in the Pleiades cluster, not only shows the same saturation effect as predicted by the models but is quantitatively in very good agreement with the predicted colors. As mentioned above, however, the $`0.20.3`$ mag offset above $`2300`$ K is still present. The COND models predict bluer colors as $`T_{\mathrm{eff}}`$ decreases, and in general more flux in the $`V,R,I`$ bands at a given $`T_{\mathrm{eff}}`$ than the DUSTY models. This stems from (i) the absence of grain opacity in COND models and (ii) the backwarming effect in the DUSTY models (see Figure 1) which destroys IR-absorbing species such as e.g. H<sub>2</sub>, H<sub>2</sub>O, allowing more flux to escape at longer wavelengths.
Although the L3.5 dwarf 2MASS-WJ0036159+182 ($`M_\mathrm{V}`$=21.15, $`M_\mathrm{R}`$=18.15, $`M_\mathrm{I}`$=15.92, $`R`$-$`I`$=2.23)(Reid et al. 1999) seems to be lying on the COND track, it is probably misleading and reveals probable errors in the photometry. Indeed, the near-IR colors of this object are not reproduced by the COND models, but rather by the DUSTY ones (see Figure 6).
### 3.2 Infra-red CMD
Figure 6 displays the same comparisons in near-IR colors $`(JK`$). As already mentioned by Leggett et al. (1998), the DUSTY models give the best fit to IR photometry of L- and very late M-dwarfs. They explain the very red colors characteristic of these objects in terms of the backwarming of the atmosphere due to grain absorption, yielding a reduction of the band strengths of the main IR absorbers, like e.g. H<sub>2</sub>O, as proposed initially by Tsuji et al. (1996). Such red colors are impossible to reproduce with dust-free opacity models (NextGen and COND). Note that for $`M_\mathrm{K}`$ brighter than $``$ 11, the DUSTY models are slightly bluer in $`(JK)`$, by $``$0.1-0.2 mag, than the NextGen models. This arises from the use of different H<sub>2</sub>O line lists: Miller et al. (1994) in NextGen and the recent AMES list (Partridge and Schwenke, 1997) in DUSTY models. Although less complete for the higher energy transitions than the AMES line list, the Miller et al. (1994) list yields better agreement between models and IR photometric observations for $`T_{\mathrm{eff}}>\mathrm{\hspace{0.17em}2200}`$-$`2300`$ K (where dust does not affect the colors), since it relies on better potential surfaces for high temperatures (see Allard et al. 2000b for details). Below $`T_{\mathrm{eff}}2000`$ K, these transitions are no longer important and the more complete AMES water linelist is to be used. Below this temperature, however, dust opacity is the main ingredient which shapes the IR colors, and the choice of the water line list is less crucial for photometric analysis. This illustrates again the remaining shortcomings in the present theory arising from still partially inaccurate molecular opacities, and stresses the need for future improvements in these inputs.
\*** FIGURE 6 \***
When dust opacity is neglected (COND and NextGen), near-IR colors become bluer due to strong H<sub>2</sub>, H<sub>2</sub>O, CO and eventually CH<sub>4</sub> absorption in the $`H`$ and $`K`$-bands. This yields very blue near-IR colors, with $`JK`$$``$$`JH<\mathrm{\hspace{0.17em}0.5}`$, as illustrated by the good agreement of Gl229B IR colors with the COND track (see also Burrows et al. 1997). The effect is more pronounced for the COND models because of the lack of TiO absorption, which extends into the J-band at these very low effective temperatures. However, the COND models overestimate the flux at and shortward of $`1\mu m`$ ($`V,R,I`$ bands) for Gl229B (Golimowski et al. 1998). The spectrum of GL229B indeed does not show strong indication for dust in its IR spectrum from 1 to 5 $`\mu m`$ (Allard et al. 1996; Marley et al. 1996; Tsuji et al. 1996; Oppenheimer et al. 1998), but reveals discrepancies in the optical with dust-free models, as mentioned above. Several solutions have been proposed for this missing source of opacity in the optical, such as a specific haze opacity due to the presence of the nearby star (Griffith et al. 1998) or warm dust layers (Tsuji et al. 1999). These suggestions, however, are based on ad-hoc hypothesis required to fit Gl229B spectrum and lack robust physical grounds. The methane-dwarf SDSS 1624 (Strauss et al. 1999), for example, does not have a hot companion, yet shows no flux blueward of 0.8 $`\mu `$m, indicating a strong absorption in the optical, as does Gl229B. This suggests that the source of absorption is a common atmospheric BD property, precluding the haze opacity scenario. Strong alkali resonance lines (K I, Na I) have been proposed recently as an alternative, or complementary solution (Tsuji et al. 1999; Burrows et al. 1999). Clearly this question is not settled yet and requires futher observational tests and consistent calculations.
For $`(JK)`$ 1.0-1.9, good agreement is found between the data and the DUSTY models (see Figure 6), except for the L0-dwarf 2MASS-WJ0746425+200032 (Reid et al. 1999) ($`M_\mathrm{K}`$=9.70, $`J`$-$`K`$=1.25), which lies significantly above both the other data and the models, making it a probable unresolved binary, as suggested by these authors.
An interesting feature of this diagram is a hint that the redder L-dwarfs of Dahn et al. (2000), and the slope of the observed L-dwarf sequence itself, seem to deviate from the DUSTY tracks for $`M_\mathrm{K}>\mathrm{\hspace{0.17em}11.5}`$ and $`(JK)>\mathrm{\hspace{0.17em}1.5}`$, the reddest observed objects having a color $`(JK)1.9`$, i.e. $`T_{\mathrm{eff}}`$ 1800 K. If the parallax and the photometry of these objects are confirmed, this may very well illustrate gravitational settling of some grains below the photosphere for these objects, as mentioned in §2.1. This will yield less backwarming and thus progressively cooler atmospheres (see Figure 1), increasing H<sub>2</sub>O, CO and H$`2`$ absorption in the IR. Once CH<sub>4</sub> forms near the photosphere, the integrated flux over the K-band is strongly reduced, although the flux still peaks around 1 $`\mu `$m, leading to Gl229B-like colors. Whether this transition occurs sharply or smoothly remains to be elucidated. Note that the CO-CH<sub>4</sub> transition occurs gradually with some of the two elements being present in the stability field of the other so that cool stars do contain some limited amount of methane and CO has been detected in Gl229B (Noll et al. 1997; Oppenheimer et al. 1998). The recently discovered methane-dwarfs (Burgasser et al. 1999), unfortunately, do not bring information about this transition because of their color selection criterium, $`JK<0.4`$. In case of an abrupt change in color between ”L” and ”methane”-dwarfs, there should be indeed a ”no-BD land” in the region $`0.5<JK<\mathrm{\hspace{0.17em}1.0}`$ with a sharp transition for $`T_{\mathrm{eff}}T_{\mathrm{gas}}=13001400`$ K, the temperature of formation of methane in the $`P310`$ bar range. This corresponds to masses $`m0.0150.02M_{}`$ for t=0.1 Gyr, $`m0.0450.05M_{}`$ for t=1 Gyr and $`m0.065M_{}`$ for t=5 Gyr. In any events, this concerns objects cooler than $`T_{\mathrm{eff}}`$ 1700-1800 K, the temperature derived with the DUSTY models along the 1 Gyr isochrone for LHS102B (Goldman et al. 1999; see also Figure 6), the reddest L-dwarf with known parallax, and for 2MASS-W1632+19, the faintest L-dwarf in the Dahn et al. (2000) sample with $`M_\mathrm{K}`$ $``$ 12.75, $`(JK`$) $``$ 1.9. These temperatures agree within 100 K with the ones derived spectroscopically by Basri et al. (1999) based on alkali resonance absorption lines in the optical spectrum. Objects cooler than 2MASS-W1632+19 thus correspond to BDs with masses $`m<\mathrm{\hspace{0.17em}0.03}`$ $`M_{}`$ for ages $`t>0.1`$ Gyr, $`m<\mathrm{\hspace{0.17em}0.06}`$ $`M_{}`$ for $`t>1`$ Gyr and to essentially all BDs for $`t>\mathrm{\hspace{0.17em}2.5}`$ Gyr (see Figure 2).
The same general behavior as in ($`JK`$) is found in ($`JH`$), with the same color saturation observed for the reddest L dwarfs for $`M_\mathrm{J}>\mathrm{\hspace{0.17em}13}`$ and $`(JH)<\mathrm{\hspace{0.17em}1.2}`$. Although not shown in the present paper, a COND isochrone at $`t`$(a few)$`\times `$10<sup>8</sup> yr fits well GL229B both in $`(JH`$), in contrast to Burrows et al. (1997) models, and $`(JK`$) colors. Their models could not reproduce the $`L`$’-band of GL229B. Since then, Leggett et al. (1999) revised the flux of GL229B, with a significant modification of the $`L`$’ band, which is now reproduced by the Burrows et al. (1997) models. The reason of their poor fit in the H-band, however, compared with the present models, remains unclear. Indeed the explanation proposed in Burrows et al. (1997) was based on previous, less accurate observations of Gl229B, as mentioned above, and does not hold anymore.
Finally, Figure 7 displays $`L^{}`$ versus $`(KL^{}`$). The $`K`$ band is affected by H<sub>2</sub>, H<sub>2</sub>O and CH<sub>4</sub> absorption, whereas the $`L^{}`$ is strongly affected by methane (Allard et al. 1996; Marley et al. 1996; Tsuji et al. 1996). The trends displayed in the previous IR diagrams are replicated, with a severe reddening due to grain backwarming effects in the DUSTY models, but also a substantial reddening for the COND and grainless model, due to the strong CH<sub>4</sub> absorption in the K-band.
The magnitudes of brown dwarfs and very-low-mass stars with dusty atmospheres for several isochrones are displayed in Tables 1-5. We stress, however, that the flux for these objects peaks around 1 $`\mu m`$, so that J,Z and H remain the most favorable bands for detection.
\*** FIGURE 7 \***
## 4 Discussion and conclusion
The comparisons in various CMDs displayed in the previous section illustrate the good general agreement between the present theory of evolution of very-low-mass stars and brown dwarfs with dusty atmospheres and observations. The models which include grain formation and opacity in the atmosphere explain successfully the IR colors of late M- and L-dwarfs, as well as their blue loop in the optical, although improvement is still needed for the latter case to reach accurate quantitative agreement with observations. As shown by Kirkpatrick et al. (1999a) for GD165B and Ruiz et al. (1997) for Kelu-1, the DUSTY models improve also significantly the agreement with observed IR spectra. The present models provide a consistent description of the physics going on in the interior and the atmosphere of objects at the bottom of and below the main sequence, from M-dwarfs to L-dwarfs. They can be used with reasonable confidence in the IR bands ($`JKL^{}M`$) to calibrate the main properties, mass, age, $`T_{\mathrm{eff}}`$, $`L`$ of identified late-M and L-dwarfs. Shortcomings still appear in the optical, whatever the treatment of dust, in spite of noticeable improvement (see Figure 4). Baraffe et al. (1998) already pointed out a possible missing source of opacity in the optical to explain the observed $`(VI)`$ colors of M-dwarfs between $`T_{\mathrm{eff}}`$ 3600 K and $``$ 2300 K. Although the new TiO linelist from AMES (Schwenke, 1998) improves appreciably the situation (see Figure 4), a $`0.20.3`$ mag discrepancy still remains with observations (see Figure 5). Below $``$ 2300 K, comparison of the DUSTY models with observations in $`(VI)`$ (see Figure 4) strongly suggests either an underestimate of optical opacities or an underestimate of the depletion of molecular absorbers during grain condensation. Problems appear as well in spectroscopic analysis of atomic line profiles where the DUSTY models cannot match the optical observations (Tinney et al. 1998; Basri et al. 1999). Another remaining caveat of the present theory arises from shortcomings in the absorption coefficients of water, yielding a $`0.2`$ mag offset in $`(JK)`$ for $`T_{\mathrm{eff}}>\mathrm{\hspace{0.17em}2300}`$ K when using the most recent H<sub>2</sub>O line list (see e.g. Figure 6). Finally, the optical spectrum of GL229B remains to be described accurately. Present COND synthetic spectra overestimate the optical flux, pointing out a missing source of opacity. Work in this direction is under progress.
Although it is important to solve these problems and eliminate these shortcomings, it is obvious from the various CMDs that substellar objects must be searched in the near-infrared. Brown dwarfs around 1500 K radiate more than 90% of their energy at wavelengths longward of $`1\mu `$m. The flux peaks at 1.1 and 1.3 $`\mu `$m, and $`J,Z,H`$ are the favored broadband filters for detection. It is thus reinsuring, given the complexity of the underlying physics, that the present models of dusty substellar objects describe the observations in near-IR bands with reasonable success.
If dust formation and absorption successfully explain the transition between M- and L-dwarfs, a theoretical and observational gap still exists between L- and methane-dwarfs. Models including dust opacity reach rapidly very red IR colors ($`JK`$) $`>`$ 2, whereas observations of the reddest or faintest L dwarfs show a saturation of their ($`JK`$) and ($`JH`$) colors at 1.9 and 1.2 respectively, which corresponds to $`T_{\mathrm{eff}}`$ $``$ 1700-1800 K (see Figure 6). This very likely illustrates ongoing dynamical processes near the photosphere. As outlined in §2.1, the proximity of the top of the convection zone from the photosphere may induce turbulent mixing and prevent the grains to diffuse downward. As convection retreats to deeper layers for cooler objects, gravitational settling of condensed species may lead to cooler atmospheres, because of the decreasing opacity, with a transition to very blue near infrared colors once methane forms near the photosphere. These dynamical processes of grain formation and diffusion are not included in the theory yet and represent indeed the next major challenge. Theory definitely needs some guidance from observations through this puzzle, by getting spectra, parallaxes and photometry of the coolest L-dwarfs and by either filling or confirming the presently expected brown dwarf dearth between the coolest L-dwarfs and the hottest methane dwarfs around $`T_{\mathrm{eff}}13001400`$ K.
Note: Various isochrones for the DUSTY models from 10<sup>7</sup> yr to 10<sup>10</sup> yr, covering a range of mass from 0.01 $`M_{}`$ to 0.1 $`M_{}`$ and $`T_{\mathrm{eff}}`$ from 900K to 3000K are available by anonymous ftp (same format as Table 1):
ftp ftp.ens-lyon.fr
username: anonymous
ftp $`>`$ cd /pub/users/CRAL/ibaraffe
ftp $`>`$ get DUSTY00\_models
ftp $`>`$ quit
Note that the DUSTY models are given in the Tables down to $`T_{\mathrm{eff}}`$=900K, but this limit is clearly unrealistic for dusty atmospheres, as seen in the various figures of the paper.
Special requests for models in any particular filters can be addressed to I. Baraffe.
We are very grateful to C. Dahn and H. Harris for providing data prior to publication and to A. Potekhin for providing the conductive opacities. We are also thankful to E. Martín and to the referee, D. Saumon, for constructive suggestions and comments. Part of this work was performed at the AAO (Sidney) under the auspice of the franco-australian cooperation program in Astronomy (CNRS/ARC). G.C and I.B are very indebted to B. Boyle, C. Tinney and the staff of the AAO for their generous hospitality during this visit. This work was supported in part by a CNRS-NSF ”dark matter” grant, NSF grant AST-9720704, NASA ATP grant NAG 5-3018 and LTSA grant NAG 5-3619 to the University of Georgia, and NASA LTSA grant NAG5-3435 to Wichita State University. The computations were done at the Pôle Scientifique de Modélisation Numérique at ENS-Lyon, on the T3E of Centre d’Etudes Nucléaires de Grenoble, on the IBM SP2 of CNUSC, the SGI Origin 2000 of the UGA UCNS, on the IBM SP2 of the San Diego Supercomputer Center (SDSC) with support from the National Science Foundation, and on the Cray T3E of the NERSC with support from the DoE. We thank all these institutions for a generous allocation of computer time.
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# 1 Introduction
## 1 Introduction
In the last years many experimental efforts have been devoted to the study of heavy-ion fission at beam energies below 10–20 MeV/u accompanied with the emission of light particles. Initial experiments involved the observation of fission fragments emitted in coincidence with neutrons or charged-particles . More recently, combined coincidence neutron and charged-particle data became available . The most important result of these experiments consists in the estimation of fission times which appear to be $`10^3`$$`10^4`$ times longer than the characteristic time of the single-particle motion. This fact implies that thermal equilibrium is being established over the intrinsic degrees of freedom at each value of the fission coordinate $`Q`$ and validates the use of the Rayleigh dissipative function technique to treat the coupling of $`Q`$ with the particle degrees of freedom .
Heavy-ion fission reactions accompanied by particle emission are used as a testing ground of the dissipation mechanisms of large-scale collective motion in hot nuclei (for a review see Refs. ). The main reason is that this process is simpler than other types of deep-inelastic collisions. The simplicity originates from the fact that fission is not sensitive to the fusion stage of a binary reaction whose understanding is far from clear.
When mass and projectile energy increase, the possibility to isolate the ground-state to scission-point motion becomes problematic because of the increasing contribution of quasifission, a fission-like process occurring at angular momenta $`J`$ exceeding some value $`J_0`$ where the fission barrier disappears. Quasifission in a clear cut sense has been observed in very heavy systems which do not have a fission barrier even at $`J=0`$ $`\mathrm{}`$. These experiments provided a motivation for the creation of the code HICOL which proved to be capable of reproducing the general features of the process .
The Rayleigh function used in HICOL is built within the one-body dissipation model, proposed in Ref. . It combines the wall formula in the mononuclear regime with an improved version of the so-called ‘completed window formula’ in the dinuclear stage. Calculations show that quasifission originates from the fact that the potential of the composite system for $`J`$ $`J_0`$ is very flat. Combined with strong friction this gives the system enough time to thermalize the relative velocity of the colliding nuclei and to relax the mass-asymmetry mode.
In medium mass reaction systems, a quasifission configuration can execute several rotations before scission, so that it is impossible to disentangle quasifission from fusion-fission by using fragment angular distributions as in heavier systems. Moreover, due to decreasing in angular momentum caused by the emission of few particles, a pocket can emerge in the initially flat potential, and a quasifission trajectory will go over into a fusion-fission one.
The introduction of quasifission trajectories to fission-evaporation routines is not an easy task. It requires answers on some principal questions. For example, in the fission case it is assumed that, for all $`J`$ within the fusion-fission angular momentum window, the system moves along the bottom of the fission valley of a non-rotating nucleus. Thus a question arises on whether the quasifission trajectories go along this path. Another related question is whether the evaporation model is applicable for particle emission along the quasifission trajectories.
The main purpose of the present work is to elucidate these questions. As a result we shall get a better perspective for a uniform description of fusion-fission and quasifission reactions accompanied by particle emission. To illustrate this opportunity, we test the compatibility of the charged particle clock with the neutron clock in calculations which effectively account for fusion-fission and quasifission trajectories. The study is performed in the $`A`$ 160, $`E_x`$= 200–300 MeV region, where a quite complete set of data on multiplicities of fission-associated light particles is available .
The paper is organized as follows. In Sect. 2 we define the quasistatic path and in Sect. 3 we introduce the geometrical quantities needed to characterize the quasistatic and dynamic configurations of the system. These quantities facilitate the comparison between the dynamic and quasistatic paths in Sect. 4. In Sect. 5 we perform Monte Carlo simulations of light-particle evaporation along dynamic trajectories and compare the results with available experimental data. Our conclusions are drawn in the last section.
## 2 The Quasistatic Path
The fundamental quantity in the construction of fission models is the total energy $`\stackrel{~}{E}\{\rho \}`$ of the nucleus expressed as a functional of the spatial nucleon density $`\rho `$. The definition of this functional and its general properties are well documented (e.g., see Refs. ). Given $`\stackrel{~}{E}\{\rho \}`$ one can use the equations
$$\frac{\delta \stackrel{~}{E}\{\rho \}}{\delta \rho }=0,\rho (𝐫)𝑑𝐫=A,$$
(1)
where $`\delta /\delta \rho `$ is the functional derivative, to find the ground state and the saddle point densities $`\rho _{\mathrm{gs}}`$ and $`\rho _{\mathrm{sd}}`$. The difference between the corresponding energies, $`\stackrel{~}{E}\{\rho _{\mathrm{sd}}\}\stackrel{~}{E}\{\rho _{\mathrm{gs}}\}`$, gives the fission barrier height.
Since nuclei are leptodermous objects the terms ‘density’ and ‘shape’ are often used synonymously. As indicated first in Ref. , the so-called conditional equilibrium shapes may play an important role in nuclear dynamics. To define these shapes, one introduces a quantity $`Q`$ characterizing the elongation of the nucleus. $`Q`$ is a functional of the density and has the form
$$Q\{\rho \}=q(𝐫)\rho (𝐫)𝑑𝐫,$$
(2)
where $`q(𝐫)`$ is a known function. The density $`\rho _Q^{}`$ of the conditional equilibrium shape is the solution of equations (1) in the class of densities $`\rho `$ restricted by the condition $`Q\{\rho \}=`$ $`Q^{}`$ where $`Q^{}`$ is a constant. As pointed out in Ref. , the sequence of conditional equilibrium shapes may acquire a physical meaning, for instance it can describe the process of fission if $`Q`$ is changing adiabatically compared to all other degrees of freedom. It was conjectured in Ref. that the adiabaticity condition is likely to be satisfied in heavy systems where $`\stackrel{~}{E}\{\rho \}`$ changes from the ground state to the saddle point by a few MeV only.
Nowadays, there is a strong evidence that the dynamical equation for fission of hot nuclei should take into account the coupling of $`Q`$ with the single-particle degrees of freedom, which causes a slowing down of the fission time scale to values exceeding the single-particle times by a few orders of magnitude. This allows one to assume that at each instant of time the single-particle degrees of freedom are in thermal equilibrium under the constraint of the known $`Q`$ and the respective velocity $`\dot{Q}`$. Expressed in formal terms, this means that the entropy $`S`$ of the system is maximal for the given constraints:
$$\frac{\delta S}{\delta \rho }=0,\rho (𝐫)𝑑𝐫=A,Q=Q^{},\dot{Q}=\dot{Q}^{}.$$
(3)
Equations (3) define a state of partial thermal equilibrium , or quasistatic state, for short.
Denoting the intrinsic excitation energy by $`E_x`$ and using the Fermi gas formula for the entropy
$$S=2\sqrt{a(E_x\stackrel{~}{E}\{\rho \}+\stackrel{~}{E}\{\rho _{\mathrm{gs}}\})},$$
(4)
one can find the density $`\rho `$ of the quasistatic state by looking for the conditional minimum of $`\stackrel{~}{E}`$. This is justified by the fact that the variation of the level density parameter $`a`$ with shape is very smooth in comparison to that of $`\stackrel{~}{E}`$ . By changing $`Q^{}`$ one gets a sequence of quasistatic states which we call the quasistatic path. The fact that fission proceeds along the quasistatic path leads to significant simplifications in the formal description of this process. It also allows us to employ the statistical-model treatment of particle emission for the fissioning nucleus, because the single-particle degrees of freedom are in thermal equilibrium at each point of the quasistatic path.
Another process in which the quasistatic shapes can be useful is quasifission. The complete thermal equilibrium of the single-particle motion at a given shape may be reached during the reseparation stage of a quasifission reaction if in the fusion stage the mass asymmetry mode and the relative velocity of the two colliding nuclei have relaxed. As a result, the quasifission trajectories in the reseparation stage will get on the quasistatic path. In practical terms, quasifission reactions are described by the one-body dissipation model of heavy ion collisions implemented in the code HICOL, while the conditional equilibrium densities can be calculated in the framework of the extended Thomas–Fermi (ETF) model of non-spherical nuclei . With these two models we wish to verify whether the dynamic trajectories in quasifission reactions follow the quasistatic path.
Our ETF calculations will be confined to mass and axially symmetric (about the $`z`$ axis) prolate density distributions $`\rho (r,z)`$ normalized to the mass number $`A`$, where $`r=\sqrt{x^2+y^2}`$ and $`x`$, $`y`$, $`z`$ are the Cartesian coordinates. The energy density of the ETF model incorporates second-order gradient corrections with spin-orbit and effective mass terms , which are very important in describing the nuclear surface. We have performed the ETF calculations using a realistic Skyrme interaction, namely SkM . From the ETF–SkM functional, we obtain fully self-consistent nuclear densities by solving the associated variational Euler–Lagrange equations in cylindrical coordinates, imposing a given value of the quadrupole moment $`Q_2`$ :
$$Q_2=2\pi _0^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\left[2z^2r^2\right]\rho (r,z)r𝑑r𝑑z.$$
(5)
To account for nuclear rotation with angular momentum $`J`$ we have included a rotational energy
$$E_{\mathrm{rot}}=\frac{J^2}{2I}$$
(6)
to the ETF energy functional. Here, $`I`$ is the rigid-body moment of inertia
$$I=\pi _0^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\left[2z^2+r^2\right]m\rho (r,z)r𝑑r𝑑z,$$
(7)
where $`m`$ is the nucleon mass. It was assumed that the spin axis is directed perpendicular to the symmetry axis of the compound nucleus.
In HICOL the constant density approximation is used. Therefore, the nuclear density is completely determined by the profile function $`y(z)`$ whose rotation about the symmetry axis $`z`$ generates the nuclear surface. This poses difficulties in the comparison between dynamic and quasistatic densities, which can nevertheless be avoided by introducing some generalized characteristics of the nuclear densities.
## 3 Geometry of the Composite System
For the sake of comparison between dynamic and quasistatic paths we introduce the elongation coordinate $`D_{\mathrm{mm}}`$ and the neck coordinate $`R_{\mathrm{neck}}`$. In axially symmetric nuclei these quantities are given by
$$D_{\mathrm{mm}}=\frac{8\pi }{A}_0^{\mathrm{}}_0^{\mathrm{}}z\rho (r,z)r𝑑r𝑑z$$
(8)
and
$$R_{\mathrm{neck}}^2=\frac{2}{\rho _0}_0^{\mathrm{}}r\rho (r,z=0)𝑑r.$$
(9)
The elongation coordinate $`D_{\mathrm{mm}}`$ defines the distance between the centers of mass of the two halves of the nucleus. It was used by Strutinsky as a constraint operator in the integro-differential equation for the profile function $`y(z)`$ of leptodermous nuclei. Our definition (9) of the neck radius is obtained from the requirement that a nucleus with constant density $`\rho _0`$ and a geometrical neck radius equal to $`R_{\mathrm{neck}}`$, has the same number of particles in the cross section of its neck as the nucleus having the distributed density $`\rho `$. In the following, $`\rho _0`$ in Eq. (9) will be identified with the one used in the code HICOL, namely $`\rho _0=A/(\frac{4}{3}\pi R_0^3)`$ where $`R_0=1.18A^{1/3}`$ fm. From Eqs. (8) and (9) one finds that for a spherical nucleus with a constant density $`\rho _0`$, the elongation $`D_{\mathrm{mm}}`$ is equal to $`\frac{3}{4}R_0`$ and that the neck radius is $`R_{\mathrm{neck}}=R_0`$.
In HICOL the profile function $`y(z)`$ of the composite system is parameterized by two spheres smoothly connected by a second-order curve . For these so-called Blocki shapes we define $`D_{\mathrm{mm}}`$ as the distance between the centers of mass of the two parts of the nucleus on both sides of a plane $`z=z_\mathrm{m}`$. For $`z_\mathrm{m}`$ we take the mean value of the left and right matching points ($`z_\mathrm{m}=0`$ for symmetric shapes). We identify $`R_{\mathrm{neck}}`$ of the Blocki profile with $`y(z_\mathrm{m})`$.
For mass symmetric shapes, the Blocki profile function reads
$$y^2(z)=\{\begin{array}{cc}R_1^2(z+s/2)^2\hfill & \text{for }R_1s/2zz_1,\hfill \\ \alpha +\beta z^2\hfill & \text{for }z_1zz_1,\hfill \\ R_1^2(zs/2)^2\hfill & \text{for }z_1zs/2+R_1.\hfill \end{array}$$
(10)
Given the volume $`V_0`$ of the nucleus, the parameters $`R_1`$, $`z_1`$, $`\alpha `$ and $`\beta `$ can be expressed in terms of the two collective degrees of freedom $`s`$ and $`\sigma `$, where $`s`$ is the distance between the centers of the spheres and
$$\sigma =\frac{V_08\pi R_1^3/3}{V_0}$$
(11)
is a measure of the constriction of the system.
Indeed, from Eq. (11) one obtains
$$R_1=R_0\left(\frac{1\sigma }{2}\right)^{1/3},$$
(12)
where $`R_0`$ is the radius of a spherical nucleus of volume $`V_0`$. If we equate the values of $`y(z)`$ and its derivative on both sides of $`z=z_1`$ and require that the total volume of the shape generated by $`y(z)`$ equals $`V_0`$, we get
$`\alpha `$ $`=`$ $`R_1^2+{\displaystyle \frac{s}{2}}\left(z_1{\displaystyle \frac{s}{2}}\right),`$ (13)
$`\beta `$ $`=`$ $`{\displaystyle \frac{s}{2z_1}}1,`$ (14)
$`z_1`$ $`=`$ $`{\displaystyle \frac{s}{2}}\sqrt{1G},`$ (15)
with
$$G=\frac{3}{s_0^2}\frac{2}{s_0^3}\frac{1+\sigma }{1\sigma },s_0=\frac{s}{2R_1}.$$
(16)
For the Blocki shapes defined by Eq. (10), Eq. (8) yields
$$D_{\mathrm{mm}}=\frac{3}{4R_0^3}\left\{\frac{s^4}{48}\left[(1G)^{3/2}1\right]+\frac{1}{2}s^2R_1^2+\frac{4}{3}sR_1^3+R_1^4\right\}.$$
(17)
The value of $`R_{\mathrm{neck}}^2`$ coincides with $`y^2(0)`$ and according to (10) is equal to $`\alpha `$. Thus, using Eqs. (13) and (15) we obtain
$$R_{\mathrm{neck}}^2=R_1^2+\frac{s^2}{4}\left(\sqrt{1G}1\right).$$
(18)
Given $`R_0`$, Eqs. (17) and (18) together with Eqs. (12) and (16) allow one to express $`D_{\mathrm{mm}}`$ and $`R_{\mathrm{neck}}`$ in terms of $`s`$ and $`\sigma `$.
In the following section, instead of the quadrupole moment, we use sometimes the moment of inertia as a constraint operator. It should be noted that for mass symmetric Blocki shapes these quantities can be found analytically:
$$Q_2=4\pi \rho _0\left[\frac{h(R_1+h)^3}{4}\left(R_1\frac{h}{3}\right)\frac{h}{6}(R_1^2h^2)z_1^2\frac{2}{15}h^2z_1^3+\frac{hz_1^4}{20}\right],$$
(19)
$`I=\pi m\rho _0[{\displaystyle \frac{(R_1+h)^3}{5}}({\displaystyle \frac{8}{3}}R_1^2{\displaystyle \frac{R_1h}{2}}+{\displaystyle \frac{h^2}{6}})`$
$`+{\displaystyle \frac{h}{3}}(R_1^2h^2)z_1^2+{\displaystyle \frac{4}{15}}h^2z_1^3+{\displaystyle \frac{hz_1^4}{30}}],`$ (20)
where $`h=s/2`$.
## 4 Dynamic and Quasistatic Paths
We now describe calculations performed in order to compare the dynamic trajectories with the quasistatic path. The calculations are carried out in the $`A160`$ composite mass region, a region of continuous experimental efforts . We start with the quasifission reaction following a <sup>60</sup>Ni$`+^{100}`$Mo collision at the beam energy $`E=600`$ MeV. In Fig. 1 we display the equidensity contour plots corresponding to the sequence of the ETF densities of conditional equilibrium for <sup>160</sup>Yb, which represents the composite system in the collision. For each ETF density shown we have calculated the values of $`D_{\mathrm{mm}}`$ and $`R_{\mathrm{neck}}`$ by means of Eqs. (8) and (9). Inserting these $`D_{\mathrm{mm}}`$ and $`R_{\mathrm{neck}}`$ into Eqs. (17) and (18) and solving these equations with respect to $`s`$ and $`\sigma `$, we can prescribe the Blocki profiles to the density distributions. From Fig. 1 one can see that such profiles are close, in general, to the $`\frac{1}{2}\rho _0`$ curves of the ETF density distributions. In dinuclear configurations, however, the nascent fragments predicted by the ETF model are somewhat flattened in the $`r`$ direction compared to the spherical form.
The densities shown in Fig. 1 were calculated accounting for the rotational energy with angular momentum $`J`$= 86 $`\mathrm{}`$. The reasons for this choice of $`J`$ will become clear in the next section. The calculations indicate that the quasistatic path depends weakly on the angular momentum. When one goes from $`J`$=86 $`\mathrm{}`$ to $`J`$=0 $`\mathrm{}`$ the quasistatic $`D_{\mathrm{mm}}`$ and $`R_{\mathrm{neck}}`$ change at most by 0.1 fm. In earlier ETF calculations the energy $`\stackrel{~}{E}`$ of the rotating nucleus along the quasistatic path was found to be very close to the sum of the energy of the non-rotating nucleus and the rotational energy computed with the moment of inertia of the latter. Moreover it was shown that the level density parameter is affected by rotation in a negligible amount if $`Q_2`$ is fixed. These findings were interpreted as an indication that rotation has a small influence on the nuclear densities calculated in conditional equilibrium. Our direct calculations of shape parameters along the quasistatic path agree with this conclusion.
We repeated the calculation of the $`J`$= 86 $`\mathrm{}`$ quasistatic path with the constraint on the moment of inertia $`I`$, instead of $`Q_2`$. The path in the $`(D_{\mathrm{mm}},R_{\mathrm{neck}})`$ space turned out to be almost the same as the one we had found with the $`Q_2`$ constraint: the differences in $`R_{\mathrm{neck}}`$ are no larger than $``$1%. Imposing the $`I`$ constraint we not only obtain the same values for $`Q_2`$ as with the $`Q_2`$ constraint, but also find that the hexadecapole moment $`Q_4`$ along the path is very similar. This means that the nuclear shapes must be equivalent with the $`I`$ or $`Q_2`$ constraint, as well.
The dynamic evolution of the shape of the <sup>60</sup>Ni$`+^{100}`$Mo composite system (at $`J`$=86 $`\mathrm{}`$ and $`E`$= 600 MeV) computed with HICOL is shown in Fig. 2. In units of $`10^{21}\mathrm{s}`$, the first stage (when neck fills in) takes about 0.2, the mononucleus lives about 30 and the scission stage lasts about 5. At the end of the first stage the individual temperatures and angular velocities of the two nuclei, predicted by HICOL, become approximately equal. The time dependence of $`D_{\mathrm{mm}}`$ for different $`J`$ values is depicted in Fig. 3, which shows a clear separation of the reactions into the fusion stage taking a fraction of $`10^{21}\mathrm{s}`$ and a much longer reseparation stage. The duration time of the latter strongly changes from one $`J`$ to another. For example, this time decreases from $`45\times 10^{21}\mathrm{s}`$ to $`15\times 10^{21}\mathrm{s}`$ when $`J`$ increases from 85 $`\mathrm{}`$ to 100 $`\mathrm{}`$.
Figure 4 shows on a ($`D_{\mathrm{mm}},R_{\mathrm{neck}}`$) plot how the dynamic trajectories with $`J`$=70, 75, 80, 85, 90, 95, 100, 105 $`\mathrm{}`$ are joining the quasistatic path determined for $`J`$= 86 $`\mathrm{}`$. The systems with $`J`$ =70 and 75 $`\mathrm{}`$ terminate at the different points of this path (in the case of $`J`$=75 $`\mathrm{}`$ this happens after a slight rebound) while those with $`J`$ =80–105 $`\mathrm{}`$, having got on the quasistatic path shortly after rebound, proceed towards the scission point. In the vicinity of the scission point, they start to deviate progressively from the quasistatic path and yield a thinner neck for the same value of the elongation. We have found the dynamic paths, in the reseparation phase, to be rather stable against variations in the excitation energy and mass asymmetry. This follows from our dynamic calculations performed for the system <sup>60</sup>Ni+<sup>100</sup>Mo at the beam energy of 1200 MeV (Fig. 5), and for <sup>48</sup>Ca on <sup>112</sup>Sn at the beam energy of 480 MeV (Fig. 6).
Very small but regular deviations of the dynamic trajectories in their slowest phase from the quasistatic path are clearly observed in Figs. 4, 5 and 6. They are probably related to the fact that the dynamical and variational calculations involve different forces, and that in the ETF calculation we varied the whole density rather than just the nuclear profile. To verify this assumption we calculated the quasistatic path for $`J`$= 86 $`\mathrm{}`$, with a $`Q_2`$ constraint, using the Yukawa-plus-exponential (YPE) forces which are employed in the code HICOL. We looked for the conditional minimum of the system energy in the class of mass symmetric Blocki shapes. The result is displayed on Fig. 6 with square symbols. According to this figure, the dynamical paths in their slowest part practically coincide with the quasistatic YPE path.
It is interesting to note that the YPE quasistatic path predicts the onset of the scission stage at noticeably smaller values of $`D_{\mathrm{mm}}`$ than the Skyrme path. This is consistent with the fact that the saddle point configuration for the YPE force is more compact than for the SkM force. The deformation energy for the YPE force reaches its maximum (of about 27.8 MeV) at $`D_{\mathrm{mm}}13.6`$ fm, whereas the SkM energy reaches its maximum (about 25.1 MeV) at $`D_{\mathrm{mm}}15.6`$ fm. These deviations manifest the scale of errors introduced into variational calculations by simple parameterizations of nuclear shapes.
## 5 Particle Emission
In the preceding section we have shown that the dynamic trajectories in the reseparation stage closely follow the quasistatic path. This means that statistical models of particle emission can be applied in this stage. Below, we describe the technique for light-particle evaporation calculations along the slow phase, which will allow us to perform comparisons with experimental data.
Our procedure is based on a Monte Carlo simulation of particle decay chains in nuclei with a time-dependent shape. Given the excitation energy $`E_x`$, angular momentum $`J`$, the dynamic path and the time $`t_\mathrm{i}`$ for the beginning of the slow phase, we calculate the neutron ($`R_\mathrm{n}`$), proton ($`R_\mathrm{p}`$) and alpha particle ($`R_\alpha `$) emission rates. We assume that the emission times $`t_\mathrm{e}`$ are distributed according to the exponential law $`\mathrm{exp}(R_{\mathrm{tot}}t_\mathrm{e})`$, where $`R_{\mathrm{tot}}=R_\mathrm{n}+R_\mathrm{p}+R_\alpha `$ is the total emission rate. A specific value of $`t_\mathrm{e}`$ is chosen using a generator of exponentially distributed random numbers. Having sampled the type of the emitted particle in proportion to the weights $`R_\mathrm{n}/R_{\mathrm{tot}}`$, $`R_\mathrm{p}/R_{\mathrm{tot}}`$ and $`R_\alpha /R_{\mathrm{tot}}`$, we find the average excitation energy $`\overline{u}`$ and root mean square angular momentum $`j_{\mathrm{rms}}=\sqrt{\overline{j^2}}`$ of the corresponding daughter nucleus and take the latter as the new decaying nucleus. Assuming its shape to be the same as that of the parent nucleus at $`t_\mathrm{i}+t_\mathrm{e}`$, we simulate the next decay. The chain of decays is terminated when the emission time exceeds the scission time of the nucleus and one proceeds to the next chain.
The pre-scission multiplicities $`M_\nu (J)`$ are calculated as
$$M_\nu (J)=\frac{1}{𝒩}\underset{i=1}{\overset{𝒩}{}}𝒩_\nu (i),$$
(21)
where $`𝒩_\nu (i)`$ is the number of particles of type $`\nu `$ ($`\nu =`$ n, p, $`\alpha `$) in a decay chain classified with the index $`i`$. Here, $`𝒩`$ is the total number of decay chains for a given $`E_x`$ and $`J`$. The explicit expressions for $`R_\nu `$, $`\overline{u}`$, $`\overline{j^2}`$ are summarized in the Appendix. They were obtained in the framework of the classical statistical model of particle emission from non-spherical nuclei . As shown in Ref. , the effects of the shape distortions on particle emission are treated by these formulas more rigorously than in heuristic models .
The input parameters of these formulas are the effective separation energies $`S_\nu ^{\mathrm{eff}}`$ calculated including deformation energies , the level density parameter $`a`$, the height $`V_\mathrm{b}`$ and the radius $`R_\mathrm{b}`$ of the corresponding spherical barrier experienced by a particle. In the following calculations we employ the YPE values of $`S_\nu ^{\mathrm{eff}}`$ for $`\nu =`$ n, p, $`\alpha `$. In Fig. 7 the YPE values (solid lines) are compared with the ETF separation energies (dashed lines) on the quasistatic path in <sup>160</sup>Yb. The ratio $`I/I_0`$ of the deformed nucleus moment of inertia to that of the spherical one is used as the coordinate along the path.
The level density parameters $`a`$ to be used later have been normalized at the spherical shape to the experimental value $`A/8.8`$ MeV<sup>-1</sup> obtained in Ref. . The shape dependence of $`a`$ has been calculated with the YPE forces following the prescription of Tõke and Swiatecki . As seen from Fig. 8, these $`a`$ are close, by magnitude and shape dependence, to the $`a`$ values from the ETF method. The dependence of $`V_\mathrm{b}`$ and $`R_\mathrm{b}`$ for n, p, $`\alpha `$ on $`A`$ and $`Z`$ is parameterized in the same way as in Ref. . The energy of the emitted neutron entering the corresponding $`R_\mathrm{b}`$ is replaced with its mean value $`2\tau `$, where $`\tau `$ is the temperature of the daughter nucleus.
Recently, Lou et al measured multiplicities of light particles in fission reactions of 10 MeV/u <sup>63</sup>Cu+<sup>92,100</sup>Mo and 20 MeV/u <sup>20</sup>Ne+<sup>144,148,154</sup>Sm. Our analysis of these data along with the data of Gonin et al on <sup>60</sup>Ni+<sup>100</sup>Mo is presented in Fig. 9 and Table 1. The measurements of multiplicities of n, p, $`\alpha `$ in the latter work were performed at 9.2 and 10.9 MeV/u. The data we analyze are obtained by interpolation of the reported values to 10 MeV/u.
In the analysis we took into account that in these reaction systems an appreciable amount of particles escape from the system during the pre-equilibrium stage. The measurements of linear-momentum transfer from projectile to target allowed Lou et al to estimate that the mass and charge removed in this stage are $`(\delta A,\delta Z)`$= (8,4), (8,4), (6,3), (7,3) and (6,3) in the reactions <sup>63</sup>Cu+<sup>92,100</sup>Mo and <sup>20</sup>Ne+<sup>144,148,154</sup>Sm, respectively. Since experimental information on $`\delta A`$ and $`\delta Z`$ for the <sup>60</sup>Ni+<sup>100</sup>Mo system is not available, we used the values $`(\delta A,\delta Z)`$=(8,4) in the closest system, <sup>63</sup>Cu+<sup>100</sup>Mo. Lou et al estimated the initial excitation energies of the equilibrated compound nuclei to be $`E_x=`$ 227, 267, 289, 277 and 282 MeV in the reactions of 10 MeV/u <sup>63</sup>Cu+<sup>92,100</sup>Mo and 20 MeV/u <sup>20</sup>Ne+<sup>144,148,154</sup>Sm, respectively. For the system of 10 MeV/u <sup>60</sup>Ni+<sup>100</sup>Mo, the linear interpolation between the excitation energies $`E_x`$= 251 MeV and 293 MeV at 9.2 MeV/u and 10.9 MeV/u, respectively, obtained in Ref. , results in $`E_x`$= 271 MeV.
This information on $`\delta A`$, $`\delta Z`$ and $`E_x`$ was used as input in our pre-scission (equilibrium) multiplicity calculations. For all systems we used $`t_\mathrm{i}=1.1\times 10^{21}`$s and simulated $`𝒩=200`$ chains of decays. For each system we took only one HICOL trajectory, namely the trajectory whose contact-to-scission time $`t_{\mathrm{cs}}`$ is closest to $`40\times 10^{21}`$ s. This condition results in $`J`$= 78, 86, 86, 85, 88, 93 $`\mathrm{}`$ in the reactions of <sup>63</sup>Cu+<sup>92</sup>Mo, <sup>60</sup>Ni+<sup>100</sup>Mo, <sup>63</sup>Cu+<sup>100</sup>Mo at 10 MeV/u and <sup>20</sup>Ne+<sup>144,148,154</sup>Sm at 20 MeV/u, respectively. The slow stage of the so-chosen (‘effective’) trajectories lasts some $`35\times 10^{21}`$ s. This is consistent with the fission time scale of $`(35\pm 15)\times 10^{21}`$ s deduced from a systematic study of pre-scission neutron multiplicities in 27 fission reactions induced by <sup>16,18</sup>O, <sup>40</sup>Ar and <sup>64</sup>Ni on targets with $`A=141`$–238 .
In the post-scission emission calculations, the thermal energy of the composite system at the moment of scission was shared between the complementary fragments in proportion to their masses. The spins of the fragments (about 6 $`\mathrm{}`$ per fragment) were taken from HICOL output. To find $`A`$ and $`Z`$ of the primary fragments, we used the calculated n, p and $`\alpha `$ pre-scission multiplicities (see Table 1). The calculated post-scission multiplicities of neutrons in all reactions but <sup>63</sup>Cu+<sup>92</sup>Mo are confined between the value of 3.7$`\pm `$0.4 measured in the system of <sup>16</sup>O+<sup>154</sup>Sm at $`E_x`$= 206 MeV and the value of 4$`\pm `$1.1 for <sup>60</sup>Ni+<sup>100</sup>Mo at 10 MeV/u which follows from interpolation of the 9.2 and 10.9 MeV/u data ($`3.6\pm 1`$ and $`4.5\pm 1.2`$, respectively) of Ref. .
With accounting for post-scission emission, which is essential, in fact, only for neutrons, the total (equilibrium) multiplicities appear to be within the likely systematic uncertainties of the experimental points. The only noticeable exceptions are the total proton multiplicities in <sup>63</sup>Cu+<sup>92</sup>Mo and <sup>60</sup>Ni+<sup>100</sup>Mo, when the measured values are smaller than the calculated ones by factors of 1.6 and 2.1, respectively. With the exception of these two data points, the overall agreement of the calculations with the rest of the data signifies a consistency of the neutron with the charged particle clock. The low proton multiplicities in <sup>60</sup>Ni+<sup>100</sup>Mo reactions observed in Ref. have motivated further experimental studies. In a recent work, Charity et al studied the nearby system <sup>64</sup>Ni+ <sup>100</sup>Mo at a similar excitation energy and found 2–3 times greater multiplicities of fusion-associated p and $`\alpha `$ compared to those of Ref. .
It should be noted that no fitting parameters were used in our analysis. The employment of the transmission coefficients for p and $`\alpha `$ from Ref. would destroy the quality of the description by strongly enhancing the $`\alpha `$ emission. For instance, the pre-scission multiplicities in the <sup>60</sup>Ni+<sup>100</sup>Mo system become 6.73, 2.29 and 1.81 for n, p and $`\alpha `$, respectively, instead of 7.27, 2.11 and 0.98 obtained with the transmission coefficients from Ref. .
## 6 Discussion and Conclusion
To shed more light on the role of quasifission trajectories in the reactions of our study, it is useful to estimate the fusion-fission $`J`$-window and compare it with available experimental information on the angular momenta, associated with the evaporation residue and fission cross sections. Nuclei emerging at the end of the evaporation cascades undergo fission if their angular momentum is confined between the angular momentum where the fission barrier height $`B_\mathrm{f}(J)`$ equals the neutron separation energy and the angular momentum where $`B_\mathrm{f}(J)`$ vanishes . To reconstruct the corresponding $`J`$-window in the beginning of the evaporation cascades, one has to account for the angular momentum removed by pre-scission particles.
In the case of <sup>60</sup>Ni+<sup>100</sup>Mo at 10 MeV/u, we find that the fusion-fission $`J`$-window at the end of the evaporation cascades is 57–82 $`\mathrm{}`$, on the average. Light particles evaporated along the effective trajectory, remove on the average about 21 $`\mathrm{}`$. Therefore, we estimate the corresponding $`J`$-window in the beginning of the equilibrium emission stage, i.e. for <sup>152</sup>Dy, to be 78–103 $`\mathrm{}`$. Since $`B_\mathrm{f}(J)`$ for <sup>152</sup>Dy vanishes at 83 $`\mathrm{}`$, trajectories within this window, with the exception of those with $`J`$=78–83 $`\mathrm{}`$, belong to quasifission.
It is interesting to note that the width of the so-defined $`J`$-window is close to the ‘total’ (fusion-fission plus quasifission) $`J`$-window widths in the closest systems to ours, where data are available. A fission $`J`$-window of 70–103 $`\mathrm{}`$ and 49–67 $`\mathrm{}`$ is implied in studies of <sup>40</sup>Ar+<sup>109</sup>Ag at 8.4 MeV/u and <sup>20</sup>Ne+<sup>159</sup>Tb at 16 MeV/u , respectively. Thus we may conclude that the majority of fission-evaporation reactions start on quasifission trajectories which after losing angular momentum end up on trajectories in the potential with a non-zero fission barrier.
To summarize, the present work was motivated by the desire to estimate the perspective for the inclusion of quasifission into fission-evaporation codes. Towards this aim we tested whether the quasifission trajectories, in the reseparation stage, follow the fission path. Our calculations were performed in the $`A`$ 160 region at the bombarding energy of 10–20 MeV/u.
For the sake of comparison between different shapes (assumed to be axially symmetric) we use the two-dimensional space of the collective variables describing elongation and constriction. The elongation is characterized by the distance between the two halves of the nucleus. The constriction is described by the neck radius defined in a way applicable for distributed densities. In this collective space, we first check the sensitivity of the quasistatic path to the angular momentum of the system and to the form of the constraint operator. The quasistatic paths found for different angular momenta (including those exceeding the critical angular momentum for fission) are practically indistinguishable from the one with $`J`$= 0 $`\mathrm{}`$. The moment of inertia as a constraint operator was found to generate a sequence of shapes which coincides with the one obtained using the constraint on the quadrupole moment.
The same space of collective variables was used to compare the sequences of shapes predicted by the HICOL code at different values of the entrance-channel angular momentum. The parts of the quasifission trajectories corresponding to the slow stages of the evolution show the eventual convergence to the quasifission trajectory with the minimal $`J`$. This latter is found to follow closely the quasistatic path obtained with ETF model. This close coincidence of dynamic trajectories with the quasistatic path was found to occur in a wide range of bombarding energies and for quite different entrance-channel mass asymmetries.
Since the dynamics of quasifission reactions is well described by quasistatic paths, at least during the slowest phase, the statistical evaporation model can be applied for the description of particle emission from such systems. This follows from the fact that at each point of the quasistatic path the system reaches thermal equilibrium.
Using single effective dynamical trajectories whose slow stage lasts some $`35\times 10^{21}`$ s, it was made possible to reproduce reasonably well experimental data on pre- and post-scission multiplicities of neutrons and total multiplicities of protons and $`\alpha `$-particles emitted from thermally equilibrated systems. This agreement was achieved without any ad hoc statistical model parameter adjustments and shows a consistency between the neutron and charged particle clock. The duration time of the slow stage of the employed dynamical trajectories was found consistent with the results of systematic studies.
## Acknowledgments
This research was supported by Grants No. PB98-1247 from the DGICYT (Spain), 1998SGR-00011 from the DGR (Catalonia), the General Secretariate of Research and Technology of the Ministry of Industry, Research and Technology of Greece ($`\mathrm{\Pi }ENE\mathrm{\Delta }`$ ’95 No. 696), and a NATO Fellowship Program for the Year 1996-97. Useful comments of Profs. W. Swiatecki and P. Schuck are greatly appreciated. We are thankful to Dr. Th. Keutgen for communicating us information on experimental data before publication.
## Appendix
In the following we outline the formalism we used for the calculation of decay rates for statistical particle emission from equilibrated compound nuclei. The excitation energy, angular momentum, deformation energy and moment of inertia of the parent nucleus are denoted as $`E_x`$, $`J`$, $`E_{\mathrm{def}}`$ and $`I_x`$, respectively. Given these quantities, the thermal energy of the parent nucleus is defined by
$$q_x=E_x\frac{J^2}{2I_x}E_{\mathrm{def}}.$$
Its reduced level density is
$$\omega _x(q_x)=\frac{1}{t_x^4(I_xa_x)^{\frac{3}{2}}}\mathrm{exp}\left[2\sqrt{a_xq_x}\right],$$
where
$$t_x=\frac{3}{4a_x}+\sqrt{\left(\frac{3}{4a_x}\right)^2+\frac{q_x}{a_x}},$$
and $`a_x`$ is the level density parameter. Similar formulas are used for the reduced level density $`\omega (q)`$ of the daughter nucleus with thermal energy $`q`$. Its moment of inertia and level density parameter are denoted by $`I`$ and $`a`$.
The mass, spin and effective separation energy of the emitted particle are denoted as $`m_\nu `$, $`s_\nu `$ and $`S_\nu ^{\mathrm{eff}}`$, respectively. Let $`z_{\mathrm{matter}}`$ be half a length of matter distribution in the deformed shape, and let $`R_\mathrm{b}`$ and $`R_{\mathrm{matter}}`$ be the barrier radius and matter radius for the spherical shape. Then half a length of the figure generated by the barrier line is postulated to be
$$z_0=z_{\mathrm{matter}}+(R_\mathrm{b}R_{\mathrm{matter}}).$$
Given the matter profile function $`y=y(z)`$ and the Coulomb potential $`\mathrm{\Phi }(z)`$ along this profile, the barrier line $`\rho (z)`$ and the potential barrier $`U(z)`$ along this line were calculated from
$$\rho (z)=Ky\left(\frac{z}{K}\right),U(z)=\frac{V_\mathrm{b}}{\mathrm{\Phi }_0}\mathrm{\Phi }\left(\frac{z}{K}\right),$$
where $`K=z_0/z_{\mathrm{matter}}`$ is the scaling factor, $`V_\mathrm{b}`$ is the $`s`$-wave potential barrier in the spherical nucleus and $`\mathrm{\Phi }_0`$ is the Coulomb potential on the edge of the spherical matter distribution.
The key characteristics of the residual nucleus are its average thermal energy $`\overline{q}`$ and temperature $`\tau `$ corresponding to this thermal energy. These quantities are calculated from
$$\overline{q}E_xS_\nu ^{\mathrm{eff}}E_{\mathrm{def}}\frac{J^2}{2I}V_\mathrm{b},$$
$$\tau =\frac{2}{a}+\sqrt{\left(\frac{2}{a}\right)^2+\frac{\overline{q}}{a}}.$$
With these definitions the particle emission rate is given by
$$R_\nu =\frac{2s_\nu +1}{2\pi }m_\nu z_0^2\tau ^2\frac{\omega (\overline{q})}{\omega _x(q_x)}\mathrm{exp}\left[\frac{1}{\tau }\left(E_xS_\nu ^{\mathrm{eff}}E_{\mathrm{def}}\frac{J^2}{2I}\overline{q}\right)\right]G(\tau ,b),$$
where
$$G(\tau ,b)=_1^1𝑑\zeta \sqrt{\eta ^2+\eta ^2\left(\frac{d\eta }{d\zeta }\right)^2}\mathrm{exp}\left[\frac{U(\zeta z_0)}{\tau }+b\zeta ^2+\frac{1}{2}b\eta ^2\right]I_0\left(\frac{1}{2}b\eta ^2\right),$$
$$b=b(\tau )=\frac{m_\nu z_0^2}{2\tau }\left(\frac{J}{I}\right)^2,\eta =\eta (\zeta )=\frac{\rho (\zeta z_0)}{z_0},$$
and $`I_0(x)`$ is a Bessel function of the first kind of imaginary argument.
The averaged square of the angular momentum of the daughter nucleus is
$$\overline{j^2}=J^24I\tau b\frac{d\mathrm{ln}G(\tau ,b)}{db}.$$
Given this quantity and $`\overline{q}`$, the average excitation energy of the daughter nucleus reads
$$\overline{u}=\overline{q}+\frac{\overline{j^2}}{2I}+u_{\mathrm{def}},$$
where $`u_{\mathrm{def}}`$ is the daughter nucleus deformation energy.
The above expressions give approximate values of $`R_\nu `$, $`\overline{j^2}`$ and $`\overline{u}`$ because they use approximate values of $`\overline{q}`$ and $`\tau `$. To correct these latter quantities we find a new value of $`\overline{q}`$ using the formula
$$\overline{q}=E_xS_\nu ^{\mathrm{eff}}E_{\mathrm{def}}\frac{J^2}{2I}2\tau \tau ^2\frac{d\mathrm{ln}G(\tau ,b(\tau ))}{d\tau }$$
at the initial value of $`\tau `$ and insert it into the formula for $`\tau `$. This procedure was iterated until convergence of $`\tau `$ was attained.
Table 1. Comparison between experimental and calculated light-particle multiplicities for the reaction systems studied in Refs. and . The neutron number $`N`$ of the composite system after the pre-equilibrium emission stage is indicated. From left to right the systems are 10 MeV/u <sup>63</sup>Cu+<sup>92</sup>Mo, <sup>60</sup>Ni+<sup>100</sup>Mo, <sup>63</sup>Cu+<sup>100</sup>Mo and 20 MeV/u <sup>20</sup>Ne+<sup>144,148,154</sup>Sm
| $`N`$ | 80 | 86 | 88 | 89 | 92 | 99 |
| --- | --- | --- | --- | --- | --- | --- |
| n<sup>pre</sup> | 3.75 | 7.27 | 7.46 | 7.65 | 8.2 | 10.87 |
| n<sup>post</sup> | 1.79 | 3.84 | 3.87 | 3.82 | 3.89 | 3.97 |
| n<sup>tot</sup> | 5.54 | 11.11 | 11.33 | 11.47 | 12.09 | 14.84 |
| n<sup>tot,exp</sup> | 6(1) | 11.25(2.7) | 10.9(1.5) | 10.3(1.3) | 12.2(1.6) | 16.9(1.2) |
| p<sup>pre</sup> | 2.97 | 2.11 | 1.97 | 2.65 | 1.85 | 0.82 |
| p<sup>post</sup> | 0.18 | 0.06 | 0.05 | 0.07 | 0.05 | 0.02 |
| p<sup>tot</sup> | 3.15 | 2.17 | 2.02 | 2.72 | 1.9 | 0.84 |
| p<sup>tot,exp</sup> | 2.02(0.3) | 1.03(0.15) | 1.6(0.3) | 2.14(0.3) | 1.83(0.5) | 0.97(0.5) |
| $`\alpha ^{\mathrm{pre}}`$ | 1.17 | 0.98 | 0.93 | 1.24 | 1. | 0.66 |
| $`\alpha ^{\mathrm{post}}`$ | 0.03 | 0.02 | 0.01 | 0.02 | 0.01 | 0.01 |
| $`\alpha ^{\mathrm{tot}}`$ | 1.2 | 1 | 0.94 | 1.26 | 1.01 | 0.66 |
| $`\alpha ^{\mathrm{tot},\mathrm{exp}}`$ | 1(0.3) | 0.77(0.11) | 1.69(0.4) | 1.3(0.3) | 0.7(0.4) | 0.63(0.4) |
## Figure Captions
Fig. 1. Equidensity contours for <sup>160</sup>Yb calculated by the ETF method at the indicated values of the quadrupole moment, with angular momentum $`J`$= 86 $`\mathrm{}`$. From outside to inside the lines represent contours of constant density $`\rho =`$ 0.1, 0.3, 0.5, 0.7, 0.9 and 1.1, in units of $`\rho _0`$. The dashed curves represent the Blocki profile with the parameters $`s`$ and $`\sigma `$ calculated in terms of $`D_{\mathrm{mm}}`$ and $`R_{\mathrm{neck}}`$ obtained from the corresponding ETF density distribution.
Fig. 2. Time evolution of the <sup>60</sup>Ni$`+^{100}`$Mo system from the code HICOL at $`E=600`$ MeV and $`J`$= 86 $`\mathrm{}`$. The time is indicated in units of $`10^{21}\mathrm{s}`$.
Fig. 3. Time dependence of $`D_{\mathrm{mm}}`$ for the reaction <sup>60</sup>Ni+<sup>100</sup>Mo at $`E=600`$ MeV and $`J`$= 79-105 $`\mathrm{}`$ predicted by the HICOL code.
Fig. 4. Dynamic (for $`J=70`$, 75, 80, 85, 90, 95, 100, 105 $`\mathrm{}`$) and quasistatic (for $`J`$= 86 $`\mathrm{}`$) paths in the <sup>60</sup>Ni+<sup>100</sup>Mo collision at 600 MeV are shown as solid and dashed lines, respectively. The spherical shape is shown as a cross. Numbers along the curves indicate the time in units of $`10^{21}\mathrm{s}`$ for the dynamic path at $`J`$=86 $`\mathrm{}`$. The insert shows the final stage of the fusion trajectories with $`J`$=70 $`\mathrm{}`$ and $`J`$=75 $`\mathrm{}`$.
Fig. 5. Dynamic paths (for $`J=70`$, 75, 80, 85, 90, 95, 100, 105 $`\mathrm{}`$) in the <sup>60</sup>Ni$`+^{100}`$Mo collision at 1200 MeV are shown as solid lines. The quasistatic path for $`J`$= 86 $`\mathrm{}`$ is represented by the dashed line. The spherical shape is shown as a cross.
Fig. 6. Dynamic paths (for $`J`$= 70, 75, 80, 85, 90, 95, 100, 105 $`\mathrm{}`$) for the collision <sup>48</sup>Ca$`+^{112}`$Sn at 480 MeV are shown as solid lines. The dashed line represents the quasistatic path corresponding to the self-consistent ETF variational calculation with the SkM force. The closed square symbols show the quasistatic path found for the YPE forces in the space of Blocki shapes.
Fig. 7. The effective separation energies of neutrons, protons and alpha particles along the quasistatic path in <sup>160</sup>Yb. The solid lines show the calculation for the dynamic shapes with YPE forces. The dashed lines indicate the ETF effective separation energies.
Fig. 8. The level density parameter along the quasistatic path. The solid line represents the $`a`$ values normalized to $`A/8.8`$ MeV<sup>-1</sup> for the spherical shape with the Tõke-Swiatecki shape-dependent factor based on YPE forces. The dashed line indicates the ETF calculation.
Fig. 9. Equilibrium total multiplicities of neutrons, protons and alpha particles (pre-scission plus post-scission) as a function of the neutron number N of the emitting system (after the pre-equilibrium stage). Experimental data points from Refs. and are shown by filled squares and crosses, respectively. The short-dashed lines represent the calculated pre-scission (equilibrium) multiplicities. The solid lines show the total calculated multiplicities. In the top (neutron) panel, the open circles represent the total measured multiplicities accounting for pre-equilibrium emission . From left to right the systems are 10 MeV/u <sup>63</sup>Cu+<sup>92</sup>Mo, <sup>60</sup>Ni+<sup>100</sup>Mo, <sup>63</sup>Cu+<sup>100</sup>Mo and 20 MeV/u <sup>20</sup>Ne+<sup>144,148,154</sup>Sm.
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# Black holes and closed trapped surfaces: a revision of a classic theorem
## 1. Introduction
Weakly asymptotically simple and empty (WASE) space-times provide a setting for the analysis of isolated gravitational objects in classical general relativity. There is sufficient freedom and generality to allow for the presence of singularities and black holes, but there is also sufficient structure, in the form of a well-behaved asymptotic region, to permit the proof of significant results. Examples of WASE space-times include the Schwarzschild, Kerr and Reissner-Nordström space-times.
Despite the evident physical importance of WASE space-times, it is often the case that scant attention is paid to their precise definition. Indeed, as will be seen, the definition proposed and used by Hawking & Ellis , henceforth H&E, is inadequate for their purposes. A more restrictive definition was subsequently proposed in as a foundation for a class of censorship theorems (see also ). One purpose of the present paper is to provide basic results for this definition further to those given in .
As has been well-known since , the concept of a closed trapped surface is central to the understanding of black holes. In the context of WASE space-times it is a standard assertion, stated as Proposition 9.2.1 in H&E, that subject to a form of weak cosmic censorship and a suitable energy condition, it is not possible for closed trapped surfaces to be seen from future null infinity. In itself this is of no direct physical significance because there need be nothing remarkable about the geometry at a space-time point which lies on a closed trapped surface. The assertion does however imply, again subject to weak cosmic censorship and an energy condition, that the presence of a closed trapped surface in a WASE space-time implies the presence of a black hole. The assertion also leads to an elementary censorship theorem according to which, in a WASE satisfying the same energy condition, if every future singularity is sufficiently strong as to be preceded by a closed trapped surface, then weak cosmic censorship must hold. The second purpose of this paper is to obtain a rigorous proof of a modified statement of H&E Proposition 9.2.1 on the basis of the definition of a WASE space-time given in .
## 2. Preliminary concepts
The following two basic definitions are taken from .
###### Definition 2.1.
An asymptote of a space-time $`(M,𝐠)`$ is a quadruple $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$, where $`(\overline{M}\text{},\overline{𝐠}\text{})`$ is a space-time-with-boundary, $`\mathrm{\Omega }`$ is a $`C^{\mathrm{}}`$ real-valued function on $`\overline{M}\text{}`$, and $`\psi :M\overline{M}\text{}`$ is a $`C^{\mathrm{}}`$ embedding such that
1. $`\psi (M)=\overline{M}\text{}\overline{M}\text{}`$;
2. $`(\psi ^{}\mathrm{\Omega })^2𝐠=\psi ^{}\overline{𝐠}\text{}`$;
3. one has $`\mathrm{\Omega }(p)=0`$ and $`𝐝\mathrm{\Omega }(p)0`$ for all $`p\overline{M}\text{}`$.
###### Definition 2.2.
An asymptote $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi ^{})`$ of a space-time $`(M^{},𝐠^{})`$ is asymptotically simple and empty (ASE) if
1. $`\text{ }\stackrel{~}{\mathrm{M}}\text{}supp(\psi _{}^{}Ricc(𝐠^{}))\text{ }\stackrel{~}{\mathrm{M}}\text{}`$;
2. $`(\text{ }\stackrel{~}{\mathrm{M}}\text{},\stackrel{~}{𝐠}\text{})`$ is strongly causal;
3. every inextendible null geodesic $`\gamma ^{}`$ of $`(M^{},𝐠^{})`$ is such that $`\psi ^{}\gamma ^{}`$ has two endpoints in $`\stackrel{~}{M}`$, both of which lie in $`\text{ }\stackrel{~}{\mathrm{M}}\text{}`$.
A space-time is asymptotically simple and empty (ASE) if it admits an asymptotically simple and empty asymptote.
Standard arguments give that any ASE space-time $`(M,𝐠)`$ is globally hyperbolic, and that for any ASE asymptote $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi ^{})`$ of $`(M^{},𝐠^{})`$, the boundary $`\text{ }\stackrel{~}{M}\text{}`$ of $`\stackrel{~}{M}`$ is the union of two disjoint connected null hypersurfaces $`\text{ }\stackrel{~}{}\text{}^+:=I^+(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{M}\text{}`$ and $`\text{ }\stackrel{~}{}\text{}^{}:=I^{}(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{M}\text{}`$. By means of condition (ii) of Definition 2.2 one can show that $`\text{ }\stackrel{~}{}\text{}^+`$ and $`\text{ }\stackrel{~}{}\text{}^{}`$ are diffeomorphic to $`𝕊^2\times `$.
In the terminology of , a slice of $`\text{ }\stackrel{~}{}\text{}^+`$ is a non-empty locally acausal compact connected topological 2-submanifold of $`\text{ }\stackrel{~}{}\text{}^+`$. Theorem 5.1 of gives that every slice of $`\text{ }\stackrel{~}{}\text{}^+`$ is homeomorphic to $`𝕊^2`$. Since, by assumption here, strong causality holds at every point of $`\text{ }\stackrel{~}{}\text{}^+`$ in $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$, Proposition 7.1 of gives that every null geodesic generator of $`\text{ }\stackrel{~}{}\text{}^+`$ cuts every slice of $`\text{ }\stackrel{~}{}\text{}^+`$, and Theorem 7.4 of gives that $`\text{ }\stackrel{~}{}\text{}^+`$ is acausal in $`(\overline{M}\text{},\overline{𝐠}\text{})`$. Any given slice of $`\text{ }\stackrel{~}{}\text{}^+`$ may be mapped along the generators of $`\text{ }\stackrel{~}{}\text{}^+`$ to yield a foliation of $`\text{ }\stackrel{~}{}\text{}^+`$ by slices of $`\text{ }\stackrel{~}{}\text{}^+`$. Similar assertions apply to $`\text{ }\stackrel{~}{}\text{}^{}`$.
Note that $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ need not be causally simple. For example, let $`(M^{},𝐠^{})`$ be Minkowski space. Let $`\stackrel{~}{\mu }\text{}:\text{ }\stackrel{~}{}\text{}^{}`$ be a future-directed null geodesic generator of $`\text{ }\stackrel{~}{}\text{}^{}`$ and let $`\stackrel{~}{\nu }\text{}:\text{ }\stackrel{~}{}\text{}^+`$ be the antipodal future-directed null geodesic generator of $`\text{ }\stackrel{~}{}\text{}^+`$. Let $`\stackrel{~}{p}\text{}|\stackrel{~}{\mu }\text{}|`$ and let $`\stackrel{~}{q}\text{}:=\dot{J}\text{}^+(\stackrel{~}{p}\text{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})|\stackrel{~}{\nu }\text{}|\text{ }\stackrel{~}{}\text{}^+`$. One then has $`J^+(\stackrel{~}{p}\text{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{}\text{}^+=\text{ }\stackrel{~}{}\text{}^+|\stackrel{~}{\lambda }\text{}_{\stackrel{~}{q}\text{}}|`$ for $`\stackrel{~}{\lambda }\text{}_{\stackrel{~}{q}\text{}}:=\stackrel{~}{\nu }\text{}|(\mathrm{},a)`$, for $`a:=\stackrel{~}{\nu }\text{}^1(\stackrel{~}{q}\text{})`$. Since $`\text{ }\stackrel{~}{}\text{}^+|\stackrel{~}{\lambda }\text{}_{\stackrel{~}{q}\text{}}|`$ is not relatively closed in $`\text{ }\stackrel{~}{}\text{}^+`$, the set $`J^+(\stackrel{~}{p}\text{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ cannot be closed in $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$.
Condition (ii) of Definition 2.2 may be decomposed into two parts, first that $`(M^{},𝐠^{})`$ is strongly causal and second that $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ is strongly causal at every point of $`\text{ }\stackrel{~}{}\text{}^+`$. Since the physical interpretation of the latter is unclear one is led, as in , to the more general concept of a “simple” space-time for which only the chronology condition is imposed on $`(M^{},𝐠^{})`$, with no additional causality conditions imposed on $`\text{ }\stackrel{~}{}\text{}^+`$ or $`\text{ }\stackrel{~}{}\text{}^{}`$ in $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$. Simple space-times are globally hyperbolic with Cauchy surfaces which, subject to the truth of the Poincaré conjecture, are diffeomorphic to $`^3`$. But the topological and causal structure of $`\text{ }\stackrel{~}{}\text{}^+`$ and $`\text{ }\stackrel{~}{}\text{}^{}`$ may exhibit new complications and subtleties. Despite the possible interest of this additional generality, $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ will for present purposes be assumed to be strongly causal as expressed in condition (ii) if Definition 2.2.
In H&E, a weakly asymptotically simple and empty (WASE) space-time is introduced with a definition which, when re-expressed in terms of asymptotes, assumes the following form.
###### Provisional Definition 2.3 (c.f. H&E p.225).
A space-time $`(M,𝐠)`$ is weakly asymptotically simple and empty (WASE) if there exists an asymptote $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ of $`(M,𝐠)`$, a space-time $`(M^{},𝐠^{})`$, an asymptote $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi ^{})`$ of $`(M^{},𝐠^{})`$ and open sets $`𝒰`$ and $`𝒰^{}`$ of $`M`$ and $`M^{}`$ respectively such that
1. $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi ^{})`$ is ASE;
2. $`\psi (𝒰)\overline{M}\text{}`$ is an open neighbourhood of $`\overline{M}\text{}`$ in $`\overline{M}\text{}`$;
3. $`\psi ^{}(𝒰^{})\text{ }\stackrel{~}{M}\text{}`$ is an open neighbourhood of $`\text{ }\stackrel{~}{M}\text{}`$ in $`\stackrel{~}{M}`$;
4. $`(𝒰,𝐠|𝒰)`$ and $`(𝒰^{},𝐠^{}|𝒰^{})`$ are globally isometric.
This definition of a weakly asymptotically simple and empty space-time is not as restrictive as its authors seem to have intended. In particular, it allows the future and past null infinities $`^+:=I^+(\overline{M}\text{},\overline{𝐠}\text{};\overline{M}\text{})\overline{M}\text{}`$ and $`^{}:=I^{}(\overline{M}\text{},\overline{𝐠}\text{};\overline{M}\text{})\overline{M}\text{}`$ of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ to be separated from one another and not join up at a spatial infinity (see Figure 1). (To add an assumption that $`𝒰`$ is connected would not help.) The deficiency undermines several of their results, and in particular their proposed proof (Proposition 9.2.1) that closed trapped surfaces are necessarily confined to black holes.
In order to overcome the problems in Definition 2.3, one may adopt the following, more restrictive definition of a WASE space-time proposed in .
###### Definition 2.4.
An asymptote $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ of a space-time $`(M,𝐠)`$ is weakly asymptotically simple and empty (WASE) if there exists an open set $`𝒰`$ of $`M`$, an extension $`(M^{},𝐠^{})`$ of $`(𝒰,𝐠|𝒰)`$, an asymptote $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi ^{})`$ of $`(M^{},𝐠^{})`$ and a topological embedding $`\xi :\overline{\psi (𝒰)\overline{M}\text{}}\text{ }\stackrel{~}{M}\text{}`$ such that
1. $`(\text{ }\stackrel{~}{\mathrm{M}}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi ^{})`$ is ASE;
2. for every $`p^{}M^{}`$, the set $`M^{}(𝒰I(p^{},𝐠^{};M^{}))`$ is compact;
3. one has $`\xi (\overline{M}\text{})=\text{ }\stackrel{~}{\mathrm{M}}\text{}`$ and $`\xi \psi |𝒰=\psi ^{}|𝒰`$;
4. for all $`q𝒰`$ and all future-pointing timelike vectors $`𝐯T_qM`$ of $`(M,𝐠)`$, the vectors $`\psi _{}𝐯`$, $`𝐯`$ and $`\psi _{}^{}𝐯`$ are future-pointing in $`(\overline{M}\text{},\overline{𝐠}\text{})`$, $`(M^{},𝐠^{})`$ and $`(\text{ }\stackrel{~}{\mathrm{M}}\text{},\stackrel{~}{𝐠}\text{})`$ respectively.
A space-time is weakly asymptotically simple and empty (WASE) if it admits a weakly asymptotically simple and empty asymptote.
###### Remark.
For any Cauchy surface $`𝒮^{}`$ of $`(M^{},𝐠^{})`$ one has $`𝒮^{}𝒰M^{}(𝒰I(p^{},𝐠^{};M^{}))`$ for any $`p^{}𝒮^{}`$ and hence that $`𝒮^{}𝒰`$ is compact.
The future and past null infinities of a WASE asymptote $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ are defined by $`^+:=I^+(\overline{M}\text{},\overline{𝐠}\text{};\overline{M}\text{})\overline{M}\text{}`$ and $`^{}:=I^{}(\overline{M}\text{},\overline{𝐠}\text{};\overline{M}\text{})\overline{M}\text{}`$ respectively. It is clear that $`\overline{M}\text{}`$ is the disjoint union of $`^+`$ and $`^{}`$.
The mappings involved in Definition 2.4 are shown in Figure 2. In Figure 3 one can see how the pathologies inherent in Definition 2.3 are eliminated by condition (ii) of Definition 2.4. This condition may be regarded as a way to require that $`𝒰`$ is a neighbourhood of spatial infinity without reference to the geometrical structure of spatial infinity.
The following two lemmas are basic in the analysis of WASE space-times. Their proofs are given in .
###### Lemma 2.5.
Within the context of Definition 2.4 one has
1. $`\psi ^{}(𝒰)\text{ }\stackrel{~}{\mathrm{M}}\text{}`$ is open in $`\stackrel{~}{M}`$
2. $`\psi (𝒰)\overline{M}\text{}`$ is open in $`\overline{M}\text{}`$. ∎
###### Lemma 2.6.
Within the context of Definition 2.4 one has
1. a subset of $`𝒰`$ is open in $`M`$ iff it is open in $`M^{}`$;
2. a subset of $`𝒰`$ is compact in $`M`$ iff it is compact in $`M^{}`$;
3. a subset of $`𝒰`$ is closed in $`M`$ if (but not only if) it is closed in $`M^{}`$. ∎
Since the mappings $`\psi :M\overline{M}\text{}`$ and $`\psi ^{}:M^{}\text{ }\stackrel{~}{M}\text{}`$ in Definition 2.1 are both diffeomorphisms onto their images, it is clear from Figure 2 that the mapping $`\xi :\overline{\psi (𝒰)\overline{M}\text{}}\text{ }\stackrel{~}{M}\text{}`$ in Definition 2.4 is such that $`\xi |\psi (𝒰)`$ is a diffeomorphism onto $`\psi ^{}(𝒰)`$. Moreover $`\xi |\psi (𝒰)`$ is a conformal isometry onto $`\psi ^{}(𝒰)`$ in the sense of $`(\xi |\psi (𝒰))^{}\stackrel{~}{𝐠}\text{}=(\mathrm{\Omega }/\mathrm{\Omega }^{})^2\overline{𝐠}\text{}|\psi (𝒰)`$. Note however that $`\xi `$ need not be differentiable at points of $`\overline{M}\text{}`$.
###### Proposition 2.7.
One has
1. $`\xi (^+)=\text{ }\stackrel{~}{}\text{}^+`$ and $`\xi (^{})=\text{ }\stackrel{~}{}\text{}^{}`$;
2. $`\xi `$ maps the null geodesic generators of $`^+`$ and $`^{}`$ onto null geodesic generators of $`\text{ }\stackrel{~}{}\text{}^+`$ and $`\text{ }\stackrel{~}{}\text{}^{}`$ respectively.
Proof. Let $`p^+`$. Lemma 2.5 gives that $`\psi (𝒰)\overline{M}\text{}`$ is an open neighbourhood of $`p`$ in $`\overline{M}\text{}`$. Hence, by the definition of $`^+`$, there exists a smooth timelike curve $`\alpha :(0,1)\psi (𝒰)`$ of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ having a future endpoint at $`p`$ in $`(\overline{M}\text{},\overline{𝐠}\text{})`$. One has $`\xi (p)\text{ }\stackrel{~}{M}\text{}`$ by condition (iii) of Definition 2.4. Since $`\xi \alpha `$ is a smooth timelike curve of $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ with a future endpoint at $`\xi (p)`$ one therefore has $`\xi (p)\text{ }\stackrel{~}{}\text{}^+`$ by the definition of $`\text{ }\stackrel{~}{}\text{}^+`$. There follows $`\xi (^+)\text{ }\stackrel{~}{}\text{}^+`$. A similar argument gives $`\xi ^1(\text{ }\stackrel{~}{}\text{}^+)^+`$ which implies $`\text{ }\stackrel{~}{}\text{}^+\xi (^+)`$. Hence one has $`\xi (^+)=\text{ }\stackrel{~}{}\text{}^+`$ and similarly $`\xi (^{})=\text{ }\stackrel{~}{}\text{}^{}`$. This establishes (i).
Let $`\gamma :I^+`$ be a future-directed null geodesic generator of $`^+`$ and let $`\stackrel{~}{\gamma }\text{}:=\xi \gamma `$. Let $`[a,b]I`$ for $`a<b`$ and let $`\mu :=\gamma |[a,b)`$. One has $`|\mu |^+`$. By Lemma 2.5 and the definition of $`^+`$ there exists a smooth future-directed timelike curve $`\mu _0:[a,b)\psi (𝒰)`$ of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ with a future endpoint at $`\gamma (b)`$. Indeed there exists a sequence of smooth future-directed timelike curves $`\mu _i:[a,b)\psi (𝒰)`$ of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ converging pointwise to $`\mu `$ in $`\overline{M}\text{}`$, each with a future endpoint at $`\gamma (b)`$. Since the smooth timelike curves $`\stackrel{~}{\mu }\text{}_i:=\xi \mu _i:[a,b)\psi ^{}(𝒰)`$ of $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ converge pointwise to $`\stackrel{~}{\mu }\text{}:=\xi \mu `$ in $`\stackrel{~}{M}`$ one has that $`\stackrel{~}{\mu }\text{}`$ is a causal curve of $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ with a future endpoint at $`\stackrel{~}{\gamma }\text{}(b)`$. Since part (i) gives $`|\stackrel{~}{\mu }\text{}|=\xi (|\mu |)\text{ }\stackrel{~}{}\text{}^+`$ and since $`\text{ }\stackrel{~}{}\text{}^+`$ is a null hypersurface of $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ it follows that $`\stackrel{~}{\mu }\text{}=\xi \mu `$ is a null geodesic generating segment of $`\text{ }\stackrel{~}{}\text{}^+`$. Since $`a`$ and $`b`$ were arbitrary in $`\mu :=\gamma |[a,b)`$ it follows that $`\stackrel{~}{\gamma }\text{}=\xi \gamma `$ is a null geodesic generating segment of $`\text{ }\stackrel{~}{}\text{}^+`$. Clearly $`\stackrel{~}{\gamma }\text{}`$ cannot have a future endpoint $`q\text{ }\stackrel{~}{}\text{}^+`$ otherwise $`\gamma `$ would have a future endpoint at $`\xi ^1(q)^+`$ and so would not be a generator of $`^+`$, contrary to hypothesis. Similarly $`\stackrel{~}{\gamma }\text{}`$ cannot have a past endpoint in $`\text{ }\stackrel{~}{}\text{}^+`$. Hence $`\stackrel{~}{\gamma }\text{}`$ is a null geodesic generator of $`\text{ }\stackrel{~}{}\text{}^+`$. The corresponding result for $`^{}`$ is similar. This establishes (ii).
Proposition 2.7 shows that the structure of $`^+`$ and $`^{}`$ for a WASE asymptote $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is directly analogous to the structure of $`\text{ }\stackrel{~}{}\text{}^+`$ and $`\text{ }\stackrel{~}{}\text{}^{}`$ for an ASE asymptote $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi )`$. In particular one may, following , define a slice of $`^+`$ (respectively a slice of $`^{}`$) as a non-empty locally acausal compact connected topological 2-submanifold of $`^+`$ (respectively $`^{}`$). Then $`^+`$ and $`^{}`$ are acausal in $`(\overline{M}\text{},\overline{𝐠}\text{})`$, every null geodesic generator of $`^+`$ cuts every slice of $`^+`$ and every null geodesic generator of $`^{}`$ cuts every slice of $`^{}`$. Slices of $`^+`$ and $`^{}`$ are mapped by $`\xi `$ to slices of $`\text{ }\stackrel{~}{}\text{}^+`$ and $`\text{ }\stackrel{~}{}\text{}^{}`$ respectively. Slices of $`\text{ }\stackrel{~}{}\text{}^+`$ and $`\text{ }\stackrel{~}{}\text{}^{}`$ are mapped by $`\xi ^1`$ to slices of $`^+`$ and $`^{}`$ respectively.
The following is a useful restriction on the causal structure of a WASE space-time. It is equivalent to a definition of asymptotic simplicity in but is re-expressed here in a form more convenient for present purposes. The change of terminology seems appropriate because the term “simple” has become overworked.
###### Definition 2.8.
A WASE asymptote $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is asymptotically chronologically consistent if $`𝒰`$, $`(M^{},𝐠^{})`$ and $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi ^{})`$ in Definition 2.4 may be chosen such that for any achronal set $`\stackrel{~}{𝒜}`$ of $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ such that $`\text{ }\stackrel{~}{𝒜}\text{}\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ one has that $`\xi ^1(\text{ }\stackrel{~}{𝒜}\text{})\psi (𝒰)\overline{M}\text{}`$ is achronal in $`(\overline{M}\text{},\overline{𝐠}\text{})`$.
## 3. The main result
A form of weak cosmic censorship hypothesis will be required. The H&E concept of future asymptotic predictability is suitable for this purpose. The following definition formulates future asymptotic predictability in terms of asymptotes and provides a weaker concept of partial future asymptotic predictability that is also well-established in the literature.
###### Definition 3.1.
Let $`𝒮`$ be a closed achronal set without edge in a WASE space-time $`(M,𝐠)`$ and let $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ be a WASE asymptote of $`(M,𝐠)`$. Then $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is future asymptotically predictable from $`\psi (𝒮)`$ if one has $`^+\overline{D}\text{}^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$. One says that $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is partially future asymptotically predictable from $`\psi (𝒮)`$ if there exists a slice $`\mathrm{\Sigma }^+`$ of $`^+`$ such that $`J^{}(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})^+\overline{D}\text{}^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$.
The main result is the following:
###### Theorem 3.2.
Let $`(\overline{M}\text{},\overline{𝐠}\text{})`$ be a WASE space-time and let $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ be a WASE asymptote of $`(M,𝐠)`$. Suppose
1. there exists a closed, edgeless achronal set $`𝒮`$ in $`(M,𝐠)`$ such that $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is future asymptotically predictable from $`\psi (𝒮)`$;
2. $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is asymptotically chronologically consistent;
3. one has $`R_{ab}k^ak^b0`$ for all null vectors $`k^a`$,
then for any closed trapped surface $`𝒯`$ of $`(M,𝐠)`$ in $`\overline{I}\text{}^+(𝒮,𝐠;M)`$ one has $`\psi (𝒯)J^{}(^+,\overline{𝐠}\text{};\overline{M}\text{})=\mathrm{}`$.
Conditions (1) and (3) of Theorem 3.2 coincide with conditions in the statement of H&E Proposition 9.2.1. However condition (2), which would seem to be necessary, makes no appearance in H&E. Note also that Theorem 3.2 requires only $`𝒯\overline{I}\text{}^+(𝒮,𝐠;M)`$ whereas H&E impose the stronger condition $`𝒯D^+(𝒮,𝐠;M)`$.
The basic idea of the H&E argument in support of the statement of their Proposition 9.2.1 is to show that if $`𝒯`$ is visible from $`^+`$ then there must be a null geodesic generator of $`\dot{I}\text{}^+(𝒯,𝐠;M)`$ which reaches from $`𝒯`$ to $`^+`$ and which is therefore of infinite affine length. A contradiction then follows by means of the Raychaudhuri equation and the null convergence condition. The argument fails though because the H&E definition of a WASE space-time is not sufficiently strong. Specifically, things begin to go wrong when they claim that, in the associated ASE space-time $`(M^{},𝐠^{})`$, for a Cauchy surface $`𝒮^{}`$ chosen such that $`𝒮^{}𝒰^{}=𝒮𝒰`$ it is necessarily the case that $`𝒮^{}𝒰^{}`$ is compact. (At this point H&E are tacitly identifying $`𝒰`$ and $`𝒰^{}`$, as is explicitly done in the present formalism.) In the first place it is unclear that there need be any Cauchy surface $`𝒮^{}`$ of $`(M^{},𝐠^{})`$ such that $`𝒮^{}𝒰^{}=𝒮𝒰`$. For example $`𝒮𝒰`$ might not be achronal in $`(M^{},𝐠^{})`$. And second, since $`𝒮`$ need not even intersect $`𝒰`$ and every Cauchy surface of $`(M^{},𝐠^{})`$ is non-compact, the set $`𝒮^{}𝒰^{}`$ could be non-compact. Definition 2.4 directly overcomes the second difficulty, as was indicated in the Remark that followed Definition 2.4. In order to overcome the first it will be necessary to employ different techniques.
At a pictorial level one might seek a counterexample to H&E Proposition 9.2.1 by arranging to have $`^+`$ both separated from $`^{}`$ and contained in the chronological future of a closed trapped surface $`𝒯`$ (see Figure 4). Since none of the null geodesic generators of $`\dot{J}\text{}^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})`$ then meet $`^+`$, the central contradiction in the proof of H&E Lemma 9.2.1 is avoided. However to obtain a full counterexample to H&E Proposition 9.2.1 one also needs to arrange for the null convergence condition to be satisfied. Even though it is not clear how this might be done, it seems unlikely that H&E Proposition 9.2.1 is correct.
The following three lemmas and a proposition are the key to the proof of the revised trapped surfaces theorem.
###### Lemma 3.3.
Let $`(M,𝐠)`$ be a WASE space-time with a asymptotically chronologically consistent WASE asymptote $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ and let $`\mathrm{\Sigma }^{}`$ be a slice of $`^{}`$. Then there exists a slice $`\mathrm{\Sigma }^+`$ of $`^+`$ such that $`J^{}(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})`$ is a closed set of $`\overline{M}\text{}`$ contained in $`\psi (𝒰)\overline{M}\text{}`$ and $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ is a closed set of $`\stackrel{~}{M}`$ contained in $`\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ for $`\stackrel{~}{\mathrm{\Sigma }}\text{}^{}:=\xi (\mathrm{\Sigma }^{})`$ and $`\stackrel{~}{\mathrm{\Sigma }}\text{}^+:=\xi (\mathrm{\Sigma }^+)`$.
Proof. The time reverse of Lemma 3.6 of gives that $`\overline{M}\text{}\overline{M}\text{}`$ cannot be contained entirely in $`I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$. Let $`\stackrel{~}{p}\text{}_+\text{ }\stackrel{~}{M}\text{}(I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{M}\text{})`$ and let $`\stackrel{~}{p}\text{}I^{}(\stackrel{~}{p}\text{}_+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{M}\text{}`$. One then has $`\stackrel{~}{p}\text{}\text{ }\stackrel{~}{M}\text{}(J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{M}\text{})`$. Let $`p^{}:=(\psi ^{})^1(\stackrel{~}{p}\text{})`$. The set $`𝒬^{}:=\{p^{}\}(M^{}(𝒰I(p,𝐠^{};M^{})))`$ is compact in $`M^{}`$ by Definition 2.4 and $`\text{ }\stackrel{~}{𝒬}\text{}:=\psi ^{}(𝒬^{})`$ is compact in $`\text{ }\stackrel{~}{M}\text{}\text{ }\stackrel{~}{M}\text{}`$. So by Lemma 4.5 (III) of the set $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+:=\dot{J}\text{}^+(\text{ }\stackrel{~}{𝒬}\text{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{}\text{}^+`$ is a slice of $`\text{ }\stackrel{~}{}\text{}^+`$. The set $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})(\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{})`$ is contained in $`I^{}(\stackrel{~}{p}\text{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ which does not intersect $`J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$. Hence $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ is contained in $`\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$.
The set $`\mathrm{\Sigma }_0^+:=\xi ^1(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+)`$ is a slice of $`^+`$. Suppose there exists a point $`x(I^{}(\mathrm{\Sigma }_0^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{}))(\psi (𝒰)\overline{M}\text{})`$. Then there exists a causal curve $`\alpha `$ of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ from $`\mathrm{\Sigma }^{}`$ to $`x`$, and a timelike curve $`\beta :[0,1]\overline{M}\text{}`$ of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ from $`x`$ to $`\mathrm{\Sigma }_0^+`$. Let $`a:=sup\{t[0,1]:\beta (t)\psi (𝒰)\overline{M}\text{}\}`$. Then $`\beta (a)`$ lies in the topological boundary of $`\psi (𝒰)\overline{M}\text{}`$ in $`\overline{M}\text{}`$ and one has $`\beta (t)(\psi (𝒰)\overline{M}\text{})I^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})J^{}(\mathrm{\Sigma }_0^+,\overline{𝐠}\text{};\overline{M}\text{})`$ for all $`t(a,1]`$. For each $`t(a,1]`$ the set $`\mathrm{\Sigma }^{}\{\beta (t)\}`$ is non-achronal in $`(\overline{M}\text{},\overline{𝐠}\text{})`$ and contained in $`\psi (𝒰)\overline{M}\text{}`$. So, by the asymptotic chronological consistency of $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$, one has that $`\xi (\mathrm{\Sigma }^{}\{\beta (t)\})=\stackrel{~}{\mathrm{\Sigma }}\text{}^{}\{\xi \beta (t)\}`$ is non-achronal in $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ for each $`t(a,1]`$. Since $`\xi \beta :[a,1]\text{ }\stackrel{~}{M}\text{}`$ is a timelike curve of $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ from $`\xi \beta (a)`$ to $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+`$ one thus has $`\xi \beta (t)I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ for all $`t(a,1]`$. Lemma 4.12 (I) of gives that $`J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ is closed in $`\stackrel{~}{M}`$. Hence one has $`\xi \beta (t)J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^{}(\mathrm{\Sigma }_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ for all $`t[a,1]`$. But since $`\beta (a)`$ lies in the topological boundary of $`\psi (𝒰)\overline{M}\text{}`$ in $`\overline{M}\text{}`$ the point $`\xi \beta (a)`$ must lie in the topological boundary of $`\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ in $`\stackrel{~}{M}`$. This is impossible because $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ is contained in $`\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ which is open in $`\stackrel{~}{M}`$. One thus has $`I^{}(\mathrm{\Sigma }_0^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})\psi (𝒰)\overline{M}\text{}`$.
Now let $`\stackrel{~}{\mathrm{\Sigma }}\text{}^+`$ be a slice of $`\text{ }\stackrel{~}{}\text{}^+`$ lying strictly to the past of $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+`$ along the generators of $`\text{ }\stackrel{~}{}\text{}^+`$. One has $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$. Hence a point of $`\stackrel{~}{M}`$ lies in $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ iff it lies on a causal curve of $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ in $`\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ from $`\stackrel{~}{\mathrm{\Sigma }}\text{}^{}`$ to $`\stackrel{~}{\mathrm{\Sigma }}\text{}^+`$. Moreover, since Lemma 4.12 (I) of gives that $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ and $`J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ are both closed in $`\stackrel{~}{M}`$, the set $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ is closed in $`\stackrel{~}{M}`$.
Let $`\mathrm{\Sigma }^+:=\xi ^1(\stackrel{~}{\mathrm{\Sigma }}\text{}^+)`$. One has $`J^{}(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})(I^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{}))\overline{M}\text{}\psi (𝒰)\overline{M}\text{}`$. Hence a point of $`\overline{M}\text{}`$ lies in $`J^{}(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})`$ iff it lies on a causal curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ in $`\psi (𝒰)\overline{M}\text{}`$ from $`\mathrm{\Sigma }^{}`$ to $`\mathrm{\Sigma }^+`$. One thus has $`J^{}(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})=\xi ^1(J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\mathrm{\Sigma }^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{}))`$. Since $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ has been shown to be closed in $`\stackrel{~}{M}`$ it follows that $`J^{}(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})`$ is closed in $`\overline{M}\text{}`$.
###### Corollary.
Let $`𝒦`$ be a compact set of $`M`$. Then there exists a slice $`\mathrm{\Sigma }_1^+`$ of $`^+`$ such that $`J^{}(\mathrm{\Sigma }_1^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})`$ does not intersect $`\psi (𝒦)`$.
Proof. One may assume $`𝒰M𝒦`$ otherwise one may redefine $`𝒰`$ as $`𝒰𝒦`$. The Lemma then gives that there exists a slice $`\mathrm{\Sigma }_1^+`$ of $`^+`$ such that $`J^{}(\mathrm{\Sigma }_1^+,\overline{𝐠}\text{};\overline{M}\text{})J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})\psi (𝒰)\overline{M}\text{}\overline{M}\text{}\psi (𝒦)`$.
###### Lemma 3.4.
Let $`\mathrm{\Sigma }^{}`$ be a slice of $`^{}`$ and let $`𝒩`$ be an open neighbourhood of $`\mathrm{\Sigma }^{}`$ in $`\overline{M}\text{}`$. Then there exists a slice $`\mathrm{\Sigma }_2^+`$ of $`^+`$ such that, for every $`q\mathrm{\Sigma }_2^+`$, the set $`𝒩`$ is cut by every past endless timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ to $`q`$.
Proof. It suffices to assume $`𝒩\psi (𝒰)\overline{M}\text{}`$ since one may otherwise redefine $`𝒩`$ as $`𝒩(\psi (𝒰)\overline{M}\text{})`$.
The set $`\text{ }\stackrel{~}{𝒩}\text{}:=\xi (𝒩)`$ is open in $`\stackrel{~}{M}`$ and $`\stackrel{~}{\mathrm{\Sigma }}\text{}^{}:=\xi (\mathrm{\Sigma }^{})`$ is a slice of $`\text{ }\stackrel{~}{}\text{}^{}`$. By the time reverse of Lemma 4.12 of one has that $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ is compact in $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ and such that $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{}\text{}^{}=\stackrel{~}{\mathrm{\Sigma }}\text{}^{}`$. The time reverse of Proposition 7.2 of gives $`\text{ }\stackrel{~}{}\text{}^+I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ whereby one has $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{}\text{}^+=\mathrm{}`$. Hence one has $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{M}\text{}=\stackrel{~}{\mathrm{\Sigma }}\text{}^{}`$. Let $`\stackrel{~}{𝒳}`$ be the compact set $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{𝒩}\text{}\text{ }\stackrel{~}{M}\text{}\text{ }\stackrel{~}{M}\text{}`$ or, if this is empty, let $`\stackrel{~}{𝒳}`$ be any non-empty compact set of $`\text{ }\stackrel{~}{M}\text{}\text{ }\stackrel{~}{M}\text{}`$. The set $`\stackrel{~}{\mathrm{\Sigma }}\text{}_{\text{ }\stackrel{~}{𝒳}\text{}}^+:=\dot{J}\text{}^+(\text{ }\stackrel{~}{𝒳}\text{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{}\text{}^+`$ is a slice of $`\text{ }\stackrel{~}{}\text{}^+`$. Since $`I^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_{\text{ }\stackrel{~}{𝒳}\text{}}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ does not intersect $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{𝒩}\text{}`$ one has $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})I^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_{\text{ }\stackrel{~}{𝒳}\text{}}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{𝒩}\text{}`$.
Let $`\mathrm{\Sigma }^+`$ be as in the statement of Lemma 3.3 and let $`\stackrel{~}{\mathrm{\Sigma }}\text{}^+:=\xi (\mathrm{\Sigma }^+)`$. Let $`\stackrel{~}{\mathrm{\Sigma }}\text{}_2^+`$ be a slice of $`\text{ }\stackrel{~}{}\text{}^+`$ lying strictly to the past of both $`\stackrel{~}{\mathrm{\Sigma }}\text{}_{\text{ }\stackrel{~}{𝒳}\text{}}^+`$ and $`\stackrel{~}{\mathrm{\Sigma }}\text{}^+`$ along the generators of $`\text{ }\stackrel{~}{}\text{}^+`$. One then has $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_2^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\psi ^{}(𝒰)\overline{M}\text{}`$ and $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_2^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})(I^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_{\text{ }\stackrel{~}{𝒳}\text{}}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\text{ }\stackrel{~}{}\text{}^+)\text{ }\stackrel{~}{𝒩}\text{}`$.
The set $`\mathrm{\Sigma }_2^+:=\xi ^1(\stackrel{~}{\mathrm{\Sigma }}\text{}_2^+)`$ is a slice of $`^+`$. Let $`q\mathrm{\Sigma }_2^+`$ and let $`\sigma :(\mathrm{},0]\overline{M}\text{}\overline{M}\text{}`$ be a future-directed, past endless timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ having a future endpoint at $`\sigma (0)=q\mathrm{\Sigma }_2^+`$. Let $`\nu :(b,0]\overline{M}\text{}`$ be the maximal segment of $`\sigma `$ to $`q\mathrm{\Sigma }_2^+\psi (𝒰)\overline{M}\text{}`$ in $`\psi (𝒰)\overline{M}\text{}`$. Then $`\stackrel{~}{\nu }\text{}:=\xi \nu `$ is a timelike curve of $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ to $`\stackrel{~}{q}\text{}:=\xi (q)\stackrel{~}{\mathrm{\Sigma }}\text{}_2^+`$. One clearly has $`|\stackrel{~}{\nu }\text{}|J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_2^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$. In order to show that $`\stackrel{~}{\nu }\text{}`$ cuts $`\stackrel{~}{𝒩}`$ it therefore suffices to show that $`\stackrel{~}{\nu }\text{}`$ cuts $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$.
Suppose first that $`\stackrel{~}{\nu }\text{}`$ is past endless in $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$. The time reverse of Lemma 3.6 of gives that $`I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ cannot contain all of $`\overline{M}\text{}\text{ }\stackrel{~}{M}\text{}`$, whilst the time reverse of Lemma 4.2 of gives $`\overline{M}\text{}\overline{M}\text{}I^+(|\nu |,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$. Hence $`\stackrel{~}{\nu }\text{}`$ cuts $`\text{ }\stackrel{~}{M}\text{}I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$. Since $`\stackrel{~}{\nu }\text{}`$ is past endless and timelike in $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ it follows that $`\stackrel{~}{\nu }\text{}`$ cuts $`\text{ }\stackrel{~}{M}\text{}J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$. Because $`\stackrel{~}{\nu }\text{}`$ has a future endpoint at $`q\stackrel{~}{\mathrm{\Sigma }}\text{}^+\text{ }\stackrel{~}{}\text{}^+I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ one thus has that $`\stackrel{~}{\nu }\text{}`$ cuts both $`I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ and $`\text{ }\stackrel{~}{M}\text{}J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ and so cuts $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$.
Now suppose that $`\stackrel{~}{\nu }\text{}`$ has a past endpoint $`\stackrel{~}{z}\text{}`$ in $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$. The point $`\stackrel{~}{z}\text{}`$ must lie in the topological boundary of the open set $`\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ in $`\stackrel{~}{M}`$ otherwise it would lie in $`\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$, in which case $`\xi ^1(\stackrel{~}{z}\text{})\psi (𝒰)\overline{M}\text{}`$ would be a past endpoint to $`\nu `$ in $`(M,𝐠)`$ and $`\nu `$ would be past extendible in $`\psi (𝒰)\overline{M}\text{}`$. Because $`\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ is open in $`\stackrel{~}{M}`$ the set $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ cannot contain $`\stackrel{~}{z}\text{}`$. Since $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ is closed in $`\stackrel{~}{M}`$ the set $`\text{ }\stackrel{~}{M}\text{}(J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{}))`$ is an open neighbourhood of $`\stackrel{~}{z}\text{}`$ in $`\stackrel{~}{M}`$ and so is cut by $`\stackrel{~}{\nu }\text{}`$. In view of $`|\stackrel{~}{\nu }\text{}|J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_2^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ one thus has that $`\stackrel{~}{\nu }\text{}`$ cuts $`\text{ }\stackrel{~}{M}\text{}J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$. Since $`\stackrel{~}{\nu }\text{}`$ has a future endpoint at $`\stackrel{~}{q}\text{}\stackrel{~}{\mathrm{\Sigma }}\text{}_2^+\text{ }\stackrel{~}{}\text{}^+I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ it follows that $`\stackrel{~}{\nu }\text{}`$ cuts both $`I^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ and $`\text{ }\stackrel{~}{M}\text{}J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ and so cuts $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$.
Since $`\stackrel{~}{\nu }\text{}:=\xi \nu `$ cuts $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ and therefore cuts $`\text{ }\stackrel{~}{𝒩}\text{}:=\xi (𝒩)`$ one has that $`\nu `$ cuts $`𝒩`$. Hence $`\sigma `$ cuts $`𝒩`$.
The final lemma will require the use of the following result.
###### Proposition 3.5.
Let $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ be a WASE asymptote of a WASE space-time $`(M,𝐠)`$ and let $`𝒮`$ be a closed edgeless achronal set of $`(M,𝐠)`$. If $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is partially future asymptotically predictable from $`\psi (𝒮)`$ then $`\psi (𝒮)`$ is closed edgeless and achronal in $`(\overline{M}\text{},\overline{𝐠}\text{})`$.
Proof. Let $`(M^{},𝐠^{})`$ be the associated ASE space-time and $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi ^{})`$ an ASE asymptote of $`(M^{},𝐠^{})`$.
If $`\overline{\lambda }\text{}`$ was a timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ from $`\psi (𝒮)\overline{M}\text{}\overline{M}\text{}`$ to $`\psi (𝒮)\overline{M}\text{}\overline{M}\text{}`$ then, because $`\overline{M}\text{}`$ is a null hypersurface of $`(\overline{M}\text{},\overline{𝐠}\text{})`$, one would have $`|\overline{\lambda }\text{}|\overline{M}\text{}\overline{M}\text{}`$ so there would exist a timelike curve $`\lambda `$ of $`(M,𝐠)`$ such that $`\overline{\lambda }\text{}=\psi \lambda `$. But then $`\lambda `$ would be a timelike curve of $`(M,𝐠)`$ from $`𝒮`$ to $`𝒮`$. This would contradict the achronality of $`𝒮`$ in $`(M,𝐠)`$. Hence $`\psi (𝒮)`$ is achronal in $`(\overline{M}\text{},\overline{𝐠}\text{})`$.
Suppose there exists $`s\overline{\psi (𝒮)}^+`$. Since $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is partially future asymptotically predictable from $`\psi (𝒮)`$ there exists a slice $`\mathrm{\Sigma }^+`$ of $`^+`$ such that $`J^{}(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})^+\overline{D}\text{}^+(\psi (𝒰),\overline{𝐠}\text{};\overline{M}\text{})`$. In the case $`sJ^{}(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})^+`$, every past endless timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ to $`s`$ would cut $`\psi (𝒮)`$, and so $`I^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$ would be an open neighbourhood of $`s\overline{\psi (𝒮)}`$ and so would intersect $`\psi (𝒮)`$. This is impossible since $`\psi (𝒮)`$ is achronal in $`(\overline{M}\text{},\overline{𝐠}\text{})`$. So suppose $`sJ^+(\psi (𝒰),\overline{𝐠}\text{};\overline{M}\text{})\mathrm{\Sigma }^+`$. There exists $`r\mathrm{\Sigma }^+`$ lying strictly to the past of $`s`$ on the null geodesic generator of $`^+`$ through $`s`$. In view of $`r\overline{D}\text{}^+(\psi (𝒰),\overline{𝐠}\text{};\overline{M}\text{})`$, every past endless timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ to $`r`$ must cut $`\psi (𝒮)`$. But then $`I^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$ is an open neighbourhood in $`\overline{M}\text{}`$ of $`rJ^{}(s,\overline{𝐠}\text{};\overline{M}\text{})\{s\}`$ and therefore of $`s^+\overline{\psi (𝒮)}`$ and so must intersect $`\psi (𝒮)`$. So again one has a contradiction to the achronality of $`\psi (𝒮)`$. One thus has $`\overline{\psi (𝒮)}^+=\mathrm{}`$.
Suppose there exists $`p\overline{\psi (𝒮)}^{}`$. Let $`\mathrm{\Sigma }^{}`$ be a slice of $`^{}`$ such that $`p\mathrm{\Sigma }^{}`$. By Lemma 3.3 there exists a slice $`\mathrm{\Sigma }^+`$ of $`^+`$ such that $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ for $`\stackrel{~}{\mathrm{\Sigma }}\text{}^{}:=\xi (\mathrm{\Sigma }^{})`$ and $`\stackrel{~}{\mathrm{\Sigma }}\text{}^+:=\xi (\mathrm{\Sigma }^+)`$. Since $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is partially future asymptotically predictable from $`\psi (𝒮)`$ one may assume that $`\mathrm{\Sigma }^+`$ is taken sufficiently far to the past in $`^+`$ to give $`\mathrm{\Sigma }^+\overline{D}\text{}^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$. Let $`\stackrel{~}{p}\text{}:=\xi (p)\stackrel{~}{\mathrm{\Sigma }}\text{}^{}`$. By Lemma 7.2 of one has $`\stackrel{~}{p}\text{}\text{ }\stackrel{~}{}\text{}^{}I^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$. Hence there is a timelike curve $`\stackrel{~}{\alpha }\text{}`$ of $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{})`$ from $`\stackrel{~}{p}\text{}\stackrel{~}{\mathrm{\Sigma }}\text{}^{}`$ to some point $`\stackrel{~}{q}\text{}\stackrel{~}{\mathrm{\Sigma }}\text{}^+`$. In view of $`|\stackrel{~}{\alpha }\text{}|J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ one has that $`\alpha :=\xi ^1\stackrel{~}{\alpha }\text{}`$ is a timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ from $`p\overline{\psi (𝒮)}^{}`$ to $`q:=\xi ^1(\stackrel{~}{q}\text{})\mathrm{\Sigma }^+`$. The set $`I^+(p,\overline{𝐠}\text{};\overline{M}\text{})`$ cannot intersect $`\psi (𝒮)`$ otherwise $`I^{}(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$ would be an open neighbourhood of $`p\overline{\psi (𝒮)}`$ and so $`\psi (𝒮)`$ would not be achronal in $`(\overline{M}\text{},\overline{𝐠}\text{})`$. Since the past endless null geodesic generating segment of $`^{}`$ to $`p`$ clearly does not cut $`\psi (𝒮)`$ it follows that $`I^+(p,\overline{𝐠}\text{};\overline{M}\text{})`$ does not intersect $`D^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$. But $`I^+(p,\overline{𝐠}\text{};\overline{M}\text{})`$ is a neighbourhood of $`q\mathrm{\Sigma }^+\overline{D}\text{}^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$ in $`\overline{M}\text{}`$ so one has a contradiction.
One now has $`\overline{\psi (𝒮)}\overline{M}\text{}=\mathrm{}`$. Since $`\psi :M\overline{M}\text{}`$ is a diffeomorphism onto its image it follows that $`\psi (𝒮)`$ is relatively closed in $`\psi (M)=\overline{M}\text{}\overline{M}\text{}`$. Hence $`\psi (𝒮)`$ is closed in $`\overline{M}\text{}`$.
One has $`edge(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})\overline{M}\text{}=\mathrm{}`$ because $`\psi (𝒮)`$ is closed in $`\overline{M}\text{}`$ and does not intersect $`\overline{M}\text{}`$. And, because $`\psi :(M,𝐠)(\overline{M}\text{},\overline{𝐠}\text{})`$ is a conformal isometry onto its image, one has $`edge(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})\psi (M)=\psi (edge(𝒮,𝐠;M))=\mathrm{}`$. Hence $`edge(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$ is empty.
###### Lemma 3.6.
Let $`(\overline{M}\text{},\overline{𝐠}\text{})`$ be a WASE space-time and let $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ be a WASE asymptote of $`(M,𝐠)`$. Suppose
1. there exists a closed edgeless achronal set $`𝒮`$ in $`(M,𝐠)`$ such that $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is partially future asymptotically predictable from $`\psi (𝒮)`$;
2. $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is asymptotically chronologically consistent.
Then for any compact set $`𝒦\overline{I}\text{}^+(𝒮,𝐠;M)`$ of $`M`$ there exists a slice $`\mathrm{\Sigma }_3^+`$ of $`^+`$ such that $`\overline{I}\text{}^+(\psi (𝒦),\overline{𝐠}\text{};\overline{M}\text{})^+J^+(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})`$.
Proof. It suffices to assume $`𝒦𝒰=\mathrm{}`$ since one may otherwise redefine $`𝒰`$ as $`𝒰𝒦`$.
Let $`\mathrm{\Sigma }_0^{}`$ be a slice of $`^{}`$ and let $`\mathrm{\Sigma }^{}`$ be a slice of $`^{}`$ lying strictly to the past of $`\mathrm{\Sigma }_0^{}`$ along the null geodesic generators of $`^{}`$. Let $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^{}:=\xi (\mathrm{\Sigma }_0^{})`$ and $`\stackrel{~}{\mathrm{\Sigma }}\text{}^{}:=\xi (\mathrm{\Sigma }^{})`$. By Lemma 3.3 there exists a slice $`\mathrm{\Sigma }_0^+`$ of $`^+`$ such that $`J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})J^{}(\mathrm{\Sigma }_0^+,\overline{𝐠}\text{};\overline{M}\text{})\psi (𝒰)\overline{M}\text{}`$ and $`J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\overline{M}\text{})\psi ^{}(𝒰)\text{ }\stackrel{~}{M}\text{}`$ for $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+:=\xi (\mathrm{\Sigma }_0^+)`$. Lemma 4.12 of gives $`\dot{J}\text{}^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})^{}=\stackrel{~}{\mathrm{\Sigma }}\text{}^{}`$ which, since $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^{}`$ lies strictly to the future of $`\stackrel{~}{\mathrm{\Sigma }}\text{}^{}`$ along the null geodesic generators of $`\text{ }\stackrel{~}{}\text{}^{}`$, implies that $`J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ is a neighbourhood of $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^{}`$ in $`\stackrel{~}{M}`$. Proposition 7.2 of gives $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^{}\text{ }\stackrel{~}{}\text{}^{}I^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ whereby one has that $`J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})I^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ is a neighbourhood of $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^{}`$ in $`\stackrel{~}{M}`$. Thus there exists an open neighbourhood $`\text{ }\stackrel{~}{𝒩}\text{}J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})`$ of $`\stackrel{~}{\mathrm{\Sigma }}\text{}_0^{}`$ in $`\stackrel{~}{M}`$. The set $`𝒩:=\xi ^1(\text{ }\stackrel{~}{𝒩}\text{})\xi ^1(J^+(\stackrel{~}{\mathrm{\Sigma }}\text{}^{},\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{})J^{}(\stackrel{~}{\mathrm{\Sigma }}\text{}_0^+,\stackrel{~}{𝐠}\text{};\text{ }\stackrel{~}{M}\text{}))=J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})J^{}(\mathrm{\Sigma }_0^+,\overline{𝐠}\text{};\overline{M}\text{})`$ is an open neighbourhood of $`\mathrm{\Sigma }_0^{}`$ in $`\overline{M}\text{}`$. In view of Proposition 3.5 one may, by passing to a subset of $`𝒩`$ if necessary, assume $`𝒩(\psi (𝒮)^+)=\mathrm{}`$. By passing to a further subset of $`𝒩`$ if necessary, one may arrange that each point of $`𝒩^{}`$ is a future endpoint of a timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ in $`𝒩`$ from $`^{}𝒩`$. Lemma 3.4 gives that there exists a slice $`\mathrm{\Sigma }_2^+`$ of $`^+`$ such that every past endless timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ to $`\mathrm{\Sigma }_2^+`$ cuts $`𝒩`$. One may assume that $`\mathrm{\Sigma }_2^+`$ lies strictly to the past of $`\mathrm{\Sigma }_0^+`$ along the null geodesic generators of $`^+`$. Since $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ is partially future asymptotically predictable from $`\psi (𝒮)`$ there exists a slice $`\mathrm{\Sigma }_3^+`$ of $`^+`$ lying strictly to the past of $`\mathrm{\Sigma }_2^+`$ along the generators of $`^+`$ such that $`\mathrm{\Sigma }_3^+\overline{D}\text{}^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$.
Suppose there exists a timelike curve $`\alpha :[0,1]\overline{M}\text{}`$ of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ from $`\psi (𝒦)\overline{M}\text{}(\psi (𝒰)\overline{M}\text{})`$ to $`\mathrm{\Sigma }_3^+`$. In view of $`J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})J^{}(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})\psi (𝒰)\overline{M}\text{}`$ the set $`I^{}(\alpha (0),\overline{𝐠}\text{};\overline{M}\text{})`$ cannot intersect $`J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})`$. So, because every past endless timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ to $`\mathrm{\Sigma }_3^+`$ cuts $`𝒩J^+(\mathrm{\Sigma }^{},\overline{𝐠}\text{};\overline{M}\text{})`$, there exists $`a(0,1)`$ such that $`\alpha (a)𝒩^{}`$. By the construction of $`𝒩`$ there exists $`\overline{x}\text{}I^{}(\alpha (a),\overline{𝐠}\text{};𝒩)^{}`$. If $`\alpha |[a,1]`$ did not cut $`\psi (𝒮)`$ one could concatenate the past endless null geodesic generating segment of $`^{}\overline{M}\text{}\psi (𝒮)`$ to $`\overline{x}\text{}`$, a timelike curve in $`𝒩\overline{M}\text{}\psi (𝒮)`$ from $`\overline{x}\text{}`$ to $`\alpha (a)`$, and the segment $`\alpha |[a,1]`$ of $`\alpha `$ from $`\alpha (0)`$ to $`\alpha (1)`$ to obtain a past endless causal curve $`\beta `$ of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ to $`\alpha (1)\mathrm{\Sigma }_3^+`$ which did not cut $`\psi (𝒮)`$. For an open neighbourhood $`𝒪_{\alpha (1)}\overline{M}\text{}\overline{\psi (𝒮)}`$ of $`\alpha (1)\overline{M}\text{}\overline{\psi (𝒮)}`$ in $`\overline{M}\text{}`$ there would exist $`c(a,1)`$ such that $`\alpha (c)𝒪_{\alpha (1)}`$. But then $`I^+(\alpha (c),\overline{𝐠}\text{};𝒪_{\alpha (0)})`$ would be an open neighbourhood of $`\alpha (1)\mathrm{\Sigma }_3^+`$ in $`\overline{M}\text{}`$ not intersecting $`D^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$, which gives a contradiction. Thus $`\alpha |[a,1]`$ must cut $`\psi (𝒮)`$ and indeed there must exist $`b(a,1)`$ such that $`\alpha (b)\psi (𝒮)`$. One now has that $`\alpha |[0,b]`$ is a timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ from $`\alpha (0)\psi (𝒦)`$ to $`\alpha (b)\psi (𝒮)`$. Hence $`I^{}(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$ intersects $`\psi (𝒦)\overline{I}\text{}^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$ and so intersects $`I^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$. This contradicts the achronality of $`\psi (𝒮)`$ in $`(\overline{M}\text{},\overline{𝐠}\text{})`$. Hence there can be no timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ from $`\psi (𝒦)`$ to $`\mathrm{\Sigma }_3^+`$. One thus has $`\psi (𝒦)I^{}(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})=\mathrm{}`$.
Suppose there exists $`\overline{y}\text{}(\overline{I}\text{}^+(\psi (𝒦),\overline{𝐠}\text{};\overline{M}\text{})^+)J^+(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})`$. Since $`^+J^+(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})`$ is relatively open in $`^+`$ one can construct an open neighbourhood $`𝒪_{\overline{y}\text{}}`$ of $`\overline{y}\text{}`$ in $`\overline{M}\text{}`$ such that every point of $`𝒪_{\overline{y}\text{}}^+`$ is a past endpoint of a timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ in $`𝒪_{\overline{y}\text{}}`$ to $`^+J^+(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})J^{}(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})`$. Then $`I^+(\psi (𝒦),\overline{𝐠}\text{};\overline{M}\text{})`$ intersects $`𝒪_{\overline{y}\text{}}^+I^{}(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})`$ and so $`I^{}(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})`$ intersects $`\psi (𝒦)`$, which is impossible. Hence $`\overline{I}\text{}^+(\psi (𝒦),\overline{𝐠}\text{};\overline{M}\text{})`$ does not intersect $`^+J^+(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})`$. There follows $`\overline{I}\text{}^+(\psi (𝒦),\overline{𝐠}\text{};\overline{M}\text{})^+J^+(\mathrm{\Sigma }_3^+,\overline{𝐠}\text{};\overline{M}\text{})`$.
It is now possible to give the proof of Theorem 3.2.
Proof of Theorem 3.2. One may, by passing to a subset of $`𝒰`$ if necessary, assume $`𝒰𝒯=\mathrm{}`$.
Suppose, for the purpose of obtaining a contradiction, that $`\psi (𝒯)J^{}(^+,\overline{𝐠}\text{};\overline{M}\text{})`$ is non-empty. Then $`J^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})^+`$ is non-empty and so is $`I^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})^+`$. By Lemma 3.6 there exists a slice $`\mathrm{\Sigma }^+`$ of $`^+`$ such that $`\overline{I}\text{}^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})^+J^+(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})`$. Since $`J^+(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})`$ is a non-empty proper subset of $`^+`$ it follows that $`\overline{I}\text{}^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})^+`$ is a non-empty proper subset of $`^+`$. Hence there exists $`\overline{q}\text{}\dot{I}\text{}^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})^+J^+(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})`$. There exists a null geodesic generator $`\overline{\gamma }\text{}`$ of $`\dot{I}\text{}^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})`$ to $`\overline{q}\text{}`$ having either a past endpoint in $`\psi (𝒯)`$ or no past endpoint in $`\overline{M}\text{}`$. In the former case $`\overline{\gamma }\text{}`$ could not be a null geodesic generating segment of $`^+`$ because it would have a past endpoint in $`\psi (𝒯)\overline{M}\text{}\overline{M}\text{}`$. In the latter case $`\overline{\gamma }\text{}`$ could not be a null geodesic generating segment of $`^+`$ because it would then cut $`^+J^+(\mathrm{\Sigma }^+,\overline{𝐠}\text{};\overline{M}\text{})`$. Hence one has $`|\overline{\gamma }\text{}|\{\overline{q}\text{}\}\overline{M}\text{}\overline{M}\text{}`$.
Suppose $`\overline{\gamma }\text{}`$ were past endless in $`(\overline{M}\text{},\overline{𝐠}\text{})`$. Then $`\overline{\gamma }\text{}`$ would be a past endless causal curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ in $`\dot{I}\text{}^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})`$. Let $`\overline{r}\text{}^+`$ lie strictly to the future of $`\overline{q}\text{}`$ along the null geodesic generator of $`^+`$ though $`\overline{q}\text{}`$. Then one could deform $`\overline{\gamma }\text{}`$ to the future in $`(\overline{M}\text{},\overline{𝐠}\text{})`$ so as to give a past endless timelike curve $`\overline{\gamma }\text{}_+`$ of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ to $`\overline{r}\text{}`$ in $`I^+(|\overline{\gamma }\text{}|,\overline{𝐠}\text{};\overline{M}\text{})I^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})I^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$. Clearly $`\overline{\gamma }\text{}_+`$ could not intersect $`\psi (𝒮)`$ because $`\psi (𝒮)`$ is achronal in $`(\overline{M}\text{},\overline{𝐠}\text{})`$. Since $`\overline{\gamma }\text{}_+`$ is timelike curve of $`(\overline{M}\text{},\overline{𝐠}\text{})`$ to $`\overline{r}\text{}\overline{M}\text{}\overline{\psi (𝒮)}`$ in $`(\overline{M}\text{},\overline{𝐠}\text{})`$ it follows that there would exist a neighbourhood of $`\overline{r}\text{}`$ in $`\overline{M}\text{}`$ that did not intersect $`D^+(\psi (𝒮),\overline{𝐠}\text{};\overline{M}\text{})`$. This would be contrary to the future asymptotic predictability of $`(\overline{M}\text{},\overline{𝐠}\text{},\mathrm{\Omega },\psi )`$ from $`\psi (𝒮)`$. Thus $`\overline{\gamma }\text{}`$ must have a past endpoint in $`(\overline{M}\text{},\overline{𝐠}\text{})`$ at $`\psi (𝒯)`$. Consequently there exists a null geodesic $`\gamma `$ of $`(M,𝐠)`$ such that $`\psi \gamma `$ is the unique maximal segment of $`\overline{\gamma }\text{}`$ in $`\overline{M}\text{}\overline{M}\text{}`$.
For each $`p|\gamma |`$, every open neighbourhood of $`\overline{p}\text{}:=\psi (p)`$ in $`\overline{M}\text{}`$ intersects both $`I^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})\overline{M}\text{}=\psi (I^+(𝒯,𝐠;M))`$ and $`\overline{M}\text{}I^+(\psi (𝒯),\overline{𝐠}\text{};\overline{M}\text{})=\psi (MI^+(𝒯,𝐠;M))`$. One thus has $`|\gamma |\dot{I}\text{}^+(𝒯,𝐠;M)`$ and hence that $`\gamma `$ is a null geodesic generator of $`\dot{I}\text{}^+(𝒯,𝐠;M)`$. Since $`\overline{\gamma }\text{}`$ has a past endpoint at $`\psi (𝒯)`$ in $`(\overline{M}\text{},\overline{𝐠}\text{})`$ it follows that $`\gamma `$ has a past endpoint at $`𝒯`$ in $`(M,𝐠)`$. Since $`\overline{\gamma }\text{}`$ has a future endpoint at $`\overline{q}\text{}^+`$ in $`(\overline{M}\text{},\overline{𝐠}\text{})`$ it follows that $`\gamma `$ is future endless and future complete in $`(M,𝐠)`$. One may assume that $`\gamma `$ is an affine future-directed null geodesic of $`(M,𝐠)`$ of the form $`\gamma :[0,\mathrm{})M`$.
Let $`𝐤`$ and $`𝐥`$ be null normal fields to $`𝒯`$ along a relative open neighbourhood $`𝒱_{\gamma (0)}`$ of $`\gamma (0)`$ in $`𝒯`$, normalised such that $`g_{ab}k^al^b=1`$, with $`𝐤(\gamma (0))=\dot{\gamma }\text{}(0)T_{\gamma (0)}M`$. The induced metric on $`𝒯`$ is given by $`h_{ab}=g_{ab}+2k_{(a}l_{b)}`$, whilst $`\text{}^{(1)}\chi _{ab}:=h_a\text{}^ch_b\text{}^dk_{c;d}`$ and $`\text{}^{(2)}\chi _{ab}:=h_a\text{}^ch_b\text{}^dl_{c;d}`$ are null second fundamental forms of $`𝒯`$ along $`𝒱_{\gamma (0)}𝒯`$. By the definition of a closed trapped surface one has $`\text{}^{(1)}\chi ^a\text{}_a<0`$ and $`\text{}^{(2)}\chi ^a\text{}_a<0`$ along $`𝒱_{\gamma (0)}`$. The vector field $`𝐤`$ along $`𝒱_{\gamma (0)}𝒯`$ defines a congruence of future endless affine null geodesics of $`(M,𝐠)`$ from $`𝒯`$ with tangents that coincide with $`𝐤`$ along $`𝒱_{\gamma (0)}𝒯`$. Let $`𝐤`$ also denote the tangents to these null geodesics. Each tangent vector to $`𝒯`$ at $`\gamma (0)`$ may be Lie propagated along $`\gamma `$ with respect to $`𝐤`$ to yield a vector field $`𝐙`$ along $`\gamma `$. From vanishing torsion one has $`Z^a\text{}_{;b}k^b=k^a\text{}_{;b}Z^b`$ and hence that $`𝐙`$ satisfies the defining equation $`(Z^a\text{}_{;b}k^b)\text{}_{;c}k^c=R^a\text{}_{bcd}k^bk^cZ^d`$ for a Jacobi field along $`\gamma `$. Note that $`𝐙`$ is orthogonal to $`𝐤`$ at $`\gamma (0)`$ and satisfies $`(k_aZ^a)_{;b}k^b=0`$ along $`\gamma `$, and so is orthogonal to $`𝐤`$ along $`\gamma `$.
One may parallelly propagate the vector $`𝐥`$ along the integral curves of $`𝐤`$ and so define $`h_{ab}=g_{ab}+2k_{(a}l_{b)}`$ along these curves. Then $`h\text{}^a\text{}_b`$ is a projection operator such that $`h\text{}^a\text{}_{b;c}k\text{}^c=0`$. One thus has that $`\text{}^{}Z^a:=h\text{}^a\text{}_bZ^b`$ satisfies $`\text{}^{}Z^a\text{}_{;b}k^b=h\text{}^a\text{}_bk^b\text{}_{;c}\text{}^{}Z^c=\text{}^{(1)}\chi \text{}^a\text{}_b\text{}^{}Z^b`$ for $`\text{}^{(1)}\chi _{ab}:=h_a\text{}^ch_b\text{}^dk_{c;d}`$ now defined all along $`\gamma `$. One also has $`(\text{}^{}Z^a\text{}_{;b}k\text{}^b)\text{}_{;c}k\text{}^c=h\text{}^a\text{}_bR\text{}^b\text{}_{cde}k\text{}^ck\text{}^d\text{}^{}Z^e`$.
The expansion and shear tensors of the vector fields $`\text{}^{}𝐙`$ along $`\gamma `$ may be expressed as $`\vartheta _{ab}=h_a\text{}^ch_b\text{}^dk_{c;d}`$ and $`\varsigma _{ab}=\vartheta _{ab}\frac{1}{2}\vartheta h_{ab}`$ respectively, where $`\vartheta :=h^{ab}\vartheta _{ab}`$ is the scalar expansion. Then, defining $`\varsigma ^2:=\frac{1}{2}\varsigma _{ab}\varsigma ^{ab}`$, one has that $`\vartheta (\lambda )`$ satisfies the Raychaudhuri equation
$`{\displaystyle \frac{d}{d\lambda }}\vartheta (\lambda )`$ $`=R_{ab}k^ak^b2\varsigma ^2(\lambda )\frac{1}{2}\vartheta ^2(\lambda )`$ (1)
and is subject to the initial condition
$`\vartheta (0)`$ $`=\chi ^a\text{}_a(\gamma (0))<0.`$ (2)
By means of condition (3) of Theorem 3.2 one thus has
$$\frac{d}{d\lambda }\left(\frac{2}{\vartheta (\lambda )}\right)1$$
(3)
for all $`\lambda [0,\mathrm{})`$ such that $`\vartheta (\lambda )0`$. Hence there exists $`\lambda _0(0,2/(\chi ^a\text{}_a(\gamma (0)))]`$ such that $`lim_{\lambda \lambda _0}\vartheta (\lambda )=\mathrm{}`$. Thus $`\gamma (\lambda )`$ is conjugate to $`𝒯`$ at $`\lambda =\lambda _0`$ and so there exists a Jacobi field $`𝐙`$ along $`\gamma `$ which is non-zero and tangent to $`𝒯`$ at $`\gamma (0)𝒯`$ and such that $`\text{}^{}𝐙(\lambda _0)=0`$. One may, by an adaptation of the technique of H&E Proposition 4.5.12, use the vector field $`\text{}^{}𝐙`$ to construct a timelike curve of $`(M,𝐠)`$ from $`𝒯`$ to $`\gamma (\lambda )`$ for any $`\lambda >\lambda _0`$. But this is impossible because $`\gamma `$ is a generator of $`\dot{I}\text{}^+(𝒯,𝐠;M)`$. This establishes the required contradiction.
This then is the revised proof of the familiar assertion that, subject to the null convergence condition and weak cosmic censorship, closed trapped surfaces are not visible from $`^+`$. The key idea of the H&E argument in support of this assertion was evidently sound, but additional constraints and analysis have been seen to be necessary to make the detailed theory of WASE space-times match intuitive expectations.
One could, as discussed in §2, consider weakening the assumed strong causality of the reference ASE asymptote $`(\text{ }\stackrel{~}{M}\text{},\stackrel{~}{𝐠}\text{},\mathrm{\Omega }^{},\psi ^{})`$ to a the chronology condition on the underlying ASE space-time $`(M^{},𝐠^{})`$. This would lead to a more general definition of a WASE space-time. It would be of some interest to find whether or not Theorem 3.2 remains true in this more general setting.
## 4. Concluding remarks
The definition of a WASE space-time, proposed in as the foundation for certain types of cosmic censorship theorems, has been seen also to provide the basis for a rigorous proof of standard relativity folklore concerning the invisibility of closed trapped surfaces from $`^+`$. Other approaches to the definition of WASE space-time have been attempted , but a comparison will not be attempted here.
Whilst this paper was in preparation, Chruściel et al. have been reconsidering another piece of relativity folklore, namely the area theorem for black holes. This states that, subject to weak cosmic censorship and an energy condition, the area of a black hole cannot decrease. A formalised statement to this effect appears as Proposition 9.2.7 in H&E, but the proof there is flawed because it is based on unsubstantiated assumptions concerning the smoothness of the event horizon. A proof is provided in of a weaker area theorem, but with full attention to matters of differentiability. A central hypothesis is one of $``$-regularity which, in present terminology, requires that there exists a neighbourhood $`𝒪`$ of the event horizon $`:=\dot{J}\text{}^{}(^+,\overline{𝐠}\text{};\overline{M}\text{})`$ in $`\overline{M}\text{}`$ such that, for any compact set $`𝒞𝒪`$ which intersects $`I^{}(^+,\overline{𝐠}\text{};\overline{M}\text{})`$, there is a null geodesic generator of $`^+`$ which cuts both $`\overline{I}\text{}^+(𝒞,\overline{𝐠}\text{};\overline{M}\text{})`$ and $`\overline{M}\text{}\overline{I}\text{}^+(𝒞,\overline{𝐠}\text{};\overline{M}\text{})`$. The authors consider various ways to derive $``$-regularity from seemingly more natural hypotheses. However, Lemma 3.6 of the present paper shows that for space-times which are WASE in the sense of and asymptotically chronologically consistent in the sense of Definition 2.8, $``$-regularity is in fact a consequence of partial future asymptotic predictability.
## acknowledgements
I am grateful to Piotr Chruściel for having prompted me to write this article, and for his subsequent comments. Financial support was provided by the NFR (Natural Sciences Research Council) of Sweden.
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LAVAL-PHY-99-20
Generating-function method for fusion rules
L. Bégin<sup>1</sup> Work supported by NSERC (Canada)., C. Cummins<sup>♯2</sup> and P. Mathieu<sup>2</sup> Work supported by NSERC (Canada) and FCAR (Québec).
Département de Physique, Université Laval, Québec, Canada G1K 7P4
Mathematics Department, University of Concordia, Montréal Québec Canada H3G 1M8
Abstract: This is the second of two articles devoted to an exposition of the generating-function method for computing fusion rules in affine Lie algebras. The present paper focuses on fusion rules, using the machinery developed for tensor products in the companion article. Although the Kac-Walton algorithm provides a method for constructing a fusion generating function from the corresponding tensor-product generating function, we describe a more powerful approach which starts by first defining the set of fusion elementary couplings from a natural extension of the set of tensor-product elementary couplings. A set of inequalities involving the level are derived from this set using Farkas’ lemma. These inequalities, taken in conjunction with the inequalities defining the tensor products, define what we call the fusion basis. Given this basis, the machinery of our previous paper may be applied to construct the fusion generating function. New generating functions for $`\widehat{sp}(4)`$ and $`\widehat{su}(4)`$, together with a closed form expression for their threshold levels are presented.
05/99, revised 04/00 (arXiv:hepth/0005002)
1. Introduction
The basic definition of a fusion coefficient is that it gives the number of independent couplings between three different fields in conformal field theory (cf. also the introduction of ; for a review of conformal field theory, and in particular fusion rules, see ). Even in theories with a Lie group symmetry, the so-called Wess-Zumino-Witten (WZW) models, an intrinsic conformal-field theoretical characterisation is unavoidable. This is manifest in formulae for the fusion coefficients: the most fundamental one is the Verlinde formula , that expresses a fusion coefficient in terms of modular $`S`$ matrix elements:
$$𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}=\underset{\widehat{\sigma }P_+^k}{}\frac{S_{\widehat{\lambda }\widehat{\sigma }}S_{\widehat{\mu }\widehat{\sigma }}S_{\widehat{\nu }\widehat{\sigma }}^{}}{S_{0\widehat{\sigma }}}$$
$`(1.1)`$
Here we use notation appropriate to a WZW model in which primary fields are in one-to-one correspondence with the integrable representations of the spectrum-generating affine algebra at a fixed level $`k`$ (this set is denoted by $`P_+^k`$) and $`0`$ stands for the basic representation, whose finite projection is the scalar representation. Fields are not distinguished from their representation labels. The matrix $`S`$ specifies the linear modular transformation properties of the characters of the primary fields among themselves. Up to a constant fixed by unitarity, it takes the form
$$S_{\widehat{\lambda }\widehat{\mu }}\underset{wW}{}ϵ(w)\mathrm{exp}\left(\frac{2\pi i}{k+g}(w(\lambda +\rho ),\mu +\rho )\right)$$
$`(1.2)`$
where $`g`$ stands for the dual Coxeter number of the algebra under consideration, $`\rho `$ is the Weyl vector, $`\lambda `$ is the finite projection of the affine weight $`\widehat{\lambda }`$ and $`W`$ is the finite Weyl group.
The remarkable fact that the ratio of two $`S`$ matrix elements is a finite character evaluated at a special point yields a close relation between fusion and tensor-product coefficients. Indeed, since the finite character and its evaluation read
$$\chi _\lambda =\frac{\underset{wW}{}ϵ(w)e^{w(\lambda +\rho )}}{_{wW}ϵ(w)e^{w\rho }}\text{and}\chi _\lambda (\xi )=\frac{\underset{wW}{}ϵ(w)e^{(w(\lambda +\rho ),\xi )}}{_{wW}ϵ(w)e^{(w\rho ,\xi )}}$$
$`(1.3)`$
we observe that
$$\chi _\lambda (\xi )=\frac{S_{\widehat{\lambda },\widehat{\sigma }}}{S_{0,\widehat{\sigma }}}\mathrm{with}\xi =\frac{2\pi i}{k+g}(\sigma +\rho )$$
$`(1.4)`$
This leads to the Kac-Walton formula which relates the fusion and the tensor-product coefficients.
The Verlinde formula does not make manifest the basic integrality property of the fusion coefficients. The $`S`$ matrix elements being in general complex numbers, it is not even clear at first sight that the fusion coefficients are real (this follows from the unitarity property of $`S`$). The integrality is ensured by the Kac-Walton formula, but in this case the positivity is not manifest.
It is mainly with the aim of displaying manifestly non-negative formulae for fusion rules that we have looked for fusion generating functions . Although the construction of explicit generating functions has an intrinsic interest, we regard the unravelling of the concept of threshold level \- reviewed below - as being the most important outcome of this analysis. It leads to a complete characterisation of fusion coefficients in terms of the corresponding tensor-product coefficients and a set of threshold levels.
As a result, the interest has shifted from the construction of fusion generating functions to the search for threshold-level computing techniques. For $`\widehat{su}(N)`$, $`N=2,3,4`$, it has been found that the threshold level is coded in a simple way in the Berenstein-Zelevinsky triangles (cf. also section 7.1 of ) describing the various distinct couplings of a tensor product . However, these formulae are difficult to generalise to larger values of $`N`$. Moreover, this approach, based on a diagrammatic description of the tensor product, is limited to the $`\widehat{su}(N)`$ algebras.
The aim of the present paper is to apply the machinery developed in to these problems. We find new generating functions for $`\widehat{sp}(4)`$ and $`\widehat{su}(4)`$, together with a closed form expression for their threshold levels. More importantly, we introduce the concept of fusion basis, that is, the set of linear and homogeneous Diophantine inequalities that describes completely the fusion rules.
The article is organised as follows. In section 2, after introducing some notation, we present a brief review of fusion rules and show, with the example of $`\widehat{su}(2)`$, how tensor-product generating functions and the Kac-Walton algorithm can be used to construct fusion-rule generating functions. A more powerful approach to the problem is then elaborated in section 3. It relies on the conjectural existence of a linear and homogeneous set of inequalities that provides a complete description of fusion rules. Given a set of fusion elementary couplings, Farkas’ lemma is then used as a technique to extract the underlying inequalities. This is what we call a fusion basis, i.e., the basis in terms of which these fusion elementary couplings are the elementary solutions. A complete analysis of the $`\widehat{su}(3),\widehat{sp}(4)`$ and $`\widehat{su}(4)`$ cases is presented in section 4, 5 and 6 respectively. In all three cases, the general expression for the threshold levels is obtained explicitly. Various arguments (based on Giambelli-type formula and level-rank duality) supporting our results are presented in Appendix A. In Appendix B, we recall previous conjectures and clarify their relation to those formulated here.
2. Fusion rules
Let $`\widehat{g}`$ be the affine Lie algebra corresponding to the finite Lie algebra $`g`$. Quantities with hats generally refers to $`\widehat{g}`$. The fundamental weights of $`\widehat{g}`$ are denoted by $`\widehat{\omega }_i`$, $`i=0,1,\mathrm{},r`$, where $`r`$ is the rank of $`g`$. An affine weight may be written as
$$\widehat{\lambda }=\underset{i=0}{\overset{r}{}}\lambda _i\widehat{\omega }_i=[\lambda _0,\lambda _1,\mathrm{},\lambda _r]$$
$`(2.1)`$
If the Dynkin labels $`\lambda _i`$ are nonnegative, the weight $`\widehat{\lambda }`$ is the highest weight of an integrable representation of $`\widehat{g}`$ at level $`k`$, with $`k`$ defined by
$$k=\underset{i=0}{\overset{r}{}}\lambda _ia_i^{}$$
$`(2.2)`$
The $`a_i^{}`$ are the co-marks: $`a_0^{}=1`$, and the remaining $`a_i^{}`$ are the coefficients of expansion of the longest root of $`g`$ in terms of the simple coroots. The set of such weights is denoted $`P_+^k`$.
To the affine weight $`\widehat{\lambda }`$, we associate a weight $`\lambda `$ of the finite algebra $`g`$
$$\lambda =\underset{i=1}{\overset{r}{}}\lambda _i\omega _i=(\lambda _1,\mathrm{},\lambda _r)$$
$`(2.3)`$
where $`\omega _i`$ for $`(i=1,\mathrm{},r)`$ are the fundamental weights of $`g`$. $`\widehat{\lambda }`$ is thus uniquely fixed from $`\lambda `$ and $`k`$. The set of integrable finite weights is written $`P_+`$.
In the conformal field-theory context, fusion rules yield the number of independent couplings between three given primary fields. Here we are interested in fusion rules in WZW models , whose generating spectrum algebra is an affine Lie algebra at integer level.
Denote the multiplicity of the representation $`\widehat{\nu }`$ in the fusion rule $`\widehat{\lambda }\times \widehat{\mu }`$ by
$$\widehat{\lambda }\times \widehat{\mu }=\underset{\widehat{\nu }P_+^k}{}𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}\widehat{\nu }$$
$`(2.4)`$
and denote by $`𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }`$ the multiplicity of the representation $`\nu `$ in the tensor product $`\lambda \mu `$:
$$\lambda \mu =\underset{\nu P_+}{}𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }\nu $$
$`(2.5)`$
where by abuse of notation, we use the same symbol for the highest weight and the highest-weight representation. The precise relation between tensor-product and fusion-rule coefficients is given by the Kac-Walton formula :
$$𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}=\underset{\genfrac{}{}{0pt}{}{\xi P_+}{w\widehat{W},w\widehat{\xi }=\widehat{\nu }P_+^k}}{}𝒩_{\lambda \mu }^{}{}_{}{}^{\xi }ϵ(w)$$
$`(2.6)`$
$`w`$ is an element of the affine Weyl group $`\widehat{W}`$, of sign $`ϵ(w)`$, and the dot indicates the shifted action,
$$w\widehat{\lambda }=w(\widehat{\lambda }+\widehat{\rho })\widehat{\rho }\widehat{\rho }=\underset{i=0}{\overset{r}{}}\widehat{\omega }_i$$
$`(2.7)`$
The Kac-Walton formula can be transformed into a simple algorithm: one first calculates the tensor product of the corresponding finite weights and then extends every weight to its affine version at the appropriate value of $`k`$ and shift-reflects back to the integrable affine sector those weights which have negative zeroth Dynkin label. Weights that cannot be shift-reflected in the integrable sector are ignored (for example this is the case for those which have zeroth Dynkin label equal to $`1`$).
The affine extension of the weights that occur in the tensor product may not be integrable at level $`k`$ but are integrable at level $`2k`$. If we divide the weight space into domains that are mapped into each other by the application of the affine Weyl reflections, then the affine reflections which contribute to the Kac-Walton algorithm, apart from the identity, are those corresponding to the domains next to the fundamental alcove and which lies in the $`P_+`$ cone. This is a crucial property of the Kac-Walton algorithm for its application to the construction of fusion-rule generating functions. Let us denote by $`\widehat{W}_f`$ this finite subset of the affine Weyl group that need to be considered . For instance, the elements of $`\widehat{W}_f`$ for the lowest rank algebras are:
| | $`\widehat{su}(2):`$ | $`\widehat{W}_f=\{id,s_0\}`$ |
| --- | --- | --- |
| | $`\widehat{su}(3):`$ | $`\widehat{W}_f=\{id,s_0,s_1s_0,s_2s_0\}`$ |
| | $`\widehat{su}(4):`$ | $`\widehat{W}_f=\{id,s_0,s_1s_0,s_3s_0,s_2s_1s_0,s_2s_3s_0,s_1s_3s_0,s_0s_1s_3s_0\}`$ |
| | $`\widehat{sp}(4):`$ | $`\widehat{W}_f=\{id,s_0,s_1s_0,s_0s_1s_0\}`$ |
| | $`\widehat{G}_2:`$ | $`\widehat{W}_f=\{id,s_0,s_1s_0,s_2s_0,s_0s_2s_1s_0\}`$ |
$`(2.8)`$
where $`s_i`$ denotes the reflection with respect to the root $`\alpha _i`$. This set of elements $`w`$ can be characterised as follows: these are the elements $`w`$ of the affine Weyl group that satisfy the requirement:
$$w\{2\alpha _0^{}+\alpha _1^{}+\mathrm{}+\alpha _r^{},\alpha _1^{},\mathrm{},\alpha _r^{}\}\mathrm{\Delta }_+^{}$$
$`(2.9)`$
where $`\mathrm{\Delta }_+^{}`$ stands for the set of positive real coroots of the affine algebra under consideration and $`r`$ stands for its rank. This condition is adapted from as further analysed in .
Note also that (2.6) may be rewritten as:
$$\widehat{\nu }P_+^k:𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}=\underset{w\widehat{W}_f^1,w\widehat{\nu }P_+}{}𝒩_{\lambda \mu }^{}{}_{}{}^{w\widehat{\nu }}ϵ(w)$$
$`(2.10)`$
where it is understood that $`w\widehat{\nu }`$ stands for its finite part since it is an index of the tensor-product coefficient. This allows us to study in isolation the contribution of a single weight in the fusion. For instance, for $`\widehat{su}(2)`$ that reads
$$𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}=𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }𝒩_{\lambda \mu }^{}{}_{}{}^{s_0\widehat{\nu }}=𝒩_{\lambda _1\mu _1}^{}{}_{}{}^{\nu _1}𝒩_{\lambda _1\mu _1}^{}{}_{}{}^{\mathrm{\hspace{0.17em}2}k+2\nu _1}$$
$`(2.11)`$
Here is an illustrative example of the Kac-Walton algorithm that will also serve to introduce the key notion of threshold level. Take the following $`sp(4)`$ tensor product: $`(1,1)(1,1)`$. Its decomposition reads
$$(1,1)(1,1)=(0,0)(0,1)\mathrm{\hspace{0.17em}2}(2,0)(0,2)(0,3)\mathrm{\hspace{0.17em}2}(2,1)(2,2)(4,0)$$
$`(2.12)`$
The $`sp(4)`$ comarks are all equal to one so that the affine extension of a weight $`(m,n)`$ at level $`k`$ is $`[kmn,m,n]`$. At level 2, the weights $`(0,3)`$ and $`(2,1)`$ are ignored (they have $`\nu _0=1)`$ and the remaining non-integrable weights are $`[2,2,2]`$ and $`[2,4,0]`$. Since the zeroth simple root is $`\widehat{\alpha }_0=[2,2,0]`$, we have $`s_0[2,2,2]=[0,0,2]`$ and $`s_0[2,4,0]=[0,2,0]`$, so that the resulting fusion is
$$[0,1,1]\times [0,1,1]=[2,0,0][1,0,1][0,2,0]$$
$`(2.13)`$
In the above example, we see that the weights $`(0,0),(0,1),(2,0)`$ appear first at level 2. It is easily checked that they reappear at every level $`k2`$. We then say that their threshold level, usually denoted by $`k_0`$, is $`2`$. The threshold level is thus the smallest value of $`k`$ such that the fusion coefficient $`𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}^{\widehat{\nu }}`$ is non-zero. If we indicate the threshold level by a subindex, by considering the extension of the above tensor product at different levels, we find
$$\begin{array}{cc}\hfill (1,1)(1,1)=(0,0)_2& (0,1)_2(2,0)_2(2,0)_3(0,2)_3\hfill \\ & (0,3)_3\mathrm{\hspace{0.17em}2}(2,1)_3(2,2)_4(4,0)_4\hfill \end{array}$$
$`(2.14)`$
To read off a fusion at fixed level $`k`$, we only keep terms with index not greater than $`k`$. The concept of threshold level was first introduced in . It can be shown (cf. ref. ) that the existence of a threshold level is a consequence the depth rule of Gepner and Witten . The notion of threshold level implies directly that
$$𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k+1)}{}_{}{}^{\widehat{\nu }}\text{and}\underset{k\mathrm{}}{lim}𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}=𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }.$$
$`(2.15)`$
To the triplet $`(\lambda ,\mu ,\nu )`$ there corresponds $`𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }`$ distinct couplings, hence $`𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }`$ values of $`k_0`$, one for each distinct coupling. Let us denote these by $`k_0^{(i)},i=1,\mathrm{},𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }`$, implementing in this notation the natural ordering $`k_0^{(i)}k_0^{(i+1)}`$. Then
$$𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}=\{\begin{array}{cc}\hfill \mathrm{max}& (i)\text{if}kk_{0}^{}{}_{}{}^{(i)}\text{and}𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }0\hfill \\ \hfill 0& \text{if}k<k_{0}^{}{}_{}{}^{(1)}\text{or}𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }=0.\hfill \end{array}$$
$`(2.16)`$
Further variations on the idea of threshold level are presented in .
Let us finally note that the fusion coefficients are invariant under the following action of the outer-automorphism group
$$𝒩_{A\widehat{\lambda },A^{}\widehat{\mu }}^{(k)}{}_{}{}^{AA^{}\widehat{\nu }}=𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}$$
$`(2.17)`$
For example, for $`\widehat{sp}(4)`$, the non-trivial outer automorphism $`a`$ exchanges the zeroth and second root, or equivalently, it acts on weights as $`a[\lambda _0,\lambda _1,\lambda _2]=[\lambda _2,\lambda _1,\lambda _0]`$. Acting on the fusion (2.13) as
$$a[0,1,1]\times a[0,1,1]=[1,1,0]\times [1,1,0]=[2,0,0][1,0,1][0,2,0]$$
$`(2.18)`$
which is easily checked from the tensor product
$$(1,0)(1,0)=(0,0)(0,1)(2,0)$$
$`(2.19)`$
which is non-truncated at level 2. Other fusions at level 2 can be obtained from (2.13) by acting on the weights as follows
$$\begin{array}{cc}\hfill a[0,1,1]\times [0,1,1]=[1,1,0]\times [0,1,1]& =a[2,0,0]a[1,0,1]a[0,2,0]\hfill \\ & =[0,0,2][1,0,1][0,2,0]\hfill \end{array}$$
$`(2.20)`$
The algorithm underlying the Kac-Walton formula suggests a simple road to the construction of fusion-rule generating functions, that is by starting from the tensor-product calculation, but keeping track of the level and taking into account the action of the affine Weyl group. We illustrate the method for the simple $`\widehat{su}(2)`$ case. Recall that the $`su(2)`$ tensor-product generating function reads
$$G^{su(2)}(L,M,N)=\frac{1}{(1LM)(1LN)(1MN)}$$
$`(2.21)`$
We start with the generating function
$$F(d,L,M,N)=\frac{1}{(1d)(1LM)(1LN)(1MN)}.$$
$`(2.22)`$
This is just the generating function for $`su(2)`$ tensor products divided by $`(1d)`$. The exponent of $`d`$ will be identified with the level. We will proceed to the generating function for $`\widehat{su}(2)`$ fusion rules by modifying (2.22). First note that at level $`k`$ we need only consider the products of $`su(2)`$ representations $`(a)`$ with $`ak`$. The generating function (2.22) includes products of representations which violate this condition. To keep terms of the form $`d^kL^a`$ with $`ak`$ introduce a dummy variable $`x`$ (using MacMahon’s notation – cf. )
$$\underset{=}{\overset{x}{\Omega }}\frac{1}{(1x^1)}F(dx,Lx^1,M,N)$$
$`(2.23)`$
This first converts $`d^kL^a`$ to $`x^{m+ka}d^kL^a`$, with $`m0`$ and then keeps terms of degree zero in $`x`$ which corresponds to keeping the terms of $`F(d,L,M,N)`$ with $`ak`$ as required. This yields:
$$\frac{1}{(1d)(1dLM)(1dLN)(1MN)}.$$
$`(2.24)`$
Repeating this procedure with $`L`$ replaced by $`M`$ yields:
$$G(d,L,M,N)=\frac{1d^2LMN^2}{(1d)(1dLM)(1dLN)(1dMN)(1dLMN^2)}.$$
$`(2.25)`$
This is still a generating function for tensor products, but with the size of the representation Dynkin labels restricted to be less than or equal to the level.
To take into account the affine Weyl group, consider a term in the expansion of the generating function which contains $`d^kN^c`$. If $`ck+1`$ then this representation is reflected back into the fundamental region of the affine Weyl group: $`cc2(ck1)=c+2k+2`$ or $`d^kN^cd^kN^{2kc+2}`$. Since this is a reflection, the corresponding character must be subtracted. In principle other affine Weyl transformations might be necessary to obtain a weight in the fundamental domain, but, as discussed earlier, for $`\widehat{su}(2)`$ one reflection suffices. At the level of generating functions the effect is to replace
$$G(d,L,M,N)G(d,L,M,N)N^2G(dN^2,L,M,N^1)$$
$`(2.26)`$
Note that the new generating function contains terms with negative powers of $`N`$ and also terms with $`c>k`$. To obtain the final function we projected out the required terms as above. Although this calculation is somewhat long (the verification here was done on a computer), the final result is very simple:
$$G^{\widehat{su}(2)}=\frac{1}{(1d)(1dLM)(1dLN)(1dMN)}$$
$`(2.27)`$
This has first been written down in . There are thus four elementary couplings:
| | $`\widehat{E}_0:d:`$ | $`(0)(0)(0)_1`$ | $`\widehat{E}_2:dLN:`$ | $`(1)(0)(1)_1,`$ |
| --- | --- | --- | --- | --- |
| | $`\widehat{E}_1:dLM:`$ | $`(1)(1)(0)_1,`$ | $`\widehat{E}_3:dMN:`$ | $`(0)(1)(1)_1.`$ |
$`(2.28)`$
As explained above, subscripts indicate the threshold level.
The notion of a model was discussed in . A model for this generating function is $`Q[\widehat{E}_0,\widehat{E}_1,\widehat{E}_2,\widehat{E}_3]`$ with the gradings of $`\widehat{E}_0,\mathrm{},\widehat{E}_3`$ respectively given by $`(1,0,0,0),(1,1,1,0),(1,1,0,1)`$ and $`(1,0,1,1)`$ for the ordering $`X_0=d,X_1=L,X_2=M,X_3=N`$. As for the finite $`su(2)`$ case, there are no relations between the elementary couplings.
The generalisation of the above calculation to other affine Lie algebras is straightforward. Starting from the tensor-product generating function augmented by the factor $`1/(1d)`$, where $`d`$ keeps track of the level, one first enforces the integrability requirement of the first two weights (those that are fused together); one then implements all the affine reflections of the set $`\widehat{W}_f`$ on the third weight and projects the alternating sum onto the integrable sector. However, even though the strategy is clear, the computations become rapidly very complicated.
To bypass this difficulty, we have argued (cf. section 3) that the use of a direct description of tensor products in terms of a system of inequalities (e.g., the Littlewood-Richardson (LR) inequalities underlying their combinatorial description for calculating tensor products – cf. section 4) simplifies the general procedure to a very large extent in addition to allowing us to use powerful algebraic results. We now look for a similar procedure here. However, this program faces an immediate difficulty since even for $`su(N)`$, a combinatorial description of fusion rules is not known. Our method is instead to find an independent route leading to the elementary couplings. Indeed, the elementary couplings are really what we need in order to apply our Grobner basis machinery. Quite remarkably, it turns out that once elementary couplings are found, there is a method that allows us to reconstruct the underlying system of Diophantine inequalities.
3. Fusion-rule elementary couplings
The construction of this section depends upon the following:
Fundamental conjecture: There exists a fusion basis, that is, a set linear and homogeneous inequalities involving $`k`$ and containing as a subset, a tensor-product basis.
For instance, the LR basis is a set linear and homogeneous inequalities. Every solution can be expanded in terms of the elementary solutions of these inequalities. For $`\widehat{su}(N)`$, the conjecture amounts to the existence of a set of additional inequalities involving the level $`k`$ that provide the proper truncation describing the fusion rules. The relation of this conjecture to the conjectures presented in is discussed in the Appendix B.
Note that homogeneity is the key property which allows us to reconstruct the fusion basis from a set of fusion elementary couplings using Farkas’ Lemma. This condition does not necessarily hold, for example we have found that the Lie superalgebra $`osp(1,2)`$ does not have a homogeneous basis.
Given homogeneity and Farkas’ lemma, the problem is reduced to finding a set of fusion elementary couplings. The Kac-Walton algorithm is one possible approach, but a rather difficult one. Instead, we will introduce a simpler approach based on the outer-automorphism group. Unfortunately, it relies on another conjecture.
Let us start from the set of tensor-product elementary couplings $`\{E_i,iI\}`$ for some set $`I`$ fixed by the algebra under study. For each $`E_i`$, we calculate the threshold level $`k_0(E_i)`$. This information specifies the affine extension of $`E_i`$. The affine extension of a tensor-product elementary coupling is necessarily a fusion-rule elementary coupling given our hypothesis that the fusion basis contains, as a subsystem, the set of inequalities that describe tensor products. Denoting by a hat the affine extension of a tensor-product elementary coupling
$$\widehat{E}_i=d^{k_0(E_i)}E_i$$
$`(3.1)`$
we have then a partial set of fusion elementary couplings with the set $`\{\widehat{E}_i,iI\}`$. Our conjecture is that the missing fusion elementary couplings can all be generated by the action of the outer-automorphism group whenever this group is nontrivial:
The outer-automorphism completeness conjecture: The complete set of elementary couplings $`\{\widehat{E}_i,iJ\}`$ for a set $`JI`$ can be generated by the action of the outer-automorphism group on the set $`\{\widehat{E}_i,iI\}`$, i.e., the full set is contained in $`\{𝒜\widehat{E}_i\}`$:
$$\{\widehat{E}_i,iJ\}\{𝒜\widehat{E}_i,iI\}$$
$`(3.2)`$
The action of the outer-automorphism group on a coupling is defined as follows. Let the three weights in the coupling be $`\{\widehat{\lambda },\widehat{\mu };\widehat{\nu }\}`$ where $`\widehat{\nu }\widehat{\lambda }\times \widehat{\mu }`$, then
$$𝒜\{\widehat{\lambda },\widehat{\mu };\widehat{\nu }\}=\{A\widehat{\lambda },A^{}\widehat{\mu };AA^{}\widehat{\nu }\}$$
$`(3.3)`$
where $`A,A^{}`$ are arbitrary elements of the outer-automorphism group; the conjectured completeness requires the consideration of all possible pairs $`(A,A^{})`$.
It should be stressed that we do not suppose that the action of $`𝒜`$ on an elementary coupling will necessarily produce another elementary coupling. Indeed, the resulting coupling could be a product of elementary couplings. What is conjectured here is that all fusion elementary couplings can be generated in this way.
If the outer-automorphism group is trivial, we expect that there will a single extra elementary coupling, the one associated to the scalar coupling: $`\widehat{E}_0`$.
As a simple example consider $`\widehat{su}(2)`$. Start with the elementary coupling $`E_1:(1)(1)(0)`$. It is easy to show that this coupling arises at level 1. This is thus the value of its threshold level. The corresponding fusion is $`[0,1]\times [0,1][1,0]`$. We now consider all possible actions of the outer-automorphims group on it. Since this group is of order 2, there are 4 possible choices for the pair
$$(A,A^{})\{(a,a),(1,1),(1,a),(a,1)\}$$
$`(3.4)`$
with $`a[\lambda _0,\lambda _1]=[\lambda _1,\lambda _0]`$. This generates the following set of four elementary couplings found previously (cf. eq (2.28)):
| | $`\widehat{E}_0:d:`$ | $`[1,0]\times [1,0][1,0]`$ | $`\widehat{E}_2:dLN:`$ | $`[0,1]\times [1,0][0,1]`$ |
| --- | --- | --- | --- | --- |
| | $`\widehat{E}_1:dLM:`$ | $`[0,1]\times [0,1][1,0]`$ | $`\widehat{E}_3:dMN:`$ | $`[1,0]\times [0,1][0,1].`$ |
$`(3.5)`$
Let us then suppose that we have a complete set of fusion elementary couplings which are the elementary solutions of set of linear and homogeneous inequalities that we are looking for. A standard theorem in the theory of linear Diophantine equations (cf. ) states that every non-negative integer solution of a given set of homogeneous Diophantine inequalities for the variables $`x_i`$ (e.g., for $`su(N)`$, these are the $`\{\lambda _i,n_{ij}\}`$) can be generated from a non-negative combination of the fundamental solutions. Hence, given the set of elementary couplings $`\{\widehat{E}_i\}`$, any coupling can be decomposed (maybe not uniquely) in the form $`_i\widehat{E}_i^{a_i}`$. Let the grading variables representing the $`x_i`$ be denoted by $`X_i`$. To the expression of $`g(\widehat{E}_i)`$ corresponds a vector $`ϵ_i`$ of components $`ϵ_{ij}`$, which is the vector form of the elementary solutions of the Diophantine equations. In other words,
$$\widehat{E}_i:g(\widehat{E}_i)=\underset{j}{}X_j^{ϵ_{ij}}$$
$`(3.6)`$
Reading off a particular coupling means that we are interested in a specific set of non-negative integers $`\{x_i\}`$ given by
$$\underset{i}{}a_iϵ_{ij}=x_j$$
$`(3.7)`$
in terms of non-negative integers $`a_i`$. We are thus looking for the existence conditions for such a coupling. This is related to Farkas’ lemma . The standard, rational, form of the lemma is (cf. , corollary 7.1d):
Farkas’ lemma: Let $`V`$ an $`m\times n`$ matrix with rational entries and let $`xQ^m`$. Then there exists $`a0`$, $`aQ^n`$ such that $`Va=x`$ if and only if for all $`uQ^m`$, $`u^{}V0`$ implies $`u^{}x0`$.
We can relate this to our problem in the following way. First note that the condition that for all $`uQ^m`$, $`u^{}V0`$ implies $`u^{}x0`$ is equivalent to the condition that for all $`uZ^m`$, $`u^{}V0`$ implies $`u^{}x0`$. Necessity is clear and sufficiency follows since if $`uQ^m`$ and $`u^{}V0`$ then $`u=cu^{}`$ with $`cQ`$, $`c>0`$ and $`u^{}Z^m`$. Then $`u^{}V0`$, so $`u^{}x0`$ and multiplying by $`c`$ gives the required inequality.
Now consider the inequalities
$$u^{}V0,uZ^m.$$
$`(3.8)`$
By writing $`u_i=w_iv_i`$, $`w_i,v_iN`$, $`i=1\mathrm{}m`$, we obtain a new system of linear Diophantine inequalities. It is not difficult to see that every solution to (3.8) can be obtained from a solution to this new system. Moreover, the new system of linear Diophantine inequalities has a finite set of fundamental solutions. These give rise to a set of fundamental solutions to (3.8) such that every solution to (3.8) is a linear combination of these fundamental solutions with non-negative integer coefficients. Call these fundamental solutions $`s_i`$, $`i=1\mathrm{}k`$. Thus the condition that for all $`uZ^m`$, $`u^{}V0`$ implies $`u^{}x0`$ is equivalent to the condition $`s_i^{}x0`$, $`i=1\mathrm{}k`$.
Putting all this together we obtain the following variation of Farkas’ lemma:
Lemma: Let $`V`$ be an $`m\times n`$ matrix with rational entries and let $`xQ^m`$. Then there exists $`a0`$, $`aQ^n`$ such that $`Va=x`$ if and only if $`s_i^{}x0`$, $`i=1\mathrm{}k`$ where $`s_i`$, $`i=1\mathrm{}k`$ are a fundamental set of solutions of the system $`u^{}V0`$, $`uZ^m`$.
We can reformulate this Lemma over the integers in a form which is more convenient for our application:
Proposition A: Suppose $`VM_{m,n}(N)`$ and let $`xN^m,aN^n`$. Then $`Va=x`$ if and only if $`u_i^{}x=\alpha _i^{}a`$, $`i=1\mathrm{}k`$ where $`u_i,\alpha _i`$, $`i=1\mathrm{}k`$ are a fundamental set of solutions of the system $`u^{}V=\alpha ^{}`$, $`uZ^m`$, $`\alpha N^n`$.
To show this, suppose that $`Va=x`$ and that $`u^{}V=\alpha ^{}`$ with $`uZ^m`$, $`\alpha N^n`$. Then $`u^{}x=u^{}Va=\alpha ^{}a`$. In particular this is true for the fundamental solutions.
Conversely, suppose $`u_i^{}x=\alpha _i^{}a`$ for every fundamental solution. Then $`u^{}x=\alpha ^{}a`$ for every $`uZ^m`$, $`\alpha N^n`$ such that $`u^{}V=\alpha ^{}`$. Since $`VM_{m,n}(N)`$, one set of solutions to $`u^{}V=\alpha ^{}`$, $`uZ^m`$, $`\alpha N^n`$ is given by taking $`u`$ to be a suitable unit vector and $`\alpha ^{}`$ to be a row of $`V`$ which gives $`Va=x`$ as required.
To link the lemma to the situation presented above, we note that the entries $`V_{ij}`$ of the matrix $`V`$ are given here by the numbers $`ϵ_{ji}`$ appearing in (3.6). Our analogue of the relation $`Va=x`$ describes a generic coupling and our goal is to find the defining system of inequalities underlying the existence of this coupling. The equalities $`u_i^{}x=\alpha _i^{}a`$ $`i=1\mathrm{}k`$ imply that $`x`$ satisfies $`u_i^{}x0`$ $`i=1\mathrm{}k`$ since $`\alpha _i`$ and $`a`$ are non-negative. In general these inequalities have solutions which are not solutions of the former equalities for any $`a`$. For example if $`V=(2)`$, then $`Va=x`$ is $`2a=x`$ which is also the equality obtained from the second part of the Proposition A. Thus $`x`$ is a non-negative even integer. But the corresponding inequality is $`x0`$. However, we have found that for the particular systems we consider, this does not happen - as can be easily verified by computing the fundamental set of solutions to the inequalities $`u_i^{}x0`$ $`i=1\mathrm{}k`$ and verifying that they are the columns of $`V`$.
As a simple illustration of this construction, let us work out the example of $`\widehat{su}(2)`$. We use the LR variables $`\{k,\lambda _1,n_{11},n_{12}\}`$ and the corresponding grading variables $`\{d,L_1,N_{11},N_{12}\}`$ in terms of which the elementary couplings and the corresponding vectors are
| | $`\widehat{E}_0:d`$ | $`ϵ_0=(1,0,0,0)`$ |
| --- | --- | --- |
| | $`\widehat{E}_1:dL_1N_{12}`$ | $`ϵ_1=(1,1,0,1)`$ |
| | $`\widehat{E}_2:dL_1`$ | $`ϵ_2=(1,1,0,0)`$ |
| | $`\widehat{E}_3:dN_{11}`$ | $`ϵ_3=(1,0,1,0)`$ |
$`(3.9)`$
For future reference, we display the LR tableaux of the corresponding tensor-product elementary couplings
$$E_1:\begin{array}{c}\text{ 1 1 }\end{array},E_2:\begin{array}{c}\text{ 1 }\end{array},E_3:\begin{array}{c}\text{ 1 }\end{array}$$
$`(3.10)`$
To the fusion elementary couplings, we associate the vectors $`ϵ_j`$ which form the matrix $`V`$ with components $`V_{ij}=ϵ_{ji}`$:
$$V=\left(\begin{array}{cccc}1& 1& 1& 1\\ 0& 1& 1& 0\\ 0& 0& 0& 1\\ 0& 1& 0& 0\end{array}\right)$$
$`(3.11)`$
and so we have the matrix equation
$$Va=x$$
$`(3.12)`$
This equation describes a general fusion coupling. We now want to unravel the underlying system of inequalities. For this, we use Proposition A, i.e., we find the fundamental solutions of $`u^{}V0`$. This is first transformed into a set of equalities $`u^{}V=\alpha ^{}`$ by introducing new non-negative parameters $`\alpha _i`$:
| | $`u_0=\alpha _0`$ | $`u_0+u_1=\alpha _2`$ |
| --- | --- | --- |
| | $`u_0+u_1+u_3=\alpha _1`$ | $`u_0+u_2=\alpha _3`$ |
$`(3.13)`$
We next apply the vector-basis arguments (see Section 7 of ). Let us choose the $`\alpha _i`$ as our independent variables. (This example is somewhat misleading due to its simplicity: in general not all the $`\alpha _i`$ can be taken as the independent variables.) The dependent variables read then
| | $`u_0=\alpha _0`$ | $`u_2=\alpha _3\alpha _0`$ |
| --- | --- | --- |
| | $`u_1=\alpha _2\alpha _0`$ | $`u_3=\alpha _1\alpha _2`$ |
$`(3.14)`$
The 4 basis vectors are obtained by setting successively one $`\alpha _i`$ equal to 1 and all the others equal to 0. These vectors are written as $`e_i`$ and their entries are
$$e_i=(u_0(\alpha _i=1),u_1(\alpha _i=1),u_2(\alpha _i=1),u_3(\alpha _i=1);\alpha _0,\alpha _1,\alpha _2,\alpha _3)$$
$`(3.15)`$
With $`i=0,1,2,3`$, we find
| | $`e_0=(1,1,1,0;1,0,0,0)`$ | $`e_3=(0,0,1,0;0,0,0,1)`$ |
| --- | --- | --- |
| | $`e_1=(0,0,0,1;0,1,0,0)`$ | $`e_2=(0,1,0,1;0,0,1,0)`$ |
$`(3.16)`$
These $`e_i`$ are manifestly linearly independent and they are non-negative expressions in the $`\alpha _i`$. In other words, their grading re-transcription of the above vectors (with $`U_i`$ and $`𝒜_i`$ denoting the grading variables of $`u_i`$ and $`\alpha _i`$ respectively) reads
| | $`_0=U_0U_1^1U_2^1𝒜_0`$ | $`_2=U_1U_3^1𝒜_2`$ |
| --- | --- | --- |
| | $`_1=U_3𝒜_1`$ | $`_3=U_2𝒜_3`$ |
$`(3.17)`$
Here we see that all $`_i`$ contain positive powers of the $`𝒜_i`$ (this is not generic and it reflects the simplicity of the $`su(2)`$ case). Hence, all solutions are generated freely from the non-negative powers of the $`_i`$.
The corresponding linear system of Proposition A is $`e_i(x,a)^{}=0`$ with $`x=(k,\lambda _1,n_{11},n_{12})`$ and $`a=(a_1,a_2,a_3,a_4)`$ non-negative integers:
| | $`k\lambda _1n_{11}=a_1`$ | $`\lambda _1n_{12}=a_3`$ |
| --- | --- | --- |
| | $`n_{12}=a_2`$ | $`n_{11}=a_4`$ |
$`(3.18)`$
which are equivalent to the inequalities:
| | $`k\lambda _1+n_{11}`$ | $`\lambda _1n_{12}`$ |
| --- | --- | --- |
| | $`n_{12}0`$ | $`n_{11}0`$ |
$`(3.19)`$
The last three conditions define the LR basis. The first one is the additional fusion constraint.
In general, we will work the elementary solutions $`e_i`$ in their exponential version $`_i`$ to keep the notation more compact and it should be clear that the (in)equalities can be read off as easily at this level.
The construction of the $`\widehat{su}(2)`$ generating function is now straightforward: since there are no relations between the elementary couplings, the generating function is simply (2.27), that is
$$G^{\widehat{su}(2)}=\underset{i=0}{\overset{3}{}}\frac{1}{(1\widehat{E}_i)}$$
$`(3.20)`$
From the $`k`$–inequality of the $`\widehat{su}(2)`$ fusion basis, we read off the threshold level of a coupling as $`k_0=\lambda _1+n_{11}`$, that is
$$k_0=(\lambda _1+\mu _1+\nu _1)/2$$
$`(3.21)`$
The threshold level is also nicely coded in the LR tableaux: all elementary couplings have threshold level 1 and they all have a single column. We can then write directly that
$$k_0=\mathrm{\#}\mathrm{columns}=\lambda _1+n_{11}$$
$`(3.22)`$
and we recover the previous result. For an $`su(2)`$ LR tableau, it is clear that the number of columns is given by this expression. More generally, for $`su(N)`$, it is simple to check that the number of columns is simply
$$\mathrm{\#}\mathrm{columns}=(\lambda +\mu +\nu ,\omega _{N1})=\underset{i=1}{\overset{N1}{}}\lambda _i+n_{11}$$
$`(3.23)`$
where $`\omega _{N1}`$ is the $`N1`$-th fundamental weight.
4. The generating function for $`\widehat{su}(3)`$ fusion rules
The $`su(3)`$ tensor-product elementary couplings are:
$$\begin{array}{cc}\hfill E_1& =\begin{array}{c}\text{ 1 1 2 }\end{array},E_2=\begin{array}{c}\text{ 1 }\end{array},E_3=\begin{array}{c}\text{ 1 }\end{array},E_4=\begin{array}{c}\text{ 1 1 1 }\end{array}\hfill \\ \hfill E_5& =\begin{array}{c}\text{ 1 1 }\end{array},E_6=\begin{array}{c}\text{ 1 2 }\end{array},E_7=\begin{array}{c}\text{ 1 1 }\end{array},E_8=\begin{array}{c}\text{ 1 1 1 2 }\end{array}.\hfill \end{array}$$
$`(4.1)`$
Using the Kac-Walton formula, the threshold level of $`E_1`$ is 1 and the corresponding fusion reads
$$\widehat{E}_1:[0,1,0]\times [0,0,1][1,0,0]$$
$`(4.2)`$
Acting on $`\widehat{E}_1`$ with $`(a^n,a^m;a^{n+m})`$ $`n,m=0,1,2`$ yields the elementary couplings:
| $`\widehat{E}_0:[1,0,0]\times [1,0,0][1,0,0]:`$ | $`d`$ | $`(1,0,0,0,0,0,0,0)`$ |
| --- | --- | --- |
| $`\widehat{E}_1:[0,1,0]\times [0,0,1][1,0,0]:`$ | $`dL_1N_{12}N_{23}`$ | $`(1,1,0,0,1,0,0,1)`$ |
| $`\widehat{E}_2:[0,1,0]\times [1,0,0][0,1,0]:`$ | $`dL_1`$ | $`(1,1,0,0,0,0,0,0)`$ |
| $`\widehat{E}_3:[1,0,0]\times [0,1,0][0,1,0]:`$ | $`dN_{11}`$ | $`(1,0,0,1,0,0,0,0)`$ |
| $`\widehat{E}_4:[0,0,1]\times [0,1,0][1,0,0]:`$ | $`dL_2N_{13}`$ | $`(1,0,1,0,0,1,0,0)`$ |
| $`\widehat{E}_5:[0,0,1]\times [1,0,0][0,0,1]:`$ | $`dL_2`$ | $`(1,0,1,0,0,0,0,0)`$ |
| $`\widehat{E}_6:[1,0,0]\times [0,0,1][0,0,1]:`$ | $`dN_{11}N_{22}`$ | $`(1,0,0,1,0,0,1,0)`$ |
| $`\widehat{E}_7:[0,1,0]\times [0,1,0][0,0,1]:`$ | $`dL_1N_{12}`$ | $`(1,1,0,0,1,0,0,0)`$ |
| $`\widehat{E}_8:[0,0,1]\times [0,0,1][0,1,0]:`$ | $`dL_2N_{11}N_{23}`$ | $`(1,0,1,1,0,0,0,1)`$ |
$`(4.3)`$
The last column is the vector $`ϵ_i`$ with entries $`(k,\lambda _1,\lambda _2,n_{11},n_{12},n_{13},n_{22},n_{23})`$. By this procedure, we have thus recovered the affine extension of the 8 tensor-product elementary couplings and found an extra elementary coupling: $`\widehat{E}_0`$.
To derive the fusion basis, we proceed as in the $`su(2)`$ case. The set of variables here is
$$(x_0,x_1,\mathrm{},x_7)=(k,\lambda _1,\lambda _2,n_{11},n_{12},n_{13},n_{22},n_{23})$$
$`(4.4)`$
and the matrix $`V`$ (with columns written in the order $`\widehat{E}_0,\mathrm{},\widehat{E}_8`$) reads
$$V=\left(\begin{array}{ccccccccc}1& 1& 1& 1& 1& 1& 1& 1& 1\\ 0& 1& 1& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 1& 1& 0& 0& 1\\ 0& 0& 0& 1& 0& 0& 1& 0& 1\\ 0& 1& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0& 1\end{array}\right)$$
$`(4.5)`$
The reformulation of $`u^{}V0`$ in terms of equalities by the introduction of appropriate nonnegative parameters reads:
| | $`u_0=\alpha _0`$ | $`u_0+u_2=\alpha _5`$ |
| --- | --- | --- |
| | $`u_0+u_1+u_4+u_7=\alpha _1`$ | $`u_0+u_3+u_6=\alpha _6`$ |
| | $`u_0+u_1=\alpha _2`$ | $`u_0+u_1+u_4=\alpha _7`$ |
| | $`u_0+u_3=\alpha _3`$ | $`u_0+u_2+u_3+u_7=\alpha _8`$ |
| | $`u_0+u_2+u_5=\alpha _4`$ | |
$`(4.6)`$
We have 17 variables and 9 equations, hence 8 free variables. Let us choose them to be the $`\alpha _i`$ except for $`\alpha _5`$. Solving for the dependent variables leads to
| | $`u_0=\alpha _0`$ | $`u_5=\alpha _0+\alpha _1+\alpha _3+\alpha _4\alpha _7\alpha _8`$ |
| --- | --- | --- |
| | $`u_1=\alpha _0+\alpha _2`$ | $`u_6=\alpha _3+\alpha _6`$ |
| | $`u_2=\alpha _1\alpha _3+\alpha _7+\alpha _8`$ | $`u_7=\alpha _1\alpha _7`$ |
| | $`u_3=\alpha _0+\alpha _3`$ | $`\alpha _5=\alpha _0\alpha _1\alpha _3+\alpha _7+\alpha _8`$ |
| | $`u_4=\alpha _2+\alpha _7`$ | |
$`(4.7)`$
The basis vectors $`e_i`$ of this system are obtained by setting one of the $`\alpha _j=1`$ and all the others equal to 0 (with the understanding the $`\alpha _5`$ is excluded from this list of free variables). It appears more natural here to express them in their exponentiated version since a projection will be needed to extract the non-negative fundamental solutions. Denote by $`U_i`$ the grading variable associated to $`u_i`$ and by $`𝒜_i`$ those associated to $`\alpha _i`$, the exponential form of the basis vectors reads
| | $`_0:U_0U_1^1U_3^1U_5^1𝒜_0𝒜_5`$ | $`_4:U_5𝒜_4`$ |
| --- | --- | --- |
| | $`_1:U_2^1U_5U_7𝒜_1𝒜_5^1`$ | $`_5:U_6𝒜_6`$ |
| | $`_2:U_1U_4^1𝒜_2`$ | $`_6:U_2U_4U_5^1U_7^1𝒜_5𝒜_7`$ |
| | $`_3:U_2^1U_3U_5U_6^1𝒜_3𝒜_5^1`$ | $`_7:U_2U_5^1𝒜_5𝒜_8`$ |
$`(4.8)`$
To get the corresponding non-negative couplings, i.e., terms containing only non-negative powers of the $`𝒜_i`$, we must keep only the non-negative powers of the $`_i`$. But this is not sufficient since negative powers of $`𝒜_5`$ can appear: we need to project the free generators of the non-negative $`_i`$ powers
$$\underset{i=0}{\overset{7}{}}\frac{1}{1_i}$$
$`(4.9)`$
to non-negative $`𝒜_5`$ powers, using, say the MacMahon algorithm (cf. the $`\mathrm{\Omega }`$ projection described in section 3 of ). After the projection, all the variables $`𝒜_i`$ are set equal to 1. Here however, it is fairly easy to find out by inspection those non-negative combinations of the $`_i`$ that have non-negative $`𝒜_5`$ terms. These are
$$_0,_2,_4,_5,_6,_7$$
$`(4.10)`$
together with
| | $`_0_1:U_0U_1^1U_2^1U_3^1U_7𝒜_0𝒜_1`$ | $`_1_7:U_7𝒜_1𝒜_8`$ |
| --- | --- | --- |
| | $`_0_3:U_0U_1^1U_2^1U_6^1𝒜_0𝒜_3`$ | $`_3_6:U_3U_4U_6^1U_7^1𝒜_3𝒜_7`$ |
| | $`_1_6:U_4𝒜_1𝒜_7`$ | $`_3_7:U_3U_6^1𝒜_3𝒜_8`$ |
$`(4.11)`$
At this point, we set all $`𝒜_i=1`$. We have thus 12 elementary non-negative solutions and the corresponding inequalities are:
| | $`\lambda _1n_{12}`$ | $`\lambda _2n_{13}`$ | $`\lambda _2+n_{12}n_{13}+n_{23}`$ |
| --- | --- | --- | --- |
| | $`n_{11}n_{22}`$ | $`n_{11}+n_{12}n_{22}+n_{23}`$ | |
$`(4.12)`$
and $`n_{ij}0`$ (except for $`n_{11}0`$ which is implied by the others), which are the LR conditions for $`su(3)`$. There are also three inequalities involving $`k`$:
$$\begin{array}{cc}\hfill k\lambda _1\lambda _2& n_{22}\hfill \\ \hfill k\lambda _1\lambda _2& n_{11}n_{23}\hfill \\ \hfill k\lambda _1& n_{13}+n_{11}\hfill \end{array}$$
$`(4.13)`$
The set of inequalities (4.12) and (4.13) represents the $`\widehat{su}(3)`$ fusion basis.
Before we leave the analysis of the $`\widehat{su}(3)`$ case, let us return to the set of equations (4.7). The last equality gives a relation between different $`\alpha _i`$. Actually this relation signals a relation between different sums of columns of $`V`$. In other words, this signals a relation between products of elementary couplings. Indeed, to link the last equality of (4.7) with such a relation, we recall that the labelling of the $`\alpha _i`$ is that of the elementary couplings, which are the columns of $`V`$. Hence, the sought for relation is simply the product form of the equality with $`\alpha _i\widehat{E}_i`$:
$$\alpha _1+\alpha _3+\alpha _5=\alpha _0+\alpha _7+\alpha _8\widehat{E}_1\widehat{E}_3\widehat{E}_5=\widehat{E}_0\widehat{E}_7\widehat{E}_8$$
$`(4.14)`$
As there is only one relation, it is easy to find the generating function. Forbidding $`\widehat{E}_1\widehat{E}_3\widehat{E}_5`$, we get
$$\begin{array}{cc}\hfill G_1& =\left(\underset{\genfrac{}{}{0pt}{}{i=0}{i1,3,5}}{\overset{8}{}}(1\widehat{E}_i)^1\right)(\frac{1}{(1\widehat{E}_1)(1\widehat{E}_5)}\hfill \\ & +\frac{\widehat{E}_3}{(1\widehat{E}_3)(1\widehat{E}_1)}+\frac{\widehat{E}_3\widehat{E}_5}{(1\widehat{E}_5)(1\widehat{E}_3)})\hfill \end{array}$$
$`(4.15)`$
If instead, we decide to forbid $`\widehat{E}_0\widehat{E}_7\widehat{E}_8`$, we would have
$$\begin{array}{cc}\hfill G^{}& =\left(\underset{\genfrac{}{}{0pt}{}{i=0}{i0,7,8}}{\overset{8}{}}(1\widehat{E}_i)^1\right)(\frac{1}{(1\widehat{E}_0)(1\widehat{E}_7)}\hfill \\ & +\frac{\widehat{E}_8}{(1\widehat{E}_7)(1\widehat{E}_8)}+\frac{\widehat{E}_0\widehat{E}_8}{(1\widehat{E}_8)(1\widehat{E}_0)})\hfill \end{array}$$
$`(4.16)`$
and simple manipulations show that $`G_1=G^{}`$. An independent proof of this generating function is presented in Appendix A.
Given the fusion basis, we can write down directly the threshold level to be
$$k_0=\mathrm{max}(\lambda _1+\lambda _2+n_{11}n_{23},\lambda _1+\lambda _2+n_{22},\lambda _1+n_{11}+n_{13})$$
$`(4.17)`$
This can also be extracted from the generating function as follows. A generic term of the $`\widehat{su}(3)`$ generating function is ($`\widehat{E}_0=d`$)
$$d^\alpha \widehat{E}_1^a\widehat{E}_2^b\widehat{E}_3^c\widehat{E}_4^d\widehat{E}_5^e\widehat{E}_6^f\widehat{E}_7^g\widehat{E}_8^h$$
$`(4.18)`$
with either $`a=0,c=0`$ or $`e=0`$. In all cases the threshold level is simply
$$k_0=a+b+c+d+e+f+g+h$$
$`(4.19)`$
In terms of the grading variables $`L_i`$ and $`N_{ij}`$, the above generic term becomes
$$d^{\alpha +k_0}L_1^{a+b+g}L_2^{d+e+h}N_{11}^{c+f+h}N_{12}^{a+g}N_{13}^dN_{22}^fN_{23}^{a+h}$$
$`(4.20)`$
From this expression we read off the relation between the $`n_{ij}`$ and the variables $`a,\mathrm{},h`$. In each three cases (where one of $`a,c,e`$ is zero), we can then solve for the sum $`k_0=a+b+c+d+e+f+g+h`$. We find
| | $`a=0:k_0`$ | $`=\lambda _1+\lambda _2+n_{11}n_{23}`$ |
| --- | --- | --- |
| | $`c=0:k_0`$ | $`=\lambda _1+\lambda _2+n_{22}`$ |
| | $`e=0:k_0`$ | $`=\lambda _1+n_{11}+n_{13}`$ |
$`(4.21)`$
This leads to the compact expression (4.17) for the $`\widehat{su}(3)`$ threshold level. This is easily checked to be equivalent to the formula given in in terms of BZ triangle data (cf. section 7.1 of ):
$$k_0=\mathrm{max}\{m_{13}+\mu _1+\mu _2,n_{13}+\nu _1+\nu _2,l_{13}+\lambda _1+\lambda _2\}$$
$`(4.22)`$
An explicit formula for the $`\widehat{su}(3)`$ fusion coefficients is written down in .
Notice that the threshold level is also simply encoded in the LR tableaux. Indeed, every elementary couplings has threshold level 1 and it corresponds to the number of columns except for $`E_8`$. This leads directly to the following formula for the threshold level of a general LR tableau
$$kk_0\mathrm{\#}\mathrm{columns}\mathrm{\#}E_8=\mathrm{\#}\mathrm{columns}\mathrm{\#}\begin{array}{c}\text{ 1 1 1 2 }\end{array}$$
$`(4.23)`$
that is, $`k_0`$ is the number of columns minus the total number of $`E_8`$ that we can take out of the tableau while preserving its LR character. Consider for instance:
$$\begin{array}{c}\text{ 1 1 1 1 1 1 2 2 }\end{array}:\begin{array}{c}\text{ 1 1 1 1 1 1 2 2 }\end{array}\begin{array}{c}\text{ 1 1 1 2 }\end{array}=\begin{array}{c}\text{ 1 1 1 2 }\end{array}$$
$`(4.24)`$
After the subtraction of one $`E_8`$, the resulting tableau is not a LR tableau: counting from right to left, we find that a $`\begin{array}{c}\text{ 2 }\end{array}`$ precedes the first $`\begin{array}{c}\text{ 1 }\end{array}`$. Therefore, no $`E_8`$ can be removed and $`k_0`$ is given by the number of columns which is 4.
5. The $`\widehat{sp}(4)`$ generating function
We first recall some results obtained in . The appropriate basis for the description of $`sp(4)`$ tensor products reads :
| | $`\lambda _1p`$ | $`\mu _1q`$ |
| --- | --- | --- |
| | $`\lambda _2r_1/2`$ | $`\mu _1q+r_1r_2`$ |
| | $`\lambda _2r_1/2+qp`$ | $`\mu _1p+r_1r_2`$ |
| | $`\lambda _2r_2/2+qp`$ | $`\mu _2r_2/2`$ |
| | $`\nu _1=r_2r_12p+\lambda _1+\mu _1`$ | $`\nu _2=pqr_2+\lambda _2+\mu _2`$ |
$`(5.1)`$
together with $`p,qN`$ and $`r_i2N`$ for $`i=1,2`$. A proper set of variables for a complete description of a particular tensor-product coupling is thus
$$\{\lambda _1,\lambda _2,\mu _1,\mu _2,r_1,r_2,p,q\}$$
$`(5.2)`$
(notice the absence of the $`\nu _i`$ Dynkin labels). Let the corresponding grading variables be
$$\{L_1,L_2,M_1,M_2,R_1,R_2,P,Q\}$$
$`(5.3)`$
The list of elementary coupling with their grading description is:
| | $`A_1:`$ | $`(0,0)(1,0)(1,0)M_1`$ |
| --- | --- | --- |
| | $`A_2:`$ | $`(1,0)(0,0)(1,0)L_1`$ |
| | $`A_3:`$ | $`(1,0)(1,0)(0,0)L_1M_1PQ`$ |
| | $`B_1:`$ | $`(0,0)(0,1)(0,1)M_2`$ |
| | $`B_2:`$ | $`(0,1)(0,0)(0,1)L_2`$ |
| | $`B_3:`$ | $`(0,1)(0,1)(0,0)L_2M_2R_1^2R_2^2`$ |
| | $`C_1:`$ | $`(0,1)(1,0)(1,0)L_2M_1Q`$ |
| | $`C_2:`$ | $`(1,0)(0,1)(1,0)L_1M_2R_2^2P`$ |
| | $`C_3:`$ | $`(1,0)(1,0)(0,1)L_1M_1P`$ |
| | $`D_1:`$ | $`(2,0)(0,1)(0,1)L_1^2M_2R_2^2P^2`$ |
| | $`D_2:`$ | $`(0,1)(2,0)(0,1)L_2M_1^2R_1^2`$ |
| | $`D_3:`$ | $`(0,1)(0,1)(2,0)L_2M_2R_2^2`$ |
$`(5.4)`$
The relation between elementary couplings are generated by
| | $`C_1C_2=A_3D_3,`$ | $`C_2C_3=A_1D_1`$ | $`C_3C_1=A_1A_3B_2`$ |
| --- | --- | --- | --- |
| | $`D_1D_2=B_3C_3^2`$ | $`D_2D_3=A_1^2B_2B_3`$ | $`D_1D_3=B_2C_2^2`$ |
| | $`C_1D_1=A_3B_2C_2`$ | $`C_2D_2=A_1B_3C_3`$ | $`C_3D_3=A_1B_2C_2`$ |
$`(5.5)`$
To find the fusion elementary couplings, we start by computing the threshold level of $`A_1`$ by the Kac-Walton formula. It is found to be 1. The corresponding level-1 fusion, denoted $`\widehat{A}_1`$, is thus
$$[1,0,0]\times [0,1,0][0,1,0]$$
$`(5.6)`$
We can act on it with the four pairs
$$(A,A^{})=\{(1,1),(a,a),(a,1),(1,a)\}$$
$`(5.7)`$
We obtain in this way two copies of $`\widehat{A}_1`$ and two copies of $`\widehat{C}_1`$, the level-1 extension of $`C_1`$. Similarly, $`A_2,`$ and $`A_3`$ are found to have level 1 and this implies the same result for $`C_2,C_3`$. $`B_1`$ is also found to have threshold level 1. Acting on it with the above sequence of outer automorphisms leads successively to $`B_1,B_2,B_3`$ and a new coupling, $`\widehat{E}_0`$:
$$\widehat{E}_0:[1,0,0]\times [1,0,0][1,0,0]$$
$`(5.8)`$
Finally, $`D_1,D_2`$ and $`D_3`$ have threshold level 2 and they are all fixed with respect to the action of the outer-automorphism group. The set $`\{\widehat{A}_i,\widehat{B}_i,\widehat{C}_i,\widehat{D}_i,\widehat{E}_0\}`$ is thus our candidate complete set of fusion elementary couplings, whose explicit expression in terms of grading variables is read from their tensor-product relative with the addition of an appropriate factors of $`d`$.
Having obtained the fusion elementary couplings, we now work out the corresponding fusion basis. Introduce the set of variables
$$(x_0,x_1,\mathrm{},x_8)=(k,\lambda _1,\lambda _2,\mu _1,\mu _2,r_1,r_2,p,q)$$
$`(5.9)`$
This fixes the ordering of the rows of $`V`$. The matrix $`V`$ is built from the columns which form the different elementary couplings in the order $`\widehat{E}_0,\widehat{A}_1,\mathrm{},\widehat{D}_3`$:
$$V=\left(\begin{array}{ccccccccccccc}1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 2& 2& 2\\ 0& 0& 1& 1& 0& 0& 0& 0& 1& 1& 2& 0& 0\\ 0& 0& 0& 0& 0& 1& 1& 1& 0& 0& 0& 1& 1\\ 0& 1& 0& 1& 0& 0& 0& 1& 0& 1& 0& 2& 0\\ 0& 0& 0& 0& 1& 0& 1& 0& 1& 0& 1& 0& 1\\ 0& 0& 0& 0& 0& 0& 2& 0& 0& 0& 0& 2& 0\\ 0& 0& 0& 0& 0& 0& 2& 0& 2& 0& 2& 0& 2\\ 0& 0& 0& 1& 0& 0& 0& 0& 1& 1& 2& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 1& 0& 0& 0& 0& 0\end{array}\right)$$
$`(5.10)`$
The transcription of the inequalities $`u^{}V0`$ into the equalities $`u^{}V=\alpha ^{}`$ takes the following form
| | $`u_0=\alpha _0`$ | $`u_0+u_2+u_3+u_8=\alpha _7`$ |
| --- | --- | --- |
| | $`u_0+u_3=\alpha _1`$ | $`u_0+u_1+u_4+2u_6+u_7=\alpha _8`$ |
| | $`u_0+u_1=\alpha _2`$ | $`u_0+u_1+u_3+u_7=\alpha _9`$ |
| | $`u_0+u_1+u_3+u_7+u_8=\alpha _3`$ | $`2u_0+2u_1+u_4+2u_6+2u_7=\alpha _{10}`$ |
| | $`u_0+u_4=\alpha _4`$ | $`2u_0+u_2+2u_3+2u_5=\alpha _{11}`$ |
| | $`u_0+u_2=\alpha _5`$ | $`2u_0+u_2+u_4+2u_6=\alpha _{12}`$ |
| | $`u_0+u_2+u_4+2u_5+2u_6=\alpha _6`$ | |
$`(5.11)`$
Solving for the dependent variables $`u_i,\alpha _j`$, $`i=0,\mathrm{},8`$ and $`j=6,7,8,9`$ gives
| | $`u_0=\alpha _0`$ | $`u_7=\frac{1}{2}(\alpha _02\alpha _2+\alpha _5+\alpha _{10}\alpha _{12})`$ |
| --- | --- | --- |
| | $`u_1=\alpha _0+\alpha _2`$ | $`u_8=\frac{1}{2}(\alpha _02\alpha _1+2\alpha _3\alpha _5\alpha _{10}+\alpha _{12})`$ |
| | $`u_2=\alpha _0+\alpha _5`$ | $`\alpha _6=2\alpha _1\alpha _5+\alpha _{11}+\alpha _{12}`$ |
| | $`u_3=\alpha _0+\alpha _1`$ | $`\alpha _7=\frac{1}{2}(\alpha _0+2\alpha _3+\alpha _5\alpha _{10}+\alpha _{12})`$ |
| | $`u_4=\alpha _0+\alpha _4`$ | $`\alpha _8=\frac{1}{2}(\alpha _0\alpha _5+\alpha _{10}+\alpha _{12})`$ |
| | $`u_5=\frac{1}{2}(\alpha _02\alpha _1\alpha _5+\alpha _{11})`$ | $`\alpha _9=\frac{1}{2}(\alpha _0+2\alpha _1+\alpha _5+\alpha _{10}\alpha _{12})`$ |
| | $`u_6=\frac{1}{2}(\alpha _4\alpha _5+\alpha _{12})`$ | |
$`(5.12)`$
As usual, the basis vectors $`e_i`$ of this system are obtained by setting one of the $`\alpha _i=1`$ and all the others equal to 0, excluding $`\alpha _6,\mathrm{},\alpha _9`$. We will give their exponentiated version, where as before, we denote by $`U_i`$ the grading variable associated to $`u_i`$ and by $`𝒜_i`$ those associated to $`\alpha _i`$:
$$\begin{array}{cc}& _0:U_0U_1^1U_2^1U_3^1U_4^1U_5^{1/2}U_7^{1/2}U_8^{1/2}𝒜_0𝒜_7^{1/2}𝒜_8^{1/2}𝒜_9^{1/2}\hfill \\ & _1:U_3U_5^1U_8^1𝒜_1𝒜_6^2𝒜_9\hfill \\ & _2:U_1U_7^1𝒜_2\hfill \\ & _3:U_8𝒜_3𝒜_7\hfill \\ & _4:U_4U_6^{1/2}𝒜_4\hfill \\ & _5:U_2U_5^{1/2}U_6^{1/2}U_7^{1/2}U_8^{1/2}𝒜_5𝒜_6^1𝒜_7^{1/2}𝒜_8^{1/2}𝒜_9^{1/2}\hfill \\ & _6:U_7^{1/2}U_8^{1/2}𝒜_7^{1/2}𝒜_8^{1/2}𝒜_9^{1/2}𝒜_{10}\hfill \\ & _7:U_5^{1/2}𝒜_6𝒜_{11}\hfill \\ & _8:U_6^{1/2}U_7^{1/2}U_8^{1/2}𝒜_6𝒜_7^{1/2}𝒜_8^{1/2}𝒜_9^{1/2}𝒜_{12}\hfill \end{array}$$
$`(5.13)`$
Next we keep only those combinations of the $`_i`$ that contain only non-negative integer powers of the $`𝒜_i`$. This projection is not so simple to work out by inspection. We thus need to use a more systematic procedure:
Consider a general expansion of the form $`_i(1_i)^1`$ and in a generic term of the form $`_i_i^{ϵ_i^{}}`$, let us collect the number of $`𝒜_i`$ factors (denote by $`a_i`$ their exponents). Of course we are only interested in those $`𝒜_i`$ that appear with negative powers, namely $`i=6,7,8,9`$. Their powers can be read off from the $`𝒜_i`$ in (5.13) and this yields the following expressions:
$$\begin{array}{cc}& a_6=2ϵ_1^{}ϵ_5^{}+ϵ_7^{}+ϵ_8^{}0\hfill \\ & 2a_7=ϵ_0^{}+2ϵ_3^{}+ϵ_5^{}ϵ_6^{}+ϵ_8^{}0\hfill \\ & 2a_8=ϵ_0^{}ϵ_5^{}+ϵ_6^{}+ϵ_8^{}0\hfill \\ & 2a_9=ϵ_0^{}+2ϵ_1^{}+ϵ_5^{}+ϵ_6^{}ϵ_8^{}0\hfill \end{array}$$
$`(5.14)`$
(These equations should be compared with the last four of (5.12), with $`\alpha _iϵ_i^{},i5`$ and $`\alpha _ie_{i4}^{}`$ for $`i10`$). We then look for the elementary solutions of this system of inequalities. There are 4 elementary solutions with $`ϵ_0^{}0`$. Their grading reformulation reads
$$_0_1_6_8^2_0_1_7_8_0_5_6_8_0_3_6$$
$`(5.15)`$
Denote their vector-reformulation respectively as $`e_i`$ with $`i=0,1,2,3`$, then the conditions $`e_ix0`$ yield, in the above order
$$\begin{array}{cc}& k\lambda _1+\lambda _2+\mu _2+r_1/2r_2\hfill \\ & k\lambda _1+\lambda _2+\mu _2r_2/2\hfill \\ & k\lambda _1+\mu _1+\mu _2p\hfill \\ & k\lambda _1+\lambda _2+\mu _1+\mu _2pqr_1/2\hfill \end{array}$$
$`(5.16)`$
The other elementary solutions are
$$\begin{array}{cc}& _2,_3,_4,_7,_5_8,_6_8,_1_8^2,_1_7^2,\hfill \\ & _3_6^2,_5_6_7,_5_6^2_8,_1_6^2_8,_1_6_7_8\hfill \end{array}$$
$`(5.17)`$
and the resulting inequalities reproduce the whole set of BZ inequalities (5.1) with the positivity requirement on $`r_i,p`$ and $`q`$ (together with $`\mu _1q+\frac{1}{2}(r_1r_2)`$ which is implied by the other ones).
Let us return to the last four equations in (5.12). As mentioned in connection to the $`\widehat{su}(3)`$ case, they indicate the ‘basic relations’: the correspondence between the $`\alpha _i`$ and the elementary couplings being fixed by the ordering of the columns of $`V`$ (e.g., $`\alpha _3\widehat{A}_3`$ and $`\alpha _7\widehat{C}_1`$). The relations correspond then respectively to
| | $`\widehat{A}_1^2\widehat{B}_2\widehat{B}_3=\widehat{D}_2\widehat{D}_3`$ | $`\widehat{E}_0\widehat{C}_1^2\widehat{D}_1=\widehat{A}_3^2\widehat{B}_2\widehat{D}_3`$ |
| --- | --- | --- |
| | $`\widehat{E}_0\widehat{B}_2\widehat{C}_2^2=\widehat{D}_1\widehat{D}_3`$ | $`\widehat{E}_0\widehat{C}_3^2\widehat{D}_3=\widehat{A}_1^2\widehat{B}_2\widehat{D}_1`$ |
$`(5.18)`$
The first and third relations appear in the list (5.5). All other linear relations in the set (5.5) can be obtained from products of the above four, allowing for the cancellations of common factors. For instance, consider the product of the left factors of the second and third relations; equating this with the product of the right factors yields
$$\widehat{E}_0^2\widehat{B}_2\widehat{C}_1^2\widehat{C}_2^2\widehat{D}_1=\widehat{A}_3^2\widehat{B}_2\widehat{D}_1\widehat{D}_3^2$$
$`(5.19)`$
Cancelling the $`\widehat{B}_2\widehat{D}_1`$ terms and taking the square root gives
$$\widehat{E}_0\widehat{C}_1\widehat{C}_2=\widehat{A}_3\widehat{D}_3$$
$`(5.20)`$
which is the affine extension of the relation $`C_1C_2=A_3D_3`$. All other linear relations can be obtained in a similar way.
We can write the $`\widehat{sp}(4)`$ generating function in the compact form
$$\begin{array}{cc}\hfill G& =\overline{E}_0\overline{B}_1\overline{B}_2\overline{B}_3[\overline{A}_1\overline{A}_2\overline{A}_3\overline{C}_1\overline{C}_2\overline{C}_3(1\widehat{A}_1\widehat{A}_3\widehat{B}_2)+\widehat{D}_1\overline{D}_1\overline{A}_2\overline{A}_3\overline{C}_2\overline{C}_3\hfill \\ & +\widehat{D}_3\overline{D}_3\overline{A}_1\overline{A}_2\overline{C}_1\overline{C}_2+\widehat{D}_2\overline{D}_2\overline{A}_1\overline{A}_2\overline{A}_3\overline{C}_1\overline{C}_3(1\widehat{A}_1\widehat{A}_3\widehat{B}_2)]\hfill \end{array}$$
$`(5.21)`$
where $`\overline{Q}`$ is defined as
$$\overline{Q}=\frac{1}{1\widehat{Q}}$$
$`(5.22)`$
This can be re-expressed under a manifestly positive form as follows
$$\begin{array}{cc}\hfill G& =\overline{E}_0\overline{A}_1\overline{A}_2\overline{B}_1\overline{B}_2\overline{B}_3\overline{C}_1\overline{C}_2\overline{C}_3+\overline{E}_0\widehat{A}_3\overline{A}_2\overline{A}_3\overline{B}_1\overline{B}_2\overline{B}_3\overline{C}_1\overline{C}_2\overline{C}_3\hfill \\ & +\overline{E}_0\widehat{A}_1\widehat{A}_3\overline{A}_1\overline{A}_2\overline{A}_3\overline{B}_1\overline{B}_3\overline{C}_1\overline{C}_2\overline{C}_3+\overline{E}_0\widehat{D}_1\overline{A}_2\overline{A}_3\overline{B}_1\overline{B}_2\overline{B}_3\overline{C}_2\overline{C}_3\overline{D}_1\hfill \\ & +\overline{E}_0\widehat{D}_3\overline{A}_1\overline{A}_2\overline{B}_1\overline{B}_2\overline{B}_3\overline{C}_1\overline{C}_2\overline{D}_3+\overline{E}_0\widehat{D}_2\overline{A}_1\overline{A}_2\overline{B}_1\overline{B}_2\overline{B}_3\overline{C}_1\overline{C}_3\overline{D}_2\hfill \\ & +\overline{E}_0\widehat{A}_3\widehat{D}_2\overline{A}_2\overline{A}_3\overline{B}_1\overline{B}_2\overline{B}_3\overline{C}_1\overline{C}_3\overline{D}_2+\overline{E}_0\widehat{A}_1\widehat{A}_3\widehat{D}_2\overline{A}_1\overline{A}_2\overline{A}_3\overline{B}_1\overline{B}_3\overline{C}_1\overline{C}_3\overline{D}_2\hfill \end{array}$$
$`(5.23)`$
We should stress that this is essentially a new result. A generating function for $`\widehat{sp}(4)`$ fusion rules was given in ; the approach, however, was ad hoc and the result was not related to any known basis.
As before the information concerning the threshold level that can be deduced from the fusion basis inequalities (5.16) can also be obtained directly from the generating function. A generic term of the $`\widehat{sp}(4)`$ generating function (5.23) reads
$$d^\alpha \widehat{A}_1^a\widehat{A}_2^b\widehat{A}_3^c\widehat{B}_1^d\widehat{B}_2^e\widehat{B}_3^f\widehat{C}_1^g\widehat{C}_2^h\widehat{C}_3^i\widehat{D}_1^j\widehat{D}_2^k\widehat{D}_3^l$$
$`(5.24)`$
Its threshold level is (since all these factors have a single power of $`d`$ except for the three $`\widehat{D}_i=d^2D_i`$):
$$k_0=a+b+c+d+e+f+g+h+i+2j+2k+2l$$
$`(5.25)`$
Now express the elementary couplings in terms of dummy variables $`\{L_1,L_2,M_1,M_2,R_1,R_2,P,Q\}`$ whose exponent are the BZ basis data, respectively $`\{\lambda _1,\lambda _2,\mu _1,\mu _2,r_1,r_2,p,q\}`$:
| | $`A_1=M_1`$ | $`A_2=L_1`$ | $`A_3=L_1M_1PQ`$ |
| --- | --- | --- | --- |
| | $`B_1=M_2`$ | $`B_2=L_2`$ | $`B_3=L_2M_2R_1^2R_2^2`$ |
| | $`C_1=L_2M_1Q`$ | $`C_2=L_1M_2R_2^2P`$ | $`C_3=L_1M_1P`$ |
| | $`D_1=L_1^2M_2R_2^2P^2`$ | $`D_2=L_2M_1^2R_1^2`$ | $`D_3=L_2M_2R_2^2`$ |
$`(5.26)`$
Next, consider each term of the generating function (5.23) and solve for $`k_0`$ in terms of the basis variables. Surprisingly there are only four different formulas for $`k_0`$. The expressions corresponding to the different terms of (5.23) are:
| | $`\mathrm{terms}1,5,6:`$ | $`k_0=\lambda _1+\lambda _2+\mu _1+\mu _2pqr_1/2`$ |
| --- | --- | --- |
| | $`\mathrm{terms}3,8:`$ | $`k_0=\lambda _1+\mu _1+\mu _2p`$ |
| | $`\mathrm{term}7:`$ | $`k_0=\lambda _1+\lambda _2+\mu _2+r_1/2r_2`$ |
| | $`\mathrm{terms}2,4:`$ | $`k_0=\lambda _1+\lambda _2+\mu _2r_2/2`$ |
$`(5.27)`$
Therefore, the threshold formula is the maximum value of these four values or equivalently
$$k_0=\lambda _1+\lambda _2+\mu _1+\mu _2\mathrm{min}(p+q+r_1/2,\lambda _2+p,\mu _1r_1/2+r_2,\mu _1+r_2/2)$$
$`(5.28)`$
Notice that by rewriting $`kk_0`$, we recover from (5.27) the 4 inequalities (5.16).
The system of inequalities (5.1) can be transformed into a system of equations by setting $`r_1/2=s_1`$ and $`r_2/2=s_2`$ and introducing the integers $`a_i`$ (cf. section 7.5 of ):
| | $`\lambda _1=p+a_1`$ | $`\nu _2=a_4+a_8`$ |
| --- | --- | --- |
| | $`\lambda _2=s_1+a_2`$ | $`a_2+p=a_3+q`$ |
| | $`\mu _1=q+a_5`$ | $`a_3+s_1=a_4+s_2`$ |
| | $`\mu _2=s_2+a_8`$ | $`a_5+2s_2=a_6+2s_1`$ |
| | $`\nu _1=a_1+a_7`$ | $`a_6+q=a_7+p`$ |
$`(5.29)`$
As shown in , this leads to the following diamond-type graphical representation of the tensor product:
......................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................$``$$``$$``$$``$$``$$``$$``$$``$$``$$``$$``$$``$$`s_1`$$`s_2`$$`q`$$`p`$$`\lambda _1`$$`\mu _2`$$`\nu _1`$$`\nu _2`$$`\mu _1`$$`\lambda _2`$$`a_6`$$`a_5`$$`a_2`$$`a_3`$$`a_4`$$`a_7`$$`a_8`$$`a_1`$............................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................
Dotted lines relate those two points that compose the label indicated beside it and opposite continuous lines are constrained to be equal, with the length of a line being defined as the sum of its extremal points except for the lines delimited by the points $`(a_6,s_1)`$ and $`(a_5,s_2)`$ where the point $`s_i`$ is counted twice (the little bar besides $`s_1`$ and $`s_2`$ being a reminder of this particularity). For those lines, the constraint reads $`a_6+2s_1=a_5+2s_2`$. Given a triple $`sp(4)`$ product, the number of such diamonds that can be drawn with only non-negative entries gives its multiplicity.
In terms of these data, the expression for the threshold level (5.27) look somewhat more symmetrical: the four expressions in (5.27) correspond respectively to the following terms:
$$k_0=a_1+a_8+\mathrm{max}\{a_4+a_7+s_1,a_5+q+s_2,a_4+q+s_1,a_4+q+s_2\}$$
$`(5.30)`$
6. The $`\widehat{su}(4)`$ generating function
Written directly in terms of LR tableaux, the $`su(4)`$ elementary solutions are:
$$\begin{array}{cc}\hfill A_1& =\begin{array}{c}\text{ 1 2 3 }\end{array},A_2=\begin{array}{c}\text{ 1 1 1 1 }\end{array},A_3=\begin{array}{c}\text{ 1 }\end{array},B_1=\begin{array}{c}\text{ 1 2 }\end{array},B_2=\begin{array}{c}\text{ 1 1 1 2 }\end{array},B_3=\begin{array}{c}\text{ 1 1 }\end{array},\hfill \\ \hfill C_1& =\begin{array}{c}\text{ 1 }\end{array},C_2=\begin{array}{c}\text{ 1 1 2 3 }\end{array},C_3=\begin{array}{c}\text{ 1 1 1 }\end{array},D_1^{}=\begin{array}{c}\text{ 1 1 1 }\end{array},D_2^{}=\begin{array}{c}\text{ 1 1 }\end{array},D_3^{}=\begin{array}{c}\text{ 1 1 2 }\end{array},\hfill \end{array}$$
$`(6.1)`$
and
$$D_1=\begin{array}{c}\text{ 1 1 1 2 3 }\end{array},D_2=\begin{array}{c}\text{ 1 1 1 2 1 3 }\end{array},D_3=\begin{array}{c}\text{ 1 1 1 1 2 }\end{array},E_1=\begin{array}{c}\text{ 1 1 1 1 1 2 }\end{array},E_2=\begin{array}{c}\text{ 1 1 1 2 }\end{array},E_3=\begin{array}{c}\text{ 1 1 1 2 1 3 }\end{array}$$
$`(6.2)`$
The relations are :
| | $`D_j^{^{}}D_k=C_iE_i`$ | $`D_jD_k^{^{}}=B_iC_jC_k`$ | $`E_iE_j=B_kD_kD_k^{^{}}`$ |
| --- | --- | --- | --- |
| | $`D_iE_i=C_jB_kD_k`$ | $`D_i^{^{}}E_i=B_jD_j^{^{}}C_k`$ | |
$`(6.3)`$
with $`i,j,k`$ a cyclic permutation of $`1,2,3`$.
Consider now the construction of the set of fusion elementary couplings using outer-automorphism completeness. Start with $`A_1:(0,0,0)(0,0,1)(0,0,1)`$, this has threshold level 1. Acting on it with
$$(A,A^{})=(a^n,a^m),n,m=0,1,2,3$$
$`(6.4)`$
where
$$a[\lambda _0,\lambda _1\lambda _2,\lambda _3]=[\lambda _3,\lambda _0\lambda _1,\lambda _2]$$
$`(6.5)`$
we generate the affine extension of the whole set $`A_i,B_i,C_i,D_i,D_i^{}`$, which thus all have threshold level 1, together with the scalar coupling $`\widehat{E}_0=[1,0,0,0]\times [1,0,0,0][1,0,0,0]`$. Finally, the affine extension of $`E_1`$ arises first at level 2: $`[0,1,0,1]\times [1,0,1,0][1,0,1,0]`$. The three weights in this coupling are fixed under the action of $`A=a^2`$. Hence, we need only to consider
$$(A,A^{})=\{(1,1),(a,a),(a,1),(1,a)\}$$
$`(6.6)`$
and this leads respectively to $`\widehat{E}_1,\widehat{E}_2,\widehat{E}_3`$, which all have $`k_0=2`$ and a new elementary coupling $`\widehat{F}=[0,1,0,1]\times [0,1,0,1][0,1,0,1]`$ (first discovered in ). Notice that at the level of tensor-products, $`\widehat{F}`$ is a composite product $`C_1C_2C_3`$. But if it were still composite for fusions, it would necessarily have level 3 since $`k_0(C_1C_2C_3)=3`$. This is the reason why $`\widehat{F}`$ must be regarded as a new elementary coupling.
The whole set of fusion elementary couplings is:
| | $`\widehat{A}_1=[1,0,0,0]\times [0,0,0,1][0,0,0,1]`$ | $`\widehat{D}_1^{}=[0,0,1,0]\times [0,1,0,0][0,0,0,1]`$ |
| --- | --- | --- |
| | $`\widehat{A}_2=[0,0,0,1]\times [0,1,0,0][1,0,0,0]`$ | $`\widehat{D}_2^{}=[0,1,0,0]\times [0,1,0,0][0,0,1,0]`$ |
| | $`\widehat{A}_3=[0,1,0,0]\times [1,0,0,0][0,1,0,0]`$ | $`\widehat{D}_3^{}=[0,1,0,0]\times [0,0,1,0][0,0,0,1]`$ |
| | $`\widehat{B}_1=[0,0,0,0]\times [0,0,1,0][0,0,1,0]`$ | $`\widehat{D}_1=[0,0,1,0]\times [0,0,0,1][0,1,0,0]`$ |
| | $`\widehat{B}_2=[0,0,1,0]\times [0,0,1,0][1,0,0,0]`$ | $`\widehat{D}_2=[0,0,0,1]\times [0,0,0,1][0,0,1,0]`$ |
| | $`\widehat{B}_3=[0,0,1,0]\times [1,0,0,0][0,0,1,0]`$ | $`\widehat{D}_3=[0,0,0,1]\times [0,0,1,0][0,1,0,0]`$ |
| | $`\widehat{C}_1=[1,0,0,0]\times [0,1,0,0][0,1,0,0]`$ | $`\widehat{E}_1=[0,1,0,1]\times [1,0,1,0][1,0,1,0]`$ |
| | $`\widehat{C}_2=[0,1,0,0]\times [0,0,0,1][1,0,0,0]`$ | $`\widehat{E}_2=[1,0,1,0]\times [1,0,1,0][0,1,0,1]`$ |
| | $`\widehat{C}_3=[0,0,0,1]\times [1,0,0,0][0,0,0,1]`$ | $`\widehat{E}_3=[1,0,1,0]\times [0,1,0,1][1,0,1,0]`$ |
$`(6.7)`$
together with two couplings that have no elementary finite relative:
$$\widehat{E}_0=[1,0,0,0]\times [1,0,0,0][1,0,0,0]\widehat{F}=[0,1,0,1]\times [0,1,0,1][0,1,0,1]$$
$`(6.8)`$
The tensor-product relations are modified by the appropriate insertions of $`d`$ or $`\widehat{E}_0`$ factors in order to put them at the same threshold level:
| | $`\widehat{E}_0\widehat{D}_j^{^{}}\widehat{D}_k=\widehat{C}_i\widehat{E}_i`$ | $`\widehat{E}_0\widehat{D}_j\widehat{D}_k^{^{}}=\widehat{B}_i\widehat{C}_j\widehat{C}_k`$ | $`\widehat{E}_i\widehat{E}_j=\widehat{E}_0\widehat{B}_k\widehat{D}_k\widehat{D}_k^{^{}}`$ |
| --- | --- | --- | --- |
| | $`\widehat{D}_i\widehat{E}_i=\widehat{C}_j\widehat{B}_k\widehat{D}_k`$ | $`\widehat{D}_i^{^{}}\widehat{E}_i=\widehat{B}_j\widehat{D}_j^{^{}}\widehat{C}_k`$ | $`\widehat{E}_0\widehat{F}=\widehat{C}_1\widehat{C}_2\widehat{C}_3`$ |
$`(6.9)`$
with $`i,j,k`$ a cyclic permutation of $`1,2,3`$.
To get the $`\widehat{su}(4)`$ basis, we first write down the $`V`$ matrix, whose columns are the vectorial transcription of the fusion elementary couplings written in terms of the grading variables. The column ordering corresponds to $`\widehat{E}_0,\widehat{A}_i,\widehat{B}_i,\widehat{C}_i,\widehat{D}_i^{},\widehat{D}_i,\widehat{E}_i,\widehat{F}`$ with $`i=1,2,3`$. The rows are labelled by the LR variables $`(k,\lambda _1,\lambda _2,\lambda _3,n_{11},n_{12},n_{13},n_{14},n_{22},n_{23},n_{24},n_{33},n_{34})`$. The matrix $`V`$ is thus
$$V=\left(\begin{array}{ccccccccccccccccccccc}1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 2& 2& 2& 2& \\ 0& 0& 0& 1& 0& 0& 0& 0& 1& 0& 0& 1& 1& 0& 0& 0& 1& 0& 0& 1& \\ 0& 0& 0& 0& 0& 1& 1& 0& 0& 0& 1& 0& 0& 1& 0& 0& 0& 1& 1& 0& \\ 0& 0& 1& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 1& 1& 1& 0& 0& 1& \\ 0& 1& 0& 0& 1& 0& 0& 1& 0& 0& 0& 0& 0& 1& 1& 1& 0& 1& 1& 1& \\ 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 1& 1& 0& 0& 0& 1& 0& 0& 1& \\ 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 1& 0& \\ 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \\ 0& 1& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 1& 0& \\ 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 1& 1& 0& 0& 0& 1& 0& 1& \\ 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 1& 0& 0& 0& \\ 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \\ 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 1& 1& 0& 0& 0& 1& 1& \end{array}\right)$$
$`(6.10)`$
With $`u=(u_0,\mathrm{},u_{12})`$, the equations $`u^{}V0`$ can be transformed into equalities by introducing the variables $`\alpha _i`$:
| | $`u_0=\alpha _0`$ | $`u_0+u_2+u_6=\alpha _{10}`$ |
| --- | --- | --- |
| | $`u_0+u_4+u_8+u_{11}=\alpha _1`$ | $`u_0+u_1+u_5=\alpha _{11}`$ |
| | $`u_0+u_3+u_7=\alpha _2`$ | $`u_0+u_1+u_5+u_9=\alpha _{12}`$ |
| | $`u_0+u_1=\alpha _3`$ | $`u_0+u_2+u_4+u_9+u_{12}=\alpha _{13}`$ |
| | $`u_0+u_4+u_8=\alpha _4`$ | $`u_0+u_3+u_4+u_8+u_{12}=\alpha _{14}`$ |
| | $`u_0+u_2+u_6+u_{10}=\alpha _5`$ | $`u_0+u_3+u_4+u_{10}=\alpha _{15}`$ |
| | $`u_0+u_2=\alpha _6`$ | $`2u_0+u_1+u_3+u_5+u_{10}=\alpha _{16}`$ |
| | $`u_0+u_4=\alpha _7`$ | $`2u_0+u_2+u_4+u_9=\alpha _{17}`$ |
| | $`u_0+u_1+u_5+u_9+u_{12}=\alpha _8`$ | $`2u_0+u_2+u_4+u_6+u_8+u_{12}=\alpha _{18}`$ |
| | $`u_0+u_3=\alpha _9`$ | $`2u_0+u_1+u_3+u_4+u_5+u_9+u_{12}=\alpha _{19}`$ |
$`(6.11)`$
We have 13 free variables; let us choose them to be the $`\alpha _i`$ for $`i=0,\mathrm{},12`$. Solving for the dependent variables leads to
| | $`u_0=\alpha _0`$ | $`u_5=\alpha _3+\alpha _{11}`$ | $`u_{10}=\alpha _5\alpha _{10}`$ |
| --- | --- | --- | --- |
| | $`u_1=\alpha _0+\alpha _3`$ | $`u_6=\alpha _6+\alpha _{10}`$ | $`u_{11}=\alpha _1\alpha _4`$ |
| | $`u_2=\alpha _0+\alpha _6`$ | $`u_7=\alpha _2\alpha _9`$ | $`u_{12}=\alpha _8\alpha _{12}`$ |
| | $`u_3=\alpha _0+\alpha _9`$ | $`u_8=\alpha _4\alpha _7`$ | |
| | $`u_4=\alpha _0+\alpha _7`$ | $`u_9=\alpha _{11}+\alpha _{12}`$ | |
$`(6.12)`$
together with
| | $`\alpha _{13}=\alpha _0+\alpha _6+\alpha _7+\alpha _8\alpha _{11}`$ | $`\alpha _{17}=\alpha _6+\alpha _7\alpha _{11}+\alpha _{12}`$ |
| --- | --- | --- |
| | $`\alpha _{14}=\alpha _0+\alpha _4+\alpha _8+\alpha _9\alpha _{12}`$ | $`\alpha _{18}=\alpha _4+\alpha _8+\alpha _{10}\alpha _{12}`$ |
| | $`\alpha _{15}=\alpha _0+\alpha _5+\alpha _7+\alpha _9\alpha _{10}`$ | $`\alpha _{19}=\alpha _0+\alpha _7+\alpha _8+\alpha _9`$ |
| | $`\alpha _{16}=\alpha _5+\alpha _9\alpha _{10}+\alpha _{11}`$ | |
$`(6.13)`$
Now, by setting successively $`\alpha _i=1`$ for $`i=0,\mathrm{},12`$ and the others equal to 0, we generate the following set of basis vectors:
| | $`_0=dN_{11}^1L_1^1L_2^1L_3^1𝒜_0𝒜_{13}^1𝒜_{14}^1𝒜_{15}^1𝒜_{19}^1`$ | $`_6=N_{13}^1L_2𝒜_6𝒜_{13}𝒜_{17}`$ |
| --- | --- | --- |
| | $`_1=N_{33}𝒜_1`$ | $`_7=N_{11}N_{22}^1L_2𝒜_7𝒜_{13}𝒜_{15}𝒜_{17}𝒜_{19}`$ |
| | $`_2=N_{14}𝒜_2`$ | $`_8=N_{34}𝒜_8𝒜_{13}𝒜_{14}𝒜_{18}𝒜_{19}`$ |
| | $`_3=N_{12}^1L_1𝒜_3`$ | $`_9=N_{14}^1L_3𝒜_9𝒜_{14}𝒜_{15}𝒜_{16}𝒜_{19}`$ |
| | $`_4=N_{22}^1N_{33}𝒜_4𝒜_{14}𝒜_{18}`$ | $`_{10}=N_{13}N_{24}^1𝒜_{10}𝒜_{15}^1𝒜_{16}^1𝒜_{18}`$ |
| | $`_5=N_{24}𝒜_5𝒜_{15}𝒜_{16}`$ | $`_{11}=N_{12}N_{23}^1𝒜_{11}𝒜_{13}^1𝒜_{16}𝒜_{17}^1`$ |
| | | $`_{12}=N_{23}N_{34}^1𝒜_{12}𝒜_{14}^1𝒜_{17}𝒜_{18}^1`$ |
$`(6.14)`$
We must now look for those combinations that contain only non-negative powers of the $`𝒜_i`$. Since each $`_i`$ contains at least one positive power of $`𝒜_i`$, these must be obtained from positive combinations of the $`_j`$. To find them, it is convenient to proceed as in the analysis of $`\widehat{sp}(4)`$. Denote by $`a_i`$ the number of $`𝒜_i`$ factors in a general term $`_i^{ϵ_i^{}}`$ of the free expansion of the $`_i`$ in non-negative powers we get (equivalently, we can read off the $`𝒜_i`$ from (6.14)):
| | $`a_{13}=ϵ_0^{}+ϵ_6^{}+ϵ_7^{}+ϵ_8^{}ϵ_{11}^{}`$ | $`a_{17}=ϵ_6^{}+ϵ_7^{}ϵ_{11}^{}+ϵ_{12}^{}`$ |
| --- | --- | --- |
| | $`a_{14}=ϵ_0^{}+ϵ_4^{}+ϵ_8^{}+ϵ_9^{}ϵ_{12}^{}`$ | $`a_{18}=ϵ_4^{}+ϵ_8^{}+ϵ_{10}^{}ϵ_{12}^{}`$ |
| | $`a_{15}=ϵ_0^{}+ϵ_5^{}+ϵ_7^{}+ϵ_9^{}ϵ_{10}^{}`$ | $`a_{19}=ϵ_0^{}+ϵ_7^{}+ϵ_8^{}+ϵ_9^{}`$ |
| | $`a_{16}=ϵ_5^{}+ϵ_9^{}ϵ_{10}^{}+ϵ_{11}^{}`$ | |
$`(6.15)`$
These relations are to be compared with (6.13). We thus look for elementary solutions of the system $`a_i0`$ for $`ϵ_i^{}`$ non-negative integers.
The full list of composites that involve $`_0`$ – these are those that generate $`k`$-dependent constraints – is
| | $`_0_4_7,`$ | $`_0_8_9,`$ | $`_0_6_9,`$ | $`_0_9_8,`$ | $`_0_7_9,`$ | $`_0_5_8,`$ |
| --- | --- | --- | --- | --- | --- | --- |
| | $`_0_7_8_{11},`$ | $`_0_7_9_{12},`$ | $`_0_7_9_{10},`$ | $`_0_5_8_{12}_4`$ | $`_0^2_7_8_9`$ | |
$`(6.16)`$
The constraints are ($`_i`$ specifies the vector $`e_j`$ such that the inequality is $`e_jx0)`$:
| | $`_0_4_7:`$ | $`k\lambda _1+\lambda _2+\lambda _3+n_{33}`$ |
| --- | --- | --- |
| | $`_0_8_9:`$ | $`k\lambda _1+\lambda _2+n_{11}+n_{14}n_{34}`$ |
| | $`_0_6_9:`$ | $`k\lambda _1+n_{11}+n_{13}+n_{14}`$ |
| | $`_0_7_8:`$ | $`k\lambda _1+\lambda _2+\lambda _3+n_{22}n_{34}`$ |
| | $`_0_7_9:`$ | $`k\lambda _1+\lambda _2+n_{14}+n_{22}`$ |
| | $`_0_5_8:`$ | $`k\lambda _1+\lambda _2+\lambda _3+n_{11}n_{24}n_{34}`$ |
| | $`_0_7_8_{11}:`$ | $`k\lambda _1+\lambda _2+\lambda _3n_{12}+n_{22}+n_{23}n_{34}`$ |
| | $`_0_7_9_{10}:`$ | $`k\lambda _1+\lambda _2+n_{14}n_{13}+n_{22}+n_{24}`$ |
| | $`_0_8_9_{12}:`$ | $`k\lambda _1+\lambda _2+n_{11}+n_{14}n_{23}`$ |
| | $`_0^2_7_8_9:`$ | $`2k2\lambda _1+2\lambda _2+\lambda _3+n_{14}+n_{22}+n_{11}n_{34}`$ |
$`(6.17)`$
When these inequalities are re-expressed in terms of BZ triangle data, they reproduce the threshold formula presented in .
The $`_0`$-independent elementary solutions, namely
$$\begin{array}{cc}& _1,_2,_3,_4,_5,_6,_7,_8,_9\hfill \\ & _6_{11},_4_{12},_7_{11},_8_{12},_9_{10},_5_{10},\hfill \\ & _8_{11}_{12},_9_{10}_{12},_7_{10}_{11}\hfill \end{array}$$
$`(6.18)`$
yield the standard LR inequalities:
| | $`\lambda _1n_{12}`$ | $`n_{11}n_{22}`$ |
| --- | --- | --- |
| | $`\lambda _2n_{13}`$ | $`n_{11}+n_{12}n_{22}+n_{23}`$ |
| | $`\lambda _2+n_{12}n_{13}+n_{23}`$ | $`n_{11}+n_{12}+n_{13}n_{22}+n_{23}+n_{24}`$ |
| | $`\lambda _3n_{14}`$ | $`n_{22}n_{33}`$ |
| | $`\lambda _3+n_{13}n_{14}+n_{24}`$ | $`n_{22}+n_{23}n_{33}+n_{34}`$ |
| | $`\lambda _3+n_{13}+n_{23}n_{14}+n_{24}+n_{34}`$ | |
$`(6.19)`$
and $`n_{ij}0`$, except for $`n_{11}0`$ and $`n_{22}0`$ which are implied by the above equations.
As in the $`\widehat{sp}(4)`$ case, we can check that the relations (6.13) code the ‘basic linear relations’ of the model. Indeed, the 7 relations read from (6.13) are
| | $`\widehat{E}_0\widehat{D}_1\widehat{D}_2^{}=\widehat{B}_3\widehat{C}_1\widehat{C}_2`$ | $`\widehat{E}_1\widehat{D}_1^{}=\widehat{B}_2\widehat{C}_3\widehat{D}_2^{}`$ |
| --- | --- | --- |
| | $`\widehat{E}_0\widehat{D}_2\widehat{D}_3^{}=\widehat{B}_1\widehat{C}_2\widehat{C}_3`$ | $`\widehat{E}_2\widehat{D}_2^{}=\widehat{B}_3\widehat{C}_1\widehat{D}_3^{}`$ |
| | $`\widehat{E}_0\widehat{D}_3\widehat{D}_1^{}=\widehat{B}_2\widehat{C}_1\widehat{C}_3`$ | $`\widehat{E}_3\widehat{D}_3^{}=\widehat{B}_1\widehat{C}_3\widehat{D}_1^{}`$ |
| | $`\widehat{E}_0\widehat{F}=\widehat{C}_1\widehat{C}_2\widehat{C}_3`$ | |
$`(6.20)`$
and these are the generators of all the $`\widehat{su}(4)`$ linear relations.
In order to construct the $`\widehat{su}(4)`$ generating function, we must choose a term ordering. We fix the ordering as follows:
$$\begin{array}{cc}& \{L_1,L_2,L_3,N_{11},N_{12},N_{13},N_{14},N_{22},N_{23},N_{24},N_{33},N_{34},d,\hfill \\ & \widehat{E}_1,\widehat{E}_2,\widehat{E}_3,\widehat{B}_1,\widehat{B}_2,\widehat{B}_3,\widehat{C}_1,\widehat{C}_2,\widehat{C}_3,\widehat{A}_1,\widehat{A}_2,\widehat{A}_3,\widehat{D}_1,\widehat{D}_2,\widehat{D}_3,\widehat{D}_1^{^{}},\widehat{D}_2^{^{}},\widehat{D}_3^{^{}},\widehat{E}_0,\widehat{F}\}\hfill \end{array}$$
$`(6.21)`$
Grobner basis methods yield the forbidden products:
$$\begin{array}{cc}\hfill \{& \widehat{E}_i\widehat{E}_j,\widehat{D}_i^{^{}}\widehat{E}_i,\widehat{D}_i\widehat{E}_i,\widehat{C}_i\widehat{E}_i,\widehat{B}_i\widehat{C}_j\widehat{C}_k,\widehat{B}_2\widehat{C}_3\widehat{D}_1\widehat{D}_2^{^{}},\widehat{B}_1\widehat{C}_3\widehat{D}_1\widehat{D}_2^{^{}},\widehat{B}_1\widehat{C}_2\widehat{D}_3\widehat{D}_1^{^{}},\hfill \\ & \widehat{E}_0\widehat{B}_1\widehat{D}_1\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_2^{^{}},\widehat{F}\widehat{E}_i,\widehat{F}\widehat{B}_i,\widehat{C}_1\widehat{C}_2\widehat{C}_3,\widehat{F}\widehat{C}_1\widehat{C}_2\widehat{C}_3\}\hfill \end{array}.$$
$`(6.22)`$
The different terms of the generating function are fully specified by their denominator (the numerators are introduced to avoid over-counting):
| | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}}\widehat{C}_1\widehat{C}_2\widehat{F},`$ | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_2\widehat{B}_3\widehat{C}_2\widehat{C}_3\widehat{D}_2\widehat{D}_3\widehat{D}_2^{^{}}\widehat{D}_3^{^{}}\widehat{E}_1\widehat{F},`$ |
| --- | --- | --- |
| | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_2\widehat{B}_3\widehat{C}_2\widehat{D}_1\widehat{D}_2\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{E}_3\widehat{F},`$ | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_2\widehat{B}_3\widehat{C}_1\widehat{D}_1\widehat{D}_2\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{E}_3\widehat{F},`$ |
| | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_2\widehat{C}_1\widehat{C}_2\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_2\widehat{C}_1\widehat{C}_2\widehat{D}_1\widehat{D}_2\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ |
| | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_2\widehat{C}_1\widehat{C}_2\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_3\widehat{C}_1\widehat{C}_3\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ |
| | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_3\widehat{C}_1\widehat{C}_3\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_3\widehat{C}_1\widehat{C}_3\widehat{D}_1\widehat{D}_2\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ |
| | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_3\widehat{C}_1\widehat{C}_3\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_3\widehat{C}_1\widehat{C}_3\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_3^{^{}},`$ |
| | $`\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_3\widehat{C}_1\widehat{C}_3\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ | $`\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_2\widehat{B}_3\widehat{C}_1\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ |
| | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_2\widehat{B}_3\widehat{C}_1\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_2\widehat{B}_3\widehat{C}_1\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ |
| | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_2\widehat{B}_3\widehat{C}_1\widehat{D}_1\widehat{D}_2\widehat{D}_1^{^{}}\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_2\widehat{B}_3\widehat{C}_1\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_2^{^{}}\widehat{D}_3^{^{}},`$ |
| | $`\widehat{E}_0\widehat{A}_1\widehat{A}_2\widehat{A}_3\widehat{B}_1\widehat{B}_2\widehat{B}_3\widehat{C}_1\widehat{D}_1\widehat{D}_2\widehat{D}_3\widehat{D}_1^{^{}}\widehat{D}_3^{^{}},`$ | |
$`(6.23)`$
Finally, note that the expression of threshold levels in terms of tableau data, that is, the analogue of the $`\widehat{su}(2,3)`$ formulas written previously is clearly
$$k_0=\mathrm{\#}\mathrm{columns}\mathrm{max}\{\mathrm{\#}D_1+\mathrm{\#}D_2+\mathrm{\#}D_3+\mathrm{\#}C_1C_2C_3\}$$
$`(6.24)`$
since the $`D_i^{}s`$ have level 1 but two columns and $`C_1C_2C_3`$ has level 2 and three columns (corresponding to the $`\widehat{F}`$ fusion coupling.)
7. Conclusion and open problems
We have obtained the fusion generating function for $`\widehat{su}(3,4)`$ and $`\widehat{sp}(4)`$ using the conjectural existence of a fusion basis. In the $`\widehat{su}(3)`$ case a first-principle derivation (presented in Appendix A) provides an independent proof of the results, thus a partial confirmation of the conjectures and the correctness of the underlying fusion basis. Moreover, different tests of the $`\widehat{sp}(4)`$ and $`\widehat{su}(4)`$ generating functions, presented in Appendix A, also support our conjectures and the fusion basis constructions. En passant, we point out that the search for the complete $`\widehat{su}(N)`$ level-rank symmetric function introduced in Appendix A is a quest that deserves further studies.
Although the theme of this paper is the construction of fusion generating functions, our most important result is the unravelling of the fusion basis concept, for which we have provided concrete examples. The main open problem is to find a fundamental and Lie algebraic way of deriving the fusion basis (analogous to the Berenstein-Zelevinsky conjectures ). We observe that the number of $`k`$-type inequalities increases rather quickly with the rank of the algebra: 1 for $`\widehat{su}(2)`$, 3 for $`\widehat{su}(3)`$, 4 for $`\widehat{sp}(4)`$ and 10 for $`\widehat{su}(4)`$. More specifically we would like to find arguments to justify the homogeneity property (on the other hand, the linearity appears to be a generic property, a direct consequence of the Kac-Walton algorithm).
With regard to the automorphism completeness conjecture we note that for simplicity (and because the discussion is to a large extent devoted to $`\widehat{su}(N)`$ for which the outer-automorphism group is rather large) we have focused on the outer-automorphism group as the essential symmetry. It is natural to extend the conjecture to the full symmetry group of fusion coefficients. However, we should stress is that the outer-automorphism conjecture is just a convenient tool. If the conjecture (or its natural extension to the full fusion symmetry group) turns out to be wrong, there are other avenues that could yield the complete set of fusion elementary couplings.
In the present work, the only information on fusion data that has been extracted, out of the fusion basis or the fusion generating function, is the expression for the threshold level in terms of the basis variables. But there are certainly more data that can be lifted. For example, given a triple product with multiplicity $`m`$, to which there correspond $`m`$ values of the threshold levels, we could ask for the expression, in terms of the Dynkin labels, of the minimum and maximum values of $`k_0`$. It is easy to write down some explicit expressions for particular fusion coefficients.
The reformulation of the problem of computing fusion rules in terms of a fusion basis solves, in principle, the quest for a combinatorial method since it reduces a fusion computation to solving inequalities. But we expect that we have not found an optimal solution to the quest for an efficient combinatorial description.
Acknowledgement
L.B. would like to thank P. Dargis for his crucial, albeit involuntary, rôle in bringing Farkas’ lemma to his attention.
Appendix A. Independent verifications of the fusion generating functions
The $`\widehat{su}(3)`$ fusion generating function is not presented here for the first time; it appeared originally in . A sketch of its proof was presented in without details. In this section we present a complete proof of the $`\widehat{su}(3)`$ generating function for fusion rules; in addition, we describe some independent checks confirming the validity of the $`\widehat{su}(4)`$ and $`\widehat{sp}(4)`$ fusion generating functions given in sections 6.3 and 7.3. The first check that we present uses Giambelli-type formulas. These can be viewed as equalities of corresponding expressions in the character rings. Since the $`\widehat{su}(n)`$ and $`\widehat{sp}(n)`$ fusion rings are quotients of the classical character rings (see and references therein), these formulas continue to hold for fusion products. For $`\widehat{su}(4)`$, we present another non-trivial check based on a level-rank duality argument.
A.1. Determinantal formula and the ‘composition’ method: deriving the $`su(3)`$ generating function for tensor products
The Giambelli formula, or more generally, determinantal formulae which give expressions for group characters as determinants, provide another method for calculating fusion generating functions in terms of simpler generating functions. This uses the technique of ‘composition’ of generating functions described previously in section 2.3 of .
The $`su(3)`$ Giambelli formula expresses a general representation in terms of a difference of products of representations with a single non-zero Dynkin label, i.e.,
$$(\lambda _1,\lambda _2)=(\lambda _1+\lambda _2,0)(\lambda _2,0)(\lambda _1+\lambda _2+1,0)(\lambda _21,0)$$
$`(\text{A.}1)`$
This can be rewritten in determinantal form as follows
$$(\lambda _1,\lambda _2)=det\left(\begin{array}{cc}(\lambda _1+\lambda _2,0)& (\lambda _21,0)\\ (\lambda _1+\lambda _2+1,0)& (\lambda _2,0)\end{array}\right)$$
$`(\text{A.}2)`$
Consider first the generating function $`G_1(L_1,L_2,M_1,R_1,R_2)`$ which is the generating function for products of the form: $`(\lambda _1,\lambda _2)(\mu _1,0)`$. Its explicit form is
$$\begin{array}{cc}\hfill G_1=\frac{1}{(1L_1N_1)(1L_2N_2)(1L_2M_1)(1M_1N_1)(1L_1M_1N_2)}& \end{array}$$
$`(\text{A.}3)`$
It is obtained by setting $`M_2=0`$ in the complete tensor-product generating function (cf. section 2.5 in ). Our point here is not to re-derive $`G_1`$ from first principles but simply to show how we can reconstruct the complete generating function out of the partial information contained in $`G_1`$. In the fusion case, we will indicate how the analogue of $`G_1`$ can be obtained, preventing the argument from being circular.
From two copies of $`G_1`$ we form the composite generating function $`G_2`$:
$$G_2(L_1,L_2,M_1,M_2,N_1,N_2)=\underset{=}{\overset{R}{\Omega }}G_1(L_1,L_2,M_1,R_1,R_2)G_1(R_1^1,R_2^1,M_2,N_1,N_2)$$
$`(\text{A.}4)`$
which is the generating function for products of the form
$$(\lambda _1,\lambda _2)(\mu _1,0)(\mu _2,0)$$
$`(\text{A.}5)`$
Note that the generating function for products
$$(\lambda _1,\lambda _2)(\mu _1+1,0)(\mu _21,0)$$
$`(\text{A.}6)`$
is $`M_2M_1^1G_2`$ and so, by (A.1), the generating function for products $`(\lambda _1,\lambda _2)(\mu _1,\mu _2)`$ is:
$$G_3=\underset{}{\overset{M_1}{\Omega }}(G_2M_2M_1^1G_2)$$
$`(\text{A.}7)`$
The coefficient of $`M_1^{\mu _1}M_2^{\mu _2}`$ is the multiplicity of the representation with Dynkin labels $`(\mu _1\mu _2,\mu _2)`$ in the product
$$(\lambda _1,\lambda _2)\left[(\mu _1,0)(\mu _2,0)(\mu _1+1,0)(\mu _21,0)\right]$$
$`(\text{A.}8)`$
To change to variables which carry the Dynkin labels we make the substitution $`M_2M_2M_1^1`$, so that $`M_1`$ now carries the first Dynkin label. This introduces negative powers of $`M_1`$, corresponding to products (A.8) with $`\mu _1<\mu _2`$, which are not required. So we must keep only non-negative degree terms in $`M_1`$ to obtain the final generating function. Denote the resulting expression as $`G_4(L_1,L_2,M_1,M_2,N_1,N_2)`$; it reads
$$\begin{array}{cc}\hfill G_4=& \frac{(1L_1L_2M_1M_2N_1N_2)}{(1L_1N_1)(1L_1M_2)(1L_2M_1)(1L_2N_2)}\hfill \\ & \times \frac{1}{(1M_2N_2)(1M_1N_1)(1L_1M_1N_2)(1L_2M_2N_1)}\hfill \end{array}$$
$`(\text{A.}9)`$
which is the usual form of the $`su(3)`$ generating function (cf. section 2.5 of ).
A.2. Extension of the determinantal formula methods to fusion rules: the $`\widehat{su}(3)`$ case
The starting point for the derivation of the $`\widehat{su}(3)`$ fusion generating function is the generating function for fusions of the form
$$[k\lambda _1\lambda _2,\lambda _1,\lambda _2]\times [k\mu _1,\mu _1,0]$$
$`(\text{A.}10)`$
These fusions are known explicitly and the information on their fusion coefficients can be lifted to the following generating function
$$\begin{array}{cc}\hfill F_1(& d,L_1,L_2,M_1,N_1,N_2)=\hfill \\ & \frac{1}{(1d)(1dL_1N_1)(1dL_2N_2)(1dL_2M_1)(1dM_1N_1)(1dL_1M_1N_2)}\hfill \end{array}$$
$`(\text{A.}11)`$
As explained in the previous subsection, the generating function for products
$$[k\lambda _1\lambda _2,\lambda _1,\lambda _2]\times [k\mu _1\mu _2,\mu _1,0]\times [k\mu _2,\mu _2,0]$$
$`(\text{A.}12)`$
is given by
$$\begin{array}{cc}\hfill F_2(d,L_1,L_2,& M_1,M_2,N_1,N_2)=\hfill \\ & \underset{=}{\overset{z}{\Omega }}\underset{=}{\overset{R}{\Omega }}F_1(z^1d,L_1,L_2,M_1,R_1^1,R_2^1)F_1(z,R_1,R_2,M_2,N_1,N_2).\hfill \end{array}$$
$`(\text{A.}13)`$
Here the variable $`z`$ is introduced in order to keep the level fixed in the composition. By the determinantal formula, the generating function is essentially
$$F_3(d,L_1,L_2,M_1,M_2,N_1,N_2)=\underset{}{\overset{M_1}{\Omega }}(F_2M_2M_1^1F_2)$$
$`(\text{A.}14)`$
except that the coefficient of $`M_1^{\mu _1}M_2^{\mu _2}`$ is the multiplicity of $`(\mu _1\mu _2,\mu _2)`$. Thus the final generating function is
$$F_4=\underset{}{\overset{M_1}{\Omega }}F_3(d,L_1,L_2,M_1,M_2M_1^1,N_1,N_2)$$
$`(\text{A.}15)`$
This reproduces the generating function given in and re-derived above.
A.3. Determinantal formula methods applied to the $`\widehat{sp}(4)`$ and $`\widehat{su}(4)`$ cases
In principle, the above procedure can be used to calculate the fusion rule generating functions for $`\widehat{su}(4)`$ and $`\widehat{sp}(4)`$. Unfortunately, the intermediate expressions are too large to be manageable, even when manipulated with computer assistance. However, it is possible to calculate the specialisation of these generating functions with all but one variable, the level-grading variable, set equal to 1. For example, in the above calculation for $`\widehat{su}(3)`$ we could have set $`L_1=L_2=N_1=N_2=1`$ at the start of the calculation since they are not needed at any intermediate steps. Similarly we can set $`M_2=M_1^1`$ at the last step which has the effect of setting $`M_2=1`$ in the final generating function. If we set all variables equal to 1, except the one that keeps track of the level, then the resulting generating function $`G(d)`$ counts the number of independent couplings at each level. The $`\widehat{su}(4)`$ and $`\widehat{sp}(4)`$ specialised generating functions have been calculated in this way and the results are:
$$G^{\widehat{su}(4)}(d)=\frac{d^6+4d^5+13d^4+16d^3+13d^2+4d+1}{(1d)^{12}(1d^2)}$$
$`(\text{A.}16)`$
and
$$G^{\widehat{sp}(4)}(d)=\frac{d^4+2d^3+5d^2+2d+1}{\left(1d\right)^9\left(1+d\right)}$$
$`(\text{A.}17)`$
These expressions agree with the specialisation of the generating functions found in sections 6 and 7; this thus provides a very strong independent verification of these results. In particular, it corroborates the closure of our set of fusion elementary couplings.
Although we will not present the details of this derivation, we would like to draw attention to some technical issues. There are potentially two problems which could arise in using the determinantal expansions. The first problem is that the determinant may contain terms which have level higher than the initial representation. For example in $`\widehat{su}(3)`$ at level 1 the determinantal expansion of the representation $`(0,1)`$ is
$$(0,1)=det\left(\begin{array}{cc}(1,0)& (0,0)\\ (2,0)& (1,0)\end{array}\right)$$
$`(\text{A.}18)`$
The representation $`(2,0)`$ is integrable only at level 2 and greater. However it can be shown, using the modification rules of , that all such terms in the determinant vanish identically in the $`\widehat{sp}(2n)`$ and $`\widehat{su}(n)`$ fusion rings. Thus, when computing with the determinantal expansions at a given level, we need only consider terms corresponding to representations which exist at that level.
The second complication which can arise is in a sense the converse of the first. There are representations which occur only at levels strictly greater than $`k`$, but which have determinantal expansions which contain products which are defined at level $`k`$. This does not occur for the $`su(n)`$ determinants. However for $`sp(4)`$ this problem can happen. The determinant formula for $`sp(4)`$ is
$$(\lambda _1,\lambda _2)=det\left(\begin{array}{cc}(\lambda _1+\lambda _2,0)& (\lambda _21,0)\\ (\lambda _1+\lambda _2+1,0)+(\lambda _1+\lambda _21,0)& (\lambda _2,0)+(\lambda _22,0)\end{array}\right).$$
$`(\text{A.}19)`$
Take for instance the representation $`(0,2)`$ which does not exist for level 1. However the determinant formula yields
$$(0,2)=(2,0)(2,0)(3,0)(1,0)(1,0)(1,0)+(2,0)(0,0)$$
$`(\text{A.}20)`$
The only product which is defined at level 1 is $`(1,0)(1,0)=(0,1)`$. Thus the above determinant yields the following modification rule: $`(0,2)=(0,1)`$ for $`\widehat{sp}(4)`$ at level $`1`$ (see for more details). Therefore, before converting the exponent of $`M_1`$ into a Dynkin label, we must ensure that it is less than or equal to the exponent of $`d`$. This can be achieved by replacing $`M_1`$ by $`M_1y^1`$ and $`d`$ by $`dy`$ and then projecting onto non-negative powers of $`y`$ and finally setting $`y=1`$.
A.4. Duality
As described in and references therein, there is a duality between fusion rules for $`\widehat{su}(n)`$ at level $`k`$ and fusion rules for $`\widehat{su}(k)`$ at level $`n`$. This duality is somewhat involved when using standard Young tableaux. However, it can be clearly seen using contravariant tableaux. This duality can be used to provide a very nice nontrivial check of the $`\widehat{su}(4)`$ generating function.
As discussed above, if all the grading variables in the $`\widehat{su}(4)`$ fusion generating function are set equal to 1, except for the one associated to the level, we obtain (A.16). However in order to use a duality argument to compare this expression with other generating functions, it needs some modifications. Duality maps Young tableaux to conjugate Young tableaux. For example $`\widehat{su}(3)`$ at level 4 has
$$\begin{array}{c}\text{ 1 }\text{ 1 }\text{ 1 }\text{ 1 }\end{array}$$
$`(\text{A.}21)`$
as a possible tableau and this maps to
$$\begin{array}{c}\text{ 1 1 1 1 }\end{array}$$
$`(\text{A.}22)`$
in $`\widehat{su}(4)`$ at level 3. In other words we have to include in the generating functions the terms corresponding to tableaux which have columns of length $`n`$ in $`\widehat{su}(n)`$. If $`_n(d,L_1,\mathrm{})`$ stands for the original $`\widehat{su}(n)`$ fusion generating function, then the procedure for incorporating tableaux augmented by columns of length $`n`$ – while maintaining the first row smaller or equal to $`k`$ — amounts to calculate
$$g_n(d)\frac{^2}{xy}xy_n(dxy,x^1L_1,x^1L_2,\mathrm{},y^1M_1,y^1M_2,\mathrm{},N_1,N_2,\mathrm{})|_{x=y=1}.$$
$`(\text{A.}23)`$
The effect of this operation is to multiply
$$d^kL_1^{\lambda _1}L_2^{\lambda _2}\mathrm{}M_1^{\mu _1}M_2^{\mu _2}\mathrm{}N_1^{\nu _1}N_2^{\nu _2}\mathrm{by}(k\lambda _1\lambda _2\mathrm{}+1)(k\mu _1\mu _2\mathrm{}+1)$$
$`(\text{A.}24)`$
which is the factor needed to add in all the Young tableaux with all allowed numbers of columns of length $`n`$. In other words, the $`\widehat{su}(3)`$ tableau $`\begin{array}{c}\text{ 1 1 1 }\end{array}`$ at level 5 should appear in following equivalent forms:
$$\begin{array}{c}\text{ 1 1 1 }\end{array},\begin{array}{c}\text{ 1 1 1 1 1 1 }\end{array},\begin{array}{c}\text{ 1 1 1 1 1 1 1 1 1 }\end{array},\begin{array}{c}\text{ 1 1 1 1 1 1 1 1 1 1 1 1 }\end{array}$$
$`(\text{A.}25)`$
that is, it should be counted four times. Doing this and setting all Dynkin-grading variables equal to 1 leads to the following generating functions:
$$\begin{array}{cc}& g_0=\frac{1}{1d}g_1=\frac{1+d}{(1d)^3}g_2=\frac{1+3d+d^2}{(1d)^6}g_3=\frac{d^4+6d^3+10d^2+6d+1}{(1d)^{10}}\hfill \\ & g_4=\frac{d^{10}+13d^9+78d^8+257d^7+513d^6+642d^5+513d^4+257d^3+78d^2+13d+1}{(1d)^{12}(1d^2)^3}\hfill \end{array}$$
$`(\text{A.}26)`$
The first two functions above correspond to the limiting algebras $`\widehat{su}(0)`$ and $`\widehat{su}(1)`$. For $`\widehat{su}(0)`$, there is only the trivial representation and it occurs at any level. Therefore, there is a single coupling at every level and there are no correction factors: $`g_0(d)=_kd^k`$. The function $`g_1`$ can be constructed by duality. We start with the generating function for $`\widehat{su}(k)`$ fusions at level 1. At level 1, we can ignore all relations between the elementary couplings; moreover, we can keep track only of those elementary couplings that occur at level 1: these are the various products involving the fundamental and the scalar representations. The truncated generating function then reads
$$\frac{1}{(1d)_i[(1dL_iN_i)(1dM_iN_i)]_{i,j}(1dL_iM_jN_{i+j})}$$
$`(\text{A.}27)`$
where in the last series of term, the summation is defined modulo $`k`$ with the understanding that $`N_k=1`$. In this function, we replace $`ddxy,L_iL_i/x,M_iM_i/y`$, multiply the result by $`xy`$, differentiate with respect to $`x,y,d`$ and set $`x=y=L_i=M_i=N_i=1,d=0`$ (to keep only the linear term in $`d`$). This gives $`(k+1)^2`$. Hence we have
$$g_1(d)=\underset{k=1}{\overset{\mathrm{}}{}}(k+1)^2d^k=\frac{1+d}{(1d)^3}$$
$`(\text{A.}28)`$
These functions $`g_n(d)`$ display very nice properties:
1- the factor $`(1d)`$ occurs to the power $`(n+2)(n+1)/2`$ in the denominator;
2- the numerator polynomial $`p_n(d)`$ satisfies $`p_n(1/d)d^{\mathrm{deg}(p_n)}=p_n(d)`$;
3- $`p_n(d)`$ has positive coefficients;
4- the difference between the degree of the numerator and denominator is $`2n`$.
The mere fact that $`g_4(d)`$ shares the generic properties of the previous $`g_n`$ functions is supporting evidence for the correctness of the $`\widehat{su}(4)`$ generating function.
The Taylor expansions of the $`g_n(d)`$ functions read:
$$\begin{array}{ccccccccccccc}g_0(d)=& 1& +& d& +& d^2& +& d^3& +& d^4& +& d^5& +\mathrm{}\\ g_1(d)=& 1& +& 4d& +& 9d^2& +& 16d^3& +& 25d^4& +& 36d^5& +\mathrm{}\\ g_2(d)=& 1& +& 9d& +& 40d^2& +& 125d^3& +& 315d^4& +& 686d^5& +\mathrm{}\\ g_3(d)=& 1& +& 16d& +& 125d^2& +& 656d^3& +& 2646d^4& +& 8832d^5& +\mathrm{}\\ g_4(d)=& 1& +& 25d& +& 315d^2& +& 2646d^3& +& 16720d^4& +& 85212d^5& +\mathrm{}\end{array}$$
$`(\text{A.}29)`$
from which duality (i.e. horizontal versus vertical) is completely manifest. (We stress that the ‘built-in duality’ for obtaining $`g_1`$ concerns only the second row and the second column.) In particular the first 4 terms 1, 25, 315 and 2646 of the $`\widehat{su}(4)`$ function match the coefficients of the 5-th column. This again provides independent evidence for the correctness of the $`\widehat{su}(4)`$ fusion generating function out of which the function $`g_4`$ has been constructed. In particular, this is a decisive test of the necessity of the extra elementary coupling $`\widehat{F}`$ and an evidence for the absence of further additional elementary couplings.
From the above functions $`g_n(d)`$ we can construct the sum
$$f(r,d)=g_0(d)+g_1(d)r+g_2(d)r^2+g_3(d)r^3+\mathrm{}$$
$`(\text{A.}30)`$
where $`r`$ is the grading variable associated to the rank +1 (i.e., its exponent is the value of $`n`$ for $`\widehat{su}(n)`$). It satisfies $`f(r,d)=f(d,r)`$ by duality. We speculate that other symmetry properties might be used to provide an explicit formula for $`f(r,d)`$.
We can illustrate this dual symmetry in a particular example. Consider the function $`\stackrel{~}{g}_n(d)`$ that counts the number of couplings of the representation $`[k1,1,0,0,\mathrm{},0]`$ with anything in $`\widehat{su}(n)`$ at level $`k`$. Since the Young tableau of $`(1,0,0,\mathrm{},0)`$ is invariant under a duality transformation exchanging $`k`$ and $`n`$, by summing up the resulting functions multiplied by $`r^n`$, one should produce an expression $`\stackrel{~}{f}(r,d)`$ symmetric in the interchange of $`r`$ and $`d`$. The function $`\stackrel{~}{g}_n(d)`$ is calculated as follows in terms of the original $`\widehat{su}(n)`$ fusion generating function $`_n`$:
$$\stackrel{~}{g}_n(d)\frac{^2}{M_1x}x_n(dx,x^1L_1,x^1L_2,\mathrm{},M_1,1,\mathrm{},1,1\mathrm{})|_{x=y=1,M_1=0,L_1=\mathrm{}=1}$$
$`(\text{A.}31)`$
As explained above, the differentiation with respect to $`x`$ is required in order to take into account all contributing Young diagrams associated to the first representation ($`\lambda )`$. The second representation being fixed to be $`(1,0,0,\mathrm{},0)`$, does not require an adjusting multiplication factor. Setting the variable $`M_1=0`$, after having differentiated with respect to it, simply serves to select the term linear in $`M_1`$. Since the representation $`(1,0,0,\mathrm{},0)`$ does not exist for $`\widehat{su}(0)`$, $`\stackrel{~}{g}_0(d)=0`$. The function $`\stackrel{~}{g}_1(d)`$ is found by duality as explained previously. The first few $`\widehat{su}(n)`$ functions $`\stackrel{~}{g}_n(d)`$ are found to be:
$$\begin{array}{cc}& \stackrel{~}{g}_1(d)=\frac{2dd^2}{(1d)^2}\stackrel{~}{g}_2(d)=\frac{3dd^2}{(1d)^3}\hfill \\ & \stackrel{~}{g}_3(d)=\frac{4dd^2}{(1d)^4}\stackrel{~}{g}_4(d)=\frac{5dd^2}{(1d)^5}\hfill \end{array}$$
$`(\text{A.}32)`$
Fortunately, the general pattern is clear: the expression of $`\stackrel{~}{g}_n`$ is easily guessed to be:
$$\stackrel{~}{g}_n(d)=\frac{(n+1)dd^2}{(1d)^{n+1}}n1$$
$`(\text{A.}33)`$
From this exact form of $`g_n(d)`$, we can write down readily the exact expression for the sum
$$\stackrel{~}{f}(r,d)=\underset{n=1}{\overset{\mathrm{}}{}}\stackrel{~}{g}_n(d)r^n=\frac{dr(2dr)}{(1dr)^2}$$
$`(\text{A.}34)`$
The result is manifestly invariant under the duality transformation that interchanges $`r`$ and $`d`$.
Appendix B. Status of previous conjectures
In this appendix, we would like to clarify the relation between the present work and our previous ones and state precisely in what sense our previous conjectures are either embodied in the present reformulation of the problem or have been proved.
A general approach to the construction of generating functions for fusion rules was proposed in . It was based on the following two conjectures:
1) Every coupling is characterised by a threshold level $`k_0`$. The multiplicity of a triple product at level $`k`$ is given by the number of couplings with threshold levels $`k`$.
2) There is a choice of forbidden couplings such that the threshold level of a coupling is the sum of the threshold levels of its components.
As already mentioned, it can be shown that conjecture 1 is a consequence a sharpened formulation of the depth rule . This leaves us with a single conjecture which we rename:
Conjecture I: There is a choice of forbidden couplings such that the threshold level of a coupling is obtained from the sum of the threshold levels of the elementary couplings that appear in its decomposition.
In the formulation of conjecture I, the element of ‘choice’ refers to the fact that both sides of a tensor-product relation do not always have the same threshold level and which one is taken as the forbidden coupling makes a difference in the generating function for fusion rules. With the notion of a set of elementary fusion couplings, which includes the scalar one (this is a new feature of the present work), all relations acquire equal threshold levels and this choice becomes immaterial. This suggests the following modification of conjecture I:
Conjecture I’: The threshold level of a fusion coupling is read off from its decomposition into the elementary fusion couplings.
A interesting aspect of this reformulation of the conjecture is that it embodies an observation that was presented as a conjecture in , namely that the level is always minimised. More precisely, in the choice of forbidden couplings, we should always forbid the one with higher threshold level. This ‘minimal level’ prescription is automatically taken into consideration here since the relations have identical levels. If one of the products appears in the relation with a factor $`\widehat{E}_0`$, it means that the product without this $`\widehat{E}_0`$ factor occurs at a lower level and it is not forbidden. For instance, the relation $`\widehat{E}_1\widehat{E}_3\widehat{E}_5=\widehat{E}_0\widehat{E}_7\widehat{E}_8`$ indicates that the coupling $`\widehat{E}_7\widehat{E}_8`$ appears at level 2. In the tensor-product relation $`E_1E_3E_5=E_7E_8`$, we have thus effectively forbid the higher-level term of the relation.
Once the notion of fusion elementary couplings in terms of which every coupling can be decomposed (conjecture I’) is introduced, this naturally calls for a reinterpretation in terms of a fusion basis. It is indeed plain that our conjecture (and the mere existence of threshold level) boils down the fundamental conjecture presented in the text, that is, the existence of a fusion basis.
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# Evolution in the Clustering of Galaxies for 𝑧<1.0
## 1 Introduction
One of the key problems in modern astronomy is understanding how galaxies form and evolve. Although this problem might seem rather straightforward, a considerable amount of uncertainty remains, primarily because of the difficulty in separating out the competing effects of density evolution, i.e. variation in the number of galaxies with redshift due to clustering or mergers, from luminosity evolution, i.e. the intrinsic evolution in an individual galaxy’s spectral energy distribution.
Historically, two principle techniques have been used to quantify evolution in the clustering of galaxies. The first technique is to invert the angular correlation function ($`w(\theta )`$) using the Limber equation (Limber, 1954; Peebles, 1980) and an observed or model redshift distribution to estimate the expected change in the amplitude of the angular correlation function for different magnitude intervals and/or cosmologies(e.g., Koo & Szalay 1984; Efstathiou et al. 1991; Roche et al. 1993; Infante & Pritchet 1995; Brainerd et al. 1996). The alternative approach is to compute the spatial correlation function ($`\xi (r)`$) for different epochs directly using spectroscopic redshifts (Le Fèvre et al. , 1996; Carlberg et al. , 1997, 1999; Small et al. , 1999). These two techniques, however, suffer from different limitations that have restricted their utility.
Studies which utilize different magnitude intervals are limited in that an apparent magnitude selection samples galaxies of different intrinsic luminosities at different redshifts, complicating the analysis considerably. Furthermore, the clustering of galaxies on small scales is the result of a highly non-linear, complex process; and, therefore, the actual validity of the power law model for the evolution of the spatial clustering function is not guaranteed, although it is at least a useful diagnostic. Finally, this approach is also limited by the applicability of the assumed redshift distribution, which can heavily influence the theoretical conversion from angular coordinates to spatial coordinates.
The spectroscopic approach, on the other hand, is hindered either by the size of the available samples, especially when the data is binned into distinct redshift intervals, or by the depth or width of the survey (which implies either a redshift variable, intrinsic luminosity selection effect or possible contamination from strong clustering). For example, Le Fèvre et al. (1996) analyze the spatial clustering for 591 galaxies in the CFRS with $`I<22.5`$, Small et al. (1999) analyze the spatial clustering of 831 galaxies with $`r21^m`$ to measure the evolution in the correlation length for $`0.2z0.5`$, and Carlberg et al. (1997) use a sample of 248 galaxies with $`K<21.5`$ while Carlberg et al. (1999) use 2300 high intrinsic luminosity galaxies distributed over a wide area, to determine the spatial correlation function for different epochs.
In this paper, we continue our development of a new technique which uses photometric redshifts to measure the angular correlation function in redshift shells (Brunner, 1997; Connolly et al. , 1998; Brunner et al. , 1999b). This novel approach minimizes the galaxy projection effect inherent in standard angular correlation measurements, while utilizing a significantly large sample that minimizes the effects of shot noise in our analysis. By adopting an ensemble approach, we are able to measure the evolution of clustering with both redshift and intrinsic luminosity. Unless otherwise noted, we assume $`h=1.0`$, $`\mathrm{\Omega }_M=0.3`$, and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, throughout this paper.
## 2 Data
The observations and reduction of the data used in this analysis have been extensively detailed elsewhere (Brunner, 1997; Brunner et al. , 1997, 1999a). In this section, we discuss the important points which impact the rest of our discussion.
### 2.1 Observations
The photometric data analyzed in this paper are located at 14:20, +52:30 covering approximately 0.054 square degrees. All of the photometric data were obtained using the Prime Focus CCD (PFCCD) camera on the Mayall 4 meter telescope at Kitt Peak National Observatory (KPNO). The observations were performed on the nights of March 31 – April 3, 1995, March 18 – 20, 1996, and May 14 – 16, 1996. The PFCCD uses the T2KB CCD, a $`2048^2`$ Tektronix CCD with $`24\mathrm{\mu m}`$ pixel scale, which at $`f/2.8`$ in the 4 meter results in a scale of $`0.47\mathrm{}/`$pixel and a field of view of $`16.0\mathrm{}\times 16.0\mathrm{}`$. All observations were made through the broadband filters: $`U,B,R,\&I`$.
### 2.2 Data Reduction
The photometric data were reduced in the standard fashion. The data were photometrically calibrated to the published Landolt (1992) standard star fields using a curve of growth analysis. A linear regression on the published stellar magnitude, the instrumental magnitude, the airmass, and a color term was performed, and the result translated to a one second standard exposure. We transformed our magnitude system to the AB system (Oke & Gunn, 1983) using published transformations (Fukugita et al. , 1995).
Source detection and photometry were performed using SExtractor version 2.0.8 (Bertin & Arnouts, 1996) with the appropriate correction for the background estimation bug applied (Bertin, 1998). SExtractor was chosen for its ability to perform matched aperture photometry, using the same detection image for each program image. Our detection image was constructed from the $`U,B,R,\&I`$ images using an optimal $`\chi ^2`$ process (Szalay et al. , 1998).
The astrometric solution for our data was determined by matching against a pre-release version of the HST Guide Star Catalog II (Lasker, 1996). The residuals of the final geometric transformation to the GSCII for the reference stars were all less than $`0.15`$ pixels, or equivalently, less than $`0.07\mathrm{}`$.
Our completeness limits were calculated by adding artificially generated galaxies to the final stacked image. The iterative, Monte-Carlo approach we used produced a completeness curve, from which both a $`90\%`$ and $`50\%`$ completeness limits in all four bands were extracted. The $`2\%`$ and $`10\%`$ photometric error magnitude limits were calculated as the mean of all valid detections in the master catalog which had a measured photometric error that was approximately the same as the target photometric error ($`0.1`$ magnitudes for $`10\%`$ photometry and $`0.02`$ magnitudes for $`2\%`$ photometry).
We empirically determined the stellar locus in each band separately using several complementary techniques: the ratio of a core to total magnitude, objects which were classified as stars on overlapping HST images, and objects which were spectroscopically identified as stars. In addition, all objects with $`I<20^m`$ were visually inspected and classified as stellar or non-stellar. The final classification was constructed by taking the union of the four separate classifications in each band, resulting in 505 stellar objects. The number-magnitude distribution of stellar objects agrees with model predictions (Bahcall and Soneira, 1980). The spatial distribution of the stellar objects is fairly random, with the possible minor exception of the image corners where the PSF increases due to focal degradations. As a result, these areas were not utilized during the actual correlation analysis.
### 2.3 Photometric Redshifts
Many cosmological tests are more sensitive to the sample size (i.e. Poisson Noise) than small errors in distance, which makes them perfect candidates for utilizing a photometric redshift catalog, including quantifying the evolution of galaxy clustering. As a result, we have developed an empirical photometric redshift technique (Connolly et al. , 1995; Brunner et al. , 1997; Brunner, 1997; Brunner et al. , 1999a), which is not designed to accurately predict the redshift for a given galaxy (Baum, 1962) or locate high redshift objects (Steidel et al. , 1996). Instead, it is designed to provide distance indicators which are statistically accurate for the entire sample, along with corresponding redshift error estimates.
The calibration data for implementing this technique was taken from overlapping spectroscopic surveys including data from both the Canada France Redshift survey (Lilly et al. , 1995), and the Deep Extragalactic Evolutionary Probe survey (Mould, 1993). The accuracy of any empirically derived relationship is predominantly dependent on the quality of the data used in the analysis — photometric redshifts being no exception. As a result, we imposed several restrictions on the calibrating data in order to minimize the intrinsic dispersion within the photometric redshift relationship. After imposing our quality assurance conditions, we were left with 190 galaxies which formed our calibration sample.
The 190 calibrating redshifts were, therefore, used to derive an iterative piecewise polynomial fit to the galaxy distribution in the $`U,B,R,\&I`$ flux space. This iterative approach utilizes a global fit to determine a rough estimate of the galaxy’s redshift, after which a more accurate local polynomial fit, corresponding to the appropriate redshift interval, was applied (Brunner et al. , 1999a). For each derived polynomial fit, the degrees of freedom remained a substantial fraction of the original data (a second order fit in four variables requires 15 parameters while a third order fit in four variables requires 35 parameters).
To estimate the error in a photometric redshift for the full photometric sample, we adopt a bootstrap with replacement algorithm, in which galaxies are randomly selected from the calibration sample and, once selected, are not removed from the set of calibrating galaxies. Thus, at the extremes, one galaxy could be selected 190 consecutive times or, alternatively, each redshift could be selected exactly once (each of these realizations has the same probability). This approach is designed to emphasize any incompleteness in the sampling of the true distribution of galaxies in the four dimensional space $`U,B,R,\&I`$ by the calibration redshifts. In order to fully account for potential sources of error in the redshift estimation, the magnitudes of the calibrating sample were drawn from a Gaussian probability distribution function with mean given by the measured magnitude and sigma by the magnitude error.
The actual photometric redshift error was calculated from 100 different realizations using this algorithm. For each different realization, a photometric redshift was calculated for every galaxy in the photometric redshift catalog. The actual error was optimally determined to be given by the normalized difference between the fifth and second quantiles of the estimated redshift distribution for each individual galaxy (Brunner et al. , 1999a). As expected the average estimated error is the largest at the upper and lower redshift limits where the incompleteness in the calibrating sample is most evident. The majority of the rest of the objects with extremely large redshift errors are blended in one or more bands, which causes these objects to be isolated from the high density surface delineated by the majority of galaxies in the four flux space $`U,B,R,\&I`$. The effect these objects impart on any subsequent analysis, however, is minimized by the inclusion of their photometric error, which causes them to be non-localized in redshift space. As a result, these objects provide a minimal contribution to many “redshift bins” rather than strongly biasing only a few bins.
A subtle, and often overlooked, effect in any photometric redshift analysis is the requirement for reliable photometry in all program filters. Ideally we would obtain accurate redshift estimates for all galaxies; however, since we need accurate, multi-band photometry in order to reliably estimate redshifts, we must place restrictions on the photometric catalog used in the analysis. In particular, we restrict the full sample of detected sources to those objects which have both $`I_{AB}<24.0`$ and measured magnitude errors $`<0.25`$ in $`U,B,\&R`$. This minimizes any selection bias to only faint early-type galaxies at high redshifts, or very high redshift drop-out objects (see Brunner et al. 1999a for more discussion), neither of which significantly affect the rest of our analysis.
## 3 Ensemble Approach
Due to the lower precision of photometric redshift determination as compared to spectroscopic redshifts (roughly a factor of 20), we have developed a new, statistical approach to quantifying the evolution of galaxies (Brunner, 1997; Connolly et al. , 1998; Brunner et al. , 1999a). This approach capitalizes on our ability to reliably generate not only a redshift estimate for a galaxy using broadband photometry, but also a reliable redshift error estimate. As a result, we define the probability density function, $`P(z)`$, for an individual galaxy’s redshift to be a Gaussian probability distribution function with mean ($`\mu `$) given by the estimated photometric redshift and standard deviation ($`\sigma `$) defined by the estimated error in the photometric redshift.
$$P(z)=\frac{1}{\sigma \sqrt{2\pi }}e^{\left(\frac{(z\mu )^2}{2\sigma ^2}\right)}$$
In order to measure an interesting cosmological quantity, we generate multiple ensembles (or realizations) of the relevant properties of the galaxy distribution using the statistical redshift and redshift error estimates (see Brunner et al. 1999a for an application of this technique to the number-redshift distribution). The quantity of interest is determined as the mean of the multiple realizations, and the associated error is given by the corresponding standard deviation.
In order to divide our sample by intrinsic luminosity, we determined the absolute magnitude distribution of the galaxies in our catalog in an ensemble approach. First, we created different realizations of our galaxy catalog. In order to minimize any systematic errors, we selected the apparent magnitude of each galaxy from a Gaussian probability distribution function with mean and sigma given from the original photometric catalog measurements. Similarly, the redshift of each galaxy was drawn from a separate Gaussian probability distribution function with mean and sigma given from the photometric redshift and corresponding redshift error estimate. The $`k`$correction was determined using the spectral classification which was part of the original redshift estimation procedure. For galaxies with large photometric redshifts, occasional discordant redshifts were calculated (i.e. outside the range of our calibration sample — $`z<0`$, or $`z>1.2`$) in which case the galaxy was dropped from that particular realization.
Together, these quantities were used to determine the absolute magnitude for each galaxy in 100 different ensemble distributions. The absolute magnitude for each galaxy was calculated as the mean over the different realizations, appropriately normalized to account for possible discordant redshifts as discussed above. The resultant distributions for the $`U`$ and $`B`$ bands are displayed in Figure Evolution in the Clustering of Galaxies for $`z<1.0`$.
## 4 Analysis
Before computing the angular correlation function, we quantified our efficiency in detecting galaxies as a function of pixel location. The primary areas where this effect is important are around bright stars, in charge transfer trails, and near the edge of the frame due to edge effects or focus degradations. We, therefore, defined bounding boxes, for each of the four stacked images, which contained all of the observable flux for the saturated stars within the image. In the end, a total of 15 regions were masked out in the frame, 17 regions were masked out in the frame, 45 regions were masked out in the frame, and 36 regions were masked out in the frame. We also masked both the edge and corners of each frame in order to reduce the effects of PSF variations on our object detection efficiency. These four mask files were concatenated to produce a total mask file which was used for the calculation of the angular correlation function in different redshift or absolute magnitude intervals.
We used the optimal estimator Landy & Szalay (1993) $`(DD2DR+RR)/RR`$, where $`D`$ stands for data and $`R`$ stands for random, to determine the angular correlation function. This required counting the number of observed pairs (that were not within masked areas), which was done in 10 bins of constant width $`\mathrm{\Delta }\mathrm{lg}(\theta )=0.25`$, centered at $`\theta =4.3\mathrm{}`$, to $`\theta =759.6\mathrm{}`$. One thousand objects were randomly placed in the non masked areas within the image, and the data-random (DR) and random-random (RR) correlation functions were calculated for the same angular bins used for the data-data (DD) auto-correlation function. Before applying the estimator, each of the correlation measurements were scaled by the appropriate number of pairs. This process was repeated ten separate times, and the results averaged to minimize any systematic effects.
This estimator uses the calculated number density of galaxies within the CCD frame to estimate the true mean density of galaxies. The small angular area of our images introduces a bias in the estimate, commonly referred to as the “Integral Constraint” (Peebles, 1980). We estimated a correction for this bias following the prescription of Landy & Szalay (1993), which is subtracted from the estimated value.
The error in the estimation of the angular correlation function was assumed to be Poisson in nature. As a result, we calculated the error in the angular correlation estimation as the square root of the number of random-random pairs in each angular bin, scaled by the relevant number of data points.
The angular correlation function is generally parameterized in the following fashion:
$$w(\theta )=A_w\theta ^\delta $$
where the exponent has previously been shown to be $`\delta =0.8`$. In general, the amplitude of the correlation function ($`A_w`$) is calculated (assuming the previous value for $`\delta `$) by minimizing the $`\chi ^2`$ with respect to $`A_w`$ (Press et al. , 1992):
$$A_w=\left(\underset{i=1}{\overset{N}{}}\mathrm{lg}(w(\theta )_i)+0.8\underset{i=1}{\overset{N}{}}\mathrm{lg}(\theta _i)\right)/N$$
This calculation can be significantly affected by outliers; and, as a result, we adopt the more robust technique of minimization of the absolute deviations to determine the angular correlation amplitude. For our simple case, this technique reduces to finding the median of the amplitudes at each angle (Press et al. , 1992).
### 4.1 The Angular Correlation Function: $`w(\theta )`$
Although not the primary aim of this paper, we measured the angular correlation function in different apparent magnitude intervals to compare with previous work. For each program filter, the absolute upper magnitude limit for the data used in the estimation of the angular correlation function was set at the $`90\%`$ catalog completeness limit and the lower magnitude limit was always set to $`15^m`$. The angular correlation function was then determined in four different magnitude ranges (each offset from the next by $`0.5^m`$). The upper magnitude limits, and the corresponding number of objects are listed in Table 1. Each estimation was repeated 100 times.
The amplitude of the different correlation functions for each band at $`\theta =1.0\mathrm{}`$ are listed in Table 2. For brevity, we only display the results for the B band. Thus, in Figure Evolution in the Clustering of Galaxies for $`z<1.0`$ the actual correlation measurements for the different magnitude bins are displayed, while in Figure Evolution in the Clustering of Galaxies for $`z<1.0`$ the measured correlation amplitudes are compared to comparable published results (e.g., Koo & Szalay 1984; Roche et al. 1993; Infante & Pritchet 1995), showing remarkably good agreement. The error bars were calculated from the one sigma upper and lower measurements of the amplitude of the angular correlation function.
### 4.2 The Multi-Variate Angular Correlation Function: $`w(\theta ,z)`$
Using the 3052 objects in the photometric redshift–template SED catalog, the multivariate angular correlation function $`w(\theta ,z_P)`$ was determined for four different redshifts by binning the data in non-overlapping redshift bins of width $`\mathrm{\Delta }z_P=0.2`$ centered at $`z_P=0.4`$ to $`z_P=1.0`$. The four different functions were calculated for the 445, 573, 946, and 582 objects in the different respective redshift bins. We display, in Figure Evolution in the Clustering of Galaxies for $`z<1.0`$, the measured change in the amplitude of the angular correlation function with redshift (i.e. $`A_w`$), and hence the strength of the angular correlation function at a fixed angular separation. The error in $`A_w(z)`$ was calculated by estimating the amplitude in each redshift interval for the one sigma upper and lower values.
In order to compare this result with theoretical expectations (e.g., semi-analytic theory), which are determined in spatial coordinates, it is necessary to convert between angular correlation measurements and spatial correlation quantities. The standard technique for determining this transformation is to assume a power law model for the spatial clustering (Peebles, 1980; Le Fèvre et al. , 1996),
$$\xi (r,z)=\left(\frac{r}{r_0(z)}\right)^\gamma =\left(\frac{r}{r_0(0)}\right)^\gamma (1+z)^{(3+ϵ)}$$
where $`ϵ`$ represents a parameterization of the evolution of the spatial correlation function, and the correlation length is measured in physical units. The conversion, for small angular separations, between angular and spatial coordinates is accomplished via the relativistic form of Limber’s equation (Limber, 1954; Peebles, 1980). Specifically, this results in the following conversion for the amplitude of the angular correlation function within the redshift region of interest:
$$A_w=\left(\frac{_0^{\mathrm{}}G(z)B(\gamma )(\frac{dN}{dz})^2𝑑z}{\left[_0^{\mathrm{}}\frac{dN}{dz}𝑑z\right]^2}\right)r_0^\gamma $$
(1)
where,
$$G(z)=(1+z)^{(3+ϵ)}\sqrt{1+\mathrm{\Omega }_0z}x(z)^{(1\gamma )},$$
$$B(\gamma )=\left(\frac{3600.0180.0}{\pi }\right)^{(\gamma 1)}\left(\frac{H_0}{c}\right)^\gamma \frac{\mathrm{\Gamma }(\frac{1}{2})\mathrm{\Gamma }(\frac{(\gamma 1)}{2})}{\mathrm{\Gamma }(\frac{\gamma }{2})}$$
is a constant quantity,
$$x(z)=2\frac{((\mathrm{\Omega }_02)(\sqrt{1+\mathrm{\Omega }_0z}+2\mathrm{\Omega }_0+\mathrm{\Omega }_0z)}{\mathrm{\Omega }_0^2(1+z)^2},$$
is the angular diameter distance (Weinberg, 1972), and $`dN/dz`$ is the number of galaxies per unit redshift. Local spectroscopic surveys have determined that $`\gamma 1.8`$ and $`r_05.0`$ Mpc in co-moving coordinates (Carlberg et al. , 1999).
By assuming a uniform redshift distribution, we have calculated the tracks of different evolutionary models for $`\mathrm{\Omega }_0=1.0`$ and $`\mathrm{\Omega }_0=0.1`$, which are displayed, along with the measured values of the amplitude of the angular correlation function in Figure Evolution in the Clustering of Galaxies for $`z<1.0`$. Of the three different scenarios, fixed clustering in co-moving coordinates ($`ϵ=1.2`$) are the least consistent with our data, independent of the value of $`\mathrm{\Omega }_0`$. The results for clustering fixed in proper coordinates ($`ϵ=0.0`$) are good, independent of the value of $`\mathrm{\Omega }_0`$, while the predictions of linear theory ($`ϵ=0.8`$) are in agreement for higher values of $`\mathrm{\Omega }_0`$.
While this result is interesting, we can take an additional step in order to directly compare our measurements with published results from spectroscopic surveys. Using Equation 1 and an ensemble averaged empirical redshift distribution, we can actually transform our measurement of the angular correlation amplitude into a determination of the spatial correlation scale length within a given redshift interval, since we are able to empirically compute our observed redshift distribution (see Brunner et al. 1999a for more details). Essentially, this only involves adding a top-hat window function (corresponding to the appropriate redshift interval) to the integrands in Equation 1, since the transformation will be applied individually to each redshift bin in which we calculated the angular correlation amplitude.
The process is complicated, however, by the fact that our $`dN/dz`$ is determined using photometric redshifts, while the transformation assumes a spectroscopic redshift interval. From Figure Evolution in the Clustering of Galaxies for $`z<1.0`$, although small, the dispersion between spectroscopic redshifts and photometric redshifts is not zero. The conversion between a spectroscopic interval and a photometric redshift interval accounts to a broadening of the top-hat function, which we accomplish using two Gaussians centered on the endpoints of the redshift interval, so that the new window function is given by
$$W(z)=\{\begin{array}{cc}e^{\frac{(zz_1)^2}{\sigma ^2}}\hfill & 0z<z_1\hfill \\ 1\hfill & z_1zz_2\hfill \\ e^{\frac{(zz_2)^2}{\sigma ^2}}\hfill & z_1<z<\mathrm{}\hfill \end{array},$$
where the desired redshift interval is given by $`z[z_1,z_2]`$, and $`\sigma =0.061`$ is the measured dispersion in the photometric redshift relation.
The results of the transformation determine the correlation length within the given redshift bin (i.e. $`r_0(z)`$), and are displayed in Figure Evolution in the Clustering of Galaxies for $`z<1.0`$ and also tabulated in Table 3 for $`\mathrm{\Omega }_0=0.1`$, and the three canonical values of $`ϵ`$. In addition, the values of the correlation length extrapolated to $`z=0`$ (i.e. $`r_0(0)`$) are tabulated in Table 4 for three different values of $`\mathrm{\Omega }_0`$. The strong dependence of the correlation scale length on the value of the evolutionary parameter ($`ϵ`$) is a direct result of the Cosmological term ($`G(z)`$) in Equation 1. Relative to the predictions of linear theory ($`ϵ=0.8`$), the results for fixed clustering in co-moving coordinates ($`ϵ=1.2`$), are suppressed by an additional factor of $`(1+z)`$, or roughly a factor of $`1.5`$$`2`$ at the redshifts of interest (recall that $`r_0(z)G(z)^\gamma `$).
When comparing our results to previous spectroscopic results, we clearly show excellent agreement with recent measurements (Carlberg et al. , 1999; Small et al. , 1999) when the predictions of linear theory are used to quantify the evolution of clustering. This is extremely encouraging for our technique, as we uniformly sample a larger redshift range, showing, for the first time within the same dataset, the slight decrease in the correlation strength for $`z<1`$, as predicted by semi-analytic models of galaxy formation (Baugh et al., 1999). On the other hand, the measurement of the correlation scale length from the CFRS data (Le Fèvre et al. , 1996) are only in agreement with our results for fixed clustering in co-moving coordinates, which disagrees with the hierarchical growth of dark matter halos (Baugh et al., 1999). The discrepancy between the CFRS and other measurements is most likely due to their relatively small sample size, their small fields, and their neglect of the redshift evolution of the Luminosity function (Carlberg et al. , 1999).
### 4.3 The Multi-Variate Angular Correlation Function: $`w(\theta ,z,M)`$
While important, the evolution of the angular correlation function with redshift smoothes over the galaxy luminosity function, ignoring variations in clustering between galaxies of different intrinsic luminosity. As a result, we subdivided our sample into three redshift intervals ($`0.2z0.6`$, $`0.4z0.8`$, $`0.6z1.0`$), and measured the angular correlation function as a function of both $`U`$ and $`B`$ absolute magnitude in intervals of $`2.0^m`$, from $`22^m`$ to $`12^m`$. In order to improve the number of galaxies in the faint end of our analysis, we rebinned the data so that the faint bin was four magnitudes wide (i.e. $`16^m`$$`12^m`$). In the end, we obtained twelve different measurements of the multivariate angular correlation function ($`w(\theta ,z,M)`$) as a function of both $`U`$ and $`B`$ absolute magnitudes.
The number of subdivisions used in this particular analysis reintroduced one of the problems our new technique was designed to avoid, namely the effects of small sample size. We, therefore, only used a joint redshift-absolute magnitude bin when the number of objects in the bin exceeded one hundred galaxies. From Figures Evolution in the Clustering of Galaxies for $`z<1.0`$Evolution in the Clustering of Galaxies for $`z<1.0`$ ($`U`$ and $`B`$ respectively), it is clear that there is no obvious evolution within a given redshift interval with intrinsic luminosity (although this could be a result of too few galaxies). Between redshift intervals, however, there is strong evolution in the amplitude of the angular correlation function ($`A_w`$), which, given the wider redshift bins, is completely consistent with the results of the previous section.
In order to quantify this evolution, we fit a line of zero slope to the points (i.e. a mean value for each redshift interval). Between the first two redshift bins (roughly $`z0.4`$ and $`z0.6`$), $`A_w`$ drops by approximately a factor of two. For the $`U`$ band measurement, the drop between the second and third redshift intervals (roughly $`z0.6`$ and $`z0.8`$) is approximately $`20\%`$, while the $`B`$ band shows another factor of approximately two decline. This result is not surprising in the context of hierarchical structure formation where one expects objects to be more strongly cluster with decreasing redshift for $`z<1`$ (e.g.Kauffmann et al. 1999).
Unfortunately, we do not have enough galaxies to definitely test for clustering evolution with redshift for a given population with a fixed intrinsic luminosity. On the other hand, our results do seem to indicate that the evolution is not strongly dependent on intrinsic luminosity as there appears to be no clear evidence for variation in the clustering amplitude within a given redshift interval. The difference in the drop in the amplitude between the middle and high redshift intervals is most likely due to the lower, average intrinsic luminosity of the galaxies in the $`B`$ band as compared to the $`U`$ band. These points will need to be addressed with future datasets.
## 5 Conclusions
In this paper, we have presented several calculations of the angular correlation function, as a function of different apparent magnitude intervals, as a function of redshift, and as a function of both redshift and absolute magnitude. The technique we have demonstrated is less sensitive to redshift distortions than the spatial correlation approach due to the width of our redshift bins. Furthermore, our technique does not require model predictions for the redshift distribution of galaxies as does the apparent magnitude interval approach. Future work in this area will incorporate spectroscopic redshifts into the calculation in order to provide limited distance information (cf. Phillipps 1985), as well as witness the application of these techniques to larger surveys.
While not the main point of this paper, the variation of the amplitude of the angular correlation with apparent magnitude is in good agreement with previously published results, which strengthens the rest of our analysis. Furthermore, we demonstrated, for the first time from within a single dataset, the slight evolution in both the amplitude of the angular correlation function, and the spatial correlation scale length with redshift for $`z<1`$, as predicted by semi-analytic models of structure formation (Baugh et al., 1999; Kauffmann et al. , 1999). These results suggest low values for $`\mathrm{\Omega }_0`$, and allow either fixed clustering in proper coordinates, or the predictions of linear theory.
Finally, we measured the evolution of the amplitude of the angular correlation function with both redshift and intrinsic luminosity. The amplitude of the angular correlation function drops dramatically with redshift. Interestingly enough, however, we do not see significant variation in the strength of clustering within a given redshift interval as a function of intrinsic luminosity. This type of evolution might be naively expected if the luminosity of a galaxy uniquely mapped to the mass of the dark matter halo at the given redshift in which it resides (i.e. more luminous galaxies cluster more strongly within a given redshift interval). Most likely, either we have too small of a sample to place significant limits on the variation of clustering with luminosity, or else the relatively constant clustering amplitude as a function of luminosity is indicative of luminosity evolution complicating the analysis.
Future surveys, both photometric and spectroscopic (e.g., SDSS) will provide extremely useful datasets with which we can explore these ideas in greater detail. In the near future, this area will witness a merging of observations, semi-analytic theory, and N-body simulations, finally providing hope that we will be able to unambiguously quantify the clustering evolution of galaxies.
First we wish to acknowledge Gyula Szokoly for assistance in obtaining the data. We also would like to thank Barry Lasker, Gretchen Greene, and Brian McLean for allowing us access to an early version of the GSC II. We also wish to acknowledge useful discussions with Pat Cote, Rich Kron, Lori Lubin, and Ray Weymann. We thank the anonymous referee for valuable suggestions on improving this work. This research has made use of NASA’s Astrophysical Data System Abstract Service. AS acknowledges support from NASA LTSA (NAG53503) and HST Grant (GO-07817-04-96A), AJC acknowledges partial support from HST (GO-07817-02-96A) and LTSA (NRA-98-03-LTSA-039).
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# Twisted Family Structure and Neutrino Large Mixing
## Abstract
I demonstrate that neutrino large mixing between $`\nu _\mu `$ and $`\nu _\tau `$ are naturally reproduced using a novel mechanism called ‘E-twisting’in a supersymmetric $`E_6`$ grand unification model. This model explains all the characteristic features of the quark/lepton Dirac masses as well as the neutrino’s Majorana masses despite the fact that all the members in 27 of each generation are assigned a common family charge. Most remarkably, this model yields a novel relation which gives the 2-3 lepton mixing angle $`\theta _{\mu \tau }`$ in terms of quark masses and CKM mixing: $`\mathrm{tan}\theta _{\mu \tau }=(m_b/m_s)V_{cb}`$, which is a kind of $`SO(10)`$ GUT relation similar to the celebrated $`SU(5)`$ bottom-tau mass ratio. This relation is a result of a common ‘twisted $`SO(10)`$’ structure <sup>1</sup><sup>1</sup>1Talk given at the Internatinal Workshop on Neutrino Oscillation and their Origin, Fujiyoshida, Japan, February 11-13, 2000.
A remarkable fact observed in SuperKamiokande is the very large lepton mixing, which is in a sharp contrast to the quark sector where the CKM mixings are all small. Why can such a large difference occurs between the quark and lepton sectors? This is a challenging question for any particle physicist who tries to find grand unified theories (GUTs). Clearly any GUT which treats the three families of quarks and leptons as a mere repetition no longer works. We need some new mechanism of family structure. Lots of proposals have been made on the origin of this large lepton mixing angles . On the other hand, the SK results indicate larger unification groups than $`SU(5)`$ including left-right symmetry, in which large neutrino mixings seem unnatural, since in such larger GUT groups, for example, $`SO(10)`$ GUT, all the fermions of a family are combined into a single representation, and the most natural prediction would be that the neutrino mixing is also very small with hierarchical masses.
Recently we have constructed a supersymmetric $`E_6`$ unified model with an extra $`U(1)`$ family charge. There we have shown that E-twisted family structure can reproduce all the characteristic features of the fermion masses, not only the quark/lepton Dirac masses but also the neutrino Majorana masses. Despite the fact that a common $`U(1)`$ charge is assigned to all the members in a fundamental representation 27 of $`E_6`$ for each family, the model well explains all the qualitative feature of different mass hierarchies among families and between up and down quark sectors, as well as the mixing angles. In this scenario we have found that in the framework of a supersymmetric $`E_6`$ grand unified model , the twisted family structure yields a novel relation
$$\mathrm{tan}\theta _{\mu \tau }=\frac{m_b}{m_s}V_{cb},$$
(1)
which I shall compare with the experiments later. Leaving the details in the papers , I here explain the essence of the model. In the supersymmetric $`E_6`$ model, in addition to $`E_6`$ gauge vector multiplet, we introduce chiral matter multiplets corresponding to the three families, ($`\mathrm{\Psi }_i`$ $`(i=1,2,3)`$) and a pair of Higgs fields, which is introduced mainly for the electroweak symmetry breaking,($`H`$, $`\overline{H}`$) <sup>2</sup><sup>2</sup>2In order to give all the unwanted fermions to get heavy masses. we need another Higgs pair, ($`\mathrm{\Phi }`$, $`\overline{\mathrm{\Phi }}`$) and which are responsible for realizing the E-twisted family structure. Also we have to add a chiral Higgs multiplet $`\varphi (\mathrm{𝟕𝟖})`$ in order to break the GUT to the standard gauge group. Here we neglect those fields and start with the low energy fermions realized in twisted family structure.. In table 1, we summarize all the fields we need in our model.
The $`E_6`$ singlet field $`\mathrm{\Theta }`$ with $`U(1)`$ charge $`1`$ plays an important role that its suitable powers compensate the mismatch of the $`U(1)`$ charge in the superpotential interaction terms. The $`U(1)`$ flavor symmetry discriminates families and induces hierarchy between them. Note that all the quarks and leptons in one generation have a common $`U(1)`$ quantum number.
The following Yukawa superpotentials which are invariant under $`R`$ parity, $`U(1)`$ and $`E_6`$ will give masses of matter superfields $`\mathrm{\Psi }_i(\mathrm{𝟐𝟕})`$ <sup>3</sup><sup>3</sup>3 There are other superpotentials including the Higgs fields, $`\varphi (\mathrm{𝟕𝟖})`$ and $`\mathrm{\Phi }(\mathrm{𝟐𝟕})`$, whose family charges are not zero and will contribute to the Yukawa terms of the 2nd and 1st families.,
$$W_Y(H)=y_{ij}\mathrm{\Psi }_i(\mathrm{𝟐𝟕})\mathrm{\Psi }_j(\mathrm{𝟐𝟕})H(\mathrm{𝟐𝟕})\left(\frac{\mathrm{\Theta }}{M_P}\right)^{f_i+f_j},$$
(2)
where $`f_i`$ denotes the $`U(1)`$ charge of $`i`$-th family. With the coupling constants $`y`$ of order 1, the effective Yukawa coupling constants are associated with additional powers of $`\lambda =\mathrm{\Theta }/M_P`$, which we assume is of the order of the Cabibbo angle $`\lambda 0.22`$. We also suppose that only the $`SU(2)`$ doublet components of $`H`$ can have the electroweak scale vacuum expectation value (VEV).
An interesting fact is that there are two $`\mathrm{𝟓}^{}`$’s of $`SU(5)`$ in each $`\mathrm{𝟐𝟕}`$, i.e., $`\mathrm{𝟓}^{}`$ of $`\mathrm{𝟏𝟔}`$ of $`SO(10)`$ (($`\mathrm{𝟏𝟔},\mathrm{𝟓}^{}`$)) and $`\mathrm{𝟓}^{}`$ of $`\mathrm{𝟏𝟎}`$ (($`\mathrm{𝟏𝟎},\mathrm{𝟓}^{}`$)). Those may be called ‘E-parity’ doublet. It is this doubling that we have a freedom to choose the low-energy $`\mathrm{𝟓}^{}`$ candidates. This actually implies that the embedding of $`SO(10)`$ into $`E_6`$ such that $`SU(5)_{\mathrm{GG}}SO(10)E_6`$ with Georgi-Glashow $`SU(5)_{\mathrm{GG}}`$, possesses a freedom of rotation of $`SU(2)_R`$. The doubling of $`\mathrm{𝟓}^{}`$’s in each $`\mathrm{𝟐𝟕}`$ also provides the low-energy surviving down-type Higgs field with the freedom of mixing parameter between two $`\mathrm{𝟓}^{}`$’s in $`H(\mathrm{𝟐𝟕})`$:
$`H(\mathrm{𝟓}^{})`$ $`=`$ $`H(\mathrm{𝟏𝟎},\mathrm{𝟓}^{})\mathrm{cos}\theta +H(\mathrm{𝟏𝟔},\mathrm{𝟓}^{})\mathrm{sin}\theta .`$ (3)
Now we pick up the low-energy matter fields among the three $`\mathrm{\Psi }_i(\mathrm{𝟐𝟕})`$ of the above. The up-quark sector is unique since $`\mathrm{𝟏𝟎}`$ and $`\mathrm{𝟓}`$ of $`SU(5)`$ appear only once in each $`\mathrm{𝟐𝟕}`$. As for the down quarks, there is a freedom for choosing three from six $`\mathrm{𝟓}^{}`$s in three $`\mathrm{\Psi }_i(\mathrm{𝟐𝟕})`$s. We have classified possible typical scenarios in Ref. ; (i) Parallel family structure, (ii) Non-parallel family structure, and (iii) E-twisted structure. Among these three possibilities, we here take the simplest and most attractive option, namely the E-twisted family:
$`(\mathrm{𝟓}_1^{},\mathbf{5}_2^{},\mathbf{5}_3^{})`$ $`=`$ $`(\mathrm{\Psi }_1(\mathrm{𝟏𝟔},\mathrm{𝟓}^{}),\mathrm{\Psi }_2(\mathrm{𝟏𝟔},\mathrm{𝟓}^{}),\mathrm{\Psi }_3(\mathrm{𝟏𝟎},\mathrm{𝟓}^{})).`$ (4)
This structure implies that the third family $`\mathrm{𝟓}^{}`$ belongs to $`\mathrm{𝟏𝟎}`$ of $`SO(10)`$. This twisting is realized by the suitable VEVs of the Higgs fields( the details will be found in Ref. ).
Let us here concentrate ourselves on the 2nd and 3rd families and see what happens to their masses and mixings. The $`2\times 2`$ mass matirces for the up-quark, down-quark and charged-lepton, $`M_u`$, $`M_d`$ and $`M_e`$, are expressed as,
$$M_u=\begin{array}{ccccc}& & u_2^c& u_3^c& \\ u_2& (\mathrm{}& y_{22}& y_{23}& )\mathrm{}\\ u_3& y_{32}& y_{33}\end{array}v\mathrm{sin}\beta =\begin{array}{ccccc}& & u_2^c& u_3^c& \\ u_2& (\mathrm{}& & f\lambda ^2& )\mathrm{}\\ u_3& & 1\end{array}vy_{33}\mathrm{sin}\beta ,$$
(5)
$$M_d=\begin{array}{ccccc}& & d_2^c& D_3^c& \\ d_2& (\mathrm{}& z_{22}^d\mathrm{cos}\theta & y_{23}\mathrm{sin}\theta & )\mathrm{}\\ d_3& z_{32}\mathrm{cos}\theta & y_{33}\mathrm{sin}\theta \end{array}v\mathrm{cos}\beta =\begin{array}{ccccc}& & d_2^c& D_3^c& \\ d_2& (\mathrm{}& e\lambda ^2& f\lambda ^2& )\mathrm{}\\ d_3& h& 1\end{array}vy_{33}\mathrm{sin}\theta \mathrm{cos}\beta .$$
(6)
$$M_e^\mathrm{T}=\begin{array}{ccccc}& & e_2^{}& E_3& \\ e_2^c& (\mathrm{}& z_{22}^e\mathrm{cos}\theta & y_{23}^e\mathrm{sin}\theta & )\mathrm{}\\ e_3^c& z_{32}q\mathrm{cos}\theta & y_{33}\mathrm{sin}\theta \end{array}v\mathrm{cos}\beta =\begin{array}{ccccc}& & e_2^{}& E_3& \\ e_2^c& (\mathrm{}& & & )\mathrm{}\\ e_3^c& h& 1\end{array}vy_{33}\mathrm{sin}\theta \mathrm{cos}\beta .$$
(7)
where $`\mathrm{tan}\beta `$ is the mixing angle of two light Higgs doublets and $`v`$ is the VEV of the standard Higgs field $`H(\mathrm{𝟐𝟕})`$. We rewrite by using simple notations in the third terms with $``$ being irrelevant for our present discussions <sup>4</sup><sup>4</sup>4We have assumed that the main contribution comes only from the Higgs field $`H(\mathrm{𝟐𝟕})`$ at least for the quark mass matrix of $`33`$ and $`23`$ elements in the quark mass matrices, $`M_u`$ and $`M_d`$.. Note that by taking the mixing $`\mathrm{sin}\theta `$ of the Higgs field $`H(\mathrm{𝟐𝟕})`$ of Eq.(3), to be of order $`\lambda ^2`$, $`h`$ becomes of order $`1`$. It is easy to obtain $`h`$ from the bottom and strange quark masses and mixing angle, $`m_b,m_s,V_{cb}`$ from Eqs.(5) and ((6)). Noting that $`h`$ gives directly the lepton mixing angle $`\mathrm{tan}\theta _{\mu \tau }=h`$ <sup>5</sup><sup>5</sup>5We can confirm that the right handed Majorana mass term indicates very small mixing and neutrino mixing mainly comes from lepton mixing., we can get the novel relation Eq.(1), or equivalently,
$`\mathrm{sin}^22\theta _{\mu \tau }`$ $`=`$ $`{\displaystyle \frac{4V_{cb}^2\left({\displaystyle \frac{m_s}{m_b}}\right)^2}{\left[V_{cb}^2+\left({\displaystyle \frac{m_s}{m_b}}\right)^2\right]^2}}.`$ (8)
By taking the experimental value of $`x=V_{cb}m_b/m_S`$, $`1x1.68`$, we can calculate the left hand side of Eq.(8), namely,
$$0.78\mathrm{sin}^22\theta _{\mu \tau }1.$$
(9)
which is remarkably in good agreement with the large lepton mixing recently observed.
This relation can be obtained from more general framework using a kind of $`SO(10)`$ GUT and may be called the second q-l relation, similar to the first q-l relation, i.e., the celebrated bottom-tau mass ratio of $`SU(5)`$. It is interesting that this relation can be most easily obtained from the twisted $`E_6)`$ model. We would like to remark that our $`E_6`$ twisted model can also explain why the bottom quark mass is smaller by almost $`\lambda ^2`$ than that of the top quark. This also comes from the Higgs mixing factor $`\mathrm{sin}\theta `$ in $`M_d`$.
To conclude, we have found that the twisted $`E_6`$ model can explain the up-down mystery ($`m_tm_b`$), as well as the down-lepton mystery ($`\theta _{\mu \tau }\theta _{cb}`$). It is well known that $`E_6`$ gauge symmetry is naturally obtained from the 10 dimensional $`E_8\times E_8`$ heterotic string theory by the Calabi-Yau compactification into 4 dimensions. Our results are interesting and encouraging and may open the door for finding out more fundamental stringy GUTs including gravity. This report is based on the works in collaboration with T. Kugo, K. Yoshioka. I would like to thank to T. Yanagida, Y. Nomura, M. Yamaguchi and many other members for their stimulating discussions. This work has been done during the Summer Institute 98 and 99 held at Yamanashi, Japan organized by T. Kugo and T. Eguchi. I am supported in part by the Grants-in-Aid for Scientific Research No. 9161 from the Ministry of Education, Science, Sports and Culture, Japan.
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# TIFR/TH/99/52, IMSc -99/10/35 Three-Manifold Invariants from Chern-Simons Field Theory with Arbitrary Semi-Simple Gauge Groups
## 1 Introduction
In recent times topological quantum theories have proved to be a very powerful tool for the study of geometry and topology of low dimensional manifolds. An example of such theories is the Chern-Simons gauge field theory which provides a general framework for knots and links in three dimensions . Vacuum expectation values of Wilson link operators in this theory yield a class of polynomial link invariants. It was E. Witten who in his pioneering paper about ten years ago developed this framework and also demonstrated that the famous Jones polynomial was related to the expectation value of a Wilson loop operator (in spin $`1/2`$ representation) in an $`SU(2)`$ Chern-Simons field theory. Since then many attempts have been made to obtain exact and explicit non-perturbative solutions to such field theories . Powerful methods for completely analytical and non-perturbative computations of the expectation values of Wilson link operators have been developed. One such method in its complete manifestation has been presented in ref..
The power of field theoretic framework through Chern-Simons theories is indeed so deep that it allows us to study knots and links not only in simple manifolds such as a three-sphere but also in any arbitrary three-manifold. For example, the link invariants obtained in these field theories can be used to construct three-manifold invariants. One such construction involves an application of Lickorish-Wallace surgery presentation of three-manifolds in terms of unoriented framed links embedded in $`S^3`$. Surgery on more than one framed knot or link can yield the same manifold. However, the rules of equivalence of framed links which yield the same three-manifold on surgery are given by theorems of Kirby, Fenn and Rourke and are known as Kirby moves. Thus a three-manifold can be characterized by an appropriate combination of invariants of the associated framed knots and links which is unchanged under Kirby moves. A three-manifold invariant of this type has been recently constructed from invariants for framed links in an $`SU(2)`$ Chern-Simons theory in a three-sphere . The algebraic formula for the invariant so obtained is rather easy to compute for an arbitrary three-manifold. The construction developed is general enough to yield other three-manifold invariants from the link invariants of Chern-Simons gauge theories based on other semi-simple gauge groups. This extension is what will be presented in the present paper.
Other three-manifold invariants have also been constructed in recent years. For example, exploiting the surgery presentations of three-manifolds in terms of unoriented framed links, Lickorish had earlier obtained a manifold invariant using bracket polynomials of cables . Evaluation of this invariant involves a tedious calculation through recursion relations. Using representation theory of composite braids , it has been possible to directly evaluate the bracket polynomials for cables without going through the recursion relations. This direct calculation has been used to demonstrate the equivalence of the invariant obtained from $`SU(2)`$ Chern-Simons theory in ref. to the Lickorish’s three-manifold invariant up to a variable redefinition . Further Lickorish’s invariant is considered to be a reformulation of Reshetikhin-Turaev invariant, which in turn is known to be equivalent to the partition function of $`SU(2)`$ Chern-Simons theory, known as Witten invariant. Thus this establishes, by an indirect method, that the field theoretic three-manifold invariant obtained from link invariants in $`SU(2)`$ Chern-Simons gauge theory using theorems of Lickorish and Wallace, Kirby, Fenn and Rourke is actually partition function of $`SU(2)`$ Chern-Simons theory, a fact already noticed for many three-manifolds in ref.. This equivalence is up to an over all normalization.
The link invariants in general depend on the framing convention used. The frame of a knot is an associated closed curve going along the length of knot and wrapping around it certain number of times. In the field theoretic language, framing has to do with the regularization prescription used to define the coincident loop correlators . In one such framing convention known as standard framing the self-linking number of every knot (i.e., linking number of the knot and its framing curve) is zero. This convention was used in obtaining the link invariants in ref. . The invariants so obtained are ambient isotopic invariant, that is, these are unchanged under all the three Reidemeister moves. However, in our present discussion, we are interested in framed links which are only regular isotopic, that is, two framed links are equivalent if and only if they are related by two of the Reidmeister moves (excluding the one that changes the writhe). The framing convention for describing such framed links is vertical framing. Here the frame is to be just vertically above the strands of every knot projected on to a plane. The link invariants in this framing exhibit only regular isotopy invariance. These framed link invariants are in general sensitive to the relative orientations of component knots of a link. Reversing the orientation in a knot component changes the representation living on the associated Wilson loop operator to its conjugate representation. We construct an appropriate linear combination of these invariants for different group representations on the framed link which is unchanged under Kirby calculus. This combination then characterises the manifold related to the given link by surgery. Though the individual terms in this combination in general depend on the relative orientations of the knots, the combination does not. This is consistent with Lickorish-Wallace theorem for surgery presentation of three-manifolds which involves only unoriented links.
The plan of the paper is as follows: In the next section, we shall briefly discuss Chern-Simons theory based on any arbitrary semi-simple gauge group. Methods of computing the expectation value of Wilson loop operators for framed knots and links will be outlined. These are generalizations of the methods for $`SU(2)`$ Chern-Simons theory presented in ref. . In Sec.3, we shall present a theorem of Lickorish and Wallace and Kirby calculus which are the necessary ingredients in the construction of three-manifolds by surgery. Using the field theoretic framed link invariants, we derive an algebraic formula for a three-manifold invariant. Sec. 4 will contain some concluding remarks.
## 2 Chern-Simons field theory and link invariants
In this section, we shall present some of the salient features of Chern-Simons theory on $`S^3`$ based on arbitrary semi-simple gauge group $`𝒢`$ and the invariants of framed links embedded in $`S^3`$. These framed link invariants will be used in the construction of three-manifold invariants in the next section.
For a matrix valued connection one-form $`A`$ of the gauge group $`𝒢`$, the Chern-Simons action $`S`$ on $`S^3`$ is given by
$$kS=\frac{k}{4\pi }_{S^3}tr(AdA+\frac{2}{3}A^3).$$
(1)
The coupling constant $`k`$ takes integer values. Clearly action (1) does not have any metric of $`S^3`$ in it. The topological operators are the metric independent Wilson loop (knot) operators defined as
$$W_R[C]=tr_RPexp_CA_R$$
(2)
for a knot $`C`$ carrying representation $`R`$; $`A_R`$ is the connection field in representation $`R`$ of the group. Reversing the orientation of a knot corresponds to placing conjugate representation $`R^{}`$ on it.
For a link $`L`$ made up of component knots $`C_1,C_2,\mathrm{}C_r`$ carrying $`R_1,R_2,\mathrm{}R_r`$ representations respectively, we have the Wilson link operator defined as
$$W_{R_1R_2\mathrm{}R_r}[L]=\underset{\mathrm{}=1}{\overset{r}{}}W_R_{\mathrm{}}[C_{\mathrm{}}].$$
(3)
We are interested in the functional averages of these link operators:
$`V[L;𝐟;R_1,R_2\mathrm{}R_r]V_{R_1R_2\mathrm{}R_r}[D_L]`$ $`=`$ $`Z^1{\displaystyle _{S^3}}[dA]W_{R_1R_2\mathrm{}R_r}[L]e^{ikS},`$ (4)
$`\mathrm{where}Z={\displaystyle _{S^3}}[dA]e^{ikS},`$
and $`D_L`$ denotes link diagram corresponding to framed link $`[L,f]`$. Link diagrams are obtained from a regular projection of a link in a plane with transverse double points (refered to as crossings) as the only self-intersections. We work in the vertical frame where the frame curve is vertically above the plane of the diagram. Hence, in vertical framing, $`𝐟=(f_1,f_2,\mathrm{},f_r)`$ on link $`L`$ is a set of integers denoting the sum of the crossing signs in the part of the diagram representing the components. These expectation values are the generalized regular isotopy invariants of framed links. These can be exactly evaluated using the method developed in ref. . The method makes use of the following two inputs:
1. Chern-Simons functional integral (containing Wilson lines) on a three-manifold with $`n`$-punctures on its boundary corresponds to a state in the space of $`n`$-correlator conformal blocks in the corresponding Wess-Zumino conformal field theory on that boundary .
2. Knots and Links can be obtained by closure of braids (Alexander theorem) or equivalently platting of braids (Birman theorem) .
Consider a manifold $`S^3`$ from which two non-intersecting three-balls are removed. This manifold has two boundaries, each an $`S^2`$. We place $`2n`$ Wilson line-integrals over lines connecting these two boundaries through a weaving pattern $`𝐁`$ as shown in the Figure (a) below. This is a $`2n`$braid placed in this manifold. The strands are specified on the upper boundary by giving $`2n`$ assignments $`(\widehat{R}_1^{},\widehat{R}_1,\widehat{R}_2^{},\widehat{R}_2,\mathrm{},\widehat{R}_n^{},\widehat{R}_n)`$. Here $`\widehat{R}=(R,ϵ)`$ denotes representation $`R`$ and orientation $`ϵ`$ ($`ϵ=\pm 1`$ for a strand going into or away from the boundary) and conjugate assignment $`\widehat{R}^{}=(\overline{R},ϵ)`$ indicates reversal of the orientation. Similar specifications are done with respect to the lower boundary by the representation assignments $`(\widehat{\mathrm{}}_1,\widehat{\mathrm{}}_1^{},\widehat{\mathrm{}}_2,\widehat{\mathrm{}}_2^{},\mathrm{}.\widehat{\mathrm{}}_n,\widehat{\mathrm{}}_n^{})`$. Then the assignments $`\{\widehat{\mathrm{}}_i\}`$ are just a permutation of $`\{\widehat{R}_i^{}\}`$. Chern-Simons functional integral over this manifold is a state in the tensor product of the Hilbert spaces associated with the two boundaries, $`_1_2`$. This state can be expanded in terms of some convenient basis. These bases are given by the conformal blocks for $`2n`$-point correlators of the associated $`𝒢_k`$ Wess-Zumino conformal field theory on each of the $`S^2`$ boundaries.
An arbitrary braid can be generated by a sequence of elementary braidings. The eigenvalues of these elementary braids are given by conformal field theory. The braiding eigenvalues depend on the framing. In vertical framing, the eigenvalues for a right-handed half-twist between two parallely oriented strands carrying representation $`R,R^{}`$ ($`ϵϵ^{}=1`$) and for anti-parallely oriented strands ($`ϵϵ^{}=1`$), as shown in the figure below, are respectively :
$`\lambda _t^{(+)}(R,R^{})`$ $`=`$ $`()^{ϵ_R+ϵ_R^{}ϵ_t}q^{(C_R+C_R^{})/2+C_t/2}`$ (5)
$`\lambda _t^{()}(R,R^{})`$ $`=`$ $`()^{|ϵ_Rϵ_R^{}|ϵ_t}q^{(C_R+C_R^{})/2C_t/2},`$ (6)
where $`t`$ takes values allowed in the product of representations of $`R`$ and $`R^{}`$ given by the fusion rules of $`𝒢_k`$ Wess-Zumino conformal field theory, $`q`$-independent phases $`()^{ϵ_R+ϵ_R^{}ϵ_t}=\pm 1`$, $`()^{|ϵ_Rϵ_R^{}|ϵ_t}=\pm 1`$, and $`C_R,C_R^{}`$ are the quadratic Casimir of representations $`R,R^{}`$ respectively. The variable $`q`$ in the above equation is related to the coupling constant $`k`$ in Chern-Simons theory as
$$q=exp(\frac{2\pi i}{k+C_v}),$$
(7)
where $`C_v`$ is quadratic Casimir of adjoint representation.
As mentioned earlier, the link invariants in vertical framing are only regular isotopy invariants. That is, these invariants do not remain unchanged when a writhe is smoothed out, but instead pick up a phase:
$$\text{}=q^{C_R}\text{},$$
$$\mathrm{and}\text{}=q^{C_R}\text{},$$
where $`\pm `$ indicate the right-handedness and left-handedness of the writhe.
Writing the weaving pattern $`𝐁`$ in Figure (a) above in terms of elementary braids, the Chern-Simons functional integral over this manifold is given by a matrix $`𝐁(\{R_i\},\{\mathrm{}_i\})`$ in $`_1_2`$. To plat this braid, we consider two balls with Wilson lines as shown in Figures (b) and (c) above. We glue these respectively from above and below onto the manifold of Figure (a). This yields a link in $`S^3`$.
The Chern-Simons functional integral over the ball (c) is given by a vector in the Hilbert space associated with its $`S^2`$ boundary. This vector $`|\psi (\{\mathrm{}_i\})`$ can again be written in terms of a convenient basis of this Hilbert space. Similarly, the functional integral over the ball of Figure (b) above is a dual vector $`\psi (\{R_i\})|`$ in the associated dual Hilbert space. Gluing these two balls on to each other (along their oppositely oriented boundaries) gives $`n`$ disjoint unknots carrying representations $`R_1,R_2,\mathrm{}.R_n`$ in $`S^3`$. Their invariant is simply the product of invariants for $`n`$ individual unknots, due to the factorization property of invariants for disjoint knots. Thus the inner product of vectors representing the functional integrals over manifolds (b) and (c) is given by
$$\psi (\{R_i\})|\psi (\{R_i\})=\underset{i=1}{\overset{n}{}}V_{R_i}[U],$$
(8)
where the invariant for unknot carrying representation $`R_i`$ is given by $`V_{R_i}[U]=dim_qR_i`$ which is the quantum dimension for representation $`R_i`$ defined in terms of highest weight $`\mathrm{\Lambda }_{R_i}`$, Weyl vector $`\rho `$ and positive roots $`\alpha _+`$ :
$$\mathrm{dim}_qR_i=\underset{\alpha _{}\mathrm{\Delta }_+}{}\frac{[(\mathrm{\Lambda }_{R_i},\alpha _+)+(\rho ,\alpha _+)]}{[(\rho ,\alpha _+)]}=\frac{S_{0\mathrm{\Lambda }_{R_i}}}{S_{00}}.$$
(9)
Here the square brackets denote $`q`$number defined as:
$$[x]=\frac{(q^{\frac{x}{2}}q^{\frac{x}{2}})}{(q^{\frac{1}{2}}q^{\frac{1}{2}})},$$
(10)
with $`q`$ as given in eqn. (7). Matrix $`S`$ represents the generator of modular transformation $`\tau 1/\tau `$ on the characters of associated $`𝒢_k`$ Wess-Zumino conformal field theory. Its form for a group $`𝒢`$ of rank $`r`$ and dimension $`d`$ is given by :
$`S_{\mathrm{\Lambda }_{R_1}\mathrm{\Lambda }_{R_2}}=(i)^{\frac{dr}{2}}|{\displaystyle \frac{L_\omega }{L}}|^{\frac{1}{2}}\left(k+C_v\right)^{\frac{1}{2}}{\displaystyle \underset{\omega W}{}}ϵ(\omega )\mathrm{exp}\left({\displaystyle \frac{2\pi i}{k+C_v}}(\omega (\mathrm{\Lambda }_{R_1}+\rho ),\mathrm{\Lambda }_{R_2}+\rho )\right),`$ (11)
where $`W`$ denotes the Weyl group and its elements $`\omega `$ are words constructed using the generator $`s_{\alpha _i}`$ – that is, $`\omega =_is_{\alpha _i}`$ and $`ϵ(\omega )=(1)^{\mathrm{}(\omega )}`$ with $`\mathrm{}(\omega )`$ as length of the word. The action of the Weyl generator $`s_\alpha `$ on a weight $`\mathrm{\Lambda }_R`$ is
$$s_\alpha (\mathrm{\Lambda }_R)=\mathrm{\Lambda }_R2\alpha \frac{(\mathrm{\Lambda }_R,\alpha )}{(\alpha ,\alpha )},$$
(12)
and $`|L_\omega /L|`$ is the ratio of weight and co-root lattices (equal to the determinant of the cartan matrix for simply laced algebras).
Having determined the state corresponding to functional integrals over three-manifolds as drawn in Figs. (b), (c) above, we shall now obtain the invariant for a link obtained by gluing the two balls (b) and (c) on to the manifold of Figure (a). The link invariant is equal to the matrix element of matrix $`𝐁`$ between these two vectors. This can be calculated by generalizing the method of ref.() for arbitrary semi-simple groups through following proposition:
Proposition 1: Expectation value of a Wilson operator for an arbitrary $`n`$ component framed link $`[L,𝐟]`$ with a plat representation in terms of a braid $`𝐁(\{R_i\},\{\mathrm{}_i\})`$ generated as a word in terms of the braid generators is given by
$$V[L;𝐟;R_1,R_2,\mathrm{}R_n]=\psi (\{R_i\})|𝐁(\{R_i\},\{\mathrm{}_i\})|\psi (\{\mathrm{}_i\})$$
(13)
Thus these invariants for any arbitrary framed link can be evaluated.
Examples
a) Unknot with framing number +1 or -1 is related to the unknot in zero framing as:
$`=`$ $`q^{C_R}\text{}=q^{C_R}dim_qR`$ (14)
$`\mathrm{and}\text{}`$ $`=`$ $`q^{C_R}\text{}=q^{C_R}dim_qR.`$ (15)
b) The invariant for a Hopf link carrying representation $`R_1`$ and $`R_2`$ on the component knots in vertical framing can be obtained in two equivalent ways using the braiding and inverse braiding (parallel strands):
$`=`$ $`{\displaystyle \underset{\mathrm{}}{}}N_{R_1R_2}^{\mathrm{}}dim_q\mathrm{}\left(\lambda _{\mathrm{}}^{(+)}(R_1,R_2)\right)^2`$ (16)
$`=`$ $`q^{C_{R_1}C_{R_2}}{\displaystyle \underset{\mathrm{}}{}}N_{R_1R_2}^{\mathrm{}}dim_q\mathrm{}q^C_{\mathrm{}},`$
$`=`$ $`{\displaystyle \underset{\mathrm{}}{}}N_{R_1R_2}^{\mathrm{}}dim_q\mathrm{}\left(\lambda _{\mathrm{}}^{(+)}(R_1,R_2)\right)^2`$ (17)
$`=`$ $`q^{C_{R_1}+C_{R_2}}{\displaystyle \underset{\mathrm{}}{}}N_{R_1R_2}^{\mathrm{}}dim_q\mathrm{}q^C_{\mathrm{}},`$
where the summation over $`\mathrm{}`$ is over all the allowed representations in $`𝒢_k`$ conformal field theory and coefficients $`N_{R_1R_2}^{\mathrm{}}`$ are given by the fusion rules of the conformal theory:
$`N_{R_1R_2}^{\mathrm{}}`$ $`=`$ $`1\mathrm{if}\mathrm{}\mathrm{R}_1\mathrm{R}_2(\mathrm{fusion}\mathrm{rules}\mathrm{of}𝒢_\mathrm{k}\mathrm{conformal}\mathrm{field}\mathrm{theory})`$
$`=`$ $`0\mathrm{otherwise}.`$ (18)
There is an explicit form for the fusion matrix in terms of elements of modular matrix $`S`$ :
$$N_{R_1R_2}^{\mathrm{}}=\underset{m}{}S_{\mathrm{\Lambda }_{R_2}\mathrm{\Lambda }_m}\left(\frac{S_{\mathrm{\Lambda }_{R_1}\mathrm{\Lambda }_m}}{S_{0\mathrm{\Lambda }_m}}\right)S_{\mathrm{\Lambda }_m\mathrm{\Lambda }_{\mathrm{}}}^{},$$
(19)
where $`\left(\frac{S_{\mathrm{\Lambda }_{R_1}\mathrm{\Lambda }_m}}{S_{0\mathrm{\Lambda }_m}}\right)`$ can be shown to be the eigenvalues of fusion matrix $`(N_{R_1})_{R_2}^{\mathrm{}}N_{R_1R_2}^{\mathrm{}}`$. We can show the topological equivalence of Hopf links (16, 17) by exploiting the properties of generators $`S`$ (11) and $`T`$ ($`\tau \tau +1`$) representing modular transformations on the characters of associated $`𝒢_k`$ Wess-Zumino conformal field theory:
$`S^2`$ $`=`$ $`C,`$ (20)
$`(ST)^3`$ $`=`$ $`1,`$ (21)
where $`C_{\mathrm{\Lambda }\mathrm{\Lambda }^{}}=\delta _{\mathrm{\Lambda }^{}\overline{\mathrm{\Lambda }}}`$ is the charge conjugation matrix and $`S^{}=S^{}=S^1`$, $`S^{}=CS=SC`$; $`S_{\mathrm{\Lambda }\mathrm{\Lambda }^{}}^{}=S_{\overline{\mathrm{\Lambda }}\mathrm{\Lambda }^{}}=S_{\mathrm{\Lambda }\overline{\mathrm{\Lambda }^{}}}`$. Modular generator $`T`$ has a diagonal form:
$$T_{\mathrm{\Lambda }_{R_1}\mathrm{\Lambda }_{R_2}}=\mathrm{exp}(\frac{i\pi c}{12})q^{C_{\mathrm{\Lambda }_{R_1}}}\delta _{\mathrm{\Lambda }_{R_1}\mathrm{\Lambda }_{R_2}},$$
(22)
where central charge of the conformal field theory is $`c=\frac{kd}{k+C_v}`$ with $`d`$ denoting dimension of the group $`𝒢`$. Eqn. (21) can be rewritten as
$$\underset{\mathrm{\Lambda }_{\mathrm{}}}{}S_{\mathrm{\Lambda }_m\mathrm{\Lambda }_{\mathrm{}}}q^{C_\mathrm{\Lambda }_{\mathrm{}}}S_{\mathrm{\Lambda }_{\mathrm{}}\mathrm{\Lambda }_t}=\alpha q^{C_{\mathrm{\Lambda }_m}C_{\mathrm{\Lambda }_t}}S_{\mathrm{\Lambda }_m\mathrm{\Lambda }_t}^{},$$
(23)
where $`\alpha =exp(\frac{i\pi c}{4}).`$ Using eqns. (9, 19, 23), we have:
$$q^{C_{R_1}C_{R_2}}\underset{\mathrm{}}{}dim_q\mathrm{}q^C_{\mathrm{}}=q^{C_{R_1}+C_{R_2}}\underset{\mathrm{}}{}dim_q\mathrm{}q^C_{\mathrm{}}=\frac{S_{\mathrm{\Lambda }_{R_1}\mathrm{\Lambda }_{R_2}}}{S_{00}}.$$
(24)
Here summation over $`\mathrm{}`$ runs over only those irreducible representations in the product $`R_1R_2`$ which are allowed by fusion rules (eqn. 18). This confirms the equality of two expressions in eqns. (16) and (17) for the invariant for Hopf link.
(c) Next consider the Hopf link $`H(R_1,R_2)`$ with framing $`+1`$ for each of its component knots as drawn below. The framing is represented by a right-handed writhe in each of the knots.
The invariant for this link is given by
$$V[H(R_1,R_2)]=q^{C_{R_1}+C_{R_2}}\text{}=q^{C_{R_1}+C_{R_2}}\frac{S_{\mathrm{\Lambda }_{R_1}\mathrm{\Lambda }_{R_2}}}{S_{00}},$$
(25)
where the first factor $`q^{C_{R_1}+C_{R_2}}`$ comes from two writhes.
Next we shall present a discussion of how such invariants for framed links in $`S^3`$ based on any arbitrary group can be used to construct a manifold invariant for arbitrary three-manifolds. This will generalize a similar construction done earlier for $`SU(2)`$ group .
## 3 Three-manifold invariants
We shall first recapitulate the mathematical details of how three-manifolds are constructed by surgery on framed unoriented links. This will subsequently be used to derive an algebraic formula in terms of the link invariants characterizing three-manifolds so constructed. Starting step in this discussion is a theorem due to Lickorish and Wallace :
Fundamental theorem of Lickorish and Wallace: Every closed, orientable, connected three-manifold, $`M^3`$ can be obtained by surgery on an unoriented framed knot or link $`[L,f]`$ in $`S^3`$.
As pointed out earlier framing $`f`$ of a link $`L`$ is defined by associating with every component knot $`K_s`$ of the link an accompanying closed curve $`K_{sf}`$ parallel to the knot and winding $`n(s)`$ times in the right-handed direction. That is the linking number $`lk(K_s,K_{sf})`$ of the component knot and its frame is $`n(s)`$. In so called vertical framing where the frame is thought to be just vertically above the two dimensional projection of the knot as shown below, we may indicate this by putting $`n(s)`$ writhes in the knot or even by just simply writing the integer $`n(s)`$ next to the knot as shown below:
Next the surgery on a framed link $`[L,f]`$ made of component knots $`K_1,K_2,\mathrm{}.K_r`$ with framing $`f=(n(1),n(2),\mathrm{}.n(r))`$ in $`S^3`$ is performed in the following manner. Remove a small open solid torus neighbourhood $`N_s`$ of each component knot $`K_s`$, disjoint from all other such open tubular neighbourhoods associated with other component knots. In the manifold left behind $`S^3(N_1N_2\mathrm{}.N_r)`$, there are $`r`$ toral boundaries. On each such boundary, consider a simple closed curve (the frame) going $`n(s)`$ times along the meridian and once along the longitude of associated knot $`K_s`$. Now do a modular transformation on such a toral boundary such that the framing curve bounds a disc. Glue back the solid tori into the gaps. This yields a new manifold $`M^3`$. The theorem of Lickorish and Wallace assures us that every closed, orientable, connected three-manifold can be constructed in this way.
This construction of three-manifolds is not unique: surgery on more than one framed link can yield homeomorphic manifolds. But the rules of equivalence of framed links in $`S^3`$ which yield the same three-manifold on surgery are known. These rules are known as Kirby moves.
Kirby calculus on framed links in $`S^3`$: Following two elementary moves (and their inverses) generate Kirby calculus:
Move I. For a number of unlinked strands belonging to the component knots $`K_s`$ with framing $`n(s)`$ going through an unknotted circle $`C`$ with framing $`+1`$, the circle can be removed after making a complete clockwise (left-handed) twist from below in the disc enclosed by circle $`C`$:
In the process, in addition to introducing new crossings, the framing of the various resultant component knots, $`K_s^{}`$ to which the affected strands belong, change from $`n(s)`$ to $`n^{}(s)=n(s)\left(lk(K_s,C)\right)^2`$.
Move II. Drop a disjoint unknotted circle $`C`$ with framing $`1`$ without any change in rest of the link:
Two Kirby moves (I) and (II) and their inverses generate the conjugate moves:
Move $`\overline{I}`$. Here a circle $`C`$ with framing $`1`$ enclosing a number strands can be removed after making a complete anti-clockwise (right-handed) twist from below in the disc bounded by curve $`C`$:
Again, this changes the framing of the resultant knots $`K_s^{}`$ to which the enclosed strands belong from $`n(s)`$ to $`n^{}(s)=n(s)+\left(lk(K_s,C)\right)^2`$.
Move $`\overline{I}\overline{I}`$. A disjoint unknotted circle $`C`$ with framing $`+1`$ can be dropped without affecting rest of the link:
Thus Lickorish-Wallace theorem and equivalence of surgery under Kirby moves reduces the theory of closed, orientable, connected three-manifolds to the theory of framed unoriented links via a one-to-one correspondence:
$$\left(\begin{array}{c}FramedunorientedlinksinS^3modulo\\ equivalenceunderKirbymoves\end{array}\right)\left(\begin{array}{c}Closed,orientable,connectedthree\\ manifoldsmodulohomeomorphisms\end{array}\right)$$
This consequently allows us to characterize three-manifolds by the invariants of associated unoriented framed knots and links obtained from the Chern-Simons theory in $`S^3`$. This can be done by constructing an appropriate combination of the invariants of framed links which is unchanged under Kirby moves and which does not see orientations of the framed link:
$$\left(\begin{array}{c}Combinationofframedlinkinvariants\\ whichdonotchangeunderKirbymoves\end{array}\right)=\left(\begin{array}{c}Invariantsofassociated\\ threemanifold\end{array}\right)$$
Using the framed link invariants presented in the previous section, we shall now construct such a three-manifold invariant which is preserved under Kirby moves. The immediate step in this direction is to construct a combination of these link invariants which would be unchanged under Kirby move I:
In order to achieve this, we have to first solve the following equation for $`\mu _{R_2}`$ and $`\beta `$ relating the invariants $`V[H(R_1,R_2)]`$ and $`V[U(R_1)]`$ for these two links respectively:
$$\underset{R_2}{}\mu _{R_2}V[H(R_1,R_2)]=\beta V[U(R_1)],$$
(26)
where summation $`R_2`$ is over all the representations (highest weight $`\mathrm{\Lambda }_{R_2}`$) with projection along the longest root $`\theta `$ as $`(\mathrm{\Lambda }_{R_2},\theta )k`$. These are all the integrable representations of $`𝒢_k`$ conformal field theory. Rewriting the framed link invariants in terms of modular transformation matrix $`S`$ (25, 9) and comparing with the identity (23), we deduce the following solution:
$$\mu _{R_2}=S_{0\mathrm{\Lambda }_{R_2}},\beta =\alpha e^{\pi ic/4}.$$
(27)
Next we will consider the following two links $`H(X;R_1,R_2)`$ and $`U(X;R_1)`$:
where $`X`$ as an arbitrary entanglement inside the box. The link $`H(X;R_1,R_2)`$ is connected sum of the link $`U(X;R_1)`$ and a framed Hopf link $`H(R_1,R_2)`$. Factorization properies of invariants of such a connected sum of links yields:
$$dim_qR_1V[H(X;R_1,R_2)]=V[U(X;R_1)]V[H(R_1,R_2)].$$
(28)
Using eqn. (26), this further implies (note $`dim_qR_1`$ is the invariant $`V[U(R_1)]`$ for unknot in representation $`R_1`$ and with zero framing):
$$\underset{R_2}{}\mu _{R_2}V[H(X;R_1,R_2)]=\alpha V[U(X;R_1)].$$
(29)
We can generalize this relation for the following links $`H(X;R_1,R_2,\mathrm{}.R_n;R)`$ and $`U(X;R_1,R_2,\mathrm{}.R_n)`$,
as
Proposition 2: The invariants for these two links are related as:
$$\underset{R}{}\mu _RV[H(X;R_1,R_2,\mathrm{}R_n;R)]=\alpha V[U(X;R_1,R_2,\mathrm{}.R_n)].$$
(30)
Thus this proposition provides for equivalence of the two links under Kirby move I up to a phase factor $`\alpha `$ on the right-hand side. Lets us now outline the proof of this proposition.
Proof: In the following figure, we have redrawn the link $`H(X;R_1,R_2,\mathrm{}R_n;R)`$ in $`S^3`$ by gluing four three-manifold: two three-balls (each with $`S^2`$ boundary) and two three-manifolds with two $`S^2`$ boundaries each. The various boundaries have been glued together along the dotted lines as indicated. This allows us to use the Proposition 1 above to evaluate the invariant for this link.
Instead of plat diagram described in Section 2, we have the closure of a braid here. The states corresponding to Chern-Simons functional integral over two three-balls with $`S^2`$ boundaries ‘1’ and ‘3’ will be represented by the vector $`|\psi (\{R_i\},R)`$ (where $`i[1,n]`$) and its dual. The matrix element corresponding to the braiding inside the three-manifold with two $`S^2`$ boundaries ‘1’ and ‘2’ can be explicitly computed in a convenient basis. We shall work with a basis represented by the following conformal block of associated Wess-Zumino theory:
where $`l_1=R_nR_{n1},\mathrm{}l_{n1}=l_{n2}R_1`$ and $`s_1=R_1R_2,\mathrm{}s_{n1}=s_{n2}R_n,s=l_{n1}R`$. In this basis the matrix corresponding to the three-manifold with boundaries marked ‘1’ and ‘2’ turns out to be:
$$\nu _1=\underset{l_1,\mathrm{}l_{n1},s_1,\mathrm{}s_{n1},s}{}q^{C_s}|\varphi _{(l_1,\mathrm{}l_{n1},s_1,\mathrm{}s_{n1},s)}^{(2)}\varphi _{(l_1,\mathrm{}l_{n1},s_1,\mathrm{}s_{n1},s)}^{(1)}|$$
(31)
where the superscripts $`(1)`$ and $`(2)`$ inside the basis states refer to two $`S^2`$ boundaries containing the three-manifold. This result involves properties of the braiding and duality matrices which we present for $`n=2`$ in the Appendix. These can readily be generalized to arbitrary $`n`$.
The matrix $`\nu _2(X)`$ representing the other three-manifold containing entanglement $`X`$ between two $`S^2`$ boundaries indicated by dotted lines ‘$`2`$’ and ‘$`3`$’ in the figure above can similarly be evaluated. In addition, we also need to write down the states $`|\psi ^{(1)}`$ and $`\psi ^{(3)}|`$ corresponding to the two three-balls with boundaries indicated by dotted lines ‘$`1`$’ and ‘$`3`$’. All these in the above basis can be written as :
$`\nu _2(X)`$ $`=`$ $`{\displaystyle \underset{\{l_i\},s,\{s_i\},\{s_i^{}\}}{}}X(\{s_i\},\{s_i^{}\})|\varphi _{(l_1,\mathrm{}l_{n1},s_1,\mathrm{}s_{n1},s)}^{(3)}\varphi _{(l_1,\mathrm{}l_{n1},s_1^{},\mathrm{}s_{n1}^{},s)}^{(2)}|`$ (32)
$`|\psi ^{(1)}`$ $`=`$ $`{\displaystyle \underset{l_1,\mathrm{}l_{n1},,s}{}}\sqrt{dim_qs}|\varphi _{(l_1,\mathrm{}l_{n1},l_1,\mathrm{}l_{n1},s)}^{(1)}`$ (33)
$`\psi ^{(3)}|`$ $`=`$ $`{\displaystyle \underset{l_1,\mathrm{}l_{n1},s}{}}\sqrt{dim_qs}\varphi _{(l_1,\mathrm{}l_{n1},l_1,\mathrm{}l_{n1},s)}^{(3)}|`$ (34)
Gluing these four three-manifolds along the oppositely oriented $`S^2`$ boundaries, we get the link $`H(X;R_1,R_2,\mathrm{}R_n;R)`$ whose invariant can now be written as:
$$V[H(X;R_1,R_2\mathrm{}R_n;R)]=\underset{\{l_i\}}{}X(\{l_i\},\{l_i\})\left[\underset{sl_{n1}R}{}(dim_qs)q^{C_s}\right].$$
(35)
Clearly the term in square bracket is the Hopf link invariant $`H(l_{n1},R)`$ of eqn. (25). Similarly, we can compute the link invariant for $`U(X;R_1,R_2\mathrm{}R_n)`$ as
$$V[U(X;R_1,R_2\mathrm{}R_n)]=\underset{\{l_i\}}{}(dim_ql_{n1})X(\{l_i\},\{l_i\}).$$
(36)
Now using eqn.(26), it is easy to prove Proposition 2:
$`{\displaystyle \underset{R}{}}\mu _RV[H(X;R_1,R_2\mathrm{}R_n;R)]`$ $`=`$ $`{\displaystyle \underset{\{l_i\}}{}}X(\{l_i\},\{l_i\}){\displaystyle \underset{R}{}}\mu _RV[H(l_{n1},R)]`$ (37)
$`=`$ $`\alpha V[U(X;R_1,R_2\mathrm{}R_n)].`$ (38)
This completes our discussion of Kirby move I. For Kirby move II, we note that for a link containing a disjoint unknot with framing $`1`$, we have:
$$\underset{\mathrm{}}{}\mu _{\mathrm{}}\text{}=\alpha ^{}\text{}$$
(39)
This follows readily due to the exact factorizations of invariants of disjoint links into those of individual links and use of the identity: $`_{\mathrm{}}S_0\mathrm{}q^C_{\mathrm{}}S_\mathrm{}0=e^{\pi ic/4}S_{00}`$.
Clearly the Eqns. (30) and (39) respectively correspond to the two generators of Kirby calculus. Presence of phases $`\alpha `$ in these equations reflects change of three-framing of the associated manifold under Kirby moves. Three-manifolds constructed by surgery on links equivalent under Kirby moves though topologically same, yet may differ in terms of their three-framing. The dependence on three-framing ($`\alpha `$ factors in above equations representing Kirby moves) can be rotated away by use of a nice property of the linking matrix under Kirby moves. For a framed link $`[L,f]`$ whose component knots $`K_1,K_2,\mathrm{}.K_r`$ have framings (self-linking numbers) as $`n_1,n_2,\mathrm{}.n_r`$ respectively, the linking matrix is defined as
$$W[L,f]=\left(\begin{array}{ccccc}n_1& lk(K_1,K_2)& lk(K_1,K_3)& \mathrm{}..& lk(K_1,K_r)\\ lk(K_2,K_1)& n_2& lk(K_2,K_3)& \mathrm{}..& lk(K_2,K_r)\\ ..& ..& n_3& \mathrm{}..& ..\\ ..& ..& ..& \mathrm{}..& ..\\ lk(K_r,K_1)& ..& ..& \mathrm{}..& n_r\end{array}\right)$$
where $`lk(K_i,K_j)`$ is the linking number of knots $`K_i`$ and $`K_j`$. The signature of linking matrix is given by
$$\sigma [L,f]=(\mathrm{no}.\mathrm{of}+\mathrm{ve}\mathrm{eigenvalues}\mathrm{of}W)(\mathrm{no}.\mathrm{of}\mathrm{ve}\mathrm{eigenvalues}\mathrm{of}W)$$
Then this signature for the framed link $`[L,f]`$ and those for the links $`[L^{},f^{}]`$ obtained by transformation under two elementary generators of Kirby calculus are related in a simple fashion:
$`\mathrm{Kirby}\mathrm{move}\mathrm{I}:\sigma [L,f]`$ $`=`$ $`\sigma [L^{},f^{}]+1;`$
$`\mathrm{Kirby}\mathrm{move}\mathrm{II}:\sigma [L,f]`$ $`=`$ $`\sigma [L^{},f^{}]1.`$ (40)
Notice, though the sign of linking numbers $`lk(K_i,K_j)`$ for distinct knots does depend on the relative orientations of knots $`K_i`$ and $`K_j`$, the signature of linking matrix does not depend on the relative orientations of component knots.
Now collecting the properties of framed link (30, 39) and the signature of linking matrix under the Kirby moves (3), we may state our main result:
Proposition 3: For a framed link $`[L,f]`$ with component knots, $`K_1,K_2,\mathrm{}.K_r`$ and their framings respectively as $`n_1,n_2,\mathrm{}.n_r`$, the quantity
$$\widehat{F}[L,f]=\alpha ^{\sigma [L,f]}\underset{\{R_i\}}{}\mu _{R_1}\mu _{R_2}\mathrm{}.\mu _{R_r}V[L;n_1,n_2,\mathrm{}n_r;R_1,R_2,\mathrm{}.R_r]$$
(41)
constructed from invariants $`V`$ (in vertical framing) of the framed link, is an invariant of the associated three-manifold obtained by surgery on that link.
Notice individual link invariants $`V[L;n_1,n_2,\mathrm{}n_r;R_1,\mathrm{}R_r]`$ do in general depend on the relative orientations of component knots. Reversal of orientation on a particular knot changes the group representation living on it to its conjugate. Since all representations are summed for each component knot, the resultant combination (41) is an invariant of unoriented link.
The combination $`\widehat{F}[L,f]`$ of link invariants so constructed is exactly unchanged under Kirby calculus. Because of the factor depending on signature of linking matrix in front, there are no extra factors of $`\alpha `$ generated by Kirby moves.
This generalizes a similar proposition obtained earlier for $`SU(2)`$ theory to any arbitrary semi-simple gauge group.
Explicit examples: Now let us give the value of this invariant for some simple three-manifolds.
1) The surgery descriptions of manifolds $`S^3`$, $`S^2\times S^1`$, $`RP^3`$ and Lens spaces of the type $`(p,\pm 1)`$ are given by an unknot with framing $`+1,0`$ , $`+2`$ and $`\pm p`$ respectively. As indicated above the knot invariant for an unknot with zero framing carrying representation $`R`$ is $`dim_qR=S_{\mathrm{\Lambda }_R0}/S_{00}`$. Thus the invariant for $`S^3`$ is:
$$\widehat{F}[S^3]=\alpha ^1\underset{\mathrm{\Lambda }_R}{}S_{0\mathrm{\Lambda }_R}q^{C_{\mathrm{\Lambda }_R}}\frac{S_{\mathrm{\Lambda }_R0}}{S_{00}},$$
where the factor $`q^{C_{\mathrm{\Lambda }_R}}`$ comes from framing $`+1`$ (one right-handed writhe). Use of identity (23), $`_{\mathrm{\Lambda }_R}S_{0\mathrm{\Lambda }_R}q^{C_{\mathrm{\Lambda }_R}}S_{\mathrm{\Lambda }_R0}=\alpha S_{00}`$ immediately yields the invariant simply to be:
$$\widehat{F}[S^3]=1$$
(42)
Next for three-manifold $`S^2\times S^1`$, we have
$$\widehat{F}[S^2\times S^1]=\underset{\mathrm{\Lambda }_R}{}S_{\mathrm{\Lambda }_R0}\frac{S_{\mathrm{\Lambda }_R0}}{S_{00}}=\frac{1}{S_{00}},$$
(43)
where orthogonality of the $`S`$-matrix has been used. Finally, for $`RP^3`$ and more generally Lens spaces $`(p,\pm 1)`$, we have
$`\widehat{F}[RP^3]`$ $`=`$ $`\alpha ^1{\displaystyle \underset{R}{}}{\displaystyle \frac{S_{0\mathrm{\Lambda }_R}q^{2C_{\mathrm{\Lambda }_R}}S_{\mathrm{\Lambda }_R0}}{S_{00}}},`$ (44)
$`\widehat{F}[(p,\pm 1)]`$ $`=`$ $`\alpha ^1{\displaystyle \underset{R}{}}{\displaystyle \frac{S_{0\mathrm{\Lambda }_R}q^{\pm pC_{\mathrm{\Lambda }_R}}S_{\mathrm{\Lambda }_R0}}{S_{00}}}.`$ (45)
2) A more general example is the whole class of Lens spaces $`(p,q)`$; above manifolds are special cases of this class. These are obtained by surgery on a framed link made of successively linked unknots with framing given by integers $`a_1,a_2,\mathrm{}\mathrm{}a_n`$:
where these framing integers provide a continued fraction representation for the ratio of two integers $`p,q`$:
$$\frac{p}{q}=a_n\frac{1}{a_{n1}\frac{1}{\mathrm{}\mathrm{}.a_3\frac{1}{a_2\frac{1}{a_1}}}}.$$
The invariant for these manifolds can readily be evaluated. The relevant link $`[L,f]`$ above is just a connected sum of framed $`n1`$ Hopf links so that its link invariant is obtained by the factorization property of invariants for such a connected sum. Placing representations $`R_1`$, $`R_2`$, ….. $`R_n`$ on the component knots, this link invariant is:
$$V[L,f;R_1,R_2,\mathrm{}.R_n]=\frac{q^{_{i=1}^na_iC_{R_i}}_{i=1}^{n1}S_{\mathrm{\Lambda }_{R_i}\mathrm{\Lambda }_{R_{i+1}}}}{S_{00}_{i=2}^{n1}S_{\mathrm{\Lambda }_{R_i}0}},$$
where the factor $`q^{_{i=1}^na_iC_{R_i}}`$ is due to the framing $`f=(a_1,a_2,\mathrm{}.a_n)`$ of knots. This finally yields a simple formula for the three-manifold invariant:
$$\widehat{F}[(p,q)]=\alpha ^{\sigma [L,f]}\alpha ^{({\scriptscriptstyle a_i})/3}\frac{(SM^{(p,q)})_{00}}{S_{00}},$$
(46)
where matrix $`M^{(p,q)}`$ is given in terms of the modular matrices $`S`$ and $`T`$:
$$M^{(p,q)}=T^{a_n}ST^{a_{n1}}S\mathrm{}\mathrm{}T^{a_2}ST^{a_1}S.$$
(47)
The corresponding expression for the partition function for $`SU(2)`$ Chern-Simons field theory in Lens spaces has also been obtained earlier in refs. .
3) Another example we take up is the Poincare manifold $`P^3`$ (also known as dodecahedral space or Dehn’s homology sphere). It is given by surgery on a right-handed trefoil knot with framing $`+1`$:
Notice, each right-handed crossing of the trefoil introduces $`+1`$ linking number between the knot and its vertical framing, and each of the two left-handed writhes contributes $`1`$ so that the total frame number is $`+1`$. The knot invariant for a right-handed trefoil (in vertical framing with no extra writhes) carrying representation $`\mathrm{}`$ is:
$$V[T;\mathrm{}]=\underset{R}{}N_{\mathrm{}\mathrm{}}^Rdim_qR()^{6ϵ_{\mathrm{}}3ϵ_R}q^{3C_{\mathrm{}}+\frac{3}{2}C_R}$$
Using this trefoil invariant and the Proposition 3, the three-manifold invariant for Poincare manifold turns out to be:
$$\widehat{F}[P^3]=\alpha ^1\underset{m,\mathrm{},R}{}\frac{()^{6ϵ_{\mathrm{}}3ϵ_R}S_{0\mathrm{\Lambda }_{\mathrm{}}}S_{0\mathrm{\Lambda }_R}S_{\mathrm{\Lambda }_{\mathrm{}}\mathrm{\Lambda }_m}S_{\mathrm{\Lambda }_{\mathrm{}}\mathrm{\Lambda }_m}S_{\mathrm{\Lambda }_R\mathrm{\Lambda }_m}^{}q^{5C_{\mathrm{}}+\frac{3}{2}C_R}}{S_{00}S_{0\mathrm{\Lambda }_m}}.$$
(48)
4) Similarly, surgery on a right-handed trefoil $`T`$ with framing number $`1`$ (that is, with four left-handed writhes in vertical framing) yields another homology three-sphere (with fundamental group presented by $`\alpha ,\beta :`$ $`(\alpha \beta )^2=\alpha ^3=\beta ^7`$ ). Its invariant is
$$\widehat{F}[T,1]=\alpha \underset{m,\mathrm{},R}{}\frac{()^{6ϵ_{\mathrm{}}3ϵ_R}S_{0\mathrm{\Lambda }_{\mathrm{}}}S_{0\mathrm{\Lambda }_R}S_{\mathrm{\Lambda }_{\mathrm{}}\mathrm{\Lambda }_m}S_{\mathrm{\Lambda }_{\mathrm{}}\mathrm{\Lambda }_m}S_{\mathrm{\Lambda }_R\mathrm{\Lambda }_m}^{}q^{7C_{\mathrm{}}+\frac{3}{2}C_R}}{S_{00}S_{0\mathrm{\Lambda }_m}}.$$
(49)
5) The surgery on a right-handed trefoil with framing number $`+3`$ yields a coset manifold $`S^3/T^{}`$ where $`T^{}`$ is the binary tetrahedral group generated by two different $`2\pi /3`$ rotations $`\alpha ,\beta `$ about two different vertices of a tetrahedron with $`\alpha ^3=\beta ^3=(\alpha \beta )^2=1.`$ The three-manifold invariant for this coset manifold is
$$\widehat{F}[S^3/T^{}]=\alpha ^1\underset{m,\mathrm{},R}{}\frac{()^{6ϵ_{\mathrm{}}3ϵ_R}S_{0\mathrm{\Lambda }_{\mathrm{}}}S_{0\mathrm{\Lambda }_R}S_{\mathrm{\Lambda }_{\mathrm{}}\mathrm{\Lambda }_m}S_{\mathrm{\Lambda }_{\mathrm{}}\mathrm{\Lambda }_m}S_{\mathrm{\Lambda }_R\mathrm{\Lambda }_m}^{}q^{3C_{\mathrm{}}+\frac{3}{2}C_R}}{S_{00}S_{0\mathrm{\Lambda }_m}}.$$
(50)
## 4 Conclusions
We have presented here a construction of a class of three-manifold invariants, one each for any arbitrary semi-simple gauge group. The construction exploits the one-to-one correspondence between framed unoriented links in $`S^3`$ modulo equivalence under Kirby moves to closed orientable connected three-manifolds modulo homeomorphisms. Three-manifolds are characterized by an appropriate combination of invariants of the associated links. This combinations of link invariants is obtained from Chern-Simons theory in $`S^3`$ and is unchanged by Kirby moves. The construction is a direct generalization of that developed for an $`SU(2)`$ Chern-Simons theory earlier. The manifold invariant obtained from $`SU(2)`$ theory has been shown to be related to partition function of Chern-Simons theory on that manifold . The generalized three-manifold invariants $`\widehat{F}(M)`$ constructed here are also related to the partition function $`Z(M)`$ of Chern-Simons theory by an overall normalization:
$$\widehat{F}(M)=\frac{Z(M)}{S_{00}}.$$
(51)
Thus, given a surgery presentation of a three-manifold, this provides a simple method of computing partition function of Chern-Simons field theory based on an arbitrary gauge group in that three-manifold.
Acknowledgements: P.R would like to thank TIFR for hospitality where part of this work got done. The work of P.R is supported by a CSIR Grant.
Appendix:
In this appendix, we shall use the properties of braiding and duality matrices to derive the result (31). for the case $`n=2`$.
A convenient basis state for $`n=2`$ is provided by the conformal block of the Wess-Zumino theory pictorially represented as:
The matrix $`\nu _1`$ for $`n=2`$ corresponds to a functional integral over a three-manifold with two $`S^2`$ boundaries, each with six punctures carrying representations $`R_2`$, $`R_1`$, $`R_1`$, $`R_2`$, $`R`$ and $`R`$.
Between these two boundaries we have braiding $`b_4b_3^2b_4b_3^2`$ and a $`+1`$ writhe on each of the three strands carrying representations $`\overline{R_1}`$, $`\overline{R_2}`$, $`\overline{R}`$. Using the explicit representation for braid generators in terms of their eigenvalues and duality matrices $`a_{sp}\left[\begin{array}{cc}R_1& R_2\\ R_3& R_4\end{array}\right]`$, we obtain $`\nu _1`$ to be:
$`\nu _1`$ $`=`$ $`q^{C_{R_1}+C_{R_2}+C_R}b_4b_3^2b_4b_3^2|\varphi _{l_1,s_1,s}^{(1)}\varphi _{l_1,s_1,s}^{(2)}|`$
$`=`$ $`q^{C_{R_1}+C_{R_2}+C_R}(\lambda _{s_1}(R_1,R_2))^2{\displaystyle \underset{p_1s_1^{},r,s_1^{\prime \prime }}{}}a_{s_1p}\left[\begin{array}{cc}\overline{R_1}& \overline{R_2}\\ \overline{R}& s\end{array}\right]\lambda _p(R_2,R)`$
$`a_{s_1^{}p}\left[\begin{array}{cc}\overline{R_1}& \overline{R}\\ \overline{R_2}& s\end{array}\right](\lambda _{s_1^{}}(R_1,R))^2a_{s_1^{}r}\left[\begin{array}{cc}\overline{R_1}& \overline{R}\\ \overline{R_2}& s\end{array}\right]\lambda _r(R,R_2)a_{s_1^{\prime \prime }r}\left[\begin{array}{cc}\overline{R_1}& \overline{R_2}\\ \overline{R}& s\end{array}\right]`$
$`|\varphi _{l_1,s_1^{\prime \prime },s}^{(1)}\varphi _{l_1,s_1,l}^{(2)}|`$ (52)
Using the following property of the duality matrices,
$`{\displaystyle \underset{l}{}}(1)^{ϵ_l}a_{pl}\left[\begin{array}{c}R_1R_4\\ R_3R_2\end{array}\right]q^{\frac{C_l}{2}}a_{p^{}l}\left[\begin{array}{c}R_1R_3\\ R_4R_2\end{array}\right]`$ $`=`$ $`(1)^{ϵ_{R_1}+ϵ_{R_2}+ϵ_{R_3}+ϵ_{R_4}ϵ_pϵ_p^{}}`$ (53)
$`a_{pp^{}}\left[\begin{array}{c}R_3R_2\\ R_4R_1\end{array}\right]q^{\frac{C_pC_p^{}+C_{R_1}+C_{R_2}+C_{R_3}+C_{R_4}}{2}}`$
and the orthogonality relation, the above equation can be simplified to give
$$\nu _1=\underset{l_1,s_1,s}{}q^{C_s}|\varphi _{l_1,s_1,s}^{(2)}\varphi _{l_1,s_1,s}^{(1)}|$$
(54)
Generalization of this result for arbitrary $`n`$ is straight forward.
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# Overcharging of a macroion by an oppositely charged polyelectrolyte
## I Introduction
Electrostatic interactions play an important role in aqueous solutions of biological and synthetic polyelectrolytes (PE). They result in the aggregation and complexation of oppositely charged macroions in solutions. For example, in the chromatin, negative DNA winds around a positive histone octamer to form a complex known as the nucleosome. The nucleosome was found to have negative net charge $`Q^{}`$ whose absolute value is as large as 15% of the bare positive charge of the protein, $`Q`$. This counterintuitive phenomenon is called the charge inversion and can be characterized by the charge inversion ratio, $`|Q^{}|/Q`$. For PE-micelle systems, charge inversion has been predicted by Monte-Carlo simulations and observed experimentally .
These and other examples have recently stimulated several theoretical studies of charge inversion accompanying the complexation of a flexible PE with a rigid spherical or cylindrical macroion of opposite sign (for more extensive bibliography on this subject see Ref. ). All these authors arrive at the charge inversion for such a complexation. It was also shown that if the PE molecule is not totally adsorbed at the surface, its remaining part is repelled by the inverted charge of the macroion and forms an almost straight radial tail (see Fig. 1). However, all these papers use different models and seemingly deal with charge inversion of different nature. Surprisingly, both Refs. show that the inverted charge of a macroion $`Q^{}`$ does not depend on the value of the bare charge $`Q`$.
In this paper we present a new theory of complexation of a flexible PE with an oppositely charged rigid sphere. We consider here only the case of a weakly charged PE which does not create Onsager-Manning condensation. We show that both in salt free and salty solutions the charge inversion by such PE is driven by repulsive correlations of PE turns at the macroion surface. Such correlations make an almost equidistant solenoid (see Fig. 1), which locally resembles one-dimensional Wigner crystal along the direction perpendicular to PE. In the absence of salt, the charge inversion ratio is smaller than 100%. In a salty solution, it grows with the salt concentration. When the Debye-Hückel screening radius $`r_s`$ becomes smaller than the distance between neighboring turns $`A`$, the charge inversion ratio can be larger than 100%.
The charge inversion of a macroion due to complexation with one PE molecule can be explained in the way similar to Refs. , which dealt with the charge inversion of a macroion screened by many rigid multivalent counterions ($`Z`$-ions). The tail repels adsorbed PE and creates correlation hole or, in other words, its positively charged image. This image in the already adsorbed layer of PE is responsible for the additional correlation attraction to the surface, which leads to the charge inversion.
We show that smearing of charge PE on the surface of the sphere employed in Ref. is a good approximation only at $`Aa`$. If $`Aa`$ smearing of charge at the surface of sphere is a rough approximation and leads to anomalously strong inversion of charge and to the unphysical independence of the inverted charge $`Q^{}`$ on $`Q`$. The reason of this phenomenon is easy to understand. Smearing means that the PE solenoid is assumed to behave as a perfect metal. A neutral metal surface can adsorb a charged PE due to image forces, making the charge inversion ratio infinite. In reality, for an insulating macroion, an image of a point charge in the PE coil can not be smaller than $`A`$ and the energy of attraction to it vanishes at growing $`A`$. Only a macroion with a finite charge $`Q`$ adsorbs a PE coil with a finite $`A`$. Therefore, $`Q^{}`$ depends on $`Q`$ and the charge inversion ratio is always finite.
Our analytic theory is followed by Monte-Carlo simulations. They demonstrate good agreement with the theory.
## II An analytical theory.
For a quantitative calculation, consider the complexation of a negative PE with linear charge density $`\eta `$ and length $`L`$, with a spherical macroion with radius $`R`$ and positive charge $`Q`$. We assume that the PE is weakly charged, i. e. $`\eta \eta _c`$, where $`\eta _c=k_BTD/e`$ is Onsager-Manning critical linear density, $`T`$ is the temperature, $`k_B`$ is the Boltzmann constant and $`D`$ is the dielectric constant of water. In this case, there is no Onsager-Manning condensation of counterions and one can use linear theory of screening. Because we are interested in the charge inversion of the complex, we assume that the PE length $`L`$ is greater than the neutralizing length $`=Q/\eta `$. In this case, a finite length $`L_1`$ of the PE is tightly wound around the macroion due to the electrostatic attraction. The rest of the PE with length $`L_2=LL_1`$ can be arranged into two possible configurations: one tail with length $`L_2`$ or two tails with length $`L_2/2`$ going in opposite directions radially outwards from the center of the macroion. In both cases, the tails are straight to minimize its electrostatic self-energy. We assume that $`R`$, so that there are many turns of the PE around the sphere. Our goal is to calculate the net charge of the complex $`Q^{}=QL_1\eta =(L_1)\eta `$ and the charge inversion ratio $`|Q^{}|/Q`$. We show that, in the most common configuration with one tail, this net charge is negative: more PE winds around the macroion than it is necessary to neutralize it.
Let us start from the salt free solution in which all Coulomb interactions are not screened. For simplicity, we assume that the PE has no intrinsic rigidity, but its linear charge density is large so that it has a rod-like configuration in solution due to Coulomb repulsion between monomers. When PE winds around the macroion, the strong Coulomb repulsion between the neighboring PE turns keeps them parallel to each other and establishes an almost constant distance $`A`$ between them (Fig. 1). The total energy of the macroion with the PE solenoid wound around it, $`F_1`$, can be written as a sum of the Coulomb energy of its net charge plus the self-energy of PE:
$$F_1=(L_1)^2/2R+L_1\mathrm{ln}(A/a).$$
(1)
Here and below we write all energies in units of $`\eta ^2/D`$, where $`D`$ is dielectric constant of water (thus, all energies have the dimensionality of length.) The second term in Eq. (1) deserves special attention. The self-energy of a straight PE of length $`L_1`$ in the solution is $`L_1\mathrm{ln}(L_1/a)`$. However, when it winds around the macroion, every turn is effectively screened by the neighboring turns at the distance $`A`$. This screening brings the self-energy down to $`L_1\mathrm{ln}(A/a)`$. At length scale greater than $`A`$, the surface charge density of the spherical complex is uniform and the excess charge $`L_1`$ is taken into account by the first term in Eq. (1). In other words, one can interpret Eq. (1) thinking about our system as the superposition of a uniformly charged sphere with charge $`(L_1)`$ and a neutral complex consisting of the solenoid on a neutralizing spherical background. The total energy of these two objects is additive. Indeed, the energy of interaction between them vanishes because the first one creates a constant potential on the second neutral one.
One can also rewrite the energy of solenoid on the neutralizing background as
$$L_1\mathrm{ln}(A/a)=L_1\mathrm{ln}(R/a)L_1\mathrm{ln}(R/A).$$
(2)
Here the first term is the self-energy of the PE with length $`L_1`$ whose turns are randomly positioned on the macroion. (Indeed, for a strongly charged PE, each PE turn is straight up to a distance of the order of $`R`$ due to its electrostatic rigidity. If we keep a PE turn fixed and average over random positions of all other turns we find our turn on the uniform spherical background of opposite charge. The absolute value of the background charge is of the order $`R`$, the energy of interaction of our turn with it is of the order $`R`$ and is negligible compared to the turn’s self-energy $`R\mathrm{ln}(R/a)`$ or $`\mathrm{ln}(R/a)`$ per unit length.) Now it is easy to identify the second term of Eq. (2) as the correlation energy. It represents the lowering of the system’s energy by forming an equidistant coil from the random one. This correlation energy, $`E_{cor}`$, is of the order of the interaction of the PE turn with its background (a stripe of of the length $`R`$ and the width $`A`$ of the surface charge of the macroion) because all other turns lie at the distance $`A`$ and beyond. Estimating $`AR^2/L_1`$, we can write
$$E_{cor}L_1\mathrm{ln}(R/A)L_1\mathrm{ln}(L_1/R).$$
(3)
Substituting Eqs. (3) and (2) into Eq. (1) for the total energy of the spherical complex, we obtain
$$F_1=L_1\mathrm{ln}(R/a)L_1\mathrm{ln}(L_1/R)+(L_1)^2/2R.$$
(4)
To take into account the PE tails, let us consider each tail configuration separately.
One tail configuration. In this case, the total free energy of the system is the sum of that of the spherical complex, the self-energy of the tail and their interaction. This gives:
$$F=F_1+L_2\mathrm{ln}(L_2/a)+(L_1)\mathrm{ln}\left[(L_2+R)/R\right].$$
(5)
To find the optimum value of the length $`L_1`$ one has to minimize $`F`$ with respect to $`L_1`$. Using Eq. (5) and the relation $`L_2=LL_1`$, we obtain
$$(L_1)\left[R^1(LL_1+R)^1\right]=\mathrm{ln}(/R),$$
(6)
where we neglected terms of the order of unity and took into account that $`L_2R`$ (as shown below, Eq. (8)). The physical meaning of Eq. (6) is transparent: The left side is the energy of the Coulomb repulsion of the net charge of the spherical complex which has to be overcome in order to bring an unit length of the PE from the tail to the sphere. The right hand side (in which, $`L_1`$ has been approximated by $``$) is the absolute value of the correlation energy gained at the sphere which helps to overcome this repulsion (See Eq. (3)). Equilibrium is reached when these two forces are equal. From Eq. (6), one can easily see that $`L_1`$ is positive, indicating a charge inversion scenario: more PE collapses on the macroion than it is necessary to neutralize it. Eq. (6) also clearly shows that correlations are the driving force of charge inversion.
To understand how the length $`L_1`$ varies for different PE length $`L`$, it is instructive to solve Eq. (6) graphically. One can see the following behavior (Fig. 2):
(a) When $`L`$ is small, Eq. (6) has no solutions, $`F/L_1`$ is always negative. The free energy is a monotonically decreasing function of $`L_1`$ and is minimal when $`L_1=L`$. In this regime, the whole PE collapses on the macroion.
(b) As $`L`$ increases beyond a length $`L^{}`$, Eq. (6) acquires two solutions, which correspond to a local minimum and a local maximum in the free energy as a function of $`L_1`$. The global minimum is still at $`L_1=L`$ and the whole PE remains in the collapsed state.
(c) When $`L`$ increases further, at a length $`L=L_c`$, the local minimum in the free energy at $`L_1<L`$ becomes smaller than the minimum at $`L_1=L`$. A first order phase transition happens and a tail with a finite length $`L_2`$ appears. $`L_c`$ can be found from the requirement that the equation $`F(L_1)F(L)=0`$ has solutions at $`0<L_1<L`$. Using Eq. (5), one gets
$$L_c+R\mathrm{ln}(/R)+R\sqrt{\mathrm{ln}(/R)}\sqrt{\mathrm{ln}\mathrm{ln}(/R)},$$
(7)
and the tail length $`L_2`$ at this critical point is
$$L_{2,c}R\sqrt{\mathrm{ln}(/R)}\sqrt{\mathrm{ln}\mathrm{ln}(/R)}.$$
(8)
As $`L`$ continues to increase, $`L_1`$ decreases and eventually saturates at the constant value
$$L_{1,\mathrm{}}=+R\mathrm{ln}(/R)L_cL_{2,c},$$
(9)
which can be found from Eq. (6) by letting $`L\mathrm{}`$. Eq. (7) - (9) are asymptotic results valid at $`//R\mathrm{}`$. If $`//R`$ is not very large one can find $`L_1(L)`$ minimizing Eq. (5) numerically. In Fig. 3 we present results for the case $`=25R`$, which corresponds to $`25/2\pi 4`$ turns. In this case, $`L_c=35.5R`$, $`L_{2,c}=4.0R`$ and $`L_{1,\mathrm{}}=30.4R`$. It should be also noted that, as Fig. 3 and Eqs. (7), (8) and (9) suggest, $`L_1`$ is almost equal $`L_{1,\mathrm{}}`$ after the phase transition.
At $`R`$, the charge inversion ratio $`|Q^{}|/Q=(L_1)/`$ can be calculated from Eqs. (7) and (9): $`|Q^{}|/Q=(R/)\mathrm{ln}(/R)1`$. Thus, the charge inversion ratio is only logarithmically larger than the inverse number of PE turns in the coil.
Using the insight gained above, we are now in a position to achieve better understanding of the nature of the approximation employed in Ref. . The authors of Ref. replaced the adsorbed PE by the same charge uniformly smeared at the macroion surface. Therefore, the term $`L_1\mathrm{ln}(A/a)`$ was omitted in Eq. (1), so that at $`Aa`$, the correlation energy was overestimated. This approximation replaces the right hand side of Eq. (6) by the self-energy of a unit length of the tail. Correspondingly, Eq. (6) now balances the self-energy of a unit length of the tail with the electrostatic energy of this unit length smeared at the surface of overcharged macroion. Thus we can call this mechanism of charge inversion “the elimination of the self-energy” or simply “metallization”.
As a result, the charge inversion obtained in Ref. , at $`Aa`$, is larger than that of our paper. (Our correlation mechanism can be interpreted as a partial elimination of the self-energy. The second term of Eq. (1) is what is left from the PE self-energy due to self screening of PE at the distance $`A`$.) Surprising independence of $`Q^{}`$ on $`Q`$ or, in other words, the possibility of an infinite charge inversion ratio obtained in Ref. is also related to smearing of PE on the macroion surface. This happens because when PE arrives at the macroion surface it looses all its (positive) self-energy. This brings about an energy gain which does not depend on the bare charge of the macroion.
On the other hand, at $`Aa`$, the smearing of PE is a good approximation and our results are close to that of Ref. .
Two tails configuration. The free energy of the system can be written similar to Eq. (5), keeping in mind that we have two tails instead of one, each with length $`L_2/2`$:
$`F`$ $`=`$ $`F_1+L_2\mathrm{ln}{\displaystyle \frac{L_2}{2a}}+2(L_1)\mathrm{ln}{\displaystyle \frac{L_2+2R}{2R}}+`$ (11)
$`+(L_2+2R)\mathrm{ln}{\displaystyle \frac{L_2+2R}{2R}}(L_2+4R)\mathrm{ln}{\displaystyle \frac{L_2+4R}{4R}}.`$
The last two terms describe the interaction between the tails. The optimum length $`L_1`$ can be found from the condition of a minimum in the free energy. Taking into account that, as shown below, $`L_2R`$ and ignoring terms of the order unity, one gets
$$(L_1)\left[R^1(L_2/2+R)^1\right]+\mathrm{ln}(L_2/R)=\mathrm{ln}(/R).$$
(12)
Comparing this equation to Eq. (6), one finds an additional potential energy cost $`\mathrm{ln}(L_2/R)`$ for bringing a unit length of the PE from the end of a tail to the sphere. It originates from the interaction of this segment with the other tail. When $`L`$ is not very large, $`L_2`$, one can neglect this additional term and the two tail system behaves like the one tail one. At a small $`L`$, the whole PE lies on the macroion surface and the system is overcharged. As $`L`$ increases, eventually a first order phase transition happens, where two tails with length of the order $`R\sqrt{\mathrm{ln}(/R)}`$ appear. On the other hand, when $`L`$ is very large, such that $`L_2`$, the new term dominates and the macroion becomes undercharged ($`L_1`$ is negative) with $`L_1`$ decreasing as a logarithmic function of the PE length: $`L_1R\mathrm{ln}(L/)`$. At an exponentially large value of $`L\mathrm{exp}(/R)`$, the length $`L_1`$ reaches zero and the whole PE unwinds from the macroion.
Above, we have described configurations with one tail and two tails separately. One should ask which of them is realized at a given $`L`$. Numerical calculations show that, when $`L`$ is not very large, the overcharged, one tail configuration is lower in energy. At a very large value of $`L`$, the complex undergoes a first order phase transition to a two tails configuration and becomes undercharged. The value of this critical length $`L_{cc}`$ can be estimated by equating the free energies (5) and (11) at their optimal values of $`L_1`$ which are $`+R\mathrm{ln}(/R)`$ and $`R\mathrm{ln}(L/)`$ respectively. In the limit where $`\mathrm{ln}(/R)1`$, keeping only highest order terms, we get $`L_{cc}^2/R`$, which indeed is a very large length scale. This order of appearance of one and two tail configurations is in disagreement with Ref. .
In practical situations, there is always a finite salt concentration in a water solution. One, therefore, has to take the finite screening length $`r_s`$ into account. For any reasonable $`r_s`$, $`L_{cc}r_s`$, and all Coulomb interactions responsible for the transition from one to two tails are screened out. Therefore, in a salty solution the two tail configuration disappears. Below we concentrate on the effect of screening on one tail or tail-less configurations only.
In a weak screening case, when $`r_sL_{2,c}`$, Coulomb interactions responsible for the appearance of the tail remain unscreened. Therefore, the lengths $`L_c`$ and $`L_{2,c}`$ remain almost unchanged. The large $`L`$ limit of $`L_1`$ however should be modified. At a very large tail length $`L_2`$ one should replace $`LL_1=L_2`$ by $`r_s`$ in Eq. (6) because the potential vanishes beyond the distance $`r_s`$. This gives
$$L_{1,\mathrm{}}(r_s)=+R\mathrm{ln}(/R)+(R^2/r_s)\mathrm{ln}(/R).$$
One can see that $`L_{1,\mathrm{}}`$ increases and charge inversion is stronger as $`r_s`$ decreases. This is because when $`r_s`$ decreases, the capacitance of the spherical complex increases, the self-energy of it decreases and it is easier to charge it.
When $`R<r_s<L_{2,c}`$, it is easy to show that the tail length, which appears at the phase transition, is equal to $`r_s`$ instead of $`L_{2,c}`$. This means that, before a tail is driven out at the phase transition, more PE condenses on the macroion in a salty solution than that for the salt free case. In other words, the critical point $`L_c`$ is shifted towards larger values:
$$L_c(r_s)=L_{1,\mathrm{}}(r_s)+r_s.$$
Obviously, $`L_c(r_s)>L_c`$ for $`r_s<L_{2,c}`$ and $`L_c(r_s)`$ approaches $`L_c`$ at $`r_sL_{2,c}`$. When $`r_s`$ approaches $`R`$, the critical length $`L_c(r_s)`$ reaches $`+2R\mathrm{ln}(/R)`$, so that the inverted charge is twice as large as that for the unscreened case.
At stronger screening, when $`r_s<R`$, to a first approximation, the macroion surface can be considered as a charged plane. The problem of adsorption of many rigid PE molecules on an oppositely charged plane has been studied in Ref. , where the role of Wigner crystal like correlations similar to that shown in Fig. 1 was emphasized. The large electrostatic rigidity of a strongly charged PE makes this calculation applicable to our problem as well. One can use results of Ref. in three different ranges of $`r_s`$: $`R>r_s>A`$, $`A>r_s>a`$, $`a>r_s`$. In all these ranges, the net charge $`Q^{}`$ of the macroion is proportional to $`R^2`$ instead of an almost linear dependence on $`R`$ in a salt free solution. The tail is not important for the calculation of the charge inversion ratio because it produces only a local effect near the place where the tail stems from the macroion. Inverted net charge $`Q^{}`$ grows with decreasing $`r_s`$, so that charge inversion ratio of the macroion reaches 100% at $`r_sA`$ and can become even larger at $`r_sA`$. For $`r_sA`$ , our results are in agreement with those of Ref. . One should be aware that $`|Q^{}|`$ ceases to increase at very small $`r_s`$. This is because at an extremely small $`r_s`$ such that the interaction between the macroion and one persistence length of the PE becomes less than $`k_BT`$, the PE desorbs from the macroion and the macroion becomes undercharged. Therefore, $`|Q^{}|`$ should reach a maximum at a very small $`r_s`$ and then decrease.
Finally, it should be noted that in the above discussion of the role of screening, we neglected the possibility of the condensation of the PE’s counterions on the sphere with inverted charge. This is valid for a large enough screening length because it is well known that in this case condensation does not occurs on a spherical macroion. Using $`Q^{}R\mathrm{ln}(/R)`$ and the standard condition for the condensation on a charged sphere , it is not difficult to show that the sphere is screened linearly if
$$r_s>R^{1\eta /2\eta _c}^{\eta /2\eta _c}.$$
When $`r_s<R`$, the macroion can be approximated as a charged plane and it is also known that a planar charge is linearly screened if the screening radius is small enough. Specifically, Eq. (73) of Ref. shows that screening is linear if
$$r_s<Ae^{\eta _c/\eta }\frac{R^2}{}e^{\eta _c/\eta }.$$
As we can see, when $`\eta `$ is less than $`\eta _c`$ by a logarithmic factor, i. e. when $`\eta <\eta _c/\mathrm{ln}(/R)`$, the range of $`r_s`$, where the macroion is nonlinear screened, almost vanishes. For $`\eta `$ of the order of $`\eta _c`$, however, there is a range of $`r_s`$ where counterion condensation on the charge-inverted sphere has to be taken into account and the sphere’s net charge is different from our estimate. There are two aspects of this counterion condensation phenomenon. Obviously, due to stronger nonlinear screening at the sphere surface, more PE collapses onto the sphere and the charge inversion ratio is even larger than what is predicted above in the linear screening theory. On the other hand, if one defines the net charge of the sphere as the sum of its bare charge, the charges of the collapsed PE monomers and the charges of all counterions condensed on it, the magnitude of this net charge is limited at the value given by the theory of counterion condensation on a sphere . As explained in Ref. , it is this charge that is observed in electrophoresis.
Until now we talked about a weakly charged PE with $`\eta \eta _c`$. In Ref. we studied adsorption of a strongly charged PE (for e.g., DNA) with $`\eta \eta _c`$ on positively charged plane. Such PE initiates Onsager-Manning counterion condensation both in the bulk and at the plane. The theory Ref. can be applied for the sphere at $`r_sR`$, too. It predicts a strong charge inversion which grows with decreasing $`r_s`$ and exceeds 100% at $`r_s<A`$.
## III Monte-Carlo simulations.
To verify the results of our analytical theory, we do Monte Carlo (MC) simulations. The PE is modeled as a chain of $`N`$ freely jointed hard spherical beads each with charge $`e`$ and radius $`a=0.2l_B`$ where $`l_B=7.12`$Å is the Bjerrum length at room temperature $`T_{rm}=298^\mathrm{o}\mathrm{K}`$ in water. The bond length is kept fixed and equal to $`l_B`$, so that our PE charge density $`\eta `$ is equal to the Manning condensation critical charge density $`\eta _c=k_BT_{rm}D/e`$. Due to the discrete nature of the simulated PE, in order to compare simulation results with theoretical predictions, we refer to the number of monomers $`N`$ as the PE length $`L`$ measured in units of $`l_B`$. The macroion is modeled as a sphere of radius $`4l_B`$ and with charge $`100e`$ uniformly distributed at its surface. To arrange the configuration of the PE globally, the pivot algorithm is used. In this algorithm, a part of the chain from a randomly chosen monomer to one of the chain ends is rotated by a random angle about a random axis (see Ref. and references therein). To relax the PE configuration locally, a flip algorithm is used. In this algorithm, a randomly selected monomer is rotated by a random angle about the axis connecting its two neighbors (if it is one of the end monomers, its new position is chosen randomly at a sphere of radius $`l_B`$ centered at its neighbor). The usual Metropolis algorithm is used to accept or reject the move. For a typical value of the parameters, we run about $`10^7`$ Monte Carlo steps and used the last 70% of them to obtain statistical averages (one Monte Carlo step is defined as the number of elementary moves such that, on average, every particle attempts to move once). Near the phase transition to the tail state, the number of steps is 5 times larger. The time for one run is typically 5 hours on an Athlon 1 Ghz computer. Assembler language is used to speed up the calculation time inside the inner loop of the program. Our code was checked by comparing with the results of Ref. and Ref. and some references therein.
Two different initial conformations of the PE are used to make sure that the system is in equilibrium. In the first initial conformation, the PE forms an equidistant coil around the macroion. In the second initial conformation, the PE makes a straight rod. Both initial conformations, within statistical uncertainty, give the same values for all the calculated properties of the systems such as the total energy, the end-end distance of the PE, the number of collapsed monomers and the critical length $`L_c`$.
An important aspect of the simulation is to determine the length of the tail and the amount of monomers residing at the macroion surface. In the literature, one usually defines a monomer as collapsed on the surface if it is found within a certain distance from it. This distance is arbitrarily chosen to be about two or three PE bond lengths. In the Appendix, we suggest an alternative more systematic method of determining the number of collapsed monomers.
Let us now describe the results of our Monte-Carlo simulations. We study the collapsed length $`L_1`$ as a function of $`L`$ for the case the macroion has radius $`R=4l_B`$ and charge $`Q=100e`$. This corresponds to $`/R=25`$, exactly the same value as the one used in Fig. 3. The result of our simulation is presented in Fig. 4 together with the theoretical curve of Fig. 3. The phase transition is observed at the chain length of 142 monomers and the critical tail length is about 16 monomers, which agrees very well with our predictions $`L_c=142`$ and $`L_{2c}=16`$.
We also study the case of a salty in solution. As everywhere in this paper, we assume that screening by monovalent salt can be described in the linear Debye-Hückel approximation. Therefore, in our simulation, we replace the Coulomb potential of the macroion $`Q/Dr`$ by the screened potential
$$V(r)=\frac{Qe^{R/r_s}}{1+R/r_s}\frac{e^{r/r_s}}{Dr},$$
(13)
where $`r_s`$ is the linear Debye-Hückel screening length. All PE monomers are still considered as point-like charges and Yukawa potential, $`r^1e^{r/r_s}`$, is used to describe their interaction. The result of our simulation for the case $`r_s=5l_B`$ is plotted by the solid square in Fig. 4. As predicted above, screening increases the maximum charge inversion ratio to 63%.
Simulation at $`r_s=4l_B`$ shows even bigger charge inversion with 70% ratio. This suggests that the maximum in charge inversion is located at even smaller screening radius. However, we did not try to run the simulation at smaller $`r_s`$ in order to find the maximum in charge inversion because, at smaller $`r_s`$, the identification of adsorbed monomers becomes less unambiguous.
The better-than-expected agreement between MC results and theoretical prediction of the critical length $`L_c`$ for the $`r_s=\mathrm{}`$ case is somewhat accidental because in Fig. 4 we compared a zero temperature theory with a finite temperature Monte-Carlo simulation. The temperature affects $`L_c`$ because tail monomers have smaller entropy compared to collapsed monomers. The self repulsion of the tail and the repulsion from the overcharged sphere limits the configuration space of the tail monomers, while at the macroion, the PE self-energy is screened at the distance $`A`$, so that the collapsed monomers have larger configuration space. Therefore, the free energy is gained when more monomers collapse on the macroion surface. This helps to push the critical length $`L_c`$ to a higher value than its value at zero temperature.
For clarifying the role of temperature, we carry out simulations at different $`T`$ and extrapolate $`L_c`$ to $`T=0^\mathrm{o}\mathrm{K}`$. The results are shown in Fig. 5. The extrapolated $`L_c`$ is 134 which is 6% lower than the zero temperature theoretical prediction of 142. Also from this figure, one can see that the temperature dependence of $`L_c`$ is linear. A simple analysis of the Monte-Carlo data shows that the entropy per monomers gained at the surface is about 2 at the critical point.
On Fig. 6 we show two typical snapshots of the system, one for the case $`L=141`$ (before the phase transition) and the other for $`L=143`$ (after the phase transition). They again confirm that the tail appears abruptly near $`L=L_c=142`$. One can clearly see an important aspect of the correlation effects: PE segments of different turns stay away from each other and locally, they resemble a one dimensional Wigner crystal, which helps to lower the energy of the system. Globally, however, the PE conformation resembles that of a tennis ball instead of a solenoid. This obvious difference between observed conformation and the theoretical solenoid-like ground state is also related to thermal fluctuations. Solenoid structure is subjected to low energy long range bending modes, energy of which is proportional to $`k^4`$, where $`k`$ is the wave vector of such mode. It is easy to show that at the room temperature with our parameters of the system, modes with $`kR^1`$ are strongly excited and they “melt” the solenoid. However, modes with large $`k`$ are not excited and, therefore, the short range order between PE turns is preserved. This leads to a compromised “tennis ball” conformation instead of a solenoid. The difference in energy between a “tennis ball” and a solenoid conformation, however, is small compare to the interaction between the sphere and the PE. This helps to explain the small difference between the results of the finite temperature Monte-Carlo simulation and our zero temperature theory.
Monte-Carlo results similar to Fig. 4 for unscreened case were independently obtained in Ref. . For the screened case, however, the authors of Ref. claimed that charge inversion reaches maximum when $`r_s3R`$ which is still very large, much larger than what is observed in our simulations. This is because instead of the Overbeck potential (13), the authors of Ref. use the Yukawa potential $`Qr^1e^{r/r_s}`$ for the macroion, where $`r`$ is the distance to the center of the macroion. This means that they put the net charge of the macroion at the center and screen it inside the macroion body. As a result, the apparent surface charge of the macroion becomes very small and charge inversion disappears. New simulations carried out by the same authors using the proper potential (13) are in agreement with our theory and Monte Carlo simulations.
Before concluding this paper, we would like to mention that in our simulation, counterion condensation on the sphere with inverted charge was neglected. As stated in the end of Sec. II, this is valid if $`\eta \eta _c`$. In our Monte-Carlo simulations $`\eta `$ is equal to $`\eta _c`$ therefore, in order to study the effect of screening on charge inversion, we choose to simulate the system at small $`r_sA`$ where condensation is not very important.
In conclusion, we have studied charge inversion for the complexation of a PE with a spherical macroion. We started from description of the correlated ground state configuration of PE at the macroion surface instead of smearing of the PE charge at the surface. As a result, we have eliminated the unphysical finite charge adsorption at the neutral sphere. Our Monte-Carlo simulations confirm that correlations are the driving force of charge inversion.
###### Acknowledgements.
The authors are grateful to A. Yu. Grosberg for many useful discussions and to S. Stoll and P. Chodanowski for the possibility to read their paper before the publication. This work was supported by NSF DMR-9985985.
## A The number of collapsed monomers of polyelectrolyte.
To better determine the number of collapsed monomers in Monte-Carlo simulation, we use the following procedure.
Firstly, we draw the histogram of the number of monomers found within a distance $`r`$ from the macroion surface during a simulation run. Up to a normalizing factor, this histogram is nothing but the probability $`P_r(n)`$ of finding $`n`$ monomers within a distance $`r`$ from the sphere surface. Secondly, at a given $`r`$, we define the value of $`n`$ corresponding to the maximum in this histogram as the most likely number of monomers $`n(r)`$ found inside the distance $`r`$ from the macroion surface.
Now, we show that much can be learned by plotting $`n(r)`$ as a function of $`r`$. In Fig. 7a, $`n(r)`$ is plotted for two typical cases of $`L=140`$ (before the phase transition) and $`L=150`$ (after the phase transition).
Clearly, as one goes away from the macroion surface, or $`r`$ grows, at first one see a rapid increase in the number of monomers $`n(r)`$ found. After a distance of about two bond lengths, this increase is slowed and stopped. It is easy to identify the first range of $`r`$, where one observes a rapid increase in $`n(r)`$ as the collapsed layer. For the case of $`L=140`$, as $`r`$ increases beyond this layer, $`n(r)`$ is always equal to the total number of monomers $`N=140`$. This is the indication of a collapsed state where all PE monomers lie in the collapsed layer near the macroion surface. The situation is completely different in the case of $`L=150`$ where beyond the collapsed layer one sees a linear increase in $`n(r)`$ until $`r=19`$ (not shown) where $`n(r)`$ saturates at the maximum possible value of 150. This is an indication of a tailed state. The slope of the increase in this second range also provides a valuable information on the conformation of the tailed state. As one can see, this slope is very close to unity, what clearly indicates a one-tail state. This is in agreement with our prediction that after the phase transition the complex is in one tail state and in disagreement with the conclusion of Ref. that the system fluctuates between one tail and two tail conformations (for a two tail state the slope would be 2).
A closer look at the tail part of Fig. 7 shows that the slope of the tail part of $`n(r)`$ actually is slightly larger than unity and grows with $`r`$. This could be expected. The PE tail near the overcharged macroion is strongly stressed in the electric field of the macroion’s inverted charge. Farther from the macroion, this electric field is weaker and due to the thermal motion of the monomers, more than one monomer can be found as $`r`$ increases by one bond length.
The final step in determining the number of collapsed monomers in the tailed state is accomplished by fitting the tail part of $`n(r)`$ by an empirical quadratic equation $`ax^2+bx+c`$. The intersection of this curve with the $`y`$ axis gives the number of collapsed monomers or the collapsed length $`L_1`$. For e.g., at $`L=150`$, the value of $`a`$, $`b`$ and $`c`$ are 0.01, 1.22 and 125 respectively (see Fig. 7b), so that the slope at the macroion surface $`x=0`$ is $`1.22`$ and the amount of collapsed monomers is 125. Also, as $`L`$ increase, the tail gets longer and becomes more stressed due to its self-energy, the slope of $`n(r)`$ decreases and is closer to 1. The fitted value for $`b`$ are 1.32 at $`L=145`$, 1.22 at $`L=150`$ and 1.09 at $`L=165`$.
Near the phase transition point $`L=142`$ one sees two maxima in the histogram $`P_r(n)`$ instead of one. One of these maxima behaves exactly as that of one tail configuration (linearly increases with $`r`$ after the collapsed layer). The other maximum behaves exactly as that of the collapsed state (constant and equal the total number of monomers $`N=142`$ after the collapsed layer). This is because near the phase transition the PE fluctuates between the collapsed and the tailed state.
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# A note on extended complex manifolds
## Abstract
We introduce a category of extended complex manifolds, and prove that the functor describing deformations of a classical compact complex manifold $`M`$ within this category is versally representable by (an analytic subspace in) $`𝐇^{}(M,T_M)`$.
By restricting the associated versal family of extended complex manifolds over $`𝐇^{}(M,T_M)`$ to the subspace $`𝐇^1(M,T_M)`$ one gets a correct limit to the classical picture.
0. Introduction. Recent advances in homological mirror symmetry \[Ba, BK, Ko1, Ma1, Ma2, P, PZ\] give a strong evidence for the existence of extended versions of at least three classical categories in algebraic geometry — the categories of compact complex manifolds, of coherent sheaves, and of submanifolds. One such category (of extended special Lagrangian submanifolods) has been constructed in \[Me1\]. In this note we suggest an extension of the category of complex manifolds.
Curiously, even with no appropriate notion of extended complex structure at hand, one can nevertheless say a lot about their moduli spaces \[Ba, BK, Ma1, Me2\]. This egg-before-chicken situation is due to a remarkably effective, purely algebraic paradigm of the modern deformation theory (see \[GM, Ko2\] and references therein):
* Given a mathematical structure one wishes to deform, the first step should be a search for a differential $``$-graded (dg-, for short) Lie $`k`$-algebra $`(𝔤=_i𝔤^i,d,[,])`$ which “controls” the deformations. Next one defines a deformation functor
$$\begin{array}{cccc}\hfill \mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}:& \left\{\begin{array}{c}\text{the category of Artin}\hfill \\ k\text{-local algebras}\hfill \end{array}\right\}& & \left\{\text{the category of sets}\right\}\end{array}$$
as follows
$$\mathrm{𝖣𝖾𝖿}_𝔤^0(A)=\{\mathrm{\Gamma }(𝔤m_A)^1d\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma },\mathrm{\Gamma }]=0\}/\mathrm{exp}(𝔤m_A)^0,$$
where $`m_A`$ is the maximal ideal of the Artin algebra $`A`$, the latter is viewed as a $``$-graded algebra concentrated in degree zero (so that $`(𝔤m_𝒜)^i=𝔤^im_A`$), and the quotient is taken with respect to the following representation of the gauge group $`\mathrm{exp}(𝔤m_A)^0`$,
$$\mathrm{\Gamma }\mathrm{\Gamma }^g=e^{\mathrm{a}d_g}\mathrm{\Gamma }\frac{e^{ad_g}1}{\mathrm{a}d_g}dg,g(𝔤m_A)^0,$$
where $`\mathrm{a}d`$ is just the usual internal automorphism of $`𝔤`$, $`\mathrm{a}d_g\mathrm{\Gamma }:=[g,\mathrm{\Gamma }]`$. Finally one tries to represent the deformation functor by a topological (pro-Artin) algebra $`𝒪_S`$ so that
$$\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}(A)=\text{Hom}_{\mathrm{c}ont}(𝒪_S,A).$$
This associates to the given mathematical structure a formal moduli space $`S`$ whose “ring of functions” is $`𝒪_S`$.
The tangent space, $`\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}(k[\epsilon ]/\epsilon ^2)`$, to the functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}`$ is isomorphic to the cohomology group $`𝐇^1(𝔤)`$ of the complex $`(𝔤,d)`$. If one extends in the obvious way the above deformation functor to the category of arbitrary $``$-graded $`k`$-local Artin algebras (which may not be concentrated in degree 0), one gets the functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ with the tangent space isomorphic to the full cohomology group $`𝐇^{}(𝔤)=_i𝐇^i(𝔤)`$.
The dg-Lie algebra controlling deformations of a given complex structure on a compact manifold $`M`$ is $`(𝔤:=_i\mathrm{\Gamma }(M,T_M\mathrm{\Omega }_M^{0,i}),\overline{})`$, where $`T_M`$ is the tangent holomorphic sheaf, and $`\mathrm{\Omega }_M^{0,i}`$ the sheaf of $`(0,i)`$ forms. The classical deformation functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}`$ is known to be versally representable by the Kuranishi analytic subspace in $`𝐇^1(𝔤)=𝐇^1(M,T_M)`$. Its extension, $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$, is the main technical tool for introducing and studying the moduli space of so called extended complex structures\[BK\]. Actually, the authors of \[BK\] go even further and study the extended deformation functor associated with the larger dg-Lie algebra, $`(\widehat{𝔤}=\mathrm{\Gamma }(M,^{}T_M\mathrm{\Omega }_M^{0,}),\overline{})`$. They prove that $`\mathrm{𝖣𝖾𝖿}_{\widehat{𝔤}}^{}`$ is non-obstructed provided $`M`$ is a Calabi-Yau manifold, and show that the associated extended moduli space, $`S_{\mathrm{e}xt}𝐇^{}(M,^{}T_M)`$, has an induced structure of Frobenius manifold \[BK, Ba\]. For arbitrary $`M`$, the moduli space $`S_{\mathrm{e}xt}`$ is canonically an $`F_{\mathrm{}}`$-manifold \[Me2\].
There is, however, an obvious problem with the above approach as it offers no geometric explanation of what this extended complex structure might be. It is an urgent and important problem \[Ma2\] to find a geometric embodiment of this notion, develop its deformation theory, and check that the base of the resulting versal deformation can be canonically identified with what one gets from the above purely algebraic approach. The present note offeres a realization of this programme for the extended deformation functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$. It is unlikely, however, that the same approach will yield a geometrical model for the functor $`\mathrm{𝖣𝖾𝖿}_{\widehat{𝔤}}^{}`$ — this seems to be a much more intricate object.
1. Extended complex manifolds. These will be defined in two steps. First comes the notion of pre-complex manifold.
1.1. Model pre-complex structures. Let $`U`$ be an open domain in $`^{2n|0}`$. A pre-complex structure on $`U`$ is a map
$$\varphi :U^{2n|n}$$
such that
* $`\varphi `$ is a smooth embedding, and
* at each $`xU`$, the completions of the associated stalks of the sheaves $`C^{\mathrm{}}(U)`$ and $`\varphi ^1(𝒪_{^{2n|n}}/𝗇ilpotents)`$ are isomorphic.
Here $`C^{\mathrm{}}(U)`$ stands for the sheaf of smooth complex valued functions on $`U`$, and $`𝒪_{^{2n|n}}`$ for the sheaf of holomorphic functions on the supermanifold $`^{2n|n}`$. The ringed space
$$𝒰^n=\left(U,𝒪_U^{\mathrm{p}re}:=\varphi ^1(𝒪_{^{2n|n}})\right)$$
is called a pre-complex domain of dimension$`n`$. The reduced sheaf $`𝒪_U^{\mathrm{p}re}/𝗇ilpotents`$ is denoted sometimes by $`𝒪_{U,\mathrm{r}d}^{\mathrm{p}re}`$.
A morphism of pre-complex domains $`𝒰^n𝒱^m`$ is, by definition, a pair $`(f_C^{\mathrm{}},[f])`$, consisting of a smooth map $`f_C^{\mathrm{}}:UV`$ and a germ, $`[f]`$, of holomorphic maps $`f`$ from a small open neighbourhood of $`𝖨m(\varphi _U)`$ in $`^{2n|n}`$ to a small open neighbourhood of $`𝖨m(\varphi _V)`$ in $`^{2m|m}`$ such that, for each $`xU`$,
$$f_C^{\mathrm{}}^{}\left(\overline{\overline{𝒪_{V,\mathrm{r}d,x}^{\mathrm{p}re}}}\right)=\overline{\overline{f^{}(𝒪_{V,\mathrm{r}d,x}^{\mathrm{p}re})}},$$
where $`\overline{\overline{}}`$ stands for the completion of the stalks at $`x`$ with respect to the natural ideals. We often abbreviate $`(f_C^{\mathrm{}},[f])`$ simply to $`f`$.
1.2. Fact. Any open domain $`U^{2n|0}`$ admits a pre-complex structure. Indeed, identifying $`^{2n|0}`$ with $`^{n|0}`$ one immediately gets a required map $`\varphi `$ by further identifying $`C^{n|0}`$ with the “diagonal” subspace in $`^{2n|n}`$ given, in the natural coordinates $`(z^\alpha ,z^{\dot{\alpha }},\psi ^{\dot{\alpha }})`$, $`\alpha ,\dot{\alpha }=1,\mathrm{},n`$, on $`^{2n|n}`$, by the equations
$$z^{\dot{\alpha }}=\overline{(z^\alpha )},\psi ^{\dot{\alpha }}=0,$$
where the bar denotes complex conjugation.
1.3. Definition. An n-dimensional pre-complex manifold is a ringed space $`=(M,𝒪_M^{\mathrm{p}re})`$ modeled on pre-complex domains of dimension $`n`$. It is called compact if the underlying smooth manifold $`M`$ is compact.
A holomorphic vector field $`v`$ on $``$ is, as usual, a $``$-linear automorphism of the structure sheaf,
$$v:𝒪_M^{\mathrm{p}re}𝒪_M^{\mathrm{p}re},$$
whose restriction to each stalk, $`𝒪_{M,x}^{\mathrm{p}re}`$, is a derivation of the ring $`𝒪_{M,x}^{\mathrm{p}re}`$. The sheaf of holomorphic vector fields on $``$ is denoted by $`𝒯_{}`$.
1.4. Definition. An extended complex manifold is a pair $`(,_{})`$ consisting of a compact pre-complex manifold $``$ and an odd holomorphic vector field $`_{}`$ such that $`[_{},_{}]=0`$<sup>1</sup><sup>1</sup>1Odd vector fields with this property are often called homological. and $`dim𝐇(,_{})<\mathrm{}`$. Here
$$𝐇(,_{}):=\frac{\mathrm{\Gamma }(M,𝖪er[_{},\mathrm{}])}{\mathrm{\Gamma }(M,𝖨m[_{},\mathrm{}])},$$
with the differential given by
$$\begin{array}{cccc}\hfill [_{},\mathrm{}]:& 𝒯_{}& & 𝒯_{}\\ & V& & [_{},V].\end{array}$$
Note that $`𝐇(,_{})`$ is canonically a Lie superalgebra with the brackets, $`[,]_𝐇`$, induced from the usual commutator of holomorphic vector fields.
One can also associate with $`(,_{})`$ the cohomology manifold, $`(,_{})`$, which is, by definition, a ringed space $`(M,𝖪er_{}/𝖨m_{})`$.
A morphism of extended complex manifolds,
$$f:(,_{})(𝒮,_𝒮),$$
is a morphism of the associated pre-complex manifolds, $`f:(M,𝒪_M^{\mathrm{p}re})(S,𝒪_M^{\mathrm{p}re})`$, which commutes with the homological vector fields, i.e.
$$_{}\left(f^{}(g)\right)=f^{}(_𝒮g)$$
for any $`g𝒪_S^{\mathrm{p}re}`$. One may reformulate this as $`f_{}(_{})=_𝒮`$.
1.5. Basic example. Let $`M`$ be a compact complex manifold. There is associated an extended complex manifold $`(=(M,𝒪_M^{\mathrm{p}re}),\overline{})`$ constructed as follows:
* First consider a natural embedding,
$$\begin{array}{cccc}\hfill \varphi :& M& & M\times \mathrm{\Pi }T_{\overline{M}}\\ & x& & (x,O(x))\end{array}$$
where $`T_{\overline{M}}`$ is the total space of the holomorphic vector bundle over the conjugate complex manifold $`\overline{M}`$, $`\mathrm{\Pi }`$ the parity change functor, and $`O:\overline{M}\mathrm{\Pi }T_{\overline{M}}`$ the zero section.
* Set
$$𝒪_M^{\mathrm{p}re}:=\varphi ^1\left(𝗍hestructuresheafonM\times \mathrm{\Pi }T_{\overline{M}}\right).$$
* Note that $`\mathrm{\Pi }T_{\overline{M}}`$ is canonically a dg-manifold (of dimension $`n|n`$) with the homological field being just the $`(1,0)`$-part of the de Rham differential on $`\overline{M}`$. The latter induces a homological vector field on $``$ which we denote by $`\overline{}`$.
1.5.1. Remark. Let $`T_M`$ be the holomorphic tangent sheaf on $`M`$. We make the cohomology $`𝐇^{}(M,T_M)=_k𝐇^k(M,T_M)`$ into a superspace by setting
$$𝐇^{}(M,T_M)_{\stackrel{~}{0}}=\underset{k𝗂sodd}{}𝐇^k(M,T_M),𝐇^{}(M,T_M)_{\stackrel{~}{1}}=\underset{k𝗂seven}{}𝐇^k(M,T_M).$$
This choice of $`_2`$-grading is in agreement with the classical deformation theory where $`𝐇^1(M,T_M)`$ is even. Then the natural map,
$$\begin{array}{ccc}T_M\mathrm{\Omega }_M^{0,p}\times T_M\mathrm{\Omega }_M^{0,q}& & T_M\mathrm{\Omega }_M^{0,p+q}\\ X\alpha \times Y\beta & & [X,Y](\alpha \beta )\end{array}$$
induces on $`𝐇^{}(M,T_M)`$ the structure of an odd Lie superalgebra (cf. \[Ma1\]). Reversing the parity, we make $`\mathrm{\Pi }𝐇^{}(M,T_M)`$ into a Lie superalgebra.
1.5.2. Lemma. Let $`M`$ be a complex manifold, and $`(,\overline{})`$ the associated extended complex manifold. Then
* $`(,\overline{})`$ is precisely $`M`$ with its original complex structure.
* $`𝐇(,\overline{})=\mathrm{\Pi }𝐇^{}(M,T_M)`$ as Lie superalgebras.
Proof. The statement (a) follows immediately from the Poincare $`\overline{}`$-lemma. The second statement requires a small computation. Let $`\{z^\alpha \}`$ be a local coordinate system on $`M`$, and $`\{z^{\dot{\alpha }},\psi ^{\dot{\alpha }}=dz^{\dot{\alpha }}\}`$ the associated local coordinate system on $`\mathrm{\Pi }T_{\overline{M}}`$. The collection $`\{z^\alpha ,z^{\dot{\alpha }},\psi ^{\dot{\alpha }}\}`$ gives rise to a coordinate chart on $`(M,𝒪_M^{\mathrm{p}re})`$ so that any $`V\mathrm{\Gamma }(M,𝒯_{})`$ can be locally represented as
$$V=\underset{\alpha =1}{\overset{n}{}}V^\alpha \frac{}{z^\alpha }+\underset{\dot{\alpha }=1}{\overset{n}{}}\left(V^{\dot{\alpha }}\frac{}{z^{\dot{\alpha }}}+W^{\dot{\alpha }}\frac{}{\psi ^{\dot{\alpha }}}\right),$$
for some local sections, $`V^\alpha ,V^{\dot{\alpha }},W^{\dot{\alpha }}`$, of $`𝒪_M^{\mathrm{p}re}`$. As
$$\overline{}=\underset{\dot{\alpha }=1}{\overset{n}{}}\psi ^{\dot{\alpha }}\frac{}{z^{\dot{\alpha }}},$$
we have
$$[\overline{}_0,V]=\underset{\alpha =1}{\overset{n}{}}(\overline{}V^\alpha )\frac{}{z^\alpha }+\underset{\dot{\alpha }=1}{\overset{n}{}}\left(\overline{}V^{\dot{\alpha }}(1)^{\stackrel{~}{V}}W^{\dot{\alpha }}\right)\frac{}{z^{\dot{\alpha }}}+\underset{\dot{\alpha }=1}{\overset{n}{}}(\overline{}W^{\dot{\alpha }})\frac{}{\psi ^{\dot{\alpha }}}.$$
Thus $`V\mathrm{\Gamma }(M,𝖪er[\overline{},\mathrm{}])`$ if and only if $`\overline{}V^\alpha =0`$, $`W^{\dot{\alpha }}=(1)^{\stackrel{~}{V}}\overline{}V^{\dot{\alpha }}`$ and $`V^{\dot{\alpha }}`$ arbitrary. Moreover, the equivalence class $`Vmod\mathrm{\Gamma }(M,𝖨m[\overline{},\mathrm{}])`$ always has a representative of the form
$$\underset{\alpha =1}{\overset{n}{}}V^\alpha \frac{}{z^\alpha }$$
where $`V^\alpha `$ are defined uniquely up to addition of a $`\overline{}`$-exact term. The Dolbeault theorem completes the proof. $`\mathrm{}`$
2. Deformation theory. We start with the deformation theory of dg-manifolds and then apply the technique to extended complex manifolds.
2.1. Dg-manifolds. An extended complex manifold is a particular case of a complex differential $`_2`$-graded (dg-, for short) manifold which is, by definition, a complex supermanifold equipped with a homological holomorphic vector field (cf. \[Ko2\]). Morphisms of dg-manifolds,
$$(𝒳,_𝒳)\stackrel{f}{}(𝒮,_𝒮),$$
are defined as in sect. 1.4. It is easy to see that the resulting category is closed with respect to the fibered products. Note that the fibres of $`f`$ are not, in general, dg-manifolds except over the points where $`_𝒮`$ vanishes. In this context we define a pointed dg-manifold as a triple $`(𝒮,_𝒮,)`$, where $`(𝒮,_𝒮)`$ is a (formal) dg-manifold and $``$ a point in $`𝒮`$ such that $`_𝒮I_{}I_{}^2`$, $`I_{}`$ being the ideal sheaf of $``$.
2.1.1. Remark. According to Kontsevich \[Ko2\], any Lie superalgebra structure, $`[,]`$, on a vector superspace $`𝔤`$ can be equivalently interpreted as a quadratic homological vector field on $`(\mathrm{\Pi }𝔤,0)`$ viewed as a pointed formal supermanifold. Thus $`𝐇^{}(M,T_M)`$ in example 1.5, and $`\mathrm{\Pi }𝐇(,_{})`$ in 1.4 are canonically pointed dg-manifolds.
2.1.2. Lemma. Let $`(𝒮,_𝒮,)`$ be a pointed dg-manifold. Then $`\mathrm{\Pi }𝒯_{𝒮,}`$, the tangent space at $``$ with the reversed parity, is canonically a Lie superalgebra.
Proof. Let $`v_1`$ and $`v_2`$ be two tangent vectors at $`𝒮`$, and $`V_1`$ and $`V_2`$ any two germs of holomorphic vector fields such that $`V_1|_{}=v_1`$ and $`V_2|_{}=v_2`$. The bilinear skew-symmetric operation,
$$\begin{array}{cccc}\hfill [,]_{}:& \mathrm{\Pi }𝒯_{𝒮,}\times \mathrm{\Pi }𝒯_{𝒮,}& & \mathrm{\Pi }T_{𝒮,}\\ & \mathrm{\Pi }v_1\times \mathrm{\Pi }v_2& & \mathrm{\Pi }[V_1,[_𝒮,V_2]]_{}\end{array}$$
is well defined. The identity $`[_𝒮,_𝒮]=0`$ implies the Jacobi identity for $`[,]_{}`$. $`\mathrm{}`$
2.2. Deformations. A deformation of a complex dg-manifold $`(,_{})`$ is, by definition, a morphism,
$$(𝒳,_𝒳)\stackrel{f}{}(𝒮,_𝒮,)$$
from a complex dg-manifold $`(𝒳,_𝒳)`$ to a pointed complex dg-base $`(𝒮,_𝒮,)`$ such that
* $`(f^1(),_𝒳|_{f^1()})\stackrel{i}{}(,_{})`$, and
* the associated morphism of complex supermanifolds, $`f:𝒳𝒮`$ is locally trivial in an neighbourhood of $`S`$.
Given two deformations, $`f:(𝒳,_𝒳)(𝒮,_𝒮,)`$ and $`\stackrel{~}{f}:(\stackrel{~}{𝒳},_{\stackrel{~}{𝒳}})(\stackrel{~}{𝒮},_{\stackrel{~}{𝒮}},)`$, of the same dg-manifold manifold $``$, a morphism from the first to the second, is, by definition, a commutative diagram
where $`m`$ (resp. $`s`$) is a morphism of (resp. pointed) dg-manifolds. If $`s`$ is the identity morphism and $`m`$ is an isomorphism, then the resulting morphism is called an equivalence of deformations.
As usually, one defines a deformation of a smooth dg-manifold over a germ of pointed (possibly, singular) dg-bases. From now on we consider only such deformations.
2.3. Remark. It may look puzzling that we allow in defintion 2.2 the homological vector field $`_𝒮`$ to be non-zero — the fibres of $`f`$ over the points where $`_𝒮0`$ are not, in general, dg-manifolds. A more natural definition of deformation would be a version of 2.2 with $`_𝒮0`$, where the base is simply a pointed analytic superspace (let us term these $`0`$-deformations). There are two advantages with understanding a deformation as in 2.2:
* By allowing $`_𝒮0`$, we do not loose generality; $`0`$-deformation is a deformation.
* The $`0`$-deformation theory of $`(,_{})`$ is, in general, obstructed, and its versal moduli space $`(𝒱_{𝗏ersal},)`$, when it exists, is a singular analytic space. The role of $`_𝒮`$ in 2.2 is to overcome all these obstructions producing thereby a smooth versal moduli dg-space, $`(𝒮,_𝒮,)`$, of which $`𝒱_{𝗏ersal}`$ is merely a (singular) analytic subspace given by the equations $`_𝒮=0`$ (cf. \[Me2\]). In a sense, the deformation theory 2.2 is a smooth resolution of the more natural $`0`$-deformation theory.
2.4. Proposition. Let $`f:(𝒳,_𝒳)(𝒮,_𝒮,)`$ be a deformation of a dg-manifold $`(,_{})`$. There is a canonical (even) morphism of Lie superalgebras,
$$Df:(\mathrm{\Pi }𝒯_{𝒮,},[,]_{})(𝐇(,_{}),[,]_𝐇).$$
Proof. Let is first construct an odd linear morphism,
$$df:𝒯_{𝒮,}𝐇(,_{}),$$
of vector superspaces. Fixing a local trivialization, $`\varphi :𝒳𝒮\times `$, we may decompose<sup>2</sup><sup>2</sup>2We apologize for being a bit sloppy in formulating this decomposition.,
$$_𝒳=_{}+_𝒮+\mathrm{\Gamma },$$
for some $`f`$-vertical odd holomorphic field $`\mathrm{\Gamma }`$ which vanishes at $`\times `$. Then the homology condition $`[_𝒳,_𝒳]=0`$ translates into the Maurer-Cartan(-like) equations,
$$[_{},\mathrm{\Gamma }]+[_𝒮,\mathrm{\Gamma }]+\frac{1}{2}[\mathrm{\Gamma },\mathrm{\Gamma }]=0.$$
For a $`v`$ in $`𝒯_{𝒮,}`$, we define a global holomorphic vector field on $``$,
$$df(v):=[V,\mathrm{\Gamma }]_{f^1()}$$
where $`V`$ is an arbitrary extension of $`v`$ to a germ of holomorphic vector field on $`(𝒮,)`$. The Maurer-Cartan equations and the fact that $`_𝒮`$ has zero at $``$ of second order imply
$$[_{},df(v)]=0.$$
Moreover, the cohomology class of $`df(v)`$ in $`𝐇(,_{})`$ does not depend on the choice of the trivialization $`\varphi `$ so that the map $`df`$ is well-defined.
Analogously, for any $`v_1,v_2𝒯_{𝒮,}`$ we have
$$[V_1,[V_2,[_{},\mathrm{\Gamma }]]]+[V_1,[V_2,[_𝒮,\mathrm{\Gamma }]]]+\frac{1}{2}[V_1,[V_2,[\mathrm{\Gamma },\mathrm{\Gamma }]]]=0$$
implying
$$[[V_1,[_𝒮,V_2]],\mathrm{\Gamma }]_{f^1()}+[[V_1,\mathrm{\Gamma }]_{f^1()},[V_2,\mathrm{\Gamma }]_{f^1()}]=0mod𝖨m[_{},\mathrm{}].$$
Thus $`Df:=df\mathrm{\Pi }`$ is a morphism of Lie superalgebras. $`\mathrm{}`$
2.5. Versality. If $`f:(𝒳,_𝒳)(𝒮,_𝒮,)`$ is a deformation of a dg-manifold $`(,_{})`$, and $`g:(\stackrel{~}{𝒮},_{\stackrel{~}{𝒮}},)(𝒮,_𝒮,)`$ is a morphism of germs of pointed dg-spaces, then the fibred product,
$$(𝒳,_𝒳)\times _{(𝒮,_𝒮)}(\stackrel{~}{𝒮},_{\stackrel{~}{𝒮}}),$$
gives rise to the induced deformation, $`g^{}(f)`$, of $``$ over the germ $`(\stackrel{~}{𝒮},_{\stackrel{~}{𝒮}},)`$.
A deformation $`f:(𝒳,_𝒳)(𝒮,_𝒮,)`$ is called versal if every other deformation $`f:(\stackrel{~}{𝒳},_{\stackrel{~}{𝒳}})(\stackrel{~}{𝒮},_{\stackrel{~}{𝒮}},)`$ of the same extended complex manifold is equivalent to the inverse image, $`g^{}(f)`$, of $`f`$ under some morphism of germs, $`g:(\stackrel{~}{𝒮},_{\stackrel{~}{𝒮}},)(𝒮,_𝒮,)`$, of pointed dg-spaces. A versal deformation $`f`$ is called minimal if the associated morphism of Lie superalgebras $`Df`$ (see Proposition 2.4) is an isomorphism. Any two minimal versal deformations of $`(,_{})`$ are isomorphic.
2.6. Theorem. Let $`(,_{})`$ be a dg-manifold with $`dim𝐇(,_{})<\mathrm{}`$. Then there exists a smooth minimal versal deformation of $`(,_{})`$.
Proof. This is a special case of the Smoothness Theorem 2.5.6 in \[Me2\]. $`\mathrm{}`$
2.7. Example. To any (compact) smooth manifold $`M`$ one may associate a dg-manifold,
$$(=\mathrm{\Pi }T_M,_{}=𝖽eRhamdifferential).$$
As $`𝐇(,_{})`$ always vanishes (an easy exercise), this dg-manifold is rigid, i.e. its any deformation is trivial. This is in accord with the topological nature of the example.
2.7. Deformation theory of extended complex manifolds. This is an obvious modification of the deformation theory of dg-manifolds with all the notions and results from 2.1–2.6 holding true. For example, a deformation of an extended complex manifold $`(,_{})`$ is a proper morphism of complex dg-manifolds,
$$(𝒳,_𝒳)\stackrel{f}{}(𝒮,_𝒮,)$$
such that
* for each $`t𝖹eros(_𝒮)`$ the fibre $`(f^1(t),_𝒳|_{f^1(t)})`$ is an extended complex manifold,
* $`(f^1(),_𝒳|_{f^1()})\stackrel{i}{}(,_{})`$, and
* the associated morphism of complex supermanifolds, $`f:𝒳𝒮`$ is locally trivial in an neighbourhood of $`S`$.
2.7.1. Remarks. (i) In view of 1.1(ii), the local triviality condition 2.7(c) ensures that the underlying smooth structure of $``$ keeps unchanged upon deformation. It is only the homological vector field that undergoes deformation.
(ii) The induced morphism, $`f:(𝒳,_𝒳)((𝒮,_𝒮),)`$, may not be locally trivial.
2.7.2. Theorem. Let $`M`$ be a compact complex manifold and $`(,\overline{})`$ the associated extended complex manifold. Then there exists a minimal versal deformation of the latter of the form
$$(𝒳,_𝒳)\stackrel{f}{}(𝐇^{}(M,T_M),,0).$$
Moreover,
* the homological vector field $``$ is an invariant of the original complex manifold $`M`$.
* The restriction of the family $`f`$ to the analytic subspace $`𝖹eros()𝐇^1(M,T_M)`$ is equivalent to the classical Kuranishi versal deformation of $`M`$.
* If $`M`$ is a Calabi-Yau manifold, then $`0`$ and the base of the above versal deformation can be canonically identified with the Barannikov-Kontsevich moduli space of (partially) extended complex structures.
Proof. The existence of a minimal versal deformation can be infered directly from Lemma 1.5.2 and Theorem 2.6. Nevertheless, we will give a detailed proof of this statement which makes all other claims, (i)–(iii), almost obvious. The idea of the proof is to show that, though the solution space, $`\widehat{𝖬C}`$, of Maurer-Cartan equations arising in our extended deformation theory is much larger than the solution space, $`𝖬C`$, of the Maurer-Cartan equations in the original purely algebraic $`(\mathrm{\Gamma }(M,T_M\mathrm{\Omega }_M^{0,}),\overline{})`$-approach, the gauge group turns out to be much larger as well, and, crucially, at the level of quotients,
$$\frac{\widehat{𝖬C}}{\widehat{𝗀augegroup}}=\frac{𝖬C}{𝗀augegroup},$$
we have a canonical isomorphism.
Let $`f:(𝒳,_𝒳)(𝒮,_𝒮,)`$ be a deformation of $`(,\overline{})`$. We may assume without loss of generality that $`(𝒮,_𝒮,)`$ is dual to a differential Artin superalgebra $`(A,_A)`$, cf. \[Me2\]. As in the proof of Proposition 2.4, we fix a local trivialization, $`\varphi :𝒳𝒮\times `$, and decompose<sup>2</sup>,
$$_𝒳=_{}+_𝒮+\mathrm{\Gamma },$$
where $`\mathrm{\Gamma }`$ is an $`f`$-vertical odd holomorphic field $`\mathrm{\Gamma }`$ vanishing at $`\times `$. If $`\{z^\alpha ,z^{\dot{\alpha }},\psi ^{\dot{\alpha }}\}`$ is a natural local coordinate system on $``$, $`\{t^i\}`$ a local coordinate system on $`𝒮`$ centered at $``$, then $`\varphi ^1`$ maps the cartesian product of these into a local coordinate system on $`𝒳`$ in which $`\mathrm{\Gamma }`$ can written as follows
$$\mathrm{\Gamma }=\underset{\alpha =1}{\overset{n}{}}\mathrm{\Gamma }^\alpha \frac{}{z^\alpha }+\underset{\dot{\alpha }=1}{\overset{n}{}}\left(\mathrm{\Gamma }^{\dot{\alpha }}\frac{}{z^{\dot{\alpha }}}+\gamma ^{\dot{\alpha }}\frac{}{\psi ^{\dot{\alpha }}}\right),$$
for some local functions, $`\mathrm{\Gamma }^\alpha ,\mathrm{\Gamma }^{\dot{\alpha }},\gamma ^{\dot{\alpha }}`$, of $`\{z^\alpha ,z^{\dot{\alpha }},\psi ^{\dot{\alpha }},t^i\}`$ which vanish at $`t^i=0`$. Hence, the local smooth map,
$`z^\alpha `$ $``$ $`z^\alpha ,`$
$`z^{\dot{\alpha }}`$ $``$ $`z^{\dot{\alpha }},`$
$`\psi ^{\dot{\alpha }}`$ $``$ $`\psi ^{\dot{\alpha }}+\mathrm{\Gamma }^{\dot{\alpha }}(z^\beta ,z^{\dot{\beta }},\psi ^{\dot{\beta }},t^i),`$
$`t^i`$ $`t^i,`$
is invertable in a small open neighbourhood of $`f^1()`$ in $`𝒳`$. It is easy to see that this transformation sends the component $`\mathrm{\Gamma }^{\dot{\alpha }}/z^{\dot{\alpha }}`$ in $`\mathrm{\Gamma }`$ to zero. Put another way, this component in $`\mathrm{\Gamma }`$ can always be eliminated by an appropriate choice — “gauge” — of the trivialization $`\varphi `$. Then the homology condition $`[_𝒳,_𝒳]=0`$ immediately implies that, in this gauge, $`\gamma ^{\dot{\alpha }}=0`$ as well. (In physics jargon, both unwanted fields, $`\mathrm{\Gamma }^{\dot{\alpha }}`$ and $`\gamma ^{\dot{\alpha }}`$, correspond to purely gauge degrees of freedom, and, moreover, can be eliminated by one single gauge transform as above). Of the remaining coordinate transformations preserving the gauge, only the following ones,
$`z^\alpha `$ $``$ $`z^\alpha +h^\alpha (z^\beta ,z^{\dot{\beta }},\psi ^{\dot{\beta }},t^i),h^\alpha _{t=0}=0,(G)`$
$`z^{\dot{\alpha }}`$ $``$ $`z^{\dot{\alpha }}`$
$`\psi ^\alpha `$ $``$ $`\psi ^{\dot{\alpha }}`$
$`t^i`$ $``$ $`t^i`$
do effectively change $`\mathrm{\Gamma }`$. Now it is clear that, in the gauge $`\mathrm{\Gamma }^{\dot{\alpha }}=\gamma ^{\dot{\alpha }}=0`$, the solution space, $`\widehat{𝖬C}`$, of Maurer-Cartan equations in the deformation theory of extended complex manifolds,
$$[\overline{},\mathrm{\Gamma }]+[_𝒮,\mathrm{\Gamma }]+\frac{1}{2}[\mathrm{\Gamma },\mathrm{\Gamma }]=0,$$
modulo the remaining gauge freedom $`(G)`$ can be canonically identified with the solution space of the Maurer-Cartan equations of the dg-Lie algebra $`(𝔤=\mathrm{\Gamma }(M,T_M\mathrm{\Omega }_M^{0,}),\overline{})`$,
$$𝖬C=\{\mathrm{\Gamma }(𝔤m_A)_{\stackrel{~}{1}}\overline{}\mathrm{\Gamma }+_A\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma },\mathrm{\Gamma }]=0\},$$
modulo the following gauge transformations,
$$\mathrm{\Gamma }\mathrm{\Gamma }^g=e^{\mathrm{a}d_g}\mathrm{\Gamma }\frac{e^{ad_g}1}{\mathrm{a}d_g}(d+_A)g,g(𝔤m_A)_{\stackrel{~}{0}}.$$
Here $`m_A`$ stands for the maximal ideal in the Artin superalgebra $`A`$. It is proved in \[Me2\] that the associated deformation functor
$$\begin{array}{cccc}\hfill \text{Def}_𝔤^{}:& \left\{\begin{array}{c}\text{the category of differential}\\ \text{Artin superalgebras}\end{array}\right\}& & \left\{\text{the category of sets}\right\}\\ & (A,_A)& & 𝖬C/𝗀augegroup\end{array}$$
is always unobstructed, and can be versally represented by $`(𝐇^{}(M,T_M),,0)`$, where the homological vector field $``$ is an invariant of $`𝔤`$. (The latter was called in \[Me2\] Chen’s vector field as its origin can be traced back to Chen’s power series connection in the theory of iterated integrals.)
With the established isomorphism of versal moduli spaces, the statement 2.7.2(iii) becomes an obvious corollary of Lemma 2.1 in \[BK\].
It remains to check the existence of the classical limit 2.7.2(ii). Picking up a Hermitian metric on $`M`$, we can construct the adjoint, $`\overline{}^{}`$, of the Dolbeault operator $`\overline{}`$ on $`𝔤`$, the Laplacian $`\mathrm{}=\overline{}\overline{}^{}+\overline{}^{}\overline{}`$, the Green function $`G`$, and we can identify the cohomology space $`𝐇^{}(M,T_M)`$ with the space of harmonic elements, $`𝖪er\overline{}𝖪er\overline{}^{}`$, in $`𝔤`$. Choosing next a harmonic basis, $`e_i`$, in $`𝐇^{}(M,T_M)`$, and denoting the associated linear coordinates by $`t^i`$ we can represent the Chen’s vector field, $`=_i^i(t)/t^i`$, as follows \[Me2\]
$$\underset{i}{}^i(t)e_i=\frac{1}{2}P[\mathrm{\Gamma },\mathrm{\Gamma }],$$
where $`P:𝔤𝐇^{}(M,T_M)`$ is the natural projection to the harmonic constituent, and $`\mathrm{\Gamma }=_{n=1}^{\mathrm{}}\mathrm{\Gamma }_{[n]}`$ is given by a recursive formula,
$`\mathrm{\Gamma }_{[1]}`$ $`=`$ $`{\displaystyle \underset{i}{}}t^ie_i`$
$`\mathrm{\Gamma }_{[2]}`$ $`=`$ $`{\displaystyle \frac{1}{2}}G\overline{}^{}[\mathrm{\Gamma }_{[1]}(t),\mathrm{\Gamma }_{[1]}(t)],`$
$`\mathrm{\Gamma }_{[3]}`$ $`=`$ $`{\displaystyle \frac{1}{2}}G\overline{}^{}\left([\mathrm{\Gamma }_{[1]}(t),\mathrm{\Gamma }_{[2]}(t)]+[\mathrm{\Gamma }_{[2]}(t),\mathrm{\Gamma }_{[1]}(t)]\right),`$
$`\mathrm{}`$
$`\mathrm{\Gamma }_{[n]}`$ $`=`$ $`{\displaystyle \frac{1}{2}}G\overline{}^{}\left({\displaystyle \underset{k=1}{\overset{n1}{}}}[\mathrm{\Gamma }_{[k]}(t),\mathrm{\Gamma }_{[nk]}(t)]\right)`$
$`\mathrm{}`$
It is now obvious that $`𝖹eros()𝐇^1(M,T_M)`$ is precisely the Kuranishi analytic subspace \[K\]. Moreover, the fibres, $`f^1(t)`$, of our minimal versal deformation $`f`$ over this subspace are precisely the extended complex manifolds, $`(_t,\overline{}_t)`$, associated, via the construction 1.5, to the usual complex manifolds $`M_t`$ lying over $`t`$ in the classical Kuranishi versal deformation of $`M`$. $`\mathrm{}`$
Acknowledgement. This work was done during author’s visit to the Max Planck Institute for Mathematics in Bonn. Excellent working conditions in the MPIM are gratefully acknowledged. I would like to thank Yu.I. Manin for many discussions.
| Max Planck Institute for Mathematics in Bonn, and |
| --- |
| Department of Mathematics, University of Glasgow |
| sm@maths.gla.ac.uk |
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# REFERENCES
HURWITZ THEOREM AND PARALLELIZABLE SPHERES
FROM TENSOR ANALYSIS
J. A. Nieto<sup>*</sup><sup>*</sup>*nieto@uas.uasnet.mxand L. N. Alejo-Armentanabor@uas.uasnet.mx
Facultad de Ciencias Físico-Matemáticas, Universidad Autónoma
de Sinaloa, C.P. 80010, Culiacán Sinaloa, México
Abstract
By using tensor analysis, we find a connection between normed algebras and the parallelizability of the spheres S<sup>1</sup>, S<sup>3</sup> and S$`^7.`$ In this process, we discovered the analogue of Hurwitz theorem for curved spaces and a geometrical unified formalism for the metric and the torsion. In order to achieve these goals we first develope a proof of Hurwitz theorem based in tensor analysis. It turns out that in contrast to the doubling procedure and Clifford algebra mechanism, our proof is entirely based in tensor algebra applied to the normed algebra condition. From the tersor analysis point of view our proof is straightforward and short. We also discuss a possible connection between our formalism and the Cayley-Dickson algebras and Hopf maps.
Pacs No.: 02.10. Vr; 02.10. Tq; 02.90.+p.
January, 2001
I. INTRODUCTION
It is known that normed algebras are closely related to supersymmetry<sup>1-3</sup> and super p-branes<sup>4</sup>, and that these two theories require tensor analysis for their formulation. Therefore, it may be interesting to study normed algebras from the tensor analysis point of view. Moreover, normed algebras, among other things, are physically interesting because they are division algebras and in this context there are a number of interesting connections with fundamental physics. Let us just give some few examples about this fact. It has been shown<sup>5</sup> that a generalized instantons in eight dimensions fit inside the family of gauge-theoretical solitons associated to normed algebras. There is a deep relation between division algebras and superparticles (see ref. 6, 7 and references there in) and twistor formulation of a massless particles<sup>8,9</sup>. Finite Lorentz transformations of vectors in 10-dimensional Minkowski space have been studied<sup>10</sup> by means of division algebras. Finally, division algebras seem to be deeply related to the geometric structures of M-theory<sup>11</sup>.
In this work, we show that tensor analysis can be used to give a straightforward connection between normed algebras and the paralellizability of the spheres S<sup>1</sup>, S<sup>3</sup> and S$`^7.`$ In the process of studing this connection, we discovered the analogue of Hurwitz theorem for curved spaces and a unified formalism for the metric and the torsion. Our strategy to achieve these goals was first to develope a proof of Hurwitz’s theorem<sup>12</sup> based in tensor analysis. It turns out that this proof is essentially based on the composition law rewritten in tensor notation. From the point of view of tensor analysis, such a proof is short and straightforward. In fact, we do not even require to use the doubling procedure<sup>12</sup> or the Clifford algebra mechanism<sup>13</sup>.
The plan of the article is as follows. In section II, we introduce tensor notation and a proof of Hurwitz theorem based in tensor analysis. In section III, we briefly review the Cartan-Shouten equations as presented by Gursey and Tze. In section IV, using the Gursey-Tze’s procedure, we show a connection between our proof of Hurwitz theorem and the paralellizability of the spheres S<sup>1</sup>, S<sup>3</sup> and S$`^7.`$ We also prove that such a connection leads to a generalization of Hurwitz theorem for curved spaces. In section V, we develope a unified formalism for the metric and the torsion. Finally, in section VI, we make a number of final comments and briefly outline a possible extension of the present work to the case of Cayley-Dickson algebras and Hopf maps.
II. AN ALTERNATIVE PROOF OF HURWITZ THEOREM
Let us start recalling the Hurwitz theorem:
Theorem (Hurwitz, 1898): Every normed algebra with an identity is isomorphic to one of following four algebras: the real numbers, the complex numbers, the quaternions, and the Cayley (octonion) numbers.
Proof (alternative): Consider a $`N=d+1`$dimensional algebra $`𝒜`$ over the real numbers $`R`$. Let
$$e_0,e_1,\mathrm{},e_d$$
(1)
be a basis of $`𝒜`$, and let
$$A=A^0e_0+A^1e_1+\mathrm{}+A^de_d$$
(2)
be the representation of a vector $`A`$ $`ϵ`$ $`𝒜`$ relative to this basis. Here, $`A^0,A^1,\mathrm{},A^dϵR.`$ Take the multiplication table in the form
$$\begin{array}{c}e_ie_j=C_{ij}^0e_0+C_{ij}^1e_1+\mathrm{}+C_{ij}^de_d,\\ \\ (i,j=0,1,\mathrm{},d),\end{array}$$
(3)
where $`C_{ij}^k`$ , the so-called structure constants, are real numbers (See, for instance, I. L. Kantor and A.S. Solodovnikov<sup>12</sup>, S. Okubo<sup>13</sup>, Abdel-Khalek<sup>14</sup>, J. Adem<sup>15</sup>, F. R. Cohen<sup>16</sup>, Y. A. Drozd and V. V. Kirichenko<sup>17</sup>.)
Assume that the basis (1) is orthonormal with bi-linear symmetric non-degenerate scalar product given by
$$<e_ie_j>=\delta _{ij},$$
(4)
where $`\delta _{ij}`$ is the so-called Kronecker delta, with $`\delta _{ij}=0`$ if $`ij`$ and $`\delta _{ij}=1`$ if $`i=j`$ .
Assume the Einstein summation convention: if the same index appears twice, once as superscript and once as a subscript, then the index is summed over all possible values. This convention allows to write (2) and (3) as
$$A=A^ie_i,$$
(5)
and
$$e_ie_j=C_{ij}^ke_k,$$
(6)
respectively.
We shall assume that $`e_i`$ transforms as covariant first-rank tensor
$$e_i^{}=\mathrm{\Lambda }_i^je_j,$$
(7)
where, in order to leave invariant (4), $`\mathrm{\Lambda }_i^j`$ satisfies the conditions $`det\mathrm{\Lambda }_i^j0`$ and $`\mathrm{\Lambda }_k^i\mathrm{\Lambda }_l^j\delta _{ij}=\delta _{kl}`$ and therefore $`\mathrm{\Lambda }_i^j`$ is an element of an orthogonal transformation $`O(N)`$. Since $`A`$ is an invariant quantity, from (5) and (7) we find that $`A^i`$ should transform as
$$A^i=\mathrm{\Lambda }_j^iA^j,$$
(8)
i.e. $`A^i`$ is a contravariant first-rank tensor. While from (6) and (7) we find that $`C_{ij}^k`$ transforms as
$$C_{st}^r=\mathrm{\Lambda }_k^r\mathrm{\Lambda }_s^i\mathrm{\Lambda }_t^jC_{ij}^k,$$
(9)
i.e. $`C_{ij}^k`$ is a mixed third-rank tensor (twice covariant and once contravariant).
According to the multiplication table (6) the product $`AB=D`$ for $`A,B`$ and $`Dϵ`$ $`𝒜`$ is given by
$$A^iB^jC_{ij}^k=D^k.$$
(10)
A normed algebra is an algebra in which the composition law
$$<ABAB>=<AA><BB>$$
(11)
holds for any $`A`$, $`B`$ $`ϵ`$ $`𝒜`$. It can be shown that this expression is equivalent to (see, for instance, section 3.1 of ref. 13)
$$<ABCD>+<CBAD>=2<AC><BD>,$$
(12)
where $`A`$, $`B,C,D`$ $`ϵ`$ $`𝒜`$. Choosing
$$Ae_i,Be_j,Ce_m\text{ and }De_{n\text{ }}$$
(13)
we find that (12) leads to
$$<e_ie_je_me_n>+<e_me_je_ie_n>=2<e_ie_m><e_je_n>.$$
(14)
Using (4) and (6), from (14) we obtain the key formula
$$C_{ij}^kC_{mn}^l\delta _{kl}+C_{mj}^kC_{in}^l\delta _{kl}=2\delta _{im}\delta _{jn}.$$
(15)
Note that, although at first sight it looks like, (15) is not a Clifford algebra. The reason for this is that, at this level, there are not any symmetries between the indices $`i,j`$ and $`k`$ of $`C_{ij}^k.`$ In this work, the formula (15) shall play a central role. Note that when $`D=1,`$ this equation admits the solution $`C_{00}^0=1.`$ Therefore, in what follows we shall be mainly interested in solutions of (15) when $`D1.`$
Let $`e_0`$ be the identity of the algebra $`𝒜.`$ Then, the multiplication table (6) implies
$$e_0e_j=C_{0j}^ke_k=e_j$$
(16)
and
$$e_je_0=C_{j0}^ke_k=e_j.$$
(17)
From (16) we find
$$C_{0j}^k=\delta _j^k,$$
(18)
while from (17) we obtain
$$C_{j0}^k=\delta _j^k,$$
(19)
where $`\delta _j^k`$ is also a Kronecker delta.
Let us now split the formula (15) as follows:
$$C_{0j}^kC_{0n}^l\delta _{kl}+C_{0j}^kC_{0n}^l\delta _{kl}=2\delta _{jn},$$
(20)
$$C_{i0}^kC_{m0}^l\delta _{kl}+C_{m0}^kC_{i0}^l\delta _{kl}=2\delta _{im},$$
(21)
$$C_{0j}^kC_{an}^l\delta _{kl}+C_{aj}^kC_{0n}^l\delta _{kl}=0,$$
(22)
$$C_{i0}^kC_{ma}^l\delta _{kl}+C_{m0}^kC_{ia}^l\delta _{kl}=0,$$
(23)
$$C_{ab}^0C_{cd}^0+C_{cb}^0C_{ad}^0+C_{ab}^eC_{cd}^f\delta _{ef}+C_{cb}^eC_{ad}^f\delta _{ef}=2\delta _{ac}\delta _{bd},$$
(24)
where the indices $`a,b,`$…, etc run from $`1`$ to $`d`$. Using (18) and (19) we note that the equations (20) and (21) are identities. Moreover, the expression (22) gives
$$C_{anj}+C_{ajn}=0,$$
(25)
while (23) leads to
$$C_{mai}+C_{iam}=0,$$
(26)
where $`C_{mai}=C_{ma}^l\delta _{il},`$ i.e. we raised and lowed indices with $`\delta ^{il}`$ and $`\delta _{il}`$ respectively. From (25) we obtain
$$C_{ab0}+C_{a0b}=0,$$
(27)
and
$$C_{abc}+C_{acb}=0.$$
(28)
While from (26) we get
$$C_{ba0}+C_{0ab}=0$$
(29)
and
$$C_{bac}+C_{cab}=0.$$
(30)
Thus, using (18) and (19), we have that either (27) or (30) implies that
$$C_{ab}^0=\delta _{ab},$$
(31)
while (28) and (30) mean that the quantity $`C_{abc}`$ is completely antisymmetric.
Now, by substituting (31) into (24) we obtain
$$C_{ab}^eC_{cd}^f\delta _{ef}+C_{cb}^eC_{ad}^f\delta _{ef}=2\delta _{ac}\delta _{bd}\delta _{ab}\delta _{cd}\delta _{cb}\delta _{ad}.$$
(32)
Multiplying this equation by $`\delta ^{ac}`$ and $`C_g^{ad}`$ $`=\delta ^{ae}\delta ^{af}C_{gef}`$ we find
$$\delta ^{ac}C_{ab}^eC_{cd}^f\delta _{ef}=(d1)\delta _{bd}$$
(33)
and
$$C_g^{ad}C_{ab}^eC_{cd}^f\delta _{ef}+C_g^{ad}C_{cb}^eC_{ad}^f\delta _{ef}=3C_{gcb},$$
(34)
respectively, where we used the fact that $`C_{abc}`$ is completely antisymmetric. Moreover, using again the property that $`C_{abc}`$ is completely antisymmetric, we find that (33) becomes
$$C_b^{ce}C_{dce}^{}=(d1)\delta _{bd},$$
(35)
while substituting (33) into (34) we have
$$C_{dg}^aC_{ab}^eC_{ec}^d=(d4)C_{gbc}.$$
(36)
Multiplying (35) by $`\delta ^{bd}`$ we find the formula
$$C^{abc}C_{abc}=d(d1),$$
(37)
which for $`d=0`$ and $`d=1,`$ admits the solution $`C_{abc}=0`$. Moreover, for $`d=3`$ the formula (37) admits the solution $`C_{abc}^\text{ }=\epsilon _{abc},`$ where $`\epsilon _{abc}`$ is the completely antisymmetric Levi-Civita symbol, with $`\epsilon _{123}=1.`$
Let us define
$$G_{abc}C_{da}^gC_{gb}^eC_{ec}^d.$$
(38)
Since $`C_{abc}^\text{ }`$ is completely antisymmetric, we find that $`G_{abc}`$ is also completely antisymmetric. From (36) and (38) we find that
$$G^{abc}G_{abc}=(d4)^2C^{abc}C_{abc},$$
(39)
which by virtue of (37) leads to
$$G^{abc}G_{abc}=d(d1)(d4)^2.$$
(40)
Substituting (38) into (40) we get
$$C_h^{ag}C_g^{br}C_r^{ch}C_{da}^eC_{eb}^fC_{fc}^d=d(d1)(d4)^2,$$
(41)
which can be rewritten in the form
$$C_h^{ag}C_g^{br}C_{da}^eC_{eb}^f(C_r^{ch}C_{fc}^d)=d(d1)(d4)^2.$$
(42)
So, considering (32) we find that (42) becomes
$$C_h^{ag}C_g^{br}C_{da}^eC_{eb}^f(2\delta _{rf}\delta ^{hd}\delta _r^h\delta _f^d\delta _r^d\delta _f^hC_f^{ch}C_{rc}^d)=d(d1)(d4)^2.$$
(43)
Now, using (35), (36) and (37) and the fact that $`C_{abc}`$ is completely antisymmetric, we obtain
$$\begin{array}{cc}C_h^{ag}C_g^{br}C_{da}^eC_{eb}^f(2\delta _{rf}\delta ^{hd}\delta _r^h\delta _f^d\delta _r^d\delta _f^h)=& \\ & \\ =2d(d1)(d1)d(d1)(d1)+d(d1)(d4)& \\ & \\ =d(d1)(2d5),& \end{array}$$
(44)
while, since with respect to the indices $`a`$ and $`h`$ the quantity $`C_h^{ag}`$ is antisymmetric and the tensor $`(C_{da}^eC_{eb}^fC_f^{ch}C_{rc}^d)`$ is symmetric, we get
$$C_h^{ag}C_g^{br}C_{da}^eC_{eb}^f(C_f^{ch}C_{rc}^d)=C_h^{ag}C_g^{br}(C_{da}^eC_{eb}^fC_f^{ch}C_{rc}^d)0.$$
(45)
Thus, by substituting the results (44) and (45) into (43), we discover the equation
$$d(d1)(2d5)=d(d1)(d4)^2,$$
(46)
which can be rewritten in the form
$$d(d1)(d3)(d7)=0.$$
(47)
The only solutions for this equation are $`d=0,1,3`$ and $`7.`$ Therefore, we have shown that the equation (15) has solution only for $`D=1,2,4`$ and $`8.`$ This implies that normed algebras with unit element are only possible in these dimensions.
We have yet to show that the cases $`D=1,D=2,D=4`$ and $`D=8`$ correspond to real, complex, quaternion and octonion algebras, respectively. The case $`D=1`$ is trivial since for any $`Aϵ𝒜,`$ we have $`A=A^0e_0,`$ where $`A^0ϵ`$ $`R`$. For the case $`D=2,`$ we have $`C_{abc}=0`$, $`C_{ab}^0=\delta _{ab},`$ $`C_{n0}^s=\delta _n^s`$ and $`C_{0n}^s=\delta _n^s.`$ These values of the structure constants determine the algebra of complex numbers. While, for the case $`D=4,`$ we have the solution of (32) $`C_{abc}=\epsilon _{abc}`$, $`C_{ab}^0=\delta _{ab},`$ $`C_{n0}^s=\delta _n^s`$ and $`C_{0n}^s=\delta _n^s`$. It is well known that these values of the structure constants determine the algebra of quaternions. Finally, for the case $`D=8`$ we have $`C_{ab}^0=\delta _{ab},`$ $`C_{n0}^s=\delta _n^s`$ and $`C_{0n}^s=\delta _n^s.`$ Now, take the structure constants as $`C_{abc}=\mathrm{\Xi }_{abc}`$, where $`\mathrm{\Xi }_{abc}`$ is a completely antisymmetric Levi-Civita symbol, with $`\mathrm{\Xi }_{abc}=1,`$ for the following values of the indices $`(a,b,c)`$:
$$(1,2,3),(5,1,6),(6,2,4),(4,3,5),(1,7,4),(3,7,6)\text{ and }(2,7,5).$$
(48)
In fact, these values of the structure constants determine the algebra of octonions. One can verify by straightforward, but tedious, computation that, in fact for $`d=7,`$ these values for the structure constants give a solution of (32).
It is known that by definition two $`(d+1)`$ dimensional algebras $`𝒜^{}`$ and $`𝒜`$ are said to be isomorphic if they have bases with identical multiplication table. Therefore, it remains to show that any other solution is isomorphic to one of the above four solutions corresponding to the real numbers, the complex numbers, the quaternions and the octonions. For this purpose it is convenient to set $`e_0^{}=e_0.`$ So that from the transformation rule (7) we find that $`\mathrm{\Lambda }_0^0=1`$ and $`\mathrm{\Lambda }_0^a=0.`$ Thus, from the relation $`\mathrm{\Lambda }_k^i\mathrm{\Lambda }_l^j\delta _{ij}=\delta _{kl},`$ which leave invariant the scalar product (4), we find that $`\mathrm{\Lambda }_a^0=0`$ and therefore we have now the relation $`\mathrm{\Lambda }_a^c\mathrm{\Lambda }_b^d\delta _{cd}=\delta _{ab}`$ which leaves invariant the scalar product $`<e_ae_b>=\delta _{ab}.`$ Consequently, we have that $`\mathrm{\Lambda }_a^d`$ is an element of $`O(d)=`$ $`O(D1)`$ which is a subgroup of $`O(D).`$ Note that the property $`det\mathrm{\Lambda }_j^i0`$ now becomes $`det\mathrm{\Lambda }_a^d0.`$ Clearly, the transformation $`\mathrm{\Lambda }_a^d`$ acts over elements of the sub-vector space $`𝒜_0`$ of $`𝒜`$ defined by $`𝒜_0=\{A<Ae_0>=0,`$ $`Aϵ𝒜\}`$, with Dim $`𝒜_0=d.`$ In fact, we can write $`𝒜`$ $`=\lambda e_0𝒜_0`$, with $`\lambda ϵR.`$
Thus, we find that the structure constants $`C_{abc}`$ transform according to
$$C_{abc}^{^{}}=\mathrm{\Lambda }_a^d\mathrm{\Lambda }_b^e\mathrm{\Lambda }_c^fC_{def}.$$
(49)
Note that, since $`\mathrm{\Lambda }_a^d\mathrm{\Lambda }_b^e\delta _{de}=\delta _{ab},`$ if $`C_{abc}`$ is a solution of (32), then $`C_{abc}^{^{}}`$ is also a solution.
The transformation (49) has the important property that $`C_{def}=0`$ if and only if $`C_{abc}^{^{}}=0.`$ Therefore for real numbers, as well as for complex numbers, the two algebras $`𝒜^{}`$ and $`𝒜`$ are isomorphic. For quaternions take $`C_{def}=\epsilon _{def}`$ then (49) implies that $`C_{abc}^{^{}}=\mathrm{\Lambda }\epsilon _{abc}`$, $`\mathrm{\Lambda }det\mathrm{\Lambda }_a^d.`$ Thus, if $`C_{def}=\epsilon _{def}`$ is a solution of (32) we have that $`C_{abc}^{^{}}=\mathrm{\Lambda }\epsilon _{abc}`$ is also a solution. Therefore, for $`D=4`$ any solution of (32) is isomorphic to the quaternionic solution, corresponding to $`C_{abc}=\epsilon _{abc}`$. Similarly, for octonions applying (49) to the completely antisymmetric symbol $`\mathrm{\Xi }_{abc}`$ we find that $`\mathrm{\Xi }_{abc}^{}=\mathrm{\Lambda }\mathrm{\Xi }_{abc},`$ where the values of the indices $`(a,b,c)`$ are given in (48).
Therefore, we have shown that up to isomorphism the dimensions $`D=1,2,4`$ and $`8`$ correspond to real, complex, quaternion and octonion algebras, respectively. And in this way using the mathematical tool of tensor analysis we have given an alternative proof of Hurwitz theorem. It is an interesting and remarkable fact that without using doubling procedure (see ref. 12) or Clifford algebra mechanism (see ref. 13) our proof has been based almost completely in tensor algebra applied to the formula (15).
III. CARTAN-SHOUTEN EQUATIONS
Define the metric tensor by
$$g_{ab}=\delta _{cd}h_a^{(c)}h_b^{(d)},$$
(50)
where $`h_a^{(c)}`$ = $`h_a^{(c)}(x^{b)})`$ is a vielbein field. Here, $`x^a`$ is a coordinate patch of the geometrical sphere S<sup>d</sup>. The quantities $`C_{abc}`$ can now be related to the S<sup>d</sup> torsion in the form
$$T_{abc}=r^1C_{efg}h_a^{(e)}h_b^{(f)}h_c^{(g)},$$
(51)
where $`r`$ is the radius of S$`^d.`$ Using (35), (36), (50) and (51) we find that the torsion $`T_{abc}`$ satisfies the equations:
$$T_a^{cd}T_{bcd}=(d1)r^2g_{ab},$$
(52)
and
$$T_{ea}^dT_{db}^fT_{fc}^e=(d4)r^2T_{abc.}$$
(53)
We recognize these expressions as the Cartan-Schouten equations<sup>18</sup>which as Gursey and Tze<sup>19</sup> noted, are mere septad-dressed, i.e. covariant forms of the algebraic identities (35) and (36). It is well known that these equations are closely related to the parallelizability of S<sup>1</sup>, S<sup>3</sup> and S<sup>7</sup> (see ref. 13). In fact, the equations (52) and (53) can be derived by adding to the riemannian symmetric connection $`\mathrm{\Gamma }_{ab}^c`$ the totally antisymmetric torsion tensor $`T_{ab}^c`$ and ”flattening” the space in the sense that
$$_{bcd}^a(\{\mathrm{\Omega }_{ab}^c\})=0,$$
(54)
where
$$_{bcd}^a=_c\mathrm{\Omega }_{bd}^a_d\mathrm{\Omega }_{bc}^a+\mathrm{\Omega }_{ec}^a\mathrm{\Omega }_{bd}^e\mathrm{\Omega }_{ed}^a\mathrm{\Omega }_{bc}^e,$$
(55)
with
$$\mathrm{\Omega }_{ab}^c=\mathrm{\Gamma }_{ab}^c+T_{ab}^c.$$
(56)
For our purpose it is convenient to show explicitly that in fact the equations (52) and (53) follow from (54)-(56). By substituting (56) into (54) we find
$$0=R_{bcd}^a+D_cT_{bd}^aD_dT_{bc}^a+T_{ec}^aT_{bd}^eT_{ed}^aT_{bc}^e,$$
(57)
Here, $`D_c`$ denotes a covariant derivative with $`\mathrm{\Gamma }_{ab}^c`$ as a connection and
$$R_{bcd}^a=_c\mathrm{\Gamma }_{bd}^a_d\mathrm{\Gamma }_{bc}^a+\mathrm{\Gamma }_{ec}^a\mathrm{\Gamma }_{bd}^e\mathrm{\Gamma }_{ed}^a\mathrm{\Gamma }_{bc}^e.$$
(58)
Using in (57) the cyclic identities for $`R_{bcd}^a`$ leads to
$$D_cT_{bda}=T_{e[bd}T_{a]c}^e,$$
(59)
where
$$T_{e[bd}T_{a]c}^e\frac{1}{3}\{T_{ebd}T_{ac}^e+T_{eab}T_{dc}^e+T_{eda}T_{bc}^e\}.$$
(60)
Substituting (59) into (57) we obtain
$$R_{abcd}^{}=T_{eab}T_{cd}^eT_{e[ab}T_{c]d}^e.$$
(61)
For the sphere S<sup>d</sup> we have
$$R_{abcd}^{}=\frac{1}{r^2}(g_{ac}g_{bd}g_{ad}g_{bc}).$$
(62)
and therefore we get the equation
$$\frac{1}{r^2}(g_{ac}g_{bd}g_{ad}g_{bc})=T_{eab}T_{cd}^eT_{e[ab}T_{c]d}^e.$$
(63)
Contracting in (63) with $`g^{ac}`$ leads to first Cartan-Shouten equation (52), while contracting (63) with $`T_f^{ac}`$ leads to the second Cartan-Shouten equation (53).
IV. NORMED ALGEBRAS AND PARALLELIZABILITY OF S<sup>1</sup>, S<sup>3</sup> and S<sup>7</sup>
Let us start ‘undressing’ (63). Using (50) and (51) we find
$$(\delta _{ac}\delta _{bd}\delta _{ad}\delta _{bc})=C_{eab}C_{cd}^eC_{e[ab}C_{c]d}^e.$$
(64)
We shall show that this formula is equivalent to the formula (32). For this purpose, let us rewrite formula (32) in form.
$$2\delta _{ac}\delta _{bd}\delta _{ab}\delta _{cd}\delta _{ad}\delta _{cb}=C_{ab}^eC_{cd}^f\delta _{ef}+C_{cb}^eC_{ad}^f\delta _{ef}.$$
(65)
Let us first show that (64) implies (65). Making the change of indices $`ac`$ and $`ca`$ in (64) we find
$$(\delta _{ca}\delta _{bd}\delta _{cd}\delta _{ba})=C_{ecb}C_{ad}^eC_{e[cb}C_{a]d}^e.$$
(66)
By adding (64) and (66) one easily obtains (65).
Let us now show that (65) implies (64). Let us start writing (65) in the form
$$C_{ab}^eC_{cd}^f\delta _{ef}(\delta _{ac}\delta _{bd}\delta _{ad}\delta _{cb})+C_{cb}^eC_{ad}^f\delta _{ef}(\delta _{ac}\delta _{bd}\delta _{ab}\delta _{cd})=0.$$
(67)
This expression suggests to define
$$F_{abcd}C_{ab}^eC_{cd}^f\delta _{ef}(\delta _{ac}\delta _{bd}\delta _{ad}\delta _{cb}).$$
(68)
Therefore (67) gives
$$F_{abcd}+F_{cbad}=0.$$
(69)
Thus, considering that $`C_{ab}^e`$ is completely antisymmetric, from (68) and (69) we discover that $`F_{abcd}`$ is also completely antisymmetric. Using this important cyclic property for $`F_{abcd}`$ it is not difficult to show that
$$F_{abcd}=C_{e[ab}C_{c]d}^e.$$
(70)
Substituting this result into (68) lead us back to (64). Thus, we have proved the equivalence between (64) and (65).
With this equivalence at hand we have a number of interesting observations. First, since in section II we showed that (65) (or (32)) admits solution only for dimensions $`d=1,3`$ and $`7`$ we have that (64) admits solution only in these dimensions. But, since (64) is the necessary and sufficient condition for the existence of parallelism in S$`^d,`$ this means that we have found an alternative proof of the fact that only the spheres S$`^1,`$S<sup>3</sup> and S<sup>7</sup> are parallelizables. Second, in section II we proved that (65) is a consequence of the normed condition (15) (or equivalent of (11)), while in section III we proved that (64) is a consequence of the paralizability condition (54). Therefore, we have find a new bridge between normed algebras and parallelizable spheres. This link can be more transparent if using (50) and (51) we dress (65) in the form
$$\frac{1}{r^2}(2g_{ac}g_{bd}g_{ab}g_{cd}g_{ad}g_{cb})=T_{ab}^eT_{cd}^fg_{ef}+T_{cb}^eT_{ad}^fg_{ef}.$$
(71)
Of course, the equations (63) and (71) are equivalent. So, from (65) we can derive (71) which in turn leads to the formula (63). Going backwards from (63) we get (61). Therefore, we have shown that normed algebra condition (65) implies the parallelizable condition (61). Similarly, we can show that the parallelizable condition (61) implies the composition law (65). Moreover, (31) and (71) suggest to define
$$T_{ab}^0r^1g_{ab}.$$
(72)
Thus, using (72) we find that (71) can be written in the form
$$T_{ab}^kT_{cd}^lg_{kl}+T_{cb}^kT_{ad}^lg_{kl}=\frac{2}{r^2}g_{ac}g_{bd}.$$
(73)
where we recall that the indices $`m`$ and $`n`$ run from $`0`$ to $`d.`$ Setting $`g_{00}=1`$ and $`g_{0a}=0`$ we obtain (72) from (73). If we now take $`T_{0j}^k=\delta _j^k`$ and $`T_{j0}^k=\delta _j^k,`$ then we can generalize (80) in the form
$$T_{ij}^kT_{mn}^lg_{kl}+T_{mj}^kT_{in}^lg_{kl}=\frac{2}{r^2}g_{im}g_{jn}.$$
(74)
If we now introduce a basis $`h_m`$ such that
$$<h_mh_n>=g_{mn},$$
(75)
and
$$h_mh_n=T_{mn}^kh_k,$$
(76)
we find that (74) leads to a generalization of (14)
$$<h_ih_jh_mh_n>+<h_mh_jh_ih_n>=\frac{2}{r^2}<h_ih_m><h_jh_n>.$$
(77)
Clearly, this expression implies the generalized composition law condition
$$<ABAB>=\frac{1}{r^2}<AA><BB>,$$
(78)
where $`A=A^ih_i`$.
The $`r^{2\text{ }}`$in the right hand side of (74) remind us that our construction is valid for spheres. However, the equation (74) allows an straightforward generalization. In fact, let us prove the theorem ($``$) below:
Before going into the details of the theorem let us define a ‘curved’ space as a space in which (75) and (74) hold, with $`g_{00}=1,`$ $`g_{0a}=0`$ and $`g_{ab}=g_{ab}(x^i)=\eta _{cd}h_a^{(c)}(x^i)h_b^{(d)}(x^i),`$ where the flat metric $`\eta _{cd}`$ is diagonal and has an arbitrary signature and $`T_{ij}^k=T_{ij}^k(x^i).`$
Theorem ($``$): The possible dimensions $`D`$ of any real normed algebra over a ‘curved’ space with an identity are limited to only 1, 2, 4 and 8.
Proof: Let us write the composition law as follows:
$$<h_ih_jh_mh_n>+<h_mh_jh_ih_n>=2<h_ih_m><h_jh_n>.$$
(79)
By virtue of (75) and (76) we find that (79) can be written as
$$T_{ij}^kT_{mn}^lg_{kl}+T_{mj}^kT_{in}^lg_{kl}=2g_{im}g_{jn}.$$
(80)
Taking $`h_0`$ as the identity with the properties that
$$<h_0h_0>=g_{00}=1,$$
(81)
and
$$<h_0h_a>=g_{0a}=0,$$
(82)
and following the same procedure as in section II, we find
$$T_{0j}^k=T_{j0}^k=\delta _j^k,$$
(83)
$$T_{ab}^0=g_{ab},$$
(84)
$$(2g_{ac}g_{bd}g_{ab}g_{cd}g_{ad}g_{cb})=T_{ab}^eT_{cd}^fg_{ef}+T_{cb}^eT_{ad}^fg_{ef}$$
(85)
with the property that $`T_{ab}^e`$ is completely antisymmetric. The rest of the story is similar to section II after formula (32). We find that (85) has solution only if $`d=1,3`$ and $`7.`$ Note that in this result $`g_{ab}`$ may be the metric not only for the spheres S$`^1,`$S<sup>3</sup> and S$`^7,`$ but also the metric of any curved space. Moreover, in ‘flat’ space $`g_{ab}`$ may have an arbitrary signature. In particular for $`D=4`$ we could associate to $`g_{ij}`$ the signature ($`g_{ij})=diag(1,1,1,1)`$ which correspond to Minkowski signature. Note also that $`T_{ij}^k`$ unifies the metric $`g_{ab}`$ and the torsion $`T_{ab}^e`$.
Summarizing, we have proved not only an equivalence between the Hurwitz theorem for normed algebras and Cartan-Shouten theorem for parallelizable spheres, but also the theorem ($``$).
V. UNIFIED FORMALISM OF THE METRIC AND THE TORSION
In the previous section, in the context of normed algebras, we showed that makes sense to unify the metric and the torsion in just one mathematical object: the third-rank tensor $`T_{ij}^k`$ . A natural question is to see what is the geometry induced by $`T_{ij}^k.`$ In this section we show that from the vanishing of the Riemann tensor associated to such a third-rank tensor it follows the metricity condition and the Cartan-Shouten equations for homogeneous spacetimes.
Consider the equation
$$_{jkl}^i(\mathrm{\Omega }_{jk}^i)=0,$$
(86)
where
$$_{jkl}^i=_k\mathrm{\Omega }_{jl}^i_l\mathrm{\Omega }_{jk}^i+\mathrm{\Omega }_{mk}^i\mathrm{\Omega }_{jl}^m\mathrm{\Omega }_{el}^i\mathrm{\Omega }_{jk}^e$$
(87)
and
$$\mathrm{\Omega }_{jk}^i=\mathrm{\Gamma }_{jk}^i+T_{jk}^i.$$
(88)
These equations are, of course the analogue of the parallelizability conditions (54)-(56). Let us see what are the consequences of (86)-(88). For this purpose let us assume that $`T_{jk}^i`$ satisfies (83) and (84) and let us set
$$\mathrm{\Gamma }_{jk}^0=0\text{ and }\mathrm{\Gamma }_{0k}^i=0.$$
(89)
Thus, the non-vanishing terms of $`\mathrm{\Omega }_{jk}^i`$ are
$$\mathrm{\Omega }_{ab}^c=\mathrm{\Gamma }_{ab}^c+T_{ab}^c,$$
(90)
$$\mathrm{\Omega }_{ab}^0=g_{ab},$$
(91)
$$\mathrm{\Omega }_{0b}^a=\delta _b^a$$
(92)
and
$$\mathrm{\Omega }_{00}^0=1.$$
(93)
At this stage it is important to note that (90)-(93) could also be obtained if instead of (83), (84) and (89) we set $`T_{jk}^0=0,T_{0k}^i=0,\mathrm{\Gamma }_{0k}^0=0`$, $`\mathrm{\Gamma }_{00}^0=1,`$ $`\mathrm{\Gamma }_{0b}^a=\delta _b^a`$and $`\mathrm{\Gamma }_{ab}^0=g_{ab}.`$ However, with this choice, the connection between (80) and (85) will be lost. This connection is, of course, important to make contact with the normed algebras for ‘curved’ space discussed in the previous section. It is worth mentioning that the formulae (83), (84) and (89) can be understood as an anzats in the sense of Kaluza-Klein theory.
Let us split (87) in the form
$$_{abc}^0=_b\mathrm{\Omega }_{ac}^0_c\mathrm{\Omega }_{ab}^0+\mathrm{\Omega }_{0b}^0\mathrm{\Omega }_{ac}^0+\mathrm{\Omega }_{db}^0\mathrm{\Omega }_{ac}^d\mathrm{\Omega }_{0c}^0\mathrm{\Omega }_{ab}^0\mathrm{\Omega }_{dc}^0\mathrm{\Omega }_{ab}^d,$$
(94)
$$_{a0c}^0=_0\mathrm{\Omega }_{ac}^0_c\mathrm{\Omega }_{a0}^0+\mathrm{\Omega }_{00}^0\mathrm{\Omega }_{ac}^0+\mathrm{\Omega }_{d0}^0\mathrm{\Omega }_{ac}^d\mathrm{\Omega }_{0c}^0\mathrm{\Omega }_{a0}^0\mathrm{\Omega }_{dc}^0\mathrm{\Omega }_{a0}^d,$$
(95)
$$_{a0c}^b=_0\mathrm{\Omega }_{ac}^b_c\mathrm{\Omega }_{a0}^b+\mathrm{\Omega }_{00}^b\mathrm{\Omega }_{ac}^0+\mathrm{\Omega }_{d0}^b\mathrm{\Omega }_{ac}^d\mathrm{\Omega }_{0c}^b\mathrm{\Omega }_{a0}^0\mathrm{\Omega }_{dc}^b\mathrm{\Omega }_{a0}^d$$
(96)
and
$$_{abc}^d=_b\mathrm{\Omega }_{ac}^d_c\mathrm{\Omega }_{ab}^d+\mathrm{\Omega }_{0b}^d\mathrm{\Omega }_{ac}^0+\mathrm{\Omega }_{eb}^d\mathrm{\Omega }_{ac}^e\mathrm{\Omega }_{0c}^d\mathrm{\Omega }_{ab}^0\mathrm{\Omega }_{ec}^d\mathrm{\Omega }_{ab}^e.$$
(97)
Considering (90)-(93) it is straightforward to see that these formulae are reduced to
$$_{abc}^0=_bg_{ac}+_cg_{ab}g_{db}\mathrm{\Gamma }_{ac}^dg_{db}T_{ac}^d+g_{dc}\mathrm{\Gamma }_{ab}^d+g_{dc}T_{ab}^d,$$
(98)
$$_{a0c}^0=_0g_{ac},$$
(99)
$$_{a0c}^b=_0\mathrm{\Omega }_{ac}^b$$
(100)
and
$$_{abc}^d=R_{abc}^d\delta _b^dg_{ac}+\delta _c^dg_{ab}+D_bT_{ac}^dD_cT_{ab}^d+T_{eb}^dT_{ac}^eT_{ec}^dT_{ab}^e,$$
(101)
respectively. Here, we recall that $`D_a`$ denotes a covariant derivative in terms of $`\mathrm{\Gamma }_{ac}^d.`$ The Equation (86) implies $`_0g_{ac}=0`$ and $`_0\mathrm{\Omega }_{ac}^b=0,`$ that is , $`g_{ac}`$ and $`\mathrm{\Omega }_{ac}^b`$ are independents of $`x^0.`$ This result remind us the dimensional reduction procedure in Kaluza-Klein theory.
Let us now focus in (98). Since $`T_{ac}^d`$ is completely antisymmetric, using (86) the equation (98) leads to
$$_bg_{ac}_cg_{ab}+g_{db}\mathrm{\Gamma }_{ac}^dg_{dc}\mathrm{\Gamma }_{ab}^d=2T_{cab},$$
(102)
Combining the indices in (102) we also get
$$_ag_{bc}_cg_{ba}+g_{da}\mathrm{\Gamma }_{bc}^dg_{dc}\mathrm{\Gamma }_{ba}^d=2T_{cba},$$
(103)
Thus, adding these two expressions we obtain the equation
$$_bg_{ac}+_ag_{bc}2_cg_{ab}+g_{db}\mathrm{\Gamma }_{ac}^d+g_{da}\mathrm{\Gamma }_{bc}^d2g_{dc}\mathrm{\Gamma }_{ab}^d=0,$$
(104)
whose solution is
$$\mathrm{\Gamma }_{cab}=\frac{1}{2}(_ag_{bc}+_bg_{ac}_cg_{ab}).$$
(105)
We recognize in this expression the traditional definition of Christoffel symbols. Moreover, it is well known that this expression is equivalent to the metricity condition
$$D_ag_{bc}=0,$$
(106)
Therefore, we have shown that the metricity condition follows from the equation (98).
Consider now the expression (101). Using (86) we get
$$R_{abc}^d\delta _b^dg_{ac}+\delta _c^dg_{ab}+D_bT_{ac}^dD_cT_{ab}^d+T_{eb}^dT_{ac}^eT_{ec}^dT_{ab}^e=0.$$
(107)
For a homogenous space we have
$$R_{abc}^d=\gamma (\delta _b^dg_{ac}\delta _c^dg_{ab}),$$
(108)
where $`\gamma `$ is a constant. Thus, introducing a new constant $`\gamma ^{}=\gamma 1`$ the equation (107) becomes
$$\gamma ^{}(\delta _b^dg_{ac}\delta _c^dg_{ab})+D_bT_{ac}^dD_cT_{ab}^d+T_{eb}^dT_{ac}^eT_{ec}^dT_{ab}^e=0.$$
(109)
We recognize this expression as the equation (57). Hence, it is straightforward to prove that expression (109) implies the Cartan-Shouten equations. Therefore, we have shown that the metricity condition (106) and the Cartan-Shouten equations follow from (86)-(88).
VI. COMMENTS
It is known that Hurwitz theorem is closely related to the generalized Frobenius theorem (see ref. 12 and references there in): Every alternative division algebra is isomorphic to one of the following : the algebra of real numbers, the algebra of complex numbers, the quaternions, and the Cayley numbers. In fact, using Hurwitz theorem the generalized Frobenius theorem can be proved . Therefore, our procedure also gives an alternative proof of such a generalized theorem. Let just show how our procedure can be used to clarify such a relation.
Alternative algebras can be defined by means of the associator
$$(e_i,e_j,e_k)=(e_ie_j)e_ke_i(e_je_k)_{ijk}^le_l.$$
(110)
In fact, if $`_{ijkl}=`$ $`\delta _{lm}_{ijk}^m`$ is completely antisymmetric for exchanges of any two indices then the algebra is called alternative. Using (4) and (6) one can show that (86) is equivalent to
$$_{ijkl}=C_{ij}^mC_{mkl}C_{jk}^mC_{iml}.$$
(111)
Now, in section II we showed that normed algebra with an identity implies that $`C_{0j}^m=\delta _j^m,C_{j0}^m=\delta _j^m`$ and $`C_{ab}^0=\delta _{ab}`$ and that $`C_{ab}^c`$ is a completely antisymmetric quantity satisfying (32). From these conditions it follows that
$$_{abcd}=C_{ab}^mC_{mcd}C_{bc}^mC_{amd}.$$
(112)
are the only non-vanishing components of $`_{ijkl}.`$ Using (32), it is not difficult to show that this expression leads to
$$_{abcd}=2\{C_{ab}^eC_{cd}^f\delta _{ef}(\delta _{ac}\delta _{bd}\delta _{ad}\delta _{cb})\}=2F_{abcd},$$
(113)
where $`F_{abcd}`$ has been defined in (68). In section IV, we proved that (113) can be obtained from any normed algebra. Therefore since (111) is equivalent to (113) we have shown that a normed algebra with an identity is alternative algebra. The fact that a normed algebra is a division algebra can be proved directly from the composition law $`<ABAB>=<AA><BB>.`$ Indeed, if $`AB=0`$ the composition law implies that $`<AA>=0`$ or $`<BB>=0,`$ which means that $`A=0`$ or $`B=0.`$ Thus, our procedure based in tensor analysis gives a straightforward proof of the fact that a normed algebra with an identity is an alternative division algebra.
It may be interesting to apply the procedure presented in this paper in different contexts. For instance, it may be helpful to through some light on the Blencowe-Duff conjecture<sup>4</sup>: Do the four forces in Nature correspond to the four division algebras? In fact, part of the motivation of this work arose as an effort for answering this question. It is known<sup>20</sup> that using an algebraic topology called K-theory<sup>21</sup> we find that the only dimensions for division algebras structures on Euclidean spaces are again 1, 2, 4, and 8. Therefore, it may be also interesting to relate the present work to K-theory. Moreover, it is known that Englert’s solution of eleven dimensional supergravity achieves the riemannian curvature-less but torsion-full Cartan geometries of absolute parallelism on S<sup>7</sup>. Therefore, it may be interesting to see if the present work may shed some light to clarify some aspects of eleven dimensional supergravity which, as it is known, is the low energy limit theory of M-theory<sup>22-27</sup>. It also seems interesting to see if tensor analysis may be useful to study the zero divisors of Cayley-Dickson algebras<sup>28</sup>and Hopf maps. Let briefly outline this last possibility.
The Cayley-Dickson algebras are defined by the product
$$AB=(A_1B_1\overline{A}_2B_2,B_2A_1+A_2\overline{B}_1),$$
(114)
where $`A=(A_1,A_2)`$ and $`B=(B_1,B_2)`$ are in R$`{}_{}{}^{2^n}=`$R$`{}_{}{}^{2^{n1}}\times `$ R$`^{2^{n1}}`$ and $`\overline{A}=(\overline{A}_1,A_2).`$ Let us denote an algebra with this structure by $`A_n`$. It is found that $`A_0=`$real numbers $`R`$, $`A_1=`$ complex numbers, $`A_2`$ =quaternions and $`A_3=`$ octonions. A Hopf map is defined as
$$F_n:A_n\times A_nA_n\times A_o$$
(115)
$`F_n=(2AB,<BB><AA>).`$
Consider the multiplication table
$$e_ie_j=D_{ij}^\alpha e_\alpha ,$$
(116)
where $`D_{ij}^\alpha `$ are the structure constants, with $`i,j,k=0,1,\mathrm{},2^n1`$ and $`\alpha ,\beta =0,1,\mathrm{},2^n`$. Suppose $`D_{ij}^\alpha `$ satisfies the conditions
$$D_{ij}^{2^n}=\delta _{ij}$$
(117)
and
$$D_{ij}^k=D_{ji}^k,$$
(118)
where in (117) we set $`\alpha =2^n`$. Now, take $`H^i=B^i+A^i`$ and $`G^i=B^iA^i`$ and consider the product
$$F^\alpha =H^iG^jD_{ij}^\alpha .$$
(119)
Using (117) and (118) we find
$$F^{2^n}=<BB><AA>$$
(120)
and
$$F^k=2A^iB^jD_{ij}^k.$$
(121)
Therefore, $`F^\alpha `$ defined in (119) reproduces the Hopf map. It remains to find the relation between $`D_{ij}^k`$ and the Cayley-Dickson product. At this respect, our final goal is to see if our procedure may shed some light on the Hopf maps
$$\begin{array}{c}F_0:S^1S^1,\\ \\ F_1:S^3S^2,\\ \\ F_2:S^7S^4,\\ \\ F_3:S^{15}S^8.\end{array}$$
(122)
which have a certain topological invariant , the Hopf invariant, equal to one.
Finally, it may be interesting to find the connection between the present paper and the Wolf‘s works of references 29 and 30, in which the Cartan-Shouten formalism is generalized to the case of non-Euclidean spaces. Moreover, a possible connection between our procedure and flexible Malcev-admissible algebras (see references 31, 32 and 33 and references there in) may deserve further research.
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# Numerical study of the spherically-symmetric Gross-Pitaevskii equation in two space dimensions
## I Introduction
Recent experiments of Bose-Einstein condensation (BEC) in dilute bosonic atoms (alkali and hydrogen atoms) employing magnetic traps at ultra-low temperatures have intensified theoretical investigations on various aspects of the condensate . The properties of the condensate are usually described by the nonlinear mean-field Gross-Pitaevskii (GP) equation , which properly incorporates the trap potential as well as the interaction among the atoms. The GP equation in both time-dependent and independent forms is formally similar to the Schrödinger equation with a nonlinear term. The effect of the interaction leads to the nonlinear term, which complicates the solution procedure. There have been several numerical studies of the GP equation in three space dimensions .
A Bose gas in lower dimensions $``$ one and two dimensions $``$ exhibits unusual features. For an ideal Bose gas BEC cannot occur in one and two space dimensions at a finite temperature because of thermal fluctuations . The absence of BEC in one and two space dimensions has also been established for interacting uniform systems . However, condensation can take place under the action of a trap potential both for an ideal as well as interacting Bose gas.
Although, there has been no experimental realization of BEC in two space dimensions, this is a problem of great theoretical and experimental interests. In a usual experiment of BEC in three space dimensions under the action of a magnetic trap the typical thermal energy $`k_BT_c`$ is assumed to be much larger than energy of oscillator quantum $`\mathrm{}\omega `$, where $`k_B`$ is the Boltzmann constant, $`T_c`$ the critical temperature, and $`\omega `$ the oscillator frequency. This will allow thermal oscillation in all three directions. Usually, in a typical experimental situation the oscillator frequencies in three different directions, $`x`$, $`y`$, and $`z`$, are different. It is possible to obtain a quasi-two-dimensional BEC in a real three-dimensional trap by choosing the frequency in the third direction $`\omega _z`$ to satisfy $`\mathrm{}\omega _z>k_BT_c>\mathrm{}\omega _x,\mathrm{}\omega _y`$. In that case the energy for thermal fluctuation is much smaller than the oscillator energy in the $`z`$ direction. Consequently, any motion in the $`z`$ direction will be frozen and this will lead to a realization of BEC in two space dimensions. The main features of BEC in two dimensions under the action of a harmonic trap has been discussed by Mullin recently . Also, there has been consideration of BEC in low-dimensional systems for particles confined by gravitational field or by a rotational container . Possible experimental configurations for BEC in spin-polarized hydrogen in two dimensions are currently being discussed .
Recent numerical studies of the GP equation in three space dimensions in time-independent and time-dependent forms have emphasized that extensive care in numerical integration is needed to obtain good convergence. With the viability of experimental detection of BEC in two space dimensions , here we perform a numerical study of the time-dependent and time-independent GP equation in two space dimensions for an interacting Bose gas under the action of a harmonic oscillator trap potential. The interatomic interaction is taken to be both attractive and repulsive in nature.
The nonlinear time-dependent and time-independent GP equations can be compared with the corresponding two types of the linear Schrödinger equation. The stationary states in both cases have a trivial time dependence of the form $`\mathrm{\Psi }(𝐫,t)=\mathrm{exp}(iEt/\mathrm{})\mathrm{\Psi }(𝐫)`$ where $`E`$ is the parametric energy and $`t`$ the time. As is well known the time-independent form of these equations determines the stationary function $`\mathrm{\Psi }(𝐫)`$, as in the hydrogen-atom problem. The time-dependent Schrödinger equation can also be directly solved to obtain the full time-dependent solution in the case of the stationary problems, from which the trivial time dependence $`\mathrm{exp}(iEt/\mathrm{})`$ can be separated. In fact, the time-dependent methods have been successfully used for the bound-state calculation in many areas of computational quantum chemistry . This way of extracting the stationary solution from the linear time-dependent Schrödinger equation continues as a powerful technique in the case of nonlinear time-dependent GP equations.
In this paper we solve the stationary BEC problem in two dimensions using both the time-dependent and time-independent GP equations in the cases of attractive and repulsive interatomic interactions and compare the two types of solutions. The time-independent GP equation is solved by integrating it with the Runge-Kutta rule complimented by the known boundary conditions at origin and infinity . The time-dependent GP equation is solved by discretization and Gauss elimination method with the Crank-Nicholson-type rule complimented again by the known boundary conditions . We find that both the time-dependent and time-independent approaches lead to good convergence for the stationary bound-state problem of the condensate. We also compare these solutions with the Thomas-Fermi approximation in the case of repulsive interatomic interaction.
In addition to obtaining the solution of the stationary problem the time-dependent GP equation can be used to study the intrinsic time-evolution problems with nontrivial time dependence and in this paper the time-dependent approach is also used to study some evolution problems. Specifically, we study the effect of suddenly altering the trapping energy on a preformed condensate. We find that in this case instead of executing sinusoidal oscillations between the stable initial and final configurations as in standard time-evolution problems governed by the linear Schrödinger equation, the condensate executes oscillations around the stable initial and final configurations with ever-growing amplitude.
In Sec. II we describe the time-dependent and time-independent forms of the GP equation. In Sec. III we describe the numerical method in some detail. In Sec. IV we report the numerical results and finally, in Sec. V we give a summary of our investigation.
## II Nonlinear Gross-Pitaevskii Equation
At zero temperature, the time-dependent Bose-Einstein condensate wave function $`\mathrm{\Psi }(𝐫,\tau )`$ at position $`𝐫`$ and time $`\tau `$ may be described by the self-consistent mean-field nonlinear GP equation . In the presence of a magnetic trap this equation is written as
$`[{\displaystyle \frac{\mathrm{}^2}{2m}}^2`$ $`+`$ $`{\displaystyle \frac{1}{2}}m\omega ^2r^2+gN|\mathrm{\Psi }(𝐫,\tau )|^2`$ (1)
$``$ $`i\mathrm{}{\displaystyle \frac{}{\tau }}]\mathrm{\Psi }(𝐫,\tau )=0.`$ (2)
Here $`m`$ is the mass of a single bosonic atom, $`N`$ the number of atoms in the condensate, $`m\omega ^2r^2/2`$ the attractive harmonic-oscillator trap potential, $`\omega `$ the oscillator frequency, and $`g`$ the strength of interatomic interaction. A positive $`g`$ corresponds to a repulsive interaction and a negative $`g`$ to an attractive interaction. The normalization condition of the wave function is
$$_0^{\mathrm{}}𝑑𝐫|\mathrm{\Psi }(𝐫,t)|^2=1.$$
(3)
For a stationary solution the time dependence of the wave function is given by $`\mathrm{\Psi }(𝐫,\tau )=\mathrm{exp}(i\mu \tau /\mathrm{})\mathrm{\Psi }(𝐫)`$ where $`\mu `$ is the chemical potential of the condensate. If we use this form of the wave function in Eq. (1), we obtain the following stationary nonlinear time-independent GP equation :
$`\left[{\displaystyle \frac{\mathrm{}^2}{2m}}^2+{\displaystyle \frac{1}{2}}m\omega ^2r^2+gN|\mathrm{\Psi }(𝐫)|^2\mu \right]\mathrm{\Psi }(𝐫)=0.`$ (4)
The time-dependent equation (1) is equally useful for obtaining a stationary solution with trivial time dependence as well as for studying evolution processes with explicit time dependence.
Here we shall be interested in the spherically symmetric solution $`\mathrm{\Psi }(𝐫,\tau )\phi (r,\tau )=\phi (r)\mathrm{exp}(i\mu \tau /\mathrm{})`$ to Eqs. (1) and (4), which can be written, respectively, as
$`[{\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}r{\displaystyle \frac{}{r}}`$ $`+`$ $`{\displaystyle \frac{1}{2}}m\omega ^2r^2+gN|\phi (r,t)|^2`$ (5)
$``$ $`i\mathrm{}{\displaystyle \frac{}{\tau }}]\phi (r,\tau )=0,`$ (6)
$`[{\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}r{\displaystyle \frac{d}{dr}}`$ $`+`$ $`{\displaystyle \frac{1}{2}}m\omega ^2r^2+gN|\phi (r)|^2`$ (7)
$``$ $`\mu ]\phi (r)=0.`$ (8)
The above limitation to the spherically symmetric solution (in zero angular momentum state) reduces the GP equations in two physical space dimensions to one-dimensional differential equations. We shall study numerically these one-dimensional equations in this paper.
As in Ref. , it is convenient to use dimensionless variables defined by $`x=r/a_{\text{ho}}`$, and $`t=\tau \omega /2,`$ where $`a_{\text{ho}}\sqrt{\mathrm{}/(m\omega )}`$, $`\alpha =\mu /(\mathrm{}\omega )`$, $`\psi (x)=a_{\text{ho}}\sqrt{2mgN}\phi (r)/\mathrm{}`$, and $`\psi (x,t)=a_{\text{ho}}\sqrt{2\pi }\phi (r,\tau )`$. In terms of these variables Eqs. (5) and (7) becomes, respectively,
$`\left[{\displaystyle \frac{1}{x}}{\displaystyle \frac{}{x}}x{\displaystyle \frac{}{x}}+x^2+cn|\psi (x,t)|^2i{\displaystyle \frac{}{t}}\right]\psi (x,t)=0,`$ (9)
$`\left[{\displaystyle \frac{1}{x}}{\displaystyle \frac{d}{dx}}x{\displaystyle \frac{d}{dx}}+x^2+c|\psi (x)|^22\alpha \right]\psi (x)=0.`$ (10)
where $`nmgN/(\pi \mathrm{}^2)`$ is the reduced number of particles and $`c=\pm 1`$ carries the sign of $`g`$: $`c=1`$ corresponds to a repulsive interaction and $`c=1`$ corresponds to an attractive interaction. The normalization condition (3) of the wave functions become
$$1=_0^{\mathrm{}}|\psi (x,t)|^2x𝑑x=\frac{1}{n}_0^{\mathrm{}}|\psi (x)|^2x𝑑x.$$
(11)
We shall be using these two slightly different normalizations of the time-dependent and time-independent wave functions for future numerical convenience.
An interesting property of the condensate wave function is its mean-square radius defined by
$$x^2=_0^{\mathrm{}}x^2|\psi (x,t)|^2x𝑑x=\frac{1}{n}_0^{\mathrm{}}x^2|\psi (x)|^2x𝑑x.$$
(12)
## III Numerical Method
### A Boundary Condition
Both in time-dependent and time-independent approaches we need the boundary conditions of the wave function as $`x0`$ and $`\mathrm{}`$. For a confined condensate, for a sufficiently large $`x`$, $`\psi (x)`$ must vanish asymptotically. Hence the nonlinear term proportional to $`|\psi (x)|^3`$ can eventually be neglected in the GP equation for large $`x`$ and Eq. (10) becomes
$`\left[{\displaystyle \frac{1}{x}}{\displaystyle \frac{d}{dx}}x{\displaystyle \frac{d}{dx}}+x^22\alpha \right]\psi (x)=0.`$ (13)
This is the equation for the oscillator in two space dimensions in the spherically symmetric state with solutions for $`\alpha =1,3,5,\mathrm{}`$ etc. A general classification of all the states of such an oscillator is well under control . In the present BEC problem, Eq. (13) determines only the asymptotic behavior. If we consider Eq. (13) as a mathematical equation valid for all $`\alpha `$ and large $`x`$, the asymptotic form of the physically acceptable solution is given by
$$\underset{x\mathrm{}}{lim}\psi (x)=N_C\mathrm{exp}\left[\frac{x^2}{2}+(\alpha 1)\mathrm{ln}x\right],$$
(14)
where $`N_C`$ is a normalization constant. Equation (14) leads to the following asymptotic log-derivative
$$\underset{x\mathrm{}}{lim}\frac{\psi ^{}(x)}{\psi (x)}=\left[x+\frac{\alpha 1}{x}\right],$$
(15)
which is independent of the constant $`N_C`$ and where the prime denotes derivative with respect to $`x`$.
Next we consider Eq. (10) as $`x0`$. The nonlinear term approaches a constant in this limit because of the regularity of the wave function at $`x=0`$. Then one has the following usual conditions
$$\psi (0)=\text{constant},\psi ^{}(0)=0,$$
(16)
as in the case of the harmonic oscillator problem in two space dimensions . Both the small- and large-$`x`$ behaviors of the wave function will be necessary for a numerical solution of the GP equation in time-dependent and time-independent forms.
### B Time-Dependent Approach: Evolution and Stationary Problems
First we describe the numerical method for solving the time-dependent equation (9). For a numerical solution it is convenient to make the substitution $`\psi (x,t)\varphi (x,t)/x`$ in this equation, when this equation becomes
$`[{\displaystyle \frac{^2}{x^2}}`$ $`+`$ $`{\displaystyle \frac{1}{x}}{\displaystyle \frac{}{x}}{\displaystyle \frac{1}{x^2}}+x^2+cn{\displaystyle \frac{|\varphi (x,t)|^2}{x^2}}`$ (17)
$``$ $`i{\displaystyle \frac{}{t}}]\varphi (x,t)=0.`$ (18)
A convenient way to solve Eq. (17) numerically is to discretize it in both space and time and reduce it to a set of algebraic equations which could then be solved by using the known asymptotic boundary conditions. We discretize this equation by using a space step $`h`$ and time step $`\mathrm{\Delta }`$ with a finite difference scheme using the unknown $`\varphi _j^k`$ which will be approximation of the exact solution $`\varphi (x_j,t_k)`$ where $`x_j=jh`$ and $`t_k=k\mathrm{\Delta }`$. As Eq. (17) involves both time and space variables it can be discretized in more than one way. The time derivative in Eq. (17) involves the wave function at times $`t_k`$ and $`t_k+\mathrm{\Delta }`$. As $`\mathrm{\Delta }`$ is small, the time-independent operations in this equation can be discretized by using the wave-function components at time $`t_k`$ or $`t_{k+1}t_k+\mathrm{\Delta }`$. If one uses the wave-function components at time $`t_k`$, Eq. (17) is discretized as
$`{\displaystyle \frac{i(\varphi _j^{k+1}\varphi _j^k)}{\mathrm{\Delta }}}`$ $`=`$ $`{\displaystyle \frac{1}{h^2}}\left[\varphi _{j+1}^k2\varphi _j^k+\varphi _{j1}^k\right]`$ (19)
$`+`$ $`{\displaystyle \frac{1}{2x_jh}}\left[\varphi _{j+1}^k\varphi _{j1}^k\right]`$ (20)
$`+`$ $`\left[x_j^2{\displaystyle \frac{1}{x_j^2}}+cn{\displaystyle \frac{|\varphi _j^k|^2}{x_j^2}}\right]\varphi _j^k.`$ (21)
This is an explicit differencing scheme, since, given $`\varphi `$ at $`t_k`$ it is straightforward to solve for $`\varphi `$ at $`t_{k+1}`$ . One should start with an approximately known solution at $`t_k`$ and propagate it in time until a converged solution is reached. We confirm in our study that this simple scheme leads to slow convergence and large unphysical oscillations in the solution.
One can express the derivatives on the right-hand-side of Eq. (21) in terms of the variables at time $`t_{k+1}`$ . Then the unknown $`\varphi _j^{k+1}`$ appears on both sides of the equation and one has an implicit scheme. We find that the implicit scheme improves substantially the numerical accuracy and convergence rate. However, we find after some experimentation that if the right-hand-side of Eq. (21) is averaged over times $`t_k`$ and $`t_{k+1}`$ one has the best convergence. This is a semi-implicit scheme based on the Crank-Nicholson scheme for discretization . We use the following rule to discretize the partial differential equation (17)
$`{\displaystyle \frac{i(\varphi _j^{k+1}\varphi _j^k)}{\mathrm{\Delta }}}`$ $`=`$ $`{\displaystyle \frac{1}{2h^2}}[(\varphi _{j+1}^{k+1}2\varphi _j^{k+1}+\varphi _{j1}^{k+1})`$ (22)
$`+`$ $`(\varphi _{j+1}^k2\varphi _j^k+\varphi _{j1}^k)]`$ (23)
$`+`$ $`{\displaystyle \frac{1}{4x_jh}}\left[(\varphi _{j+1}^{k+1}\varphi _{j1}^{k+1})+(\varphi _{j+1}^k\varphi _{j1}^k)\right]`$ (24)
$`+`$ $`{\displaystyle \frac{1}{2}}\left[x_j^2{\displaystyle \frac{1}{x_j^2}}+cn{\displaystyle \frac{|\varphi _j^k|^2}{x_j^2}}\right](\varphi _j^{k+1}+\varphi _j^k).`$ (25)
Similar discretization rule has been used for the solution of the GP equation in three space dimensions . The first and second space derivatives of the wave function as well as the the wave function itself have been approximated by the average over their values at the initial time $`t_k`$ and the final time $`t_{k+1}`$. This procedure leads to accurate and stable numerical results. Considering that the wave function is known at time $`t_k`$, Eq. (LABEL:f) is an equation in three unknowns $``$ $`\varphi _{j+1}^{k+1},\varphi _j^{k+1}`$ and $`\varphi _{j1}^{k+1}`$. In a lattice of $`N`$ points Eq. (LABEL:f) represents a tridiagonal set for $`j=2,3,\mathrm{},(N1)`$. This set has a unique solution if the wave functions at the two end points $`\varphi _1^{k+1}`$ and $`\varphi _N^{k+1}`$ are known. In the present problem these values at the end points are provided by the known asymptotic conditions. The tridiagonal set of equations is solved by the Gauss elimination method and back substitution using a typical space step $`h=0.0001`$ and time step $`\mathrm{\Delta }=0.03`$. Although, the iterative method should work for any value of $`\mathrm{\Delta }`$, we found the convergence to be faster with this value of $`\mathrm{\Delta }`$ and we used this value throughout the present investigation.
The time-dependent method could be used to study stationary as well as time-evolution problems. First we consider the stationary problem. For the ground and the first excited states of the condensate we start with the following analytically known wave functions of the harmonic oscillator problem (13) :
$$\varphi (x)=x\psi (x)=\sqrt{2}x\mathrm{exp}(x^2/2),$$
(27)
$$\varphi (x)=x\psi (x)=\sqrt{2}x(1x^2)\mathrm{exp}(x^2/2),$$
(28)
respectively, at an initial time $`t=0`$. We then repeatedly propagate these solutions in time using the Crank-Nicholson-type algorithm (LABEL:f). The boundary condition (16), that $`\varphi (0)=0`$, is implemented at each time step . Also, the solution at each time step will satisfy the asymptotic condition (14). Starting with $`cn=0`$, at each time step we increase or decrease the nonlinear constant $`cn`$ by an amount $`\mathrm{\Delta }_1`$ typically around 0.01. This procedure is continued until the desired final value of $`cn`$ is reached. Then the final solution is iterated several times (between 10 to 40 times) to obtain a stable converged result. The resulting solution is the ground state of the condensate corresponding to the specific nonlinear constant $`cn`$. We found the convergence to be fast for small $`|cn|`$. However, the final convergence of the scheme breaks down if $`|cn|`$ is too large. In practice these difficulties start for $`cn>20`$ for the ground state for a positive $`c`$ (repulsive interaction) in a computational analysis in double precision. For an attractive interaction there is no such problem as the GP equation does not sustain a large nonlinearity $`|cn|`$ as we comment in detail in the next section.
As the time dependence of these stationary states is trivial $``$ $`\psi (x,t)=\psi (x)\mathrm{exp}(i2\alpha t)`$ $``$ the chemical potential $`\alpha `$ can be obtained from the propagation of the converged ground-state solution at two successive times, e.g., $`\psi (x,t_k)`$ and $`\psi (x,t_{k+1})`$. From the numerically obtained ratio $`\psi (x,t_k)/\psi (x,t_{k+1})=\mathrm{exp}(i2\alpha \mathrm{\Delta })`$ $`\alpha `$ can be obtained as the time step $`\mathrm{\Delta }`$ is known.
The time-dependent method could also be used to study evolution problems. One such evolution problem describes the fate of the condensate if the trap potential is removed or altered suddenly after the formation of the condensate. As a stable condensate is formed under the action of the trap potential, after a sudden change in the trap potential, the condensate will gradually modify with time. To study the time evolution of a condensate wave function as the trap is removed or altered suddenly, we have to start the time evolution of the known precalculated wave function of the condensate with the initial trap potential and allow it to evolve in time using the time-dependent GP equation with the full nonlinearity but with the altered trap potential, which could be zero.
### C Time-Independent Approach
The time-independent GP equation (10) has the following structure
$`y^{}=G(x,\psi (y)),`$ (29)
with $`y=x\psi ^{}`$, where the prime denotes the $`x`$ derivative. With this realization, a numerical integration of Eq. (10) can be implemented using the following four-point Runge-Kutta rule in steps of $`h`$ from $`x_j`$ to $`x_{j+1}`$
$`\psi _{j+2}`$ $`=`$ $`\psi _{j+1}+h\psi _{j+1}^{},`$ (30)
$`x_{j+1}\psi _{j+1}^{}`$ $`=`$ $`x_j\psi _j^{}+{\displaystyle \frac{1}{6}}(s_0+2s_1+2s_2+s_3),`$ (31)
where
$`s_0`$ $`=`$ $`hG(x_j,\psi _j),`$ (32)
$`s_1`$ $`=`$ $`hG[x_j+{\displaystyle \frac{h}{2}},\psi _j+{\displaystyle \frac{h(x_j\psi _j^{}+s_0/2)}{2(x_j+h/2)}}],`$ (33)
$`s_2`$ $`=`$ $`hG[x_j+{\displaystyle \frac{h}{2}},\psi _i+{\displaystyle \frac{h(x_j\psi _j^{}+s_1/2)}{2(x_j+h/2)}}],`$ (34)
$`s_3`$ $`=`$ $`hG[x_j+h,\psi _j+{\displaystyle \frac{h(x_j\psi _j^{}+s_2)}{(x_j+h)}}].`$ (35)
Equation (10) is integrated numerically for a given $`\alpha `$ using this algorithm starting at the origin ($`x=0`$) with the initial boundary condition (16) with a trial $`\psi (0)`$ and a typical space step $`h=`$ 0.0001. The integration is propagated to $`x=x_{\text{max}}`$, where the asymptotic condition (15) is valid. The agreement between the numerically calculated log-derivative of the wave function and the theoretical result (15) was enforced to five significant figures. The maximum value of $`x`$, up to which we needed to integrate (10) numerically for obtaining this precision, is $`x_{\text{max}}=5`$. If for a trial $`\psi (0)`$, the agreement of the log-derivative can not be obtained, a new value of $`\psi (0)`$ is to be chosen. The proper choice of $`\psi (0)`$ was implemented by the secant method. Even with this method, sometimes it is difficult to obtain the proper value of $`\psi (0)`$ for a given $`\alpha `$. Unless the initial guess is “right” and one is sufficiently near the desired solution, the method could fail, specially, for large $`|cn|`$ and lead numerically to either the trivial solution $`\psi (x)=0`$ or an exponentially divergent nonnormalizable solution in the asymptotic region.
## IV Numerical Result
### A Stationary Problem
First we consider the ground-state solution of Eq. (10) for different $`\alpha `$ in cases of both attractive and repulsive interactions using the time-independent method. In the presence of the nonlinearity, for attractive (repulsive) interatomic interaction, the solutions of the GP equation for the ground state appear for values of chemical potential $`\alpha <1`$ ($`\alpha >1`$). The relevant parameters for the solutions $``$ the wave-function at the origin $`\psi (0)`$, reduced number $`n`$, and mean-square radii $`x^2`$ $``$ are listed in Table I. The numerical integration was performed up to $`x_{\text{max}}=5`$ with $`h=0.0001`$ where the asymptotic boundary condition (15) is implemented.
Using the known tabulated values of $`n`$ in each case we also solved the time-dependent GP equation and the wave functions and energies so calculated agree well with the respective quantities calculated with the time-independent approach. The solutions were obtained using space step $`h=0.0001`$, time step $`\mathrm{\Delta }=0.03`$ and the parameter $`\mathrm{\Delta }_10.01`$. The largest value of $`x`$ used in discretization (LABEL:f) is $`x_{\text{max}}=10`$. The wave functions for different values of $`\alpha `$ (and $`n`$) for the attractive and repulsive interparticle interactions for the cases shown in Table I are exhibited in Figs. 1(a) and 1(b), respectively, where we plot $`\psi (x)`$ versus $`x`$ using the time-dependent and time-independent approaches. The curves in Figs. 1(a) and 1(b) appear in the same order as the rows in Table I and it is easy to identify the corresponding values of $`\alpha `$ from the values of $`\psi (0)`$ of each curve. From Figs. 1(a) and (b) we find that the nature of the wave function for these two cases are quite different. However, the wave functions calculated with time-dependent and time-independent approaches agree reasonably with each other.
Fig. 1. Ground-state condensate wave function $`\psi (x)`$ versus $`x`$ for (a) attractive and (b) repulsive interparticle interactions using the time-dependent (dashed line) and time-independent (full line) approaches. The parameters for these cases are given in Table I. In the time-dependent method we used time step $`\mathrm{\Delta }=0.03,\mathrm{\Delta }_1=0.01`$, space step $`h=0.0001`$ and $`x_{\text{max}}=8`$, in the time-independent method we used space step $`h=0.0001`$ and $`x_{\text{max}}=5.`$ In the case of the repulsive interparticle interaction we also show the solution (36) corresponding to the Thomas-Fermi approximation (dashed-dotted line). The curves appear in same order as in Table I. with the lowermost curve corresponding to the first row.
In the absence of previous solutions of this problem we compare the stationary solutions in the repulsive case ($`c=1`$) with those obtained via a well-known approximation, e.g., the Thomas-Fermi approximation. In this approximation the kinetic energy term in Eq. (10) is neglected and one has the following simple approximate solution
$$\psi (x)=\sqrt{2\alpha x^2},$$
(36)
for $`x^22\alpha `$ and zero otherwise. In Fig. 1 (b) we also plot the Thomas-Fermi approximation (36). We find that as expected, for a large condensate, this approximation is a reasonable approximation. However, it turns out to be a bad approximation for a small condensate.
Table I: Parameters for the numerical solution of the GP equation (10) for $`c=\pm 1`$ for the ground state wave function. The first four columns refer to the attractive interaction $`c=1`$ and the last four columns refer to the repulsive interaction $`c=1`$.
| $`\alpha `$ | $`\psi (0)`$ | $`n`$ | $`x^2`$ | $`\alpha `$ | $`\psi (0)`$ | $`n`$ | $`x^2`$ |
| --- | --- | --- | --- | --- | --- | --- | --- |
| 1.0 | 0 | 0 | 0 | 1.0 | 0 | 0 | 0 |
| 0.8 | 0.9185 | 0.3663 | 0.9030 | 1.2 | 0.8719 | 0.4353 | 1.1027 |
| 0.4 | 1.6795 | 0.9147 | 0.7297 | 1.6 | 1.4415 | 1.5276 | 1.3219 |
| $``$0.4 | 2.8255 | 1.4798 | 0.4757 | 2.2 | 1.9276 | 3.7509 | 1.6741 |
| $``$2.0 | 4.6249 | 1.7695 | 0.2400 | 3.0 | 2.3626 | 7.8377 | 2.1679 |
| $``$4.0 | 6.3252 | 1.8319 | 0.1385 | 4.0 | 2.7786 | 14.7609 | 2.8041 |
It is appropriate to comment on the numerical accuracy of the present time-dependent and independent methods, which seems to be limited typically by the difference between the time-dependent and independent solutions in Fig. 1. When the solution can be obtained numerically, as in the cases shown in Table I, the time-independent method can yield very accurate results. This accuracy can be increased by controlling the space step $`h`$ and $`x_{\text{max}}`$. This is not so in the case of the time-dependent method, where the numerical result exhibits small periodic oscillation after iteration specially for large values of $`|cn|`$ which we detail below.
The numerical solution of the time-dependent method is independent of the space step $`h`$ provided that a typical value around $`h=0.0001`$ is employed as in the present study. No visible difference in the solution is found if $`h`$ is increased by a factor of 2 or 3. However, the solution is more sensitive to the number of time iterations, specially, for a large value of $`|cn|`$, for a fixed integration time step $`\mathrm{\Delta }`$ or the step $`\mathrm{\Delta }_1`$ by which the nonlinear constant in the GP equation is increased at each time step until the final value of $`cn`$ is reached. We show this variation in Figs. 2 (a) and (b) where we plot $`|\psi (x,T)|`$ as a function of reduced time $`Tt/0.03`$ for $`x=0`$ and 2 for different choices of $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }_1`$ in the repulsive case for $`n=3.7509`$ and $`\alpha =2.2`$ corresponding to the fourth row of Table I. The zero of reduced time $`T`$ is made to coincide with the time step $`t_k`$ at which the full nonlinear constant $`cn`$ is obtained for the first time during iteration. This choice of time will allow us to compare the fluctuations of the solution during the time propagation of the full GP equation. In Fig. 2(a) we present our results for $`\mathrm{\Delta }=0.03`$ and for $`\mathrm{\Delta }_1=`$ 0.018754 and 0.0046886. In Fig. 2(b) we present our results for $`\mathrm{\Delta }_1=0.0046886`$ and for $`\mathrm{\Delta }=0.03,`$ and 0.05. From Figs. 2 (a) and (b) we find that there is numerical oscillation of the solution with time in this approach which is independent of small variations of $`\mathrm{\Delta }`$ near 0.03 and $`\mathrm{\Delta }_1`$ around 0.01. These oscillations determine the numerical error of the time-dependent approach and become larger when we employ a $`\mathrm{\Delta }`$ very different from 0.03, or $`\mathrm{\Delta }_1`$ very different from 0.01. The oscillations can really be large if an improper value of step $`\mathrm{\Delta }`$ or $`\mathrm{\Delta }_1`$ is choosen as can be seen from Fig, 2(b) for $`\mathrm{\Delta }=0.05`$. The results remain stable if we reduce these steps up to $`\mathrm{\Delta }0.01`$ and $`\mathrm{\Delta }_10.003`$. For very small $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }`$ accumulative errors also increase. This accumulative numerical error increases as the number of iterations is very large (several thousands) and a large number of iterations is needed to cover a given time interval with a small time step $`\mathrm{\Delta }`$.
Table II: Amplitude of oscillation $`A(x,T)`$ (in units of 0.01) of $`|\psi (x,T)|`$ at different times $`T`$ for $`x=0`$ and 2 calculated with $`\mathrm{\Delta }=0.03`$ and $`\mathrm{\Delta }_1=0.0046886`$ in the repulsive case for $`cn=3.7509`$. The average value of converged $`|\psi (0,T)|=1.9310`$ and $`|\psi (2,T)|=0.6667.`$
| $`T=`$ | 0 | 167 | 294 | 406 | 533 | 645 | 791 |
| --- | --- | --- | --- | --- | --- | --- | --- |
| $`|A(0,T)|`$ | 1.54 | 1.27 | 3.32 | 2.93 | 4.13 | 4.49 | 7.35 |
| $`T=`$ | 0 | 162 | 291 | 400 | 536 | 637 | 789 |
| $`|A(2,T)|`$ | 0.77 | 0.95 | 1.13 | 1.33 | 1.43 | 2.01 | 2.77 |
Fig. 2. Ground-state condensate wave function $`|\psi (x,T)|`$ versus reduced time $`Tt/0.03`$ for $`x=0`$ and 2 in the repulsive case for the nonlinear constant $`cn=3.7509`$ for (a) $`\mathrm{\Delta }=0.03`$, and $`\mathrm{\Delta }_1=0.0046886`$ (full line), and 0.018754 (dashed line) and (b) $`\mathrm{\Delta }_1=0.0046886`$ and for $`\mathrm{\Delta }=`$ 0.03 (full line), and 0.05 (dashed line). The zero of $`T`$ is taken to be the time at which the full nonlinearity is achieved for the first time.
We show a quantitative account of the above oscillation in Table II where we plot the maximum error in $`|\psi (x,T)|`$ (amplitude of oscillation of $`|\psi (x,T)|`$) for $`x=0`$ and 2 at different times calculated with steps $`\mathrm{\Delta }=0.03`$ and $`\mathrm{\Delta }_1=0.0046886`$. We find that the error increases slowly, but not necessarily monotonically, with time. The average value of the converged $`|\psi (0,T)|`$ is 1.9310 and that for $`|\psi (2,T)|`$ is $`0.6667.`$ The maximum deviations from these values as shown in Table II do not occur at the same values of $`T`$. We find from Table II that for small $`T(0)`$ the maximum average error in $`|\psi (x,T)|`$ is about $`1\%`$. For $`T800`$ this maximum average error could be as high as $`4\%`$. As these errors are oscillating with time, at a given $`T`$ this error could be smaller or even zero. Considering that we are dealing with nonlinear equations these errors are well within the acceptable limits. The errors shown in Table II would also be the typical errors in time-evolution problems which we study in the next subsection.
For repulsive interaction, it was increasingly difficult to find the solution of the GP equation using both time-dependent and time-independent methods for larger nonlinearity than those reported in Figs. 1 (a) and (b). The inputs of the time-independent method are $`\alpha `$ and an appropriate $`\psi (0)`$. In this method it became difficult (or impossible) to find the appropriate $`\psi (0)`$ and find a solution for large $`cn(>20)`$. For large nonlinearity the secant method led to radially excited state for the appropriate $`\psi (0)`$. In the time-dependent method the only input is the value of $`cn`$. For a large $`cn`$ in the repulsive case, the numerically obtained solution for the wave function shows many oscillations and is clearly unacceptable physically. A Crank-Nicholson-type approach was also used to solve the GP equation in three space dimensions . The numerical instability also set a limit in that investigation in finding stationary ground-state solution for large values of nonlinearity.
For attractive interparticle interaction, the wave function is more sharply peaked at $`x=0`$ than in the case of the repulsive interparticle interaction and one has a smaller reduced number $`n`$ and mean square radius $`x^2`$. In this case we find from Table I that with a reduction of the chemical potential $`\alpha `$, the reduced number $`n`$ increases slowly and the mean square radius $`x^2`$ decreases rapidly, so that the density of the condensate $`\rho n/x^2`$ tends to diverge as $`n`$ tends to a maximum value $`n_{\text{max}}`$. The increase in density lowers the interaction energy. The kinetic energy of the system is responsible for the stabilization. As the central density increases further for stronger attractive interparticle interaction, kinetic energy can no longer maintain equilibrium of the system and the system collapses. Consequently, for $`n>n_{\text{max}}`$, there is no stable solution of the GP equation. Numerically, from a plot of $`n`$ versus $`1/\rho `$ we find this maximum number consistent with $`\rho ^1=0`$ to be
$$n_{\text{max}}\eta N_{\text{max}}1.88.$$
(37)
There is no such limit on $`n`$ in the repulsive case. In that case with the increase of the chemical potential $`\alpha `$ the condensate increases in size as the number of particles in the condensate increases. These behaviors of the Bose-Einstein condensate in two dimensions were also noted in three dimensions . However, in three dimensions the corresponding maximum value was $`n_{\text{max}}4N_{\text{max}}|a|/a_{\text{ho}}2.30`$ .
Both the time-dependent and time-independent approaches are equally applicable for spherically-symmetric radially excited states. For the first excited state, with one node in the wave function, we verified that the convergence was as good as in the ground-state case reported here. However, it is a routine study and we do not report the results here.
### B Evolution Problem
Fig. 3. Condensate wave function $`\psi (x)`$ in the repulsive case at different times $`T=t/0.03`$ for an expanding condensate after the trap is removed suddenly at $`T=0`$. The initial condensate has $`n=8`$, and the time evolution is performed using time step $`\mathrm{\Delta }=0.03,\mathrm{\Delta }_1=0.01`$, space step $`h=0.0001`$ and $`x_{\text{max}}=15`$.
Fig. 4. The central probability density $`|\psi (0)|^2`$ in the repulsive case at different times $`T=t/0.03`$ for an oscillating condensate with $`n=8`$ after the trap energy is suddenly reduced to half at $`T=0`$. The $`|\psi (0)|^2`$ for condensates corresponding to the initial and final traps are denoted by the two straight lines. The time evolution is performed using time step $`\mathrm{\Delta }=0.03,\mathrm{\Delta }_1=0.01`$, space step $`h=0.0001`$, and $`x_{\text{max}}=15`$.
Next we consider two time-evolution problems using the time-dependent method. We consider the ground state in the repulsive case with $`cn=8`$. In the first problem, at $`t=0`$, the trap is suddenly removed. In the second, at $`t=0`$, the trap energy is suddenly reduced to half of the starting value. In both cases we study how the system evolves with time by solving the time-dependent GP equation using time step $`\mathrm{\Delta }=0.03`$, $`\mathrm{\Delta }_10.01`$ and space step $`h=0.0001`$. Both these problems are intrinsic time-dependent problems and can be studied numerically and experimentally.
The condensate cannot exist in the absence of the trap. In the first case after the trap is removed at $`t=0`$, the radius of the condensate increases and the wave function extends over a larger region of space. We solve the time-dependent GP equation at different times. In Fig. 3 we plot the wave function at different reduced times $`T=t/0.03`$. The condensate increases in size monotonically with time and eventually disappears.
Fig. 5. Condensate wave function $`\psi (x)`$ at different times $`T=t/0.03`$ for an oscillating condensate after the trap energy is suddenly reduced to half at $`T=0`$. All parameters are the same as in Fig. 3.
In the second problem at $`t=0`$, we reduce the trap energy suddenly to half of the initial value corresponding to a stable final configuration for the condensate in the repulsive case. The system is now found to oscillate between the initial and final stationary states. In the absence of the nonlinearity, the system executes sinusoidal oscillations between the two stable configurations. However, in this nonlinear problem the system executes oscillations with evergrowing amplitude. To illustrate this oscillation we plot in Fig. 4 the the central probability density $`|\psi (0)|^2`$ versus reduced time $`T=t/0.03`$. The $`|\psi (0)|^2`$ for condensates corresponding to the initial and final trap energies are denoted by the two straight lines. We see that the oscillation increases with time. In our numerical study we find that after a very large number of iterations (several thousands) the amplitude may become very large. However, we do not know if this result makes sense physically as the cumulative numerical error of the type shown in Fig. 2 will also grow after a very large number of iterations, which will possibly invalidate our conclusion. However, the solution presented in Fig. 4 is stable numerically and is the acceptable physical solution of the problem after a small number of iterations. This interesting behavior can possibly be observed experimentally and deserves further theoretical and numerical studies. In Fig. 5 we plot the wave functions of the system at different times which have very acceptable and smooth behavior. As the number of particles of the system continues fixed, the wave functions of smaller amplitudes have larger spacial extension \[mean square radius (12)\] so that the normalization condition (11) is preserved.
## V Summary
In this paper we present a numerical study of the Gross-Pitaevskii equation for BEC in two space dimensions under the action of a harmonic oscillator trap potential for bosonic atoms with attractive and repulsive interparticle interactions using time-dependent and time-independent approaches. Both approaches are used for the study of the stationary problem. In addition some evolution problems were studied by the time-dependent approach. We derive the boundary conditions (15) and (16) of the solution of the dimensionless GP equations (9) and (10). These boundary conditions are used for the solution of the stationary problem using both the time-dependent and time-independent approaches.
The time-dependent GP equation is solved by discretizing it using a Crank-Nicholson-type scheme, whereas the time-independent GP equation is solved by numerical integration using the four-point Runge-Kutta rule. In both cases numerical difficulty appears for large nonlinearity ($`cn>20`$). For medium nonlinearity, the accuracy of the time-independent method can be increased by reducing the space step $`h`$. However, the solution of the time-dependent approach exhibits intrinsic oscillation with time iteration which is independent of space and time steps used in discretization.
The ground-state wave function is found to be sharply peaked near the origin for attractive interatomic interaction. For a repulsive interatomic interaction the wave function extends over a larger region of space. In the case of an attractive potential, the mean square radius decreases with an increase of the number of particles in the condensate. Consequently, a stable solution of the GP equation can be obtained for a maximum number of particles in the condensate as given in Eq. (37).
In addition to the stationary problem we studied two evolution problems using the time-dependent approach. A stable bound state is considered and the trap potential is suddenly removed or reduced to half at $`t=0`$. If the trap is removed suddenly, the system gradually and monotonically increases in size with time and eventually it disappears occupying the whole space with zero density. If the trap energy is suddenly reduced to half, the system oscillates around the two stationary positions. The amplitude of the oscillation continue to increase with time. This behavior is interesting and can be studied experimentally in the future.
The work is supported in part by the Conselho Nacional de Desenvolvimento Científico e Tecnolánd Fundação de Amparo à Pesquisa do Estado de São Paulo of Brazil.
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# Entanglement, thermalisation and stationarity: The computational foundations of quantum mechanics
## I Motivation and overview
Everyone knows that “classical theory is *absolutely incapable* of describing the distribution of light from a blackbody” (Feynman, , I-41-2) and just about every other result of physics, but the opinion was formed out of frustration a century ago when we did not have an understanding of classical mechanics itself as we do today. I show below how the notions of holographic imaging and computational irreversibility suffice for an elegant classical interpretation of entanglement and the quantum state space formalism, respectively, and conversely, that the quantum postulates have been largely a substitute for this insight. More particularly, the present theory is reversal of the traditional assumption that our knowledge of physics is more fundamental than the science of knowledge itself.
I consider the computational issue of physical role of the *observer’s physical data states* in the process of observation. By Landauer’s principle, these must be irreversibly altered at every observation, unlike in Heisenberg’s theory, where the concern was only with possible disturbance to the *observed* entity. I shall also show how quantum uncertainty also results from this perspective, providing a precise computational notion of smallness. The states bear representational correspondence to quantum wavefunctions, yielding, as I shall show, Schrödinger’s equation as a computational result, and the constancy of $`h`$ as a principle of scale in communication between observers.
Of particular concern, naturally, is the probabilistic nature of quantum information, which inspired Schrödinger’s cat and many-worlds interpretations, but which, I hope to show with reasonable conviction, is essentially thermal, and in that sense classical. Formal notions of information have always been based on probability arising from thermal motion, beginning with Boltzmann’s work in the kinetic theory, but the mere attribution of classical thermalisation does not mean that the underlying *mechanics* should be classical: indeed, the Fermi-Pasta-Ulam (FPU) problem, the difficulty of modelling thermal diffusion by deterministic simulation, already casts doubt on the adequacy of classical mechanisms to guarantee thermalisation. On the other hand, the notion of irreversibility in Landauer’s theory plays a key role, as I shall show, in accounting for the probabilistic aspects of quantum theory.
Although I do not attempt to dismiss quantum mechanisms *per se*, for example particle creation and annihilation, the result implies that quantum effects must be indistinguishable from those that would result from thermodynamics and computational considerations anyway. In the one instance where particulate views are weakest, viz radiation in vaccuum, the usual attribution of thermalisation to nonlinear interactions becomes logically inapplicable, bypassing the FPU issue. I have shown recently Prasad2000a that classical considerations of dynamic stationarity under wall jitter, needed to account for the thermalisation classically, exactly reproduce both Planck’s law and the harmonic oscillators, which forces me to consider applying classical electromagnetism in other quantum situations as well.
As entanglement is specially of interest in quantum computation, I demonstrate in §II that it is indeed reproduced by classical electromagnetic waves; this is distinct from *all* prior considerations, including those of Bohm Bohm1951 Bohm1952 , and Einstein, Podalsky and Rosen (EPR) EPR1935 , as they invariably assumed particles in the classical picture, and were forced into issues of superluminal communication and hidden variables. No such problems occur in the wave picture, classical or otherwise, showing that it is the particulate perception, not classical mechanics, which has been at fault. The conclusion is formally supported, as discussed in §III, by a mathematical result by BenDaniel, that representable physical information is necessarily one of continuous fields and not particles. I further demonstrate that the nature of quantum information is fundamentally holographic, which makes entanglement purely an artifact of *a posteriori* selection, consistent with a recent proof by Bennett *et al.* of the impossibility of communicating classical information through entanglement BSSTT1999 .
I then establish my main contention, in §III, that the quantum state space is itself fundamentally computational in origin, by reproducing its complex vector space form purely from considerations of the most general representation for information of any kind, as well as Schrödinger’s equation, generalising Dirac’s derivation to any such space. This does not suffice to prove the dependence on $`h`$, for which I revert, in §IV, to the principles of stationarity and antinodal equipartition already established for the cavity spectrum, and apply these to the classical wave formalism, equivalent to the de Broglie waves of matter, obtainable from Hamilton-Jacobi theory.
Lastly, I show how these principles readily yield the quantisation of matter, quantum uncertainty, and the constancy of $`h`$, the last as a transitive result of pairwise interactions among observers and observed entities. In particular, the theory suffices to prove that $`h`$ is a scale factor relating only the dimensions of energy (mass) \[M\] and time \[T\], but not of space \[L\], permitting a fundamental, relative difference of spatial scale in interactions between separated entities (§V), which directly leads to the logical foundations of relativity and shows the current notions of general relativity and cosmology to be too simplistic to have been correct Prasad2000c .
## II Entanglement: 1-bit holography
Accordingly, I shall first show that classical electromagnetic waves do exactly reproduce quantum entanglement. This generally concerns a bipartite state of two particles of the form $`|\psi ^\pm =|a_1b_2\pm |b_1a_2`$, where the subscripts denote the particles and $`a`$ and $`b`$, their measured properties. Entanglement lies in the inequality of the amplitudes $`x_1y_2|\psi ^\pm `$, which are probabilistic in quantum theory, but would be deterministic in the classical picture. Key to our result is the observation that the combined bras $`x_iy_j|`$ can be identified as the resulting combined stationary physical state of two detectors, to be explained in terms of computational principles, and that this state is really a spectral component of the overall physical state, the temporal factor $`e^{i\omega t}`$ being implicit in the notation. Also omitted is the fact that the “particles” are travelling; the corresponding factor $`e^{ikz}`$ is crucial to our classical wave interpretation.
For simplicity of argument, I shall restrict myself to a single Einstein-Podalsky-Rosen (EPR) experiment, in which circularly polarised $`\gamma `$ photons are emitted by a disintegrating source and the photons are detected by analysing their linear polarisations. This entails no loss of generality other than that already outlined, viz that the treatment is applicable only to photons, as the resulting arguments will be extensible even to the de Broglie waves of matter. We denote the detection events as $`x_1|`$, etc., $`x`$ identifying the linear polarisation and the subscript, the detector; and the initial photon states as $`|r_1`$, $`|l_2`$, etc., representing the right- and left-circular polarisations, respectively, so that the entangled photon state becomes $`|\psi ^\pm =|r_1r_2\pm |l_1l_2|r_1r_2\pm l_1l_2`$, according to the parity of the parent particle. We would describe the circular polarisation classically by an electric vector field propagating in the $`z`$-direction,
$$𝐄(z,t)=𝐄_x\mathrm{cos}(kz\omega t+\alpha )+𝐄_y\mathrm{sin}(kz\omega t+\alpha ),$$
(1)
$`𝐄_x`$ and $`𝐄_y`$ being the amplitude component in the corresponding directions, so that a detection at $`z_i`$, $`i\{1,2\}`$, means simply that the phase, and therefore the delay $`t_i`$, is determined at $`z_i`$,
$$\omega t_i=\alpha +kz_i\{\begin{array}{c}(n+1/2)\pi \hfill \\ n\pi \hfill \end{array},\alpha [0,\pi /4),$$
(2)
according to whether the detected polarisation was $`x`$ or $`y`$, respectively, i.e. *the measured bit equivalently determines whether the detector was within an even or odd half wavelength interval from the source.* With one detector, this is all we would have, but with two non-colocated detectors, we could associate the two detections with a common source if, and only if, the ranges were to coincide, or in other words, the two values represented the *same bit* of information. We can reinforce this conclusion by interpreting the quantum kets also classically as
$$\begin{array}{cc}\hfill |r\frac{1}{\sqrt{2}}[|x+|y]& \frac{1}{\sqrt{2}}[𝐄_xe^{ikz}+i𝐄_ye^{ikz}]\hfill \\ \hfill \text{and }|l\frac{1}{\sqrt{2}}[|x|y]\frac{1}{\sqrt{2}}|x& \frac{1}{\sqrt{2}}[𝐄_xe^{ikz}i𝐄_ye^{ikz}],\hfill \end{array}$$
(3)
where $``$ denotes the correspondence and $`e^{i\omega t}`$ is once again omitted for brevity. The analysed polarisations $`x|`$, $`y|`$ likewise correspond to the classical phasors $`𝐄_x^{}e^{ikz}`$ and $`𝐄_y^{}e^{ikz}`$, respectively, and lead to the classical dot products
$$\begin{array}{cc}\hfill x_1y_2|r_1r_2& _{z_1,z_2}𝐄_x^{}e^{ikz_1}𝐄_y^{}e^{ikz_2}\frac{1}{\sqrt{2}}[𝐄_xe^{ikz_1}+i𝐄_ye^{ikz_1}]\frac{1}{\sqrt{2}}[𝐄_xe^{ikz_2}+i𝐄_ye^{ikz_2}]𝑑z_1𝑑z_2\hfill \\ & =\frac{1}{\sqrt{2}}_{z_1}𝐄_x^{}e^{ikz_1}𝐄_xe^{ikz_1}𝑑z_1\frac{i}{\sqrt{2}}_{z_2}𝐄_y^{}e^{ikz_2}𝐄_ye^{ikz_2}𝑑z_2\hfill \\ & =iE^2/2,E|𝐄_𝐱|=|𝐄_𝐲|\hfill \\ \hfill x_1x_2|r_1r_2& _{z_1,z_2}𝐄_x^{}e^{ikz_1}𝐄_x^{}e^{ikz_2}\frac{1}{\sqrt{2}}[𝐄_xe^{ikz_1}+i𝐄_ye^{ikz_1}]\frac{1}{\sqrt{2}}[𝐄_xe^{ikz_2}+i𝐄_ye^{ikz_2}]𝑑z_1𝑑z_2\hfill \\ & =\frac{1}{\sqrt{2}}_{z_1}𝐄_x^{}e^{ikz_1}𝐄_xe^{ikz_1}𝑑z_1\frac{1}{\sqrt{2}}_{z_2}𝐄_x^{}e^{ikz_2}𝐄_xe^{ikz_2}𝑑z_2\hfill \\ & =E^2/2,\hfill \end{array}$$
(4)
since terms involving both $`z_1`$ and $`z_2`$ in the exponent vanish in the Césaro sum exactly as in quantum theory. Similarly,
$$\begin{array}{cc}\hfill x_1y_2|l_1l_2& _{z_1,z_2}𝐄_x^{}e^{ikz_1}𝐄_y^{}e^{ikz_2}\frac{1}{\sqrt{2}}[𝐄_xe^{ikz_1}i𝐄_ye^{ikz_1}]\frac{1}{\sqrt{2}}[𝐄_xe^{ikz_2}i𝐄_ye^{ikz_2}]𝑑z_1𝑑z_2\hfill \\ & =\frac{1}{\sqrt{2}}_{z_1}𝐄_x^{}e^{ikz_1}𝐄_xe^{ikz_1}𝑑z_1\frac{i}{\sqrt{2}}_{z_2}𝐄_y^{}e^{ikz_2}𝐄_ye^{ikz_2}𝑑z_2\hfill \\ & =iE^2/2,\hfill \\ \hfill x_1x_2|l_1l_2& _{z_1,z_2}𝐄_x^{}e^{ikz_1}𝐄_x^{}e^{ikz_2}\frac{1}{\sqrt{2}}[𝐄_xe^{ikz_1}i𝐄_ye^{ikz_1}]\frac{1}{\sqrt{2}}[𝐄_xe^{ikz_2}i𝐄_ye^{ikz_2}]𝑑z_1𝑑z_2\hfill \\ & =\frac{1}{\sqrt{2}}_{z_1}𝐄_x^{}e^{ikz_1}𝐄_xe^{ikz_1}𝑑z_1\frac{1}{\sqrt{2}}_{z_2}𝐄_x^{}e^{ikz_2}𝐄_xe^{ikz_2}𝑑z_2\hfill \\ & =E^2/2,\hfill \end{array}$$
(5)
yielding the same result as the quantum amplitudes (cf. (Feynman, , III-18-3)),
$$x_1y_2|\psi ^\pm \{\begin{array}{c}0\hfill \\ iE^2\hfill \end{array}\text{ and }x_1x_2|\psi ^\pm \{\begin{array}{c}E^2\hfill \\ 0\hfill \end{array},$$
(6)
the difference being only that the classical amplitudes on the right are exact and not probabilistic. The mystery of entanglement clearly cannot lie in in the quantumness of the waves, but only in their treatment as particles EPR1935 , which introduces the notion of interaction and communication in the first place. The particulate perception also misses our information argument altogether, viz that entanglement must signify that the detected events concern the same bit of information, representing a common source (eq. 2).
This form of information is well understood in holography, where the image is formed by coincidence of antinodes of multiple waves. In ordinary holography, a single wavelength is used for each colour in the image, and angular diversity, i.e. waves from multiple angles, is used to eliminate aliases in the image by spatial coincidence. One form of holographic radar I worked on Prasad1986 employed frequency diversity for the dealiasing: we would have exactly the same principle in our EPR experiment if our source produced photon pairs at multiple frequencies, as eq. (2) would become
$$\omega _jt_{ji}=\alpha +k_jz_i\{\begin{array}{c}(n+1/2)\pi \hfill \\ n\pi \hfill \end{array},j=1,2,\mathrm{}.$$
(7)
With frequency-selective detectors, we would have a pair of detected event bits for each frequency $`j`$, but the aliases belonging to any pair of frequencies $`\omega _1`$ and $`\omega _2`$ can only coincide within intervals corresponding to the larger of the two wavelengths {$`\lambda _1c/\omega _1`$, $`\lambda _2c/\omega _2`$}, thus diminishing the density of aliases, as well as refining the uncertainty interval to the smaller wavelength. Even with a limited number of query frequencies, it is thus possible to localise the source within a known neighbourhood. In the EPR case, we have only one bit of coincidence information, and thus an infinite number of equally significant aliases, which made the information hard to recognise.
We still need to reexamine the issue of superluminal communication, which implies the conjecture that the detected value at $`z_2`$ can be influenced by a measurement at $`z_1`$ *after* the emission of the photons from the source. A stronger result, that entanglement as such confers no ability whatsoever to communicate classical information, has in any case already been established recently from within the confines of quantum theory BSSTT1999 , but we need to be able to arrive at the same result within our imaging interpretation. One difficulty, which could have made our classical wave analysis unthinkable in the past, is that eq. (1) describes a pure, endless sinusoidal wave, not a wave packet with which we could associate a discrete detection event and time-of-flight issues. Even the quantum wavefunctions in the EPR scenario, however, do not denote wave packets at all, but merely individual spectral components, i.e. pure sinusoids, and moreover, we have already established that the multiple spectral components forming a wave packet would only yield a sharper image. Accordingly, we only need to verify that communication *per se* does not occur in the present interpretation either, and the reason may be recalled from past theory concerning relativity, viz that the information represented by the measurement at $`z_2`$ is available only *after* correlation with $`z_1`$, signifying *a posteriori* identification.
The EPR correlation is thus mysterious only in the particulate view, which must hence be wrong. Further, although I showed entanglement to be classical only in the case of photons, the imaging principle is identically applicable to any wave formalism, including quantum wavefunctions, meaning that in all such cases, entanglement simply concerns independent detections of a common source. Two other difficulties exist in this regard, first, that a similar classical wave analysis does not appear possible for particles in general, and second, that the quantum wavefunctions and amplitudes are probabilistic, which is again hard to reconcile with the classical mechanics of particles. These are the very ones that make the wave-particle duality conceptually difficult in the first place, as the wave aspect by itself is inherently common to both classical and quantum descriptions, as evidenced above.
Accordingly, I show in the remaining sections that these difficulties too vanish on applying other modern engineering principles from control, computation and thermodynamic theories, the principle in each case being based solely on precise classical reasoning. In §III, I shall first show that the quantum formalism of states is in fact simply the most general representation of information, and therefore the representation necessary and sufficient for analysing the most complex situations including particle physics and the computational principles of the brain. I shall then show, in §IV, how this formalism gets represented and linked to observed physics, in particular, how we arrive at a nonzero $`h`$ and its constancy, as well as at the probabilistic character and the uncertainty principle of quantum mechanics, also on basis of sound classical, but modern, engineering.
## III Representability: the anti-particulate principle
The contention, essentially, is that the quantum formalism of states is foremost computational, and then acquires its probabilistic character once again because of a fundamental principle of computation related to the thermodynamics of representation. The separation was unobvious in the past because the probabilistic nature suggests a closer relation to Shannon’s theory, obfuscating the representational aspect of the observer’s data states, and yet not revealing the mechanisms that cause it to be statistical in the first place. The distinction is subtle but fundamental, as the logical notion of representation is *per se* deterministic, despite the fact that any physical embodiment is bound to be subject to thermal erosion.
Representability, or definability, as such appears to be sufficient to imply inherent quantisation of fields, as recently shown by BenDaniel BenDaniel1999 , the equivalent computational argument being that the total knowledge of any finite set of observers can only be finite to the extent it is representable and communicable, using a finite number of sentences in a language with a finite alphabet, as any other knowledge would be inexpressible by definition and therefore outside the realm of science. Evidently, the premise could have been used avoid many of the measure-theoretic difficulties historically encountered in quantum theory, as suggested in Appendix A, where Cauchy’s notion of continuity is also shown to be literally equivalent to our consideration of finite representation. More particularly implied is that:
*Physical information fundamentally represents continuous fields and this it can do perfectly; conversely, particles can never be perfectly represented or known, nor directly represented.*
One may recall that in classical mechanics, a particle is considered to be a geometrical point requiring no spatial description beyond a triplet of real numbers defining its instantaneous location, with possibly a second triple specifying its velocity. The distinction is that the point coordinates do not suffice to *delimit* the object as a point particle; to delimit without contextual knowledge requires infinite bits of representation. For example, as already explained, with many bits of coincidence information, we can locate a source very closely, albeit still with infinite alias regions, so that contextual knowledge is needed to limit the localisation to one neighbourhood. With infinite information we would be able to locate the source *exactly* and with no contextual knowledge at all, but the exercise assumes that our source is a point entity. No such difficulty occurs if we accept our representative data as inherently referring to a continuous field and not a point particle. This principle formalises our contention in §II that the particulate view is as such erroneous and partly responsible for the semantic difficulties of quantum theory, as will be illustrated again with respect to the notion of intrinsic spin.
To demonstrate the computational origin of the state space formalism, we now attempt to construct the formalism purely by seeking the most general representation for knowledge of any kind. We could pick arbitrary variables, for instance, a combination of groceries, train schedules and stock prices. Despite the generality, we would still have two conditions available in any such choice, viz numerical values and, if pertaining to real phenomena, observability at various times. In the resulting multidimensional state space, the observed data would be essentially represented in the direction of state vectors, in absence of a meaningful norm for all such spaces. Differences between states could be trivially treated as vectors, and the temporal evolution of a state would be represented only as a continuing change of direction. As shown in Appendix A, we do not need to be able to measure the represented variables with infinite precision at infinitely small intervals; rather, it is our inability in both regards that limits us to finite data of finite precision, and therefore to continuous fields, representing the impossibility of point delimitation, as just explained.
The quantum equations of motion are then obtained as general computational result, by considering the most general form of temporal evolution over such a space. By the representabilty principle, this too must be limited to a finite number of coefficients of finite precision, the latter being implicit, yielding
$$\underset{i=0}{\overset{n}{}}a_n\frac{d^n}{dt^n}|\psi =|\varphi =F^{(n)}(t)|\psi ,$$
(8)
where the kets merely denote our arbitarily chosen state space, and $`F^{(n)}(t)`$ signifies a continuously varying applied “force” on the system, assuming our variables are sufficiently well behaved. The generality is the reason we may apply techniques from Laplace transform theory to solve eq. (8); in particular, the operator sum $`_{i=0}^na_nd^n/dt^n`$ transforms to $`_{i=0}^na_ns^n`$, and yields the characteristic equation $`_{i=0}^na_ns^n=0`$, so that by the fundamental theorem of algebra, the complex plane becomes necessary and sufficient to represent all possible patterns of evolution. We have thus reproduced the broad notion of state space as well as the need for complex valued representation, but without *a priori* assumption of quantisation, so that these ideas are valid for any state space, not just one of quantum theory.
The general state space notion of evolution automatically leads to Schrödinger’s equation, as shown by Dirac (Dirac, , §27), because the derivative of the incremental evolution $`|\psi (t)|\psi (t+\mathrm{\Delta }t)`$ must have the operator form
$$i\mathrm{}\frac{d}{dt}|\psi =H(t)|\psi ,$$
(9)
where the imaginary coefficient $`i`$ derives from the fact that vector magnitudes lose significance in state representation, and $`\mathrm{}`$ is merely a scale factor relating the rate of evolution operator $`H`$ to the observer’s scale of time. This also leads to spectral decomposition in quantum theory (Dirac, , §29), but it should be realised that in general, $`e^{i\omega t}`$ is simply the imaginary component of (the eigenfunctions of) the Laplace operator $`sd/dt`$. It is only by restricting to stationary states that the real parts, representing transients $`e^{\alpha t}`$, vanish from quantum theory, and the reasons both for stationarity and the constancy of $`\mathrm{}`$ in eq. (9) are also basically computational, as will be explained in §IV.
The representational freedom explains why we were able to drop $`e^{ikz}`$ in the quantum versions of eq. (3) and yet arrive at the entanglement prediction, eq. (6); we have paid for our semantic imprecision, of course, by confusing the result with superluminal communication. The quantum notion of spin, as a mysteriously intrinsic property of particles, results from a similar imprecision: if in light of the foregoing theory, we recognise a particle wavefunction $`|\psi \left(\begin{array}{c}\psi _x\hfill \\ \psi _y\hfill \end{array}\right)`$ as a continuous physical field, and remember that we are referring to a travelling wave component, so that both spatial and temporal factors $`e^{ikz}`$ and $`e^{i\omega t}`$, respectively, are required for its complete physical representation, it becomes clear that the effect of a Pauli spin matrix, say $`\left(\begin{array}{cc}0\hfill & i\hfill \\ i\hfill & 0\hfill \end{array}\right)`$, on $`|\psi `$ is to simply rotate the local direction of motion as one travels within the field, literally describing a vortex field rather than an inherently non-geometrical form of angular momentum. It is only in the particulate view, where the complex exponential factors get dropped, that we end up with representational, and consequently semantic, loss. On a related note, in ignoring the magnitudes, we had also lost the $`E^2`$ in eq. (6), and this had cost us insight from classical electromagnetism.
## IV Thermodynamic stationarity
I still need to show how $`\mathrm{}`$ acquires the special significance it does in quantum theory, when applied to mechanical information, and account for the probabilistic nature of quantum wavefunctions, else the foregoing theory would be considered no more than a coincidence.
The first part has been recently achieved Prasad2000a by going back to Planck’s theory, and proving that thermal wall jitter, which is in fact necessary to account for the thermalisation of radiation in the first place as there are no nonlinear interactions *between* photons that might suffice for this purpose, unlike the case of gas molecules in kinetic theory; is in fact sufficient to exactly reproduce both Planck’s postulates and the radiation spectrum law. The wall jitter continually changes the stationary modes of the cavity, not just the energy distribution therein; we need to consider only the fraction of thermal changes that are slower than the electromagnetic transit time across the cavity, because faster motions can only affect transient behaviour. We therefore sum over the equilibrial distribution of modes, and in the process identify families of modes that are related by whole numbers of antinodes. It is trivial to verify that within such a family, the frequencies would be necessarily related as $`\{f`$, $`2f`$, $`3f`$, …$`\}`$, and importantly, incremental wall displacements can only change the modal distribution by an integral number of antinodes, i.e. between members of the same family, so that, discounting transients, the harmonic families behave exactly as Planckian oscillators, as each family can have exactly one member energised at any instant. The energy of an antinodal lobe depends only on the amplitude and not the frequency, and in the equilibrial state, the antinodal lobes must have the same energy, reproducing Planck’s quantisation rule
$$E=hf;$$
(10)
the sums over the families then trivially reproduces the equations in Planck’s derivation. The detailed balance, needed for balancing the various contributions, from the thermal Doppler spread due to wall jitter, possible interactions with the wall atoms, nonlinearities due to imperfect vacuum within the cavity, and so on, turns out to be identical in form to the Bose-Einstein derivation. Planck’s law
$$\widehat{U}(f)=\frac{hf}{e^{hf/k_BT}1},$$
(11)
$`\widehat{U}(f)`$ denoting the equilibrial expection energy at frequency $`f`$, thus turns out to be a strictly classical result, stemming from classical electromagnetics and classical thermalisation.
More importantly, the very form of eq. (11) shows that $`h`$ relates to the antinodal lobes and the spectral domain in almost the same way as the Boltzmann constant $`k_B`$ does to gas molecules and the spatial domain, as it “exposes” $`h`$ to measurement the same way as the law of Brownian motion $`R^2=6k_BTt/\mu `$, $`R`$ being the mean radial distance travelled as a function of time $`t`$, and $`\mu `$, the mobility; historically exposed $`k_B`$, allowing it to be empirically determined from diffusion phenomena (Feynman, , I-41-4). It should also be noticed that an antinodal lobe is indeed the $`\lambda /2`$ interval of uncertainty we encountered when interpreting entanglement (§II), so that our holographic notion of information would be consistent with a stationary representation in thermodynamic equilibrium, to be explored via Hamilton-Jacobi formalism below.
The constancy of $`h`$ was also independently established using Dirac’s result that given any anti-commuting relation $`[.,.]`$ and two pairs of conjugate variables $`u_1`$, $`v_1`$ and $`u_2`$, $`v_2`$, we would get, with no assumption of dependence between the pairs, the cross-relation
$$[u_1,v_1](u_2v_2v_2u_2)=(u_1v_1v_1u_1)[u_2,v_2]$$
(12)
implying that they must have the same constant of anticommutation (Dirac, , §21),
$$\begin{array}{cc}\hfill u_1v_1v_1u_1& =K[u_1,v_1]\hfill \\ \hfill \text{and }u_2v_2v_2u_2& =K[u_2,v_2].\hfill \end{array}$$
(13)
We could, for instance, choose classical electric and magnetic wave amplitudes $`𝐄`$ and $`𝐁`$ within our cavity as $`u_1`$ and $`v_1`$, and the induced emf $`V`$ and current $`i`$ in a probe for $`u_2`$ and $`v_2`$. We should then get the same value for $`K`$ for the anticommutation of $`V`$ and $`i`$ as for $`E`$ and $`B`$. We could next reuse the symbols $`u_1`$ and $`v_1`$ for $`V`$ and $`i`$, respectively, and identify $`u_2`$ and $`v_2`$ with the dynamical variables of a third system, establishing the validity of $`K`$ for the third system, and so on, proving that every physical system that could directly or indirectly interact with our cavity would be similarly quantised with the same value of $`h`$, proving it to be a universal constant by transitivity of interaction.
For a general notion of waves and their stationarity in the context of matter, we must, as previously mentioned, turn to Hamilton-Jacobi theory, beginning with the following relation between Hamilton’s principal and characteristic functions, $`S`$ and $`W`$ respectively, for a single particle system in Cartesian coordinates (Goldstein, , §10-8),
$$S(q,P,t)=W(q,P)Et,$$
(14)
$`q`$ being the particle’s generalised coordinate, $`E`$, its total energy, and $`P`$, its momentum, which is a constant of the motion in the configuration space for which $`S`$ is a solution of the Hamilton-Jacobi equation
$$H(q_1,\mathrm{},q_n;\frac{F_2}{q_1},\mathrm{},\frac{F_2}{q_n};t)+\frac{F_2}{t}=0.$$
(15)
$`W`$ closely relates to the action-angle variables used in the early days of quantum mechanics, the angle variable $`w`$ being defined as $`W/J`$, $`J`$ being the action variable: one would first solve the classical problem using these variables, and quantisation was achieved by replacing $`J`$ with a multiple of $`h`$. The approach dropped out of vogue for two reasons, first, the classical model was not always solvable and the state space formulation proved simpler and more general, for reasons we have already established, and second, the cause of quantisation was then not known, so that one could not aspire for deeper understanding by sticking to the classical formalism. The second reason is now no longer valid, as the cause of quantisation has been uncovered in the context of radiation equilibrium. It is meaningful, therefore, to reexamine the classical approach in this light, particularly to seek a general classical formalism of travelling waves to which we could apply our antinodal holographic ideas.
Eq. (14) describes constant-$`S`$ wavefronts travelling in the same direction as the particle: the particle’s velocity is given by $`u=E/p=E/mv`$, and we also have $`𝐩=W`$, but the “waves” are not given to be periodic. We need a spectral decomposition, and since, as Goldstein points out, $`S=WEt`$ must be proportional to the phase $`\omega t`$ and $`W`$ is independent of time, $`E`$ must be proportional to $`\omega `$. Since this is similar to Planck’s quantisation rule (eq. 10), it was traditionally assumed that classical mechanics stopped just short of quantum theory, as it gave no reason for assuming the constant of proportionality to be nonzero. Indeed, a plane wave trial solution of the form $`\psi =\psi _0e^{iS/\mathrm{}}`$, corresponding to $`E=\mathrm{}\omega `$, would turn Schrödinger’s equation (9) to
$$\left[\frac{1}{2m}(S)^2+V\right]+\frac{S}{t}=\frac{i\mathrm{}}{2m}^2S,$$
(16)
via the following well-known form of eq. (9)
$$\frac{\mathrm{}^2}{2m}^2\psi V\psi =\frac{\mathrm{}}{i}\frac{\psi }{t}$$
(17)
and the derivatives
$$\frac{\psi }{t}=\frac{i}{\mathrm{}}\psi \frac{S}{t}\text{ and }\frac{\psi }{x}=\frac{i}{\mathrm{}}\psi \frac{S}{x}$$
(18)
which yield
$$^2\psi =\frac{i}{\mathrm{}}\psi ^2S\frac{\psi }{\mathrm{}^2}(S)^2$$
(19)
for the Laplacian of $`\psi `$. Eq. (16) would reduce to the Hamilton-Jacobi equation (15) if $`\mathrm{}`$ were to vanish, and this was the basis for Bohr’s Correspondence Principle (BCP), that quantum mechanics reduces to classical theory in the short-wavelength limit, as $`\mathrm{}^2S(S)^2`$ is equivalent to, for a one-particle system, $`(\lambda /p)(dp/dx)2\pi `$.
Observe, however, that for the validity of eq. (14), we could not possibly set $`E=0\omega `$, i.e. the BCP is not classically meaningful at all. The error lies in mistaking the trial solution $`\psi _0e^{iS/\mathrm{}}`$ to be a complete description, instead of as a mere spectral component of the overall dynamics. In the spectral decomposition, we would look for stationary modes for which the $`Et`$ term in eq. (14) vanishes anyway, and we would not need $`\mathrm{}`$ to vanish, because the right hand term in eq. (16) then merely refers to a Fourier component, not the whole picture.
Importantly, *stationarity is a necessary condition for both the physical data states of the observer and of the observed entities, as both the observer’s knowledge and the represented information cannot change between observations.* This is an essentially computational principle, but it forces us to consider only the stationary states of material systems as being physically relevant. The stationarity further means that the Hamiltonian becomes completely separable (cf. (Goldstein, , §10-6)), so that the generalised coordinates $`q_i`$ and momenta $`p_i`$ are related as in-phase and quadrature components, respectively, at each of the characteristic frequencies. In particular, the combined stationary states of radiation and matter must be stationary with respect to either, and we can then apply Dirac’s transitivity argument (eq. 13) to each such pair {$`q_i,p_i`$} and the {$`𝐄_i`$, $`𝐁_i`$} radiation amplitudes at the same frequencies $`f_i`$, to infer that the antinodal equipartition must be identically applicable to these stationary modes of matter.
The stationarity and equilibrium principles fully explain both quantisation and the probabilistic aspects of quantum theory. As example, recall that radiation quantisation was conceived to in order to explain the almost instantaneous nature of photoelectric emission ResnickHalliday , but by including the observer’s data states in the picture, we can readily see that the abruptness is logically unavoidable, as the observer cannot know of states intermediate to those representable within its own material embodiment. The condition of thermal equilibrium, necessary for potential long term stability of the observer’s data states, guarantees, via the spectral decomposition of the Hamilton-Jacobi wave analysis of the observer’s embodiment described above, that the data states can only again change in antinodal increments of the same energy as those of radiation. We may therefore represent the observed variables by kets $`|\psi `$, the detector values, representing the observer’s state space, by bras of the form $`\varphi |`$, and get the same amplitudes $`\varphi |\psi `$ as when $`\varphi |`$ is set to represent a second observed entity instead, interacting with the first. In either case, the amplitude must be complex because of the algebraic completeness of the state space formulation. The amplitude is consistent with the (square root of the) probability of the observed system *and* the observer *thermally* and *irreversibly* transiting to the combined state $`\varphi ,\psi `$, the irreversibility being implied by the condition that the $`\varphi |`$ must be itself equilibrial and capable of lasting till erasure or the next observation. As particularly described by Landauer Landauer1961 , the irreversibility means that the final state is attained regardless of the previous states of the observer and the observed, reproducing the apparent “collapse” of the quantum wavefunction in the process of observation. Lastly, the fact that the transitions can only occur in increments of an antinodal lobe reproduces Heisenberg’s uncertainty principle.
## V Summary and loose ends
To summarise, I have established, though not in this order,
1. that the quantumstate space formalism is no more than the most general representation of information of any kind whatsoever, which would account for its discovery in the first place, and its applicability in diverse fields ranging from particle physics to studies of the brain (§III);
2. that quantum wavefunctions merely constitute Fourier in-phase and quadrature components, as already familiar to electrical engineers working with strictly classical electromagnetism (§III, §IV), and that their quantisation and randomness are due to separate computational and thermodynamic causes (vii, viii, x);
3. that the de Broglie waves of matter are obtainable from Hamilton-Jacobi theory, and that the Correspondence Principle was erroneously conceived in the context because the computational requirements of stationarity and thermal equilibrium (viii) were not known;
4. that complex values occur in quantum mechanics because of the generality of the information represented, as a result of the fundamental theorem of algebra, which incidentally proves the computational completeness of the quantum formalism that had not been established in prior theory, supporting (i), and relates Schrödinger’s equation to control theory (§III);
5. that quantum entanglement, as best illustrated by the EPR paradox, amounts simply to 1-bit holography in the classical wave picture (§II), and that issues of hidden variables and action-at-a-distance are merely consequences of preconceived particulate view;
6. that the particulate view in fact contradicts formal considerations of the representability of physical information, which show that only continuous fields can at all be physically represented, and thus observed or communicated, without contextual bias, validating the present notions of entanglement (v), Hamilton-Jacobi waves (iii) and the classical field interpretation of spin (§III);
7. that both quantisation and the probabilistic nature of quantum wavefunctions are completely explained by the same principles of stationarity and antinodal equipartition (§IV);
8. that both stationarity and thermal equilibrium are mandated by the computational consideration of stability of the observer’s data states and the represented knowledge (§III), connecting the representational generality (i, iv, vi) with the actual physics (ii, iii, v, vii);
9. that the Hamilton-Jacobi analysis (§IV) establishes not only the sufficiency of stationarity and thermalisation for (vii), but also their necessary involvement in this role, obviating any postulate or alternative explanation; and
10. that the constancy of $`h`$ is purely a consequence of the transitivity of these constraints between observed systems as well as between physical observers, signifying the establishment of a universal scale, given that $`h`$ relates the (independent) dimensions of energy (mass) and time, by communication, again a matter of logic and information.
To illustrate the attribution of quantum probabilities to thermalisation, consider the hypothetical case of an observer frozen to absolute zero temperature. By the third law of thermodynamics (Nernst theorem), such an observer would be frozen into one state of knowledge and would be incapable of making any observations whatsoever. An interesting example of the obfuscation by notions of probability and information theory in the past is the following observation that the germs of BenDaniel’s result were already present in kinetic theory: Boltzmann’s notion of thermodynamic information as the spatial localisation of a gas molecule closely parallels our notion of particle delimitation, and the randomness of molecular motion had nothing directly to do with this measure, as with each bit of information serves to (deterministically) halve the region of uncertainty. The information and randomness aspects were thus separable in Boltzmann’s theory, but were not recognised as such because the determinism properly belongs to the computational issue of representation, and the present distinction of computational and informational aspects does not seem to have been at all made in prior literature.
I have not attempted to discuss how the present theory should be applied to other quantum issues like superconductivity, exchange coupling, *zitterbewegung*, particle creation and annihilation, broken symmetries, or the grand unified theories (GUTs), etc. These are simply too numerous to be recast or corrected by one treatise or author. The fundamental principles of quantum mechanics, which are basis for all such applications, have already been shown to be fundamentally computational and thermodynamic. The fact that the equilibrial antinodal lobe, representing the quantum of change in both radiation and matter waves as described, is independent of the wavelength $`\lambda c/f`$ means that the establishment of thermal equilibrium and quantum mechanical consistency do not depend on *a priori* equality of spatial scale between the interacting entities. The hypothesis of relative spatial scale suffices for deducing both the postulates of relativity and Maxwell’s equations of electromagnetism, and more importantly, shows that the current ideas in general relativity and cosmology are entirely too simplistic, as separately described in Prasad2000c . There again, the problem in prior theory has never been a shortage of mathematical skills, but of the computational intuition of the representability of physical information, and the cognition of fundamental limitations arising from this constraint.
###### Acknowledgements.
To my colleagues A Joseph Hoane and Daniel Oblinger for valuable discussions involving my EPR solution.
## Appendix A Finite domain calculus
The premise of representational finiteness allowed me, in 1983-84, to define a notional framework as follows:
1. Every function $`f`$ is represented by at most a finite number of symbols denoting algebraic variables and operations. The continuity of a domain $`X`$ is likewise defined as the condition that at any finite cardinality $`\mathrm{\#}(X)`$, additional points may be *physically* introduced, *by resizing, adjustment of magnification, or technological replacement, etc.* to indefinitely define new points in the neighbourhood $`N(x)`$ of every point $`xX`$.
2. Continuity of a function $`f:XY`$ is then defined by the conditions that $`X`$ and $`Y`$ are both continuous, every new point $`x^{}N_ϵ(x)`$ introduces a computed value $`y^{}=f(x^{})N_\delta (y)`$, as in Cauchy’s definition.
The point is that Cauchy’s definition involves only finite domains at any finite stage in the implied execution of the limit operations, and even the issues of divergence in quantum field theory, for instance, are really concerned with how computed behaviour varies with the precision of the referenced domain, which is implicitly finite in all of mathematics. Like the intuition of continuity given by BenDaniel’s result (§III), this is too simple to be obvious.
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# Abundance evolution of intermediate mass elements (C to Zn) in the Milky Way halo and disk
## 1 Introduction
In the past ten years or so, progress in our understanding of the chemical evolution of the Milky Way came mainly from observations concerning the composition of stars in the halo and the local disk. The seminal works of Edvardsson et al. (1993) for the disk, and Ryan et al. (1996) and McWilliam et al (1995) for the halo (along with many others) provided detailed abundance patterns that reveal, in principle, the chemical history of our Galaxy.
The interpretation of these data is not straightforward, however, since it has to be made in the framework of some appropriate model of galactic chemical evolution (GCE). Only one of the three main ingredients of GCE models can be calculated from first principles at present: the stellar yields. For the other two ingredients, i.e. the stellar initial mass function (IMF) and the star formation rate (SFR), one has to rely on empirical prescriptions.
Considerable progress in GCE studies was made possible after the publication of the yields from massive stars of Woosley and Weaver (1995, hereafter WW1995). This work made available, for the first time, yields for an extensive set of isotopes (from H to Zn), stellar masses (from 11 to 40 M) and metallicities (from Z=0 to Z=Z), making thus possible a detailed comparison of theory to observations. Only two works until now explored fully the potential of the WW1995 yields. Timmes et al. (1995) adopted a simple GCE model with infall, appropriate for the Milky Way disk but certainly not for the halo (see Sec. 3.3); in the framework of that model they made a case-by-case assessment of the strengths and weaknesses of the WW1995 yields, identifying the large yields of Fe as the main weak point. On the other hand, Samland (1998) utilised a chemo-dynamical model for the Milky Way evolution (describing, presumably, correctly the halo and the disk), but introduced several approximations on the stellar lifetimes and the metallicity dependant yields of WW1995; he evaluated then the deviation of the published yields from the “true” galactic ones, the latter being derived by a comparison of his model results with observations of the halo and disk abundance patterns.
Those two works are the only ones that utilised metallicity dependant yields and studied the full range of intermediate mass chemical elements. Several other works focused on specific elements and utilised only metallicity independant yields (e.g. Pagel and Tautvaisiene 1995; Chiappini et al. 1997, 1999; Thomas et al. 1998 etc.)
In this work we reassess the chemical evolution of the elements from C to Zn in the Milky Way, using the WW1995 yields. Our work differs in several aspects from the one of Timmes et al. (1995) and, in fact, from any other work on that topic, performed in the framework of simple GCE models: the main novelty is that we use appropriate models for both the halo and the disk, correctly reproducing the main observational constraints for those two galactic subsystems (see Sec. 4). Moreover, we adopt the Kroupa et al. (1993) IMF, which presumably describes the distribution of stellar mases better than the Salpeter IMF (adopted in Timmes et al 1995, Samland 1998, and most other studies of that kind). Also, w.r.t. the work of Timmes et al. (1995), our comparison to observations benefits from the wealth of abundance data made available after the surveys of Ryan et al. (1996), McWilliam et al (1995), Chen et al (2000) and many others (listed in Table 1). These data allow to put even stronger constraints on the stellar yields as a function of metallicity. We notice that we do not include yields from intermediate mass stars in our study, since we want to see to what extent those stars (or other sources) are required to account for the observations.
The plan of the paper is as follows: In Sec. 2 we discuss briefly the uncertainties currently affecting the yields of massive stars and present the yields of WW1995. We also present those of a recent work (Limongi et al. 2000), which compare fairly well to those of WW1995 but show interesting differences for several elements. Moreover, we present the recent yields of Iwamoto et al. (1999) for supernovae Ia, calculated for white dwarfs resulting from stars of solar and zero initial metallicities, respectively; they are slightly different from the “classical” W7 model for SNIa (Thielemann et al. 1986), and we adopt them in our study. In Sec. 3 we present our chemical evolution model, stressing the importance of adopting appropriate ingredients for the halo and the disk. In Sec. 4 we “validate” our model by comparing successfully its results to the main observational constraints. We also show that current massive star yields have difficulties in explaining the solar composition of Sc, Ti and V. In Sec. 5 we present the main result of this work, i.e. a detailed comparison of the model to observations of abundance patterns in halo and disk stars. This comparison allows to identify clearly the successes and inadequacies of the WW1995 yields; some of those inadequacies may be due to physical ingredients not as yet incorporated in “standard” stellar models (i.e. mass loss or rotationally induced mixing), but the origins of others are more difficult to identify. Since the evolution of Fe (usually adopted as “cosmic clock”) is subject to various theoretical uncertainties - Fe yields of massive stars, rate of Fe producing supernovae Ia etc - we also plot our results as a function of Ca; comparison to available observations (never performed before) gives then a fresh and instructive view of the metallicity dependence of the massive star yields. In Sec. 6, we discuss qualitatively some possible ways to interpret the recent data of Israelian et al (1998) and Boesgaard et al (1999) on oxygen vs iron; these data suggest that oxygen behaves differently from the other alpha-elements and, if confirmed, will require some important revision of current ideas on stellar nucleosynthesis. Finally, in Sec. 7 we compare the model evolution of the Mg isotopic ratios to recent observations of disk and halo stars; we find that the WW1995 yields underestimate the production of the neutron-rich Mg isotopes at low metallicities.
## 2 Yields of massive stars and supernova Ia
Massive stars are the main producers of most of the heavy isotopes in the Universe (i.e. those with mass number A$`>`$11). Elements up to Ca are mostly produced in such stars by hydrostatic burning, whereas Fe peak elements are produced by the final supernova explosion (SNII), as well as by white dwarfs exploding in binary systems as SNIa. Most of He, C, N and minor CO isotopes, as well as s-nuclei comes from intermediate mass stars (2-8 M). A detailed discussion of the yields of massive stars and their role in galactic chemical evolution studies has been presented in a recent review (Prantzos 2000); here we summarize the most important points.
Extensive calculations performed in the 90ies by a few groups with 1-D stellar codes (Woosley and Weaver 1995, Arnett 1996, Thielemann et al. 1996, Chieffi et al. 1998, Maeder 1992, Woosley et al. 1993, Aubert et al. 1996, Limongi et al. 2000) have revealed several interesting features of nucleosynthesis in massive stars. In particular, the structure and composition of the pre-supernova star reflects the combined effect of (i) the various mixing mechanisms (convection, semi-convection, rotational mixing etc.), determining the extent of the various “onion-skin” layers, (ii) the amount of mass-loss (affecting mostly the yields of the He and CNO nuclei, present in the outer layers) and (iii) the rates of the relevant nuclear ractions (determining the abundances of the various species in each layer).
On the other hand, the calculation of the Fe-core collapse supernova explosion is still one of the major challenges in stellar astrophysics. Multi-dimensional hydrodynamical simulations in the 90ies revealed the crucial role played by neutrino transport in the outcome of the explosion (e.g. Janka 1998 and references therein). In the absence of a well-defined explosion scheme, modelers of supernovae nucleosynthesis have to initiate the explosion somehow (by introducing either an “internal energy bomb”, or a “piston”, e.g. Aufderheide et al. 1991) and adjust the shock energy as to have a pre-determined final kinetic energy, usually the “classical” value of 10<sup>51</sup> ergs (after accounting for the binding energy of the ejected matter). This procedure introduces one more degree of uncertainty in the final yields. Moreover, the ejected amount of Fe-peak nuclei depends largely on the position of the mass-cut, the surface separating the material falling back onto the neutronized core from the ejected envelope. The position of this surface depends on the details of the explosion (i.e. the delay between the bounce and the neutrino-assisted explosion, during which the proto-neutron star accretes material) and cannot be evaluated currently with precision (e.g. Thielemann et al. 1999 and references therein).
In the light of the aforementioned results, intermediate mass elements produced in massive stars may be divided in three major groups:
(i) C, N, O, Ne, and Mg are mainly produced in hydrostatic burning phases. They are mostly found in layers which are not heavily processed by explosive nucleosynthesis. The yields of these elements depend on the pre-supernova model (convection criterion, mixing processes, mass loss and nuclear reaction rates).
(ii) Al, Si, S, Ar and Ca are also produced by hydrostatic burning, but their abundances are subsequently affected by the passage of the shock wave. Their yields depend on both the pre-supernova model and the shock wave energy.
(iii) Fe-peak elements as well as some isotopes of lighter elements like Ca, S and Ti are produced by the final SN explosion (SN II). Their yields depend crucially upon the explosion mechanism and the position of the ”mass-cut”.
The outcome of nucleosynthesis depends also on the initial metallicity of the star. During H-burning the initial CNO transforms into <sup>14</sup>N, which transforms mostly into <sup>22</sup>Ne during He-burning, through $`\alpha `$-captures and a $`\beta `$ decay. The surplus of neutrons in <sup>22</sup>Ne (10 protons and 12 neutrons) affects the products of subsequent burning stages, in particular those of explosive burning. This neutron surplus increases with initial metallicity and favours the production of odd nuclei (<sup>23</sup>Na, <sup>27</sup>Al, <sup>31</sup>P etc.), giving rise to the so-called ”odd-even” effect.
In the past few years, several groups have reported results of pre and post-explosive nucleosynthesis calculations in massive stars with detailed networks. Thielemann et al. (1996) used bare He cores of initial metallicity $`Z_{}`$, while Arnett (1996) simulated the evolution of He cores (with polytropic-like trajectories) and studied different initial metallicities. Full stellar models (neglecting however, rotation and mass loss ) were studied by Woosley and Weaver (1995, for masses 12, 13, 15, 18, 20, 22, 25, 30, and 40 $`M_{}`$ and metallicities Z=0, 10<sup>-4</sup>, 10<sup>-2</sup>, 10<sup>-1</sup>, and 1 Z) and Limongi, Straniero and Chieffi (2000, for masses 13, 15, 20, 25 $`M_{}`$ and metallicities Z=0, 5 10<sup>-2</sup> and 1 Z). Comparison of the various yields on a star by star basis shows that there are large discrepancies between the different authors (due to differences in the adopted physics) although for some elements, like oxygen, there is a rather good agreement. Moreover, the yields do not show a monotonic behaviour with stellar mass.
Notice that the overall yield used in chemical evolution studies depends on both the individual stellar yields and the stellar IMF. Despite a vast amount of theoretical and observational work, the exact shape of the IMF is not well known yet (Gilmore et al. 1998 and references therein). It is however clear that the IMF flattens in the low mass range and cannot be represented by a power law of a single slope (e.g Kroupa et al. 1993). The shape of the IMF introduces a further uncertainty of a factor $``$ 2 as to the absolute yield value of each isotope (Wang and Silk 1993).
In Fig. 1 we present the metallicity dependant yields of Woosley and Weaver 1995 (hereafter WW1995) and Limongi, Straniero and Chieffi 2000 (hereafter LSC2000), folded with a Kroupa et al. (1993) IMF. They are presented as overproduction factors, i.e. the yields (ejected mass of a given element) are divided by the mass of that element initially present in the part of the star that is finally ejected, i.e.
$$<F>=\frac{_{M1}^{M2}Y_i(M)\mathrm{\Phi }(M)𝑑M}{_{M1}^{M2}X_{,i}(MM_R)\mathrm{\Phi }(M)𝑑M}$$
(1)
where: $`\mathrm{\Phi }(M)`$ is the IMF, $`M1`$ and $`M2`$ the lower and upper mass limits of the stellar models (12 M and 40 M for WW1995, 13 M and 25 Mfor LSC2000, respectively), $`Y_i(M)`$ are the individual stellar yields and $`M_R`$ the mass of the stellar remnant. Adopting $`X_{,i}`$ in Eq. (1) creates a slight inconsistency with the definition of the overpoduction factor given above, but it allows to visualize the effects of metallicity in the yields of secondary and odd elements.
From Fig. 1 it can be seen that i) most of the intermediate mass elements are nicely co-produced (within a factor of 2) in both calculations of solar metallicity stars; ii) some important discrepancies (e.g. Li, B, F) can be readily understood in terms of neutrino-induced nucleosynthesis, included in the WW1995 but not in the LSC2000 calculation; iii) the odd-even effect is clearly present in both calculations, but seems to be more important in LSC2000. For solar metallicity stars most of the even Z elements are produced with similar yields in both calculations, while odd Z elements in LSC2000 are produced with systematically lower yields than in WW1995. A common feature of both calculations is the relative underproduction of C, N, Sc, V and Ti w.r.t. O. C and N clearly require another source (intermediate mass stars and/or Wolf-Rayet stars, see Prantzos et al. 1994 and Sec. 4.2). The situation is less clear for the other elements, Sc, V and Ti.
In this work we adopt the metallicity dependant yields of WW1995, keeping in mind that the use of LSC2000 yields may lead to different results for some odd elements. For illustration purposes we shall also use the WW1995 yields at constant (=solar) metallicity. There are interesting differences between the two cases (i.e. constant vs. variable metallicity yields) and this instructive comparison has never been done before. We notice that in the case of the most massive stars (M$`>`$30 M) WW1995 performed 3 calculations, making different assumptions about the kinetic energy of the supernova ejecta. We adopt here their set of models A, in which, following the explosion, most of the heavy elements in the inner core fall back to form a black hole of a few solar masses; because of the form of the IMF, these very massive stars play a negligible role in shaping the elemental abundance ratios. As stressed in the Introduction, we consider no yields from intermediate mass stars in this work; our explicit purpose is to check to what extent massive stars can account for observations of intermediate mass elements and for which elements the contribution of intermediate mass stars is mandatory.
There is a strong observational argument, suggesting that massive stars are not the sole producers of Fe peak nuclei in the solar neighbourhood : the observed decline in the \[O/Fe\] ratio (Fig. 3, lower panel) from its $``$3 times the solar value in the halo stars (\[O/Fe\] $``$0.5 for \[Fe/H\]$`<`$-1) down to solar in disk stars. This decline is usually interpreted as due to injection of Fe and Fe group elements by SN Ia. Assuming that massive stars are the only source of O and Fe in the halo phase and they produce a ratio of Fe/O$``$1/3 solar, the remaining $``$2/3 of Fe in the late disk should be produced by a late source, presumably SNIa.
The WW1995 yields lead to approximately solar abundance ratios of O/Fe (or $`\alpha `$-element/Fe). This lead Timmes et al. (1995) to suggest that the Fe yields of WW1995 are probably overestimated. Following their suggestion, we adopt here half the nominal values for the WW1995 yields of Fe-peak elements (from Cr to Zn). Taking into account the uncertainties currently affecting those yields, such a reduction is not unreasonable. Our procedure allows to reproduce the observed O/Fe, but does not alter the abundance ratios between Fe-peak elements.
To account for the additional source of Fe-peak elements we utilise the recent yields of SNIa from the exploding white dwarf models of Iwamoto et al. (1999). These are updated versions of the original W7 model of Thielemann et al. (1986). In this model, the deflagration is starting in the centre of an accreting Chandrashekhar-mass CO white dwarf, burns $``$ half of the stellar material in Nuclear Statistical Equilibrium and produces $``$ 0.7 $`M_{}`$ of <sup>56</sup>Fe ( in the form of <sup>56</sup>Ni). It also produces all other Fe-peak isotopes and in particular <sup>58</sup>Ni and <sup>54</sup>Cr. This can be seen in Fig. 2, where the overproduction factors (normalised to the one of <sup>56</sup>Fe) of the SNIa yields are plotted for two models: one calculated for a white dwarf resulting from a star with solar initial metallicity (W7) and another for a white dwarf resulting from a star of zero initial metallicity (W70). The main difference between the two model results lies in the large underproduction of odd-isotopes in the latter case. In our calculation, we use the yields of those two models, linearly interpolated as a function of metallicity.
The problem with SNIa is that, although the current rate of SNIa/SNII is constrained by observations in external spiral galaxies (Tammann et al. 1994), the past history of that rate (depending on the nature of progenitor systems) is virtually unknown. Thus, at present, it is rather a mystery why the timescale for the onset of SNIa activity (presumably producing the observed decline of O/Fe in the disk) coincides with the timescale for halo formation. An original suggestion was recently made in Kobayashi et al. (1998), whereby SNIa appear at a rate which is metallicity dependant; the interest of this scenario lies in the fact that SNIa enter the cosmic scene at just the right moment. For the purpose of this work, we shall adopt the formalism of Matteucci and Greggio (1986), adjusting it as to have SNIa appearing mostly after the first Gyr, i.e. at a time when \[Fe/H\]$``$-1.
At this point we would like to point out that two recent observations (Israelian et al. 1998 and Boesgaard et al. 1999) challenged the “traditional” view of O vs Fe evolution, by finding a trend of O/Fe constantly increasing with decreasing metallicity (open triangles in Fig. 3). This intriguing trend is not confirmed by subsequent studies (Fullbright and Kraft 1999), but the question remains largely open today. If the new findings are confirmed, some of our ideas on stellar nucleosynthesis should be revised. Some possibilities of such a revision are explored in Sec. 6.
## 3 The model of galactic chemical evolution
Models of chemical evolution for the halo and the disk of the Milky Way are constructed adopting the standard formalism (Tinsley 1980, Pagel 1997). The classical set of the equations of galactic chemical evolution is solved numerically for each zone, without the Instantaneous Recycling Approximation (IRA). At the star’s death its ejecta is assumed to be thoroughly mixed in the local interstellar medium (instantaneous mixing approximation), which is then characterized by a unique composition at a given time. Abundance scatter cannot be treated in that framework, and this constitutes an important drawback of this type of “classical” models, since observations suggest a scatter of element to element ratios which increases with decreasing metallicity (Ryan et al. 1996). The basic ingredients of the model are described below.
### 3.1 Stellar lifetimes and remnant masses
The stellar lifetimes $`\tau _M`$ as a function of stellar mass M are taken from the work of the Geneva group (Schaller et al. 1992, Charbonnel et al. 1996), where the effects of mass loss on the duration of H and He burning phases are taken into account.
Stars with mass M$`<`$9$`M_{}`$ are considered to become white dwarfs with mass $`M_R(M/M_{})`$ =0.1$`(M/M_{}`$)+0.45 (Iben and Tutukov 1984). Stars with mass M $`>`$9$`M_{}`$ explode as core collapse supernovae leaving behind a neutron star of mass $`M_R=1.4M_{}`$ (as suggested by the observations of neutron stars in binary systems, e.g. Thorsett and Chakrabarty 1999). The heaviest of those stars may form a black hole, but the mass limit for the formation of stellar black holes is not known at present and cannot be inferred from theoretical or observational arguments (e.g. Prantzos 1994), despite occasional claims to the contrary. Due to the steeply decreasing stellar Initial Mass Function in the range of massive stars (see Sec. 3.2), as far as the mass limit for stellar black hole formation is M$`{}_{BH}{}^{}>`$40$`M_{}`$ the results of chemical evolution are not expected to be significantly affected by the exact value of M<sub>BH</sub>.
We stress that in our calculations we do take into account the amount of mass returned in the interstellar medium (ISM) by stars with M$`<`$11 M and M$`>`$40 M in the form of H, He, but also of all heavier elements, up to Zn. Since no yields are available for 9-11 M and $`>`$40 M stars (and since we deliberately neglect yields for intermediate mass stars), we simply assume that those stars return at their death in the ISM their initial amount of each element, i.e. that their net yield is zero for all elements (except for deuterium, which is destroyed). In that way we do not introduce any artificial modification of the adopted yields. This procedure is crucial for a correct evaluation of the metal/H ratio at a given time, especially at late times.
### 3.2 Star formation rate and initial mass function
Observations of average SFR vs. gas surface density in spirals and starbursts (Kennicutt 1998) are compatible with a Schmidt type law
$`\mathrm{\Psi }(t)=\nu \sigma _{gas}^k(t)`$ (2)
with $`k`$=1-2. However, this concerns only the disk averaged SFR and Kennicutt (1998) points out that the local SFR may have a different behaviour. Indeed, theoretical ideas of SFR in galactic disks suggest a radial dependence of the SFR (Wyse and Silk 1989) and such a dependence is indeed required in order to explain the observed abundance, colour and gas profiles in spirals (Boissier and Prantzos 1999, Prantzos and Boissier 2000). For the purposes of this work we adopt a Schmidt law with $`k`$=1.5; when combined with the adopted infall prescription (see next section) this leads to a slowly varying star formation history in the galactic disk, compatible with various observables (see Sec. 4). For consistency, we keep the same form of the SFR in the halo model, although there is no observational hint for the SFR behaviour during this early stage.
We adopt the IMF from the work of Kroupa et al. (1993, hereafter KTG93), where the complex interdependence of several factors (like stellar binarity, ages and metallicities, as well as mass-luminosity and colour-magnitude relationships) is explicitly taken into account. It is a three-slope power-law IMF $`\mathrm{\Phi }(M)M^{(1+x)}`$; in the high mass regime it has a relatively steep slope of $`X`$=1.7 (based on Scalo 1986), while it flattens in the low-mass range ($`X`$=1.2 for 0.5$`<`$M/M$`<`$1. and $`X`$=0.3 for M$`<`$0.5 M). We adopt this IMF between 0.1 and 100 M, although we are aware that there is some debate as to the exact form of the low-mass part. Again, for consistency, we adopt the same IMF in the halo and in the disk model.
### 3.3 Gaseous flows: infall and outflow
In most models of chemical evolution of the solar neighborhood, it is implicitly assumed that the old (halo) and young (disk) stars are parts of the same physical system, differing only by age; the same model is used to describe the whole evolution, from the very low metallicity regime to the current (supersolar) one (e.g. Timmes et al. 1995).
This assumption is, of course, false. The halo and the disk are different entities; different processes dominated their evolution, as revealed by the corresponding metallicity distributions (MD). In the case of the disk, observations show that the number of metal-poor stars is much smaller than what is predicted by the simple “closed-box” model of chemical evolution (the “G-dwarf problem”); the simplest explanation of that is that the disk evolved not as a closed box, but by slowly accreting infalling gas (e.g. Pagel 1997). In the case of the halo, the observed MD suggests that metal production was inefficient in those early times; the currently accepted explanation is that a strong outflow, at a rate $``$9 times the star formation rate, has occured during the halo evolution (as initially suggested by Hartwick 1976).
It is clear, then, that a unique model is inadequate to cover the whole evolution of the solar neighborhood. Still, this is done in most cases. Only in a handful of works has this point been taken into account, by adopting different prescriptions for the halo and the disk (Prantzos et al. 1993, Ferrini et al. 1993, Pardi et al. 1995, Chiappini et al. 1997, Travaglio et al. 1999), although not always the appropriate ones. The importance of that point is twofold: First, the corresponding MDs (the strongest constraints to the models) are only reproduced when appropriate models are used. Secondly, infall and outflow modify the timescales required for the gas to reach a given metallicity. This is important when one is interested in elements produced by e.g. intermediate mass stars, which enter late the galactic scene.
Another important point, related to the first one, is that the halo and the disk are, most probably, not related by any temporal sequence. Indeed, the gas leaving the halo ended, quite probably, in the bulge of the Galaxy, not in the disk, as argued e.g. by Wyse (2000 and references therein) on the basis of angular momentum conservation arguments. The disk may well have started with primordial metallicity, but a very small amount of gas. The corresponding small number of low metallicity stars that were formed by that gas explains readily the G-dwarf problem.
In the light of these arguments, we treat then the halo and the disk as separate systems, not linked by any temporal sequence. The local disk is assumed to be built up by slow accretion of gas with primordial composition. An exponentially decreasing infall rate $`f(t)e^{t/\tau }`$ with $`\tau `$ $`>`$ 7 Gyr is adopted. Such a long timescale has been shown (Chiappini et al. 1997, Prantzos and Silk 1998) to provide a satisfactory fit to the data of Wyse and Gilmore (1995) and Rocha-Pinto and Maciel (1996). We have normalized the infall rate $`f(t,R)`$, as to obtain the local disk surface density $`\mathrm{\Sigma }_T(R)`$=55 M pc<sup>-2</sup> at an age T=13.5 Gyr. Notice that chemodynamical models also support the idea of long time scales for the disk formation (Samland et al 1997).
For the halo model, there are less constraints: neither the duration of the halo phase, nor the final gas fraction or amount of stars are known. We assume then a duration of 1 Gyr and an outflow rate $`R_{out}`$ = 9 $`\mathrm{\Psi }(t)`$, in order to reproduce the observed halo MD. For consistency, we use the same SFR law and the same IMF as in the disk.
## 4 Evolution of the halo and the disk
We run two chemical evolution models, one for the halo (with outflow, for 1 Gyr) and one for the disk (with infall, for 13.5 Gyr), starting in both cases with gas of primordial composition. The only observational constraints common for the halo and the disk are: i) the metallicity distributions of low mass long-lived stars, and ii) the element/element ratio vs. metallicity (in particular, the O vs. Fe evolution). In the case of the disk there are several more constraints (see Sec. 4.2) but we turn first to (i) and (ii).
### 4.1 Metallicity distribution and O vs. Fe in the halo and the disk
In Fig. 3 we present our results and compare them to observations. The metallicity distributions ($`f`$=dN/d\[Fe/H\]) are normalised to $`f_{max}`$=1 and presented in the upper panel of Fig. 3. The adopted prescriptions (strong outflow for the halo and slow infall for the disk) lead to a satisfactory agreement between theory and observations, as expected on the basis of the discussion in Sec. 3.3. Notice that in the case of the disk, the theoretical curve shows a low metallicity tail below \[Fe/H\]=-1. However, the number of stars in the tail is extremely small, less than 10<sup>-2</sup> of the total. Although there is no “physical” discontinuity in the disk population at \[Fe/H\]=-1, we systematically show below all our results for the disk corresponding to \[Fe/H\]$`>`$-1 with thick solid curves, in order to stress that they correspond to what is traditionally thought as the “disk phase” of the Milky Way. Results for \[Fe/H\]$`<`$-1 are shown with thin solid curves, indicating that such stars do, in principle, exist, but in very small numbers.
Because a large part of Fe in the disk comes from SNIa (at least in our models) it is not clear whether the final G-dwarf metallicity distribution is mostly shaped by infall or by the rate of SNIa. In other terms, how can one be certain that the observed “G-dwarf problem” requires indeed large infall timescales (such as those discussed in Sec. 3.1 and adopted here)? We notice that the G-dwarf problem concerns mainly the low metallicity regime i.e. around \[Fe/H\]=-1 to -0.6; it is in this metallicity range that the closed box model predicts an excess of low-mass stars w.r.t the observations. But at those early times, corresponding to the first $``$2-4 Gyr of the disk’s history, the ratio of SNIa/SNII is still small (with the adopted prescription for the SNIa rate) and most of the Fe comes from SNII. Thus, the success of the model in reproducing the G-dwarf metallicity distribution does rely on the infall prescription, and not on the SNIa rate prescription. SNIa start becoming major sources of Fe somewhat later (around \[Fe/H\]=-0.5).
In the lower panel of Fig. 3 we show the corresponding evolution of O vs. Fe. It is virtually identical in the two models, up to \[Fe/H\]$``$-1, since both elements are primaries and produced in the same site (massive, short-lived, stars); their abundance ratio is then independant of infall or outflow prescriptions. As discussed in Sec. 2, the observed decline of O/Fe in the disk is reproduced by the delayed appearance of SNIa, producing $``$2/3 of the solar Fe.
Fig. 4 presents the evolution of the halo and the disk in a more “physical” way than in Fig. 3, i.e. various quantities are plotted as a function of time; time is plotted on a logarithmic scale (on the left, so that the halo evolution can be followed) and on a linear scale (on the right). The differences between the two models can be clearly seen. In particular, at a given time, the metallicity \[Fe/H\] (middle panel) is larger in the halo than in the disk (by 0.3 dex, i.e. a factor of 2); metallicity increases more slowly in our disk model than in the halo one. It takes $``$2 Gyr to the disk to reach \[Fe/H\]$``$-1, compared to $``$1 Gyr in the case of the halo. However, as noticed already, this early disk evolution concerns only very few stars.
### 4.2 Evolution of the local disk
There are many more observational constraints for the local disk than for the halo; an extensive presentation of those constraints can be found in Boissier and Prantzos (1999, their Table 1 and references therein). Here we present only briefly those constraints. Besides the MD and the O vs. Fe evolution, a satisfactory disk model should also reproduce:
(a) The current surface densities of gas ($`\mathrm{\Sigma }_G`$), stars ($`\mathrm{\Sigma }_{}`$), the total amount of matter ($`\mathrm{\Sigma }_T`$) and the current star formation rate ($`\mathrm{\Psi }_0`$);
(b) The age-metallicity relationship, traced by the Fe abundance of long-lived, F-type stars;
(c) The abundances of various elements and isotopes at solar birth $`(X_i,)`$ and today ($`X_i,0`$);
(d) The present day mass function (PDMF), resulting from the stellar IMF and the SFR history, which gives an important consistency check for the adopted SFR and IMF.
In Fig. 5 we present our results and compare them to constraints (a) and (b). It can be seen that the adopted SFR and infall rate lead to a current gas surface density of $`\mathrm{\Sigma }_G`$ 10 $`M_{}`$ pc<sup>-2</sup> and a final stellar surface density of $`\mathrm{\Sigma }_{}`$36 M pc<sup>-2</sup>, both in good agreement with observations. A current SFR $``$3.5 M pc<sup>-2</sup> Gyr<sup>-1</sup> is obtained at T=13.5 Gyr, also in agreement with observations. The evolution of the SFR is quite smooth, its current value being about half the maximum one in the past.
The lower panel of Fig. 5 shows the disk age-metallicity relation. The existence of an age-metallicity relation (AMR) in the disk is one of the important issues in studies of chemical evolution of the solar neighborhood. Several studies in the past showed a trend of decreasing metallicity with increasing stellar ages (Twarog 1980, Meusinger et al. 1991, and Jonch-Sorensen 1995). Edvardsson et al (1993) found an AMR consistent with these results but with a considerable scatter about the mean trend. However, this scatter (difficult to interpret in the framework of conventional models), may be due to contamination of the Edvardsson et al. (1993) sample by stars from different galactic regions (Garnett and Kobulnicky 2000). Indeed, the recent survey of Rocha-Pinto et al (2000, also on Fig. 5), suggests a scatter almost half of that in Edvardsson et al. (1993). In view of the current uncertainty, we consider that the mean trend of the disk AMR obtained with our model is in satisfactory agreement with observations.
In Fig. 6 we compare our results to constrain (c), i.e to the elemental (upper panel) and isotopic (lower panel) composition of the Sun. It is assumed that the Sun’s (and solar system’s) composition is representative of the one of the local interstellar medium (ISM) 4.5 Gyr ago, but this assumption is far from been definitely established. Indeed, CNO abundances in young stars and gas in the nearby Orion nebula show that the metallicity of this young region is lower than solar (Cunha & Lambert 1994, Cardelli & Federmann 1997); this cannot be readily interpreted in conventional models of chemical evolution. On the other hand, the Fe abundance of young stars determined by Edvardsson et al. (1993) seems to be compatible with the conventional picture, while the data of Rocha-Pinto et al (2000) suggest that the Sun is indeed Fe-rich w.r.t. other stars of similar age (Fig. 5). One should keep in mind this question (of the Sun being “typical” or not) when making detailed comparison of its composition to model predictions.
An inspection of Fig. 6 shows that there is good overall agreement between theory and observations, i.e. about 80% of the elements and isotopes are co-produced within a factor of two of their solar values. One should notice the following:
\- The carbon isotopes require another source. <sup>12</sup>C may be produced either by intermediate mass stars, as usually assumed, or by Wolf-Rayet stars with metallicity dependant yields (Maeder 1992, Prantzos et al. 1994). <sup>13</sup>C is made probably in intermediate mass stars. The evolution of <sup>12</sup>C/<sup>13</sup>C in the disk and its implications for the synthesis of those isotopes is studied in Prantzos et al. (1996).
\- The nitrogen isotopes also require another source. <sup>14</sup>N has the same candidate sites as <sup>12</sup>C. Novae seem to be a viable source for <sup>15</sup>N, but current uncertainties of nova models do not allow definite conclusions.
\- Fluorine is produced by neutrino-induced nucleosynthesis in WW1995, and this is also the case for a few other rare isotopes, not shown in Fig. 6 (<sup>7</sup>Li, <sup>11</sup>B). This is an interesting “new” nucleosynthesis mechanism, but because of the many involved uncertainties (see Woosley et al. 1990) it cannot be considered as established yet. One should keep an “open eye” for other, more conventional, sites of fluorine (as well as lithium and boron) nucleosynthesis, like e.g. Wolf-Rayet stars (Meynet and Arnould 1999).
\- The obtained overabundance of <sup>40</sup>K may reflect the large uncertainty in the abundance of this isotope at solar system formation (see Anders and Grevesse 1989), as already pointed out in Timmes et al. (1995).
\- Sc, V and Ti isotopes are underproduced, indicating that all currently available models of massive stars have some problems with the synthesis of these species.
\- There is a small overproduction of Ni, due to the isotope <sup>58</sup>Ni, which is abundantly produced in the W7 and W70 models of SNIa. This is also true for <sup>54</sup>Cr, a minor isotope of Cr. The amount of those nuclides depends mostly on the central density of the exploding white dwarf and the overproduction problem may be fixed by varying this parameter. Alternatives to the W7 model have recently been calculated by Iwamoto et al. (1999). On the other hand, Brachwitz et al. (2000) have explored the effect of electron capture rates on the yields of Chandrasekhar mass models for SNIa; they showed that the problem with <sup>54</sup>Cr may disappear (depending on the ignition density) while the one with <sup>58</sup>Ni is slightly alleviated. It can be reasonably expected that in future, improved, SNIa models, the overproduction problem of those nuclei will be completely solved.
Notice that in our calculation, the Fe-peak isotopic yields of WW1995 have been reduced by a factor of two, in order to reproduce the observed O/Fe ratio in halo stars ($``$3 times solar, see Fig. 3 and Sec. 5); otherwise, the WW1995 massive stars alone can make almost the full solar abundance of Fe-peak nuclei (as shown in Timmes et al. 1995), leaving no room for SNIa. Taking into account the uncertainties in the yields, especially those of Fe-peak nuclei (see Sec. 2) our reduction imposed on the WW1995 Fe yields is not unrealistic.
The nice agreement between theory and observations in Fig. 6 comes as a pleasant surprise, in view of the many uncertainties discussed in the previous section. It certainly does not guarantee that each individual yield is correctly evaluated. It rather suggests that the various factors of uncertainty cancel out (indeed, it is improbable that they all “push” towards the same direction!) so that a good overall agreement with observations results. As stressed in Timmes et al. (1995), the abundances of the intermediate mass isotopes span a range of 8 orders of magnitude; reproducing them within a factor of two suggests that models of massive stars nucleosynthesis are, globally, satisfactory. At least to first order, currently available yields of massive stars + SNIa can indeed account for the solar composition between O and Zn (with the exceptions of Sc, Ti and V, and possibly F).
## 5 Abundance ratios in the halo and the local disk
We calculated the abundance evolution of all the isotopes between H and Zn in the framework of our halo and local disk models, by using two different sets of massive star yields: i) the yields of WW1995 at constant (=solar) metallicity (Case A in the following), and ii) the metallicity dependant yields of WW1995, by interpolating between the values given for metallicities Z/Z=0, 10<sup>-4</sup>, 10<sup>-2</sup>, 10<sup>-1</sup> and 1 (Case B in the following). Because of our neglect of the C and N yields of intermediate mass stars, total metallicity is not consistently calculated in our models; we use oxygen as metallicity indicator, in order to inerpolate in the WW1995 tables (in the WW1995 models, the initial abundances of all elements are scaled to metallicity). Obviously, Case B (also studied by Timmes et al. 1995) is the “reference” case, whereas Case A is only studied for illustration purposes. In both cases, the yields of the W7 and W70 models of Iwamoto et al (1999) for SNIa are used (interpolated as a function of metallicity), whereas no yields from intermediate mass stars are considered; our explicit purpose is to check to what extent massive stars can account for observations of intermediate mass elements and for which elements the contribution of intermediate mass stars is mandatory. We stress again that we do take into account the contibution of intermediate and low mas stars to the H and He “budget”, since this is crucial for a correct evaluation of the metal/H ratio, especially at late times (Sec. 3.1).
Since most of the available data on the composition of stars concerns elemental abundances, we computed the corresponding evolution by summing over the calculated isotopic abundances. We present our results in Fig. 7 and compare them to a large body of observational data; most of the data come from the surveys of Ryan et al. (1996) and Mc William (1997) for the halo and Edvardsson et al. (1993) and Chen et al. (2000) for the disk, but we included a large number of other works, concerning specific elements (the corresponding references are listed in Table 1). We do not attempt here any discussion on the quality of these data (this would be beyond the scope of this work), and we refer the reader to the recent review of Ryan (2000) for that. It is obvious that systematic differences between various studies introduce a scatter larger than the real one (and, perhaps, unrealistic trends in some cases). Our reference Case B is shown in thick curves (dashed for the halo and solid for the disk), while Case A is in thin curves.
Before presenting our results we notice that in our models metallicity reaches \[Fe/H\]$``$-4 at a time t$``$10<sup>7</sup> yr and \[Fe/H\]$``$-3 at a time t$``$2 10<sup>7</sup> yr; these timescales correspond to the lifetimes of stars of mass M$``$20 M and M$``$10 M, respectively. Any variations in the abundance ratios in the metallicity range -4 $`<`$ \[Fe/H\] $`<`$ -3 results then from the fact that stars of different masses (starting from 100 M and going to 10 M) enter progressively the galactic scene. The discussion of Sec. 2 shows that the yields of individual stars are very uncertain, much more than those integrated over the IMF (the latter reproduce, at least, the solar composition!). Besides, there is absolutely no guarantee that the model reproduces correctly the relation between age and metallicity at those early times. For instance, in a recent work Argast et al. (2000) find that the halo became chemically homogeneous and reached \[Fe/H\]=-3 after $``$160 Myr, a duration six times longer than in our calculations. For those reasons we consider that any abundance trends of our models at \[Fe/H\]$`<`$-3 are not significant, but we show them for completeness. Integration over the whole IMF of massive stars is only made for \[Fe/H\]$`>`$-3 and we consider that our results are significant only after that point. Finally, we notice that we have reduced the WW1995 yields of Fe-peak isotopes by a factor of two, in order to reproduce the observed $`\alpha `$/Fe ratio in the halo.
### 5.1 Carbon and Nitrogen
Observations indicate a flat \[C/Fe\]$``$0 in the halo and the disk, with a large dispersion at all metallicities. Both our cases A and B show indeed \[C/Fe\]$``$0 in the halo (since both C and Fe are primaries), and a slow decline of C/Fe in the disk due to Fe production by SNIa. As discussed in Sec. 2, a complementary source of C is required in the disk. This may be either intermediate mass stars (IMS) or Wolf-Rayet (WR) stars. However, as discussed in Prantzos et al. (1994), IMS have masses M$`>`$3 M and lifetimes $`\tau <`$5 10<sup>8</sup> yr. Such stars can certainly evolve during the halo phase (if the duration of that phase is indeed $``$1 Gyr, as assumed here) and enrich the halo with C, thus rising the C/Fe ratio at \[Fe/H\]$`<`$-1. Such a behaviour is not observed, however, suggesting either that low mass stars (M$`<`$2 M) or WR stars are the main carbon sources in the disk. The latter possibility is favoured in Prantzos et al. (1994) and Gustafsson et al. (1999)
Nitrogen behaves in a similar way as carbon, i.e. the observed \[N/Fe\]$``$0 in the halo and the disk, with a large scatter at low metallisities. Our Case A (metallicity independant yields) shows also a flat \[N/Fe\]$``$0 evolution in the halo and a decline in the disk, exactly as for carbon. However, in the realistic Case B, N behaves as secondary: \[N/Fe\] increases steadily up to \[Fe/H\]$``$-1. Its value remains $``$constant in the disk phase, because Fe production by SNIa compensates for the larger N yields of more metal rich stars. However, the final N/Fe is only $``$1/3 its solar value.
Obviously, curent massive star yields fail, qualitatively and quantitavely, to reproduce the observed evolution of N/Fe. What are the alternatives? in our view, there are two:
a) Intermediate mass stars, producing primary N through hot-bottom burning in the AGB phase, are the most often quoted candidate. Large uncertainties still affect that complex phase of stellar evolution, but recent studies (e.g. Lattanzio 1998 and references therein) find that hot-bottom burning does indeed take place in such stars. If N is indeed produced as a primary in IMS, and their N yields are metallicity independant, then the N/Fe in the disk should decline (because of SNIa). Metallicity dependant N yields from WR stars (Maeder 1992) could compensate for that, keeping the N/Fe ratio $``$constant in the disk. On the other hand, if N from massive stars is indeed secondary, at some very low metallicity level (let’s say \[Fe/H\]$`<`$-3) the N/Fe ratio should also decline; this would be an important test of IMS being the main N source in the halo. If such a decline is not observed, we are lead to the second alternative, namely
b) Massive stars, producing primary N by an as yet unidentified mechanism, obviously requiring proton mixing in He-burning zones. Such mixing does not occur in standard stellar models, but “new generation” models including rotation offer just such a possibility (Heger et al. 1999, Maeder and Meynet 2000). In that case, N is produced not by the original carbon entering the star, but by the carbon produced in He-burning; as a consequence, it is produced as a primary. In that case, massive stars could be the main source of N and C in the halo.
The discussion of this section suggests then an intriguing possibility: massive stars could well be the main source of C and N in both the halo and the disk (in the latter case, through the WR winds), leaving only a minor role to intermediate mass stars!
### 5.2 $`\alpha `$ \- elements O, Mg, Si, S, Ca, Ti
The alpha elements (O, Mg, Si, S, Ca, Ti) present a well known behaviour. The $`\alpha `$/Fe ratio is $``$constant in the halo, at \[$`\alpha `$/Fe\]$``$0.3-0.5 dex, and declines gradually in the disk. The latter feature is interpreted as due to (and constitutes the main evidence for) the contribution of SNIa to the disk composition.
This behaviour is indeed apparent in Fig. 7; despite the large scatter, all the alpha elements show the aforementioned trend. We stress here again that the recent data of Israelian et al. (1998) and Boesgard et al. (1999), also plotted in Fig. 7 (with different symbols), challenge this picture in the case of oxygen. If true, these new data should impose some revision of our ideas on massive star nucleosynthesis, probably along the lines suggested in Sec. 6.
Until the situation is clarified, we stick to the “old paradigm”. In the framework of this “paradigm”, Pagel and Tautvaisiene (1995) have shown that the $`\alpha `$/Fe evolution can be readily explained by a very simple model (with IRA), the metallicity independant yields of Thielemann et al. (1996) and SNIa during the disk phase. On the other hand, Timmes et al. (1995), using the metallicity dependant yields of WW1995 (but an inappropriate model for the halo, see Sec. 3.3), found good agreement with observations, provided that the Fe yields of WW1995 are reduced by a factor of $``$2.
Our results in Fig. 7 point to the following:
\- For O, Si, S and Ca, both Cases A and B give virtually identical results. These elements behave as true primaries, without any metallicity dependence of their yields. Moreover, after the WW1995 Fe yields are reduced by a factor of 2, a fairly good agreement with observations is obtained.
\- The situation is far less satisfactory for Mg and Ti. For both of them, the WW1995 yields at solar metallicity are larger than at lower metallicities (see Fig. 1). This is puzzling since Mg and Ti are also supposed to be primaries (in fact, more puzzling in the case of Mg, since Ti is produced close to the “mass-cut” and subject to more important uncertainties). As a result, our Case A is marginally compatible with observations of Mg/Fe; the reference Case B does not match at all the observations, despite the reduction of the Fe yields by a factor of 2. In the case of Ti, both Cases A and B fail to match the observations.
These features were also noticed in Timmes et al. (1995) and the problem with the WW1995 yields of Mg and Ti pointed out; however, no satisfactory alternative was suggested. Since the Mg yields of WW1995 are steeply increasing function of stellar mass, our use of the Kroupa et al. (1993) IMF (steeper than the Salpeter IMF used by Timmes et al. 1995) leads to a low Mg/Fe ratio, even after reduction of the Fe yields. Our Fig. 1 (lower panel) suggests that the yields of LSC2000 could match better the halo data, since the Mg/Fe and Ti/Fe ratios obtained for Z=0 are larger than solar. On the other hand, Fig. 1 shows that in both WW1995 and LSC2000, Mg and Ti have lower overproduction factors than all the other alpha elements, at all metallicities; this means that, even if the halo Mg/Fe and Ti/Fe ratios are better reproduced with the LSC2000 yields, the corresponding $`\alpha `$/Mg and $`\alpha `$/Ti ratios will certainly not match the observational data. Thus, at present, none of the two available sets of metallicity dependant yields offers a solution to the problem of Mg and Ti.
The fact that Pagel and Tautvaisiene (1995) find good agreement with observations by using the Thielemann et al. (1996) yields may suggest that this set of yields indeed solves the problem. This is also the case in Chiappini et al. (1999), who use a somewhat different prescription for SNIa rate than here, and metallicity independant yields from Thielemann et al. (1996) and WW19995. Notice, however, that metallicity independant yields (those of Thielemann et al. 1996 are for solar metallicity only) should not be used for studies of the halo, even if the problem is less severe in the case of primary elements. The equivalent set of WW1995 yields also reproduces the Mg/Fe evolution in the halo (our Case A), but it is not appropriate. We need to understand how massive stars make a $``$constant Mg/Fe and Ti/Fe ratio at all metallicities, by using stellar models with the appropriate initial metallicity.
### 5.3 Sodium and Aluminium
Na and Al are two monoisotopic, odd elements. Their theoretical yields are, in principle, affected by the “odd-even” effect (see Sec. 2). This effect seems to be stronger in the case of LSC2000 than in WW1995 (Fig. 1), at least for the adopted IMF.
The observational situation for those elements is not quite clear. Recent observations (Stephens 1999) suggest that Na/Fe decreases as one goes from \[Fe/H\]=-1 to \[Fe/H\]=-2, as expected theoretically. However, most other observations do not support this picture, showing instead a flat \[Na/Fe\]$``$0 ratio with a large scatter. Our Case A evolution of Na/Fe is similar to the $`\alpha `$/Fe evolution and, obviously, incorrect. In Case B, Na/Fe increases steadily after \[Fe/H\]$``$-2.5 and reaches a plateau after \[Fe/H\]$``$-1. Neither case matches the observations well. As we shall see in Sec. 5.7, the situation improves considerably when only the halo data of Stephens (1999) and the disk data of Edvardsson et al. (1993) and Feltzing and Gustafsson (1998) are used; then Na vs. Ca shows the behaviour of an odd element, as it should.
Ryan et al. (1996) find a steep decline of Al/Fe at low metallicities, down to “plateau” value of \[Al/Fe\]$``$-0.8, but they stress that their analysis neglects NLTE effects and underestimate the real Al/Fe ratio; for that reason we do not plot their data in Fig. 7 (Ryan et al. 1996 suggest that a NLTE correction to their data would move the “plateau” value to \[Al/Fe\]$``$-0.3, i.e. consistent with what expected for an odd-Z element). On the other hand, the NLTE analysis of the data of Baumüller and Gehren (1997, open triangles in Fig. 7) suggests a practically flat Al/Fe ratio in the halo, a rather unexpected behaviour for an “odd” element. In our model Case A, Al behaves like an $`\alpha `$ element. In Case B, the “odd-even” behaviour is manifest: a small increase of Al/Fe is obtained as metallicity increases from \[Fe/H\]$``$-2.5 to \[Fe/H\]$``$-1 (the model trend below \[Fe/H\]=-3, due to stellar mass and lifetime effects, is not significant, as stressed in the begining of Sec. 5). Once again, theory does not match observations and observations do not show the expected behaviour.
It should be noted at that point that intermediate mass stars of low metallicity could, perhaps, produce some Na and Al through the operation of the Ne-Na and Mg-Al cycles in their H-burning shells and eject them in the interstellar medium through their winds. There are indeed, indications, that in low mass, low metallicity stars of globular clusters such nucleosynthesis does take place (Kraft et al. 1998). If this turns out to be true also for intermediate mass stars of low metallicity, it might considerably modify our ideas of Na and Al nucleosynthesis in the halo.
### 5.4 Potassium, Scandium, Vanadium
K, Sc and V are three odd-Z elements produced mainly by oxygen burning. However, the first one is produced in hydrostatic burning and the other two in explosive burning, i.e. their nucleosynthesis is more uncertain. Their yields are affected in similar ways by the initial metallicity of the star, as can be seen in Fig. 1.
Currently available observations show a rather different behaviour for those elements: Sc/Fe remains $``$solar in the whole metallicity range -3$`<`$\[Fe/H\]$`<`$0. V/Fe is also $``$solar in the disk and the late halo, but appears to be supersolar in the range $`3<`$\[Fe/H\]$`<`$-2 (although the data is rather scarce for a definite conclusion). Finally, K/Fe declines in the disk, while the rare halo data point to supersolar ratio \[K/Fe\]$``$0.5, i.e. its overall behaviour is similar to that of an $`\alpha `$-element!
From the theoretical point of view, the situation is also unsatisfactory. Cases A and B produce distinctively different results for Sc and V, but not so for K. In Case B, the Sc/Fe and V/Fe ratios are subsolar in the halo, while K/Fe is supersolar. Also, in that case, K/Fe declines in the disk, Sc/Fe remains $``$constant and V/Fe increases.
This “strange” theoretical behaviour results from the interplay of several factors, which do not affect all those elements in the same way: odd-even effect, Fe yield reduction and contribution of SNIa. Thus, the metallicity dependence of the yields between Z=0.1 Z and Z=Z is stronger for V than for the other two. In fact, the V yield at metallicity Z=0.1 Z is lower than at Z=0.01 Z in WW1995, which is counterintuitive (making V/Fe to decrease between \[Fe/H\]=-2 and \[Fe/H\]=-1). Also, SNIa contribute more to the production of V than to the one of Sc or K (at least according to the W7 model). For those reasons, Sc/Fe is $``$ constant in the disk, while K/Fe declines and V/Fe increases.
Although our Case B seems to match well the available data for K, we think that this is rather fortuitous: we obtain a supersolar K/Fe in the halo because of the reduction of the Fe yields by a factor of two and of the adopted IMF (Timmes et al. 1995 obtain a solar K/Fe in the halo for the same reduced Fe yields, probably because they use the Salpeter IMF).
In our view, the evolution of those three elements is far from being well understood, either observationaly or theoretically. They do not show any sign of the expected odd-even effect (rather the opposite behaviour is observed for K!). However, if theoretical “prejudices” are put aside, the situation may not be as bad for Sc and V: indeed, they are part of the “low iron group” elements and their abundances may well follow the one of Fe, as suggested by current observations. In that case, the “odd-even” effect is overestimated in the theoretical yields adopted here or those of LSC2000 (Fig. 1). We also noticed that their solar abundances are underproduced by current nucleosynthesis models (Sec. 4.2 and Fig. 6).
### 5.5 Fe-peak elements: Cr, Mn, Co, Ni, Cu and Zn
The various isotopes of the Fe peak are produced by a variety of processes (see WW1995): isotopes with mass number A$`<`$57 are produced mainly in explosive O and Si burning and in nuclear statistical equilibrium (NSE). Isotopes with A$`>`$56 are produced in NSE (mostly in “alpha-rich freeze-out”), but also by neutron captures during hydrostatic He- or C-burning. Because of the many uncertainties involved in the calculations (sensitivity to the neutron excess, the mass-cut, the explosion energy etc.) the resulting yields are more uncertain than for the other intermediate mass nuclei.
Observations show that the abundance ratio to Fe of Cr, Co, Ni and Zn is $``$solar down to \[Fe/H\]$``$-2.5 to -3. This fact, known already in the late 80ies, suggests that those elements behave similarly to Fe (at least in this metallicity range) and, therefore, are produced in a quite similar way. However, observations in the mid-90ies (Ryan et al. 1996, McWilliam 1997) show that, as one goes to even lower metallicities, a different picture is obtained (see Fig. 7): Cr/Fe is subsolar and decreasing, while Co/Fe is supersolar and increasing; the situation is less clear for Ni/Fe, but in all cases the scatter is larger at very low metallicities than at higher ones.
For the reasons mentioned in the beginning of Sec. 5, we do not consider the trends of our models in the range \[Fe/H\]$`<`$-3 to be significant. We do not then attempt here to interpret those recent intriguing findings, which point, perhaps, to some interesting physics affecting the evolution of the first stellar generations. We simply notice that such an attempt is made in Nakamura et al. (1999), who study the sensitivity of the corresponding yields to various parameters (neutron excess, mass-cut, explosion energy). Their conclusion is that the observed Co/Fe excess cannot be explained by any modification of those parameters.
The yields of WW1995 show a mild metallicity dependence in the case of Cr and Ni and a more important one in the cases of Mn, Co, Cu and Zn. For that reason, we obtain different results for those elements between our Cases A and B (Fig. 7). The situation for each of those elements is as follows:
\- The Cr/Fe evolution is reproduced satisfactorily for \[Fe/H\]$`>`$-2.5; in the disk, Cr and Fe are produced in similar amounts by SNIa and the Cr/Fe ratio remains $``$constant.
\- Co/Fe decreases steadily as one goes to low metallicities (in Case B). This trend is not observed in the data and suggests that the “odd-even” effect for that nucleus is overestimated in WW1995; we notice that LSC2000 find a much smaller effect (Fig. 1).
\- The WW1995 yields adequately describe the Ni/Fe evolution, except at the lowest metallicities (\[Fe/H\]$`<`$-3). The LSC2000 yields would face the same problem, as can be seen in Fig. 1. The excess of Ni/Fe obtained in the disk model is due to the overproduction of <sup>58</sup>Ni by the W7 model of SNIa (see Sec. 4.2).
\- The WW1995 yields suggest a $``$constant (solar) Zn/Fe in the halo, albeit at a value lower than actually observed. On the other hand, they suggest that Zn/Fe should increase in the disk, while observations show no such increase. An inspection of the LSC2000 yields in Fig. 1 suggests that they would face the same problems.
\- Finally, the WW1995 yields offer an excellent description of the observed evolution of Mn/Fe and Cu/Fe. If the observations are correct, we have an exquisite realisation of the “odd-even” effect for Fe-peak nuclei (especially in the case of Mn), almost a “text-book” case. An inspection of the LSC2000 yields shows that they would do equally well.
### 5.6 Fluorine, Neon, Phosphorous, Chlorine, Argon
We present in Fig. 7 the evolution of those elements according to our models, although no observational data exist for them in stars; fluorine is an exception, its abundance being measured in giants and barium stars (Jorissen et al. 1992).
We recall that F is produced in WW1995 mainly by neutrino-induced nucleosynthesis (spallation of <sup>20</sup>Ne) and the corresponding yields are very uncertain. As seen in Fig. 1, the F yield of WW1995 are metallicity dependant, and this is also reflected in the evolution of the F/Fe ratio (Case A vs Case B). We notice again that F may also be produced in other sites, like in the He-burning shells of AGB stars (as suggested by the calculations of Forestini and Charbonnel 1997) or in WR stars. The recent calculations of Meynet and Arnould (2000) show that the F yields of the latter site are also metallicity dependant, but they are important only for metallicities \[Fe/H\]$`>`$-1; at lower metallicities, very few massive stars turn into WR. Obviously, if AGB and WR stars are the main producers of F, the evolution of F/Fe ratio may be quite different from the one shown in Fig. 7.
The main Ne isotope is <sup>20</sup>Ne, i.e. Ne should evolve as an $`\alpha `$-element. The evolution of Ne/Fe in Fig. 7 is similar to the one of C/Fe. The yields of WW1995 show a small metallicity dependence (reflected in Case A vs. Case B) not exhibited by the yields of LSC2000.
Like Ne, Ar is also an even-Z element. There is no metallicity dependence in the Ar yields of WW1995 (which explains the similarity between cases A and B), neither in those of LSC2000. Ar is expected to behave like Si or Ca.
P and Cl are odd-Z elements. When the WW1995 Fe yields are divided by 2, a $``$solar P/Fe and a supersolar Cl/Fe ratio is obtained for halo stars. In the disk, enhanced P production by massive stars (due to the “odd-even” effect) and by SNIa compensate for the Fe production by SNIa; as a consequence, the P/Fe ratio decreases only very slightly. On the contrary, this compensation does not occur for Cl and the Cl/Fe ratio decreases in the disk.
In the absence of observational data, the nucleosynthesis of these elements can not be put on a firm basis. Their solar abundances are relatively well reproduced with the WW1995 yields (Fig. 6), and this is quite encouraging. On the other hand, we notice that the LSC2000 yields show a more pronounced “odd-even” effect for P and Cl than WW1995.
### 5.7 Chemical evolution with respect to Ca
Traditionally, the results of galactic chemical evolution studies are presented as a function of Fe/H, i.e. Fe is assumed to play the role of “cosmic clock”. However, in view of the uncertainties on Fe production and evolution (due to mass cut and explosion energy in SNII, or to the uncertain evolution of the rate of SNIa), it has been suggested that Fe should be replaced by a “robust” $`\alpha `$ element, like e.g. O or Ca.
In view of the uncertainties currently affecting the observational status of oxygen, we choose here Ca as the reference element. Among the data listed in Table 1 (and plotted in Fig. 7) we selected those including observations of Ca abundances and we plot the element/Ca ratios in Fig. 8 as a function of Ca/H. We also plot on the same figure the corresponding model results obtained with the metallicity dependant yields of WW1995 and the W7 model for SNIa (i.e. our Case B).
Several interesting features can be noticed:
\- For O, Al, K and V, existing data concern only the disk phase and are consistent with X/Ca$``$solar. Model results show that O/Ca and K/Ca ratios are solar over the whole metallicity range; they also show clearly the “odd-even” effect for Al/Ca, V/Fe and Cu/Fe.
\- Among the $`\alpha `$-elements, the observed Mg/Ca and Si/Ca ratios are solar down to very low metallicities. In our models, we also find constant Mg/Ca and Si/Ca ratios, slightly below the observed values in the former case, and in fair agreement with the observations in the latter.
\- The observed Na/Ca evolution shows clearly the “odd-even” effect, especially with the recent data of Stephens (1999) for the metallicity range -1.5$`<`$\[Ca/H\]$`<`$-0.5 and those of Feltzing and Gustafsson (1998) for \[Ca/H\]$`>`$0. This behaviour is fairly well reproduced by the model.
\- The observed Sc/Ca and Ti/Ca ratios are slightly below their solar values in the halo, with some hint for a decrease of the latter ratio at very low metallicities. Model results are broadly compatible with those observations.
\- Cr/Ca, Fe/Ca and Mn/Ca ratios are all lower than solar in the whole metallicity range, exactly as observed. The agreement between the model results and the data is excellent for all three cases, down to the lowest metallicities; notice that the evolution of Cr w.r.t. Fe was not so well reproduced by the model at the lowest metallicities (Fig. 7).
\- Finally, the observed Co/Ca and Ni/Ca ratios decrease with decreasing Ca/H down to \[Ca/H\]$``$-2 and increase at lower metallicities. The former trend is rather well reproduced by the model, but not the latter. The problematic behaviours of Co and Ni at low metallicities do not disappear when Ca is adopted as “cosmic clock”.
## 6 Alternatives for Oxygen vs Iron
In the previous sections we treated oxygen exactly as the other $`\alpha `$-elements, i.e. by assuming that\[O/Fe\]$``$0.4$``$constant in the halo. However, the recent intriguing findings of Israelian et al. (1998) and Boesgaard et al. (1999) suggest that O/Fe continues to rise as one goes from the disk to halo stars of low metallicities (we shall call these data “new data” in this section). Although the observational status of O/Fe is not settled yet, the “new data” certainly call for alternatives to the “standard” scenario to be explored.
An obvious alternative is to assume that Fe producing SNIa enter the galactic scene as early as \[Fe/H\]$``$-3, instead of \[Fe/H\]$``$-1 in the “standard” scenario. Indeed, the first white dwarfs, resulting from the evolution of $``$8 M stars, are produced quite early on in the galactic history; if their companions are almost equally massive, their red giant winds would push rapidly the white dwarf beyond the Chandrasekhar mass, and induce a SNIa explosion. The subsequent evolution of the SNIa rate (not well known today), should then be such as to ensure a continuous, smooth decline of O/Fe with \[Fe/H\], as the “new data” suggest. Such a behaviour is indeed obtained in the calculations of Chiappini et al. (1999), which have not been adjusted as to fit the new data: it is a direct consequence of their adopted formalism for the SNIa rate.
The problem with this “alternative” is that it also affects the evolution of the other $`\alpha `$/Fe abundance ratios in the halo. Observationaly, none of the $`\alpha `$-elements shows a behaviour comparable to the one suggested by the “new data” for oxygen (see Fig. 7 for Mg, Si and Ca). The “new data” can simply not be explained in terms of SNIa only, because this would spoil the current nice agreement with the other $`\alpha `$-elements (see Fig. 9a). \[Notice: C/Fe would also decrease with metallicity quite early in that case, but this is not a serious problem, since C from intermediate mass stars could keep the C/Fe ratio close to solar, as observed (and indicated in Fig. 9)\].
A second possibility is that the O yields from massive stars are, for some reason, metallicity dependant. It is already known that this happens for the C and N yields of massive stars, for metallicities Z$`>`$0.1 Z: because of intense stellar winds, the most massive stars lose their envelope already during He-burning. This envelope is rich in H-burning products (like He and N) and later in early He-burning products (essentially C). Thus, less mass is left in the He-core to be processed into oxygen (Maeder 1992). As discussed in Sec. 5.1, this metallicity dependence of C yields from massive (WR) stars, can indeed explain the observed C/O evolution in the disk. However, Prantzos et al. (1994) have shown that the effect is clearly negligible for the evolution of oxygen in the disk, at least with Maeder’s (1992) yields. And at lower metallicities, the effect is virtually inexistent: even the most massive stars present negligible mass losses. Thus, current models suggest that metallicity dependant Oxygen yields cannot help explaining the new data.
However, the effect may have been underestimated. After all, stellar mass loss is yet poorly understood. Suppose then that, starting at \[Fe/H\]$``$-3, massive stars produce less and less oxygen as their metallicity increases, because an ever larger part of their envelope is removed. Their inner layers, producing the other $`\alpha `$-elements and Fe, are not affected by mass loss; the resulting $`\alpha `$/Fe abundance ratio is constant with metallicity, while the corresponding O/Fe is decreasing with metallicity. The problem encountered by the first alternative seems to be solved. However, in the expelled mass of those stars, the abundances of He, N and C should be particularly enhanced. The resulting N/Fe and C/Fe ratios should be steadily increasing with metallicity in the halo (see Fig. 9), which is not observed; and introducing N and C from IMS would only make things worse. Thus, several arguments suggest that metallicity dependant oxygen (and, by necessity carbon) yields of massive stars cannot explain the “new data”.
A third alternative concerns the possibility of having metallicity dependant yields of Fe and all elements heavier than oxygen (while keeping the O,N,C yields independant of metallicity below \[Fe/H\]$``$-1). In that case, the yields of $`\alpha `$-elements and Fe would decrease with decreasing metallicity at the same rate, producing a quasi-constant $`\alpha `$/Fe abundance ratio in the halo, as observed. The O/Fe and C/Fe ratios would both decrease with increasing metallicity (Fig. 9); however, in the latter case, this decrease would be compensated by C production from IMS, so that the C/Fe ratio would remain $``$constant in the halo, as observed. Thus, from the three studied alternatives, we think that only the last one cannot be at present rejected on observational grounds.
What could be the physics behind such a metallicity dependence of the yields of $`\alpha `$-elements and Fe in massive stars? First, we notice that the required effect is very small: a factor of $``$3 increase is required in the yields for a 100-fold increase in metallicity (between \[Fe/H\]=-3 and \[Fe/H\]=-1, see Fig. 9), i.e. of the same order as the “odd-even” effect in Fig. 1. Our scenario requires that the supernova layers inside the C-exhausted core (i.e. the layers containing all the elements heavier than oxygen) be well mixed during the explosion. Various instabilities could contribute to that, either in the pre-supernova stage (in the O-burning shell, Bazan and Arnett 1998) or during the explosion itself (as in SN1987A, Arnett et al. 1989). This is required in order to ensure that the $`\alpha `$/Fe ratio will be $``$ constant in the ejecta. But the main ingredient is that the structure of the star depends on metallicity, in the sense that lower metallicity cores are more compact than higher metallicity ones. Then, at the lowest metallicities (say \[Fe/H\]$``$-3), after the passage of the shock wave, a relatively large proportion of the well mixed C-exhausted core will fall back to the black hole, feeling a strong gravitational potential. At higher metallicities, the core is less compact and a larger proportion of the C-exhausted core escapes. At all metallicities, oxygen (and lighter elements as well) are located in the losely bound He-layers and manage always to escape with the same (metallicity independant) yields.
If the “new data” of Israelian et al. (1998) and Boesgaard et al. (1999) on O vs Fe are confirmed, some radical revision of our ideas on stellar nucleosynthesis will be required. At present, we think that our third alternative (schematically illustrated in the right panels of Fig. 9) is both plausible and compatible with all currently available data.
## 7 Evolution of Mg isotopic ratios
There are very few cases where observations allow to check models of isotopic abundance evolution in the Galaxy, especially concerning the early (i.e. halo) phase of that evolution. One of these rare cases concerns the Mg isotopes <sup>25</sup>Mg and <sup>26</sup>Mg.
All magnesium isotopes are mainly produced by hydrostatic burning in the carbon and neon shells of massive stars. The production of the neutron-rich isotopes <sup>25</sup>Mg and <sup>26</sup>Mg is affected by the neutron-excess (i.e. their yields increase with initial stellar metallicity) while <sup>24</sup>Mg is produced as a primary (in principle). Thus, the isotopic ratios <sup>25</sup>Mg/<sup>24</sup>Mg and <sup>26</sup>Mg/<sup>24</sup>Mg are expected to increase with metallicity.
Observational evidence of a decline of the abundances of <sup>25</sup>Mg and <sup>26</sup>Mg relative to <sup>24</sup>Mg in low metallicity stars was reported as early as 1980 (Tomkin and Lambert 1980). In a recent work Gay and Lambert (1999) derived Mg isotopic abundance ratios for 19 dwarf stars in the metallicity range -1.8$`<`$\[Fe/H\]$`<`$0, using high resolution spectra of the MgH A-X 0-0 band at 5140 Å. They compared their observations with the theoretical predictions of Timmes et al. (1995) in the solar neighbourhood and found an overall good agreement.
The evolution of Mg isotopic abundance ratios of our models is plotted as a function of \[Fe/H\] in Fig. 10. The upper panel represents the evolution of <sup>25</sup>Mg/<sup>24</sup>Mg and the lower panel the one of <sup>26</sup>Mg/<sup>24</sup>Mg. Both ratios increase slowly with \[Fe/H\]. <sup>25</sup>Mg/<sup>24</sup>Mg becomes slightly larger than the corresponding solar ratio at \[Fe/H\]$``$0, while <sup>26</sup>Mg/<sup>24</sup>Mg is 60% higher than solar at that metallicity. This is consistent with the results of Fig. 6 (lower panel), showing that <sup>26</sup>Mg is produced with its solar value at Sun’s formation, while <sup>25</sup>Mg and <sup>24</sup>Mg are slightly underproduced. We notice that Timmes et al. (1995) find also supersolar Mg isotopic ratios at \[Fe/H\]=0, but the <sup>26</sup>Mg excess is not as large as ours. We think that this difference is due to our use of the Kroupa et al. (1993) stellar IMF, favouring the <sup>26</sup>Mg yields w.r.t those of <sup>24</sup>Mg; Timmes et al. use the Salpeter IMF.
In Fig. 10 we compare our results with observations from various sources, including the recent data of Gay and Lambert (1999). The observational trends are, globally, reproduced by our model for disk stars, although the <sup>26</sup>Mg/<sup>24</sup>Mg ratio is higher than observed for stars of near solar metallicity. More interesting is the fact that the model isotopic ratios are systematically lower than observations for halo stars (below \[Fe/H\]$``$-1). This was also noticed in Timmes et al. (1995). It may well be that the WW1995 yields underestimate the importance of the neutron-excess in the production of the Mg isotopes at those metallicities. Another possibility is that there is some other source of the neutron-rich Mg isotopes in the late halo, like e.g. AGB stars with He-shells hot enough to activate the <sup>22</sup>Ne($`\alpha `$,n)<sup>25</sup>Mg neutron source. This reaction, would not only provide neutrons for the s-process in those stars, but it would also produce large amounts of <sup>25</sup>Mg and <sup>26</sup>Mg. At present, the operation of that source in AGB stars of disk-like metallicities seems improbable, but there is no evidence as to what may happen at lower metallicities.
## 8 Summary
In this work we present a comprehensive study of the evolution of the abundances of intermediate mass elements (C to Zn) in the Milky Way halo and in the local disk. We use a consistent model in order to describe the evolution of those two galactic subsystems. The model assumes strong outflow in the halo phase and slow infall in the disk, which allow to correctly reproduce the corresponding metallicity distributions; these observables constitute the strongest constraints for chemical evolution models of those regions. Also, we consider the halo and the disk to evolve independently, since there is no hint at present for a physical connection between the two (see Sec. 3.3). We note that this type of modelisation has very rarely been done before.
The second important ingredient of this study is the consistent use of metallicity dependant yields for all isotopes. We adopt the yields of WW1995 and we note that there is a remarkably good agreement between them and the more recent ones of LSC2000 (but also some important differences). Only one study of similar scope has been done before with the metallicity dependant WW1995 yields (Timmes et al. 1995), but it utilised an inconsistent model for the halo. The study of Samland (1998) used appropriate models for the halo and the disk, but made several approximations concerning the stellar lifetimes and the metallicity dependence of the yields. We note that we have divided the (uncertain, anyway) Fe-peak isotopic yields of WW1995 by a factor of 2, in order to obtain abundance ratios w.r.t Fe consistent with observations; indeed, Timmes et al. (1995) recognised the problem with the WW1995 Fe yields and presented also results for twice and half the nominal values. We also performed calculations with metallicity independant yields (at solar metallicity only) in order to illustrate the differences with the metallicity dependant ones. In all cases we used the recent yields of Iwamoto et al. (1999) for SNIa, which are also metallicity dependant (this dependence affects very little the results). We only used yields from massive stars and SNIa, in order to find out for which elements and to what extent is the contribution of other sources mandatory.
We compared our results to a large body of observational data. In Sec. 4 we “validated” our model, by showing that it reproduces in a satisfactory way all the main observational constraints for the halo and the local disk. We found that the resulting elemental and isotopic compositions at a galactic age of 9 Gyr compare fairly well to the solar one ; among the few exceptions, the most important ones concern:
a) The C and N isotopes, which are underproduced. For the major ones (<sup>12</sup>C and <sup>14</sup>N), both WR and IMS are candidate sources; for <sup>13</sup>C and <sup>15</sup>N, IMS and novae are, respectively, the main candidates.
b) The isotopes of Sc, Ti and V, for which there is no other candidate source. The fact that the corresponding LSC2000 yields are even lower than WW1995 may point to some generic problem of current nucleosynthesis models for those elements.
We consider our results for the halo evolution to be significant only above \[Fe/H\]$`>`$-3. The reason is that at lower metallicities massive stars have lifetimes comparable to the age of the halo at that point; since the yields of individual stars are very uncertain, we consider that the corresponding results have little meaning. Only when the age of the halo becomes significantly larger than the lifetime of the “lightest” massive star (and ejecta are averaged over the IMF for all massive stars) we consider our results to become significant. For that reason, we are not able to draw any conclusion on the puzzling behaviour of the Fe-peak elements (Cr, Co, Ni) observed recently below \[Fe/H\]$``$-3.
We have compiled a large number of observational data on the composition of halo stars. The main conclusions of the comparison of our results to those data (Sec. 5 and Figs. 7 and 8) are the following:
\- C and N require other sources than those studied here. For C, it could be WR or low mass stars, contributing to C production in the disk. For N, the source of primary N required in the halo could be either IMS with hot-bottom burning or rotationally induced mixing in massive stars.
\- The evolution of the $`\alpha `$-elements O, Si, S and Ca is well understood (baring the discrepant “new data” for O, see below) with the assumption that SNIa contribute most of Fe in the disk; however, the WW1995 yields underproduce Mg and Ti, and inspection of the LSC2000 yields shows that they would not be of help.
\- Similarly, the odd-Z elements Sc and V are underproduced at all metallicities by both WW1995 and LSC2000 yields; this discrepancy points to some important revision required in current models of nucleosynthesis in massive stars, at least for those elements. It is significant that observationally, neither Sc nor V show the theoretically expected behaviour of odd-Z elements, suggesting that the “odd-even” effect may be overestimated in current nucleosynthesis models.
\- Observed abundances of Na and Al also do not show the theoretically expected behaviour of odd-Z elements, when they are plotted w.r.t Fe (Fig. 7). However, other sources may be involved in the nucleosynthesis of those two elements (e.g. H-shell burning in intermediate mass stars in the red giant stage), which prevents from drawing definite conclusions. It is remarkable that, when the observed Na evolution is plotted vs. Ca (Fig. 8), Na shows indeed the expected behaviour of odd-Z element. Observations of Na vs Fe at low metallicities are necessary to establish the behaviour of this element. In the case of Al, NLTE effects play an important role in estimating its abundance at low metallicities and render difficult a meaningful comparison of observations to theory.
\- Among the Fe-peak elements, several important discrepancies between theory and observations are found when results are plotted w.r.t. Fe (Fig. 7). The theoretical trends of Cr, Co, Ni and Zn deviate from the observed ones to various extents; in the case of Ni, the adopted W7 model for SNIa largely overproduces the main isotope <sup>58</sup>Ni in the disk, as well as<sup>54</sup>Cr, a minor Cr isotope. We notice that, when results are plotted w.r.t. Ca (Fig. 8), the observed behaviour of Cr is well reproduced by the model; this might imply that it is the Fe yields that are problematic at low metallicities. We notice that Cr is produced at layers lying at larger distance from the core than Fe, and are thus less subject to the uncertainties of the mass-cut.
\- There is a remarkably good agreement between the theoretical and the observed behaviour of the odd-Z Fe-peak elements Mn and Cu, when their evolution is plotted w.r.t. Fe (or w.r.t. Ca, in the case of Mn).
The recent data of Israelian et al. (1998) and Boesgaard et al. (1999) suggest that oxygen behaves differently than the other $`\alpha `$-elements. Although this new picture of O vs Fe is not confirmed yet, we explored in this work a few alternatives to the “standard” scenario presented here. We thus showed in Sec. 6 (and Fig. 9), albeit qualitatively only, that the only “reasonable” way to accomodate the new data is by assuming that the yields of both Fe and all $`\alpha `$-elements (except O, C and He) decrease with decreasing metallicity for \[Fe/H\]$`<`$-1; we also proposed a qualitative explanation for such a behaviour.
Finally, we compared the model evolution of the Mg isotopic ratios to current observations (Sec. 7 and Fig. 10). We found that, although the WW1995 yields of Mg describe relatively well the observations in the disk, they systematically underproduce the halo data. This suggests that the “odd-even” effect for those isotopes has been underestimated at low metallicities in WW1995.
In summary, we have revisited the chemical evolution of the halo and the local disk with consistent models and metallicity dependant yields of massive stars and SNIa. We showed that current yields are remarkably successful in reproducing a large number of observations, but need revision in several important cases. For some of those cases, the inclusion of non-classical ingredients in stellar models (i.e. mass-loss for C, rotationally induced mixing for primary N) could clearly help, but for most of the others (Sc, V and Ti at all metallicities, Fe-peak elements at very low metallicities) the situation remains unclear. Finally, we explored a few alternatives that could help to explain the new O vs Fe data and concluded that viable solutions exist, but would require some important modifications of our current understanding of massive star nucleosynthesis.
Acknowledgements. Aruna Goswami acknowledges the hospitality of IAP (Paris, France) where part of the work was being carried out. We are grateful to M. Limongi, T. Beers, A. McWilliam, Y. Chen and E. Carretta for kindly providing us their data in electronic form. This work is supported by CSIR/CNRS bi-lateral co-operation programme No. 19(207)/CNRS/98-ISTAD.
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# 1 Introduction
## 1 Introduction
Vacuum fluctuations are an essential ingredient of any quantum field theory, and also in quantum gravity they play an important role. The presence in the gravitational action of a dimensional coupling of the order of $`10^{33}cm`$ – the “Planck length” – indicates that the strongest fluctuations occur at very small scale: this is the famous “spacetime foam”, first studied by Wheeler and then by Hawking and Coleman through functional integral techniques .
More recently, Ashtekar and others analysed the possible occurrence of large fluctuations in 2+1 gravity coupled to matter. In this case the theory is classically solvable and admits a standard Fock-space quantization. In 3+1 dimensions, however, Einstein quantum gravity is a notoriously intractable theory, where everything (states, transition amplitudes, time…) is highly non-trivial.
The non-renormalizable UV divergences of the perturbative expansion may indicate that quantum gravity is not a fundamental microscopic theory, but an effective low-energy limit , and will be eventually replaced by a theory of strings or branes. On the other hand, it is known from particle physics that the Einstein lagrangian can be obtained, without any geometrodynamical assumption, as the only one which correctly accounts for a gravitational force mediated by helicity-2 particles . For this reason, it is important to investigate – besides the standard perturbative expansion – all the basic properties of the Einstein lagrangian. In the past years we took an interest into Wilson loops , vacuum correlations at geodesic distance , and the expression of the static potential through correlations between particles worldlines .
In this work we study a set of gravitational field configurations, called “dipolar zero modes”, which were not considered earlier in the literature. They give an exactly null contribution to the Einstein action, being thus candidates to become large fluctuations in the quantized theory. We give an explicit expression, to leading order in $`G`$, for some of the field configurations of this (actually quite large) set. We also give an estimate of possible suppression effects following the addition to the pure Einstein action of cosmological or $`R^2`$ terms.
Our zero modes have two peculiar features, which make them relatively easy to compute: (i) they are formally solutions of the Einstein equations with auxiliary virtual sources; (ii) their typical length scale is such that they can be treated in the weak field approximation. We shall see that these fluctuations can be large even on a “macroscopic” scale. There are some, for instance, which last $`1s`$ or more and correspond to the field generated by a virtual source with size $`1cm`$ and mass $`10^6g`$. This seems paradoxical, for several reasons, both theoretical and phenomenological. We have therefore been looking for possible suppression processes. Our conclusion is that a vacuum energy term $`(\mathrm{\Lambda }/8\pi G)d^4x\sqrt{g(x)}`$ in the action could do the job, provided it was scale-dependent and larger, at laboratory scale, than its observed cosmological value. This is at present only a speculative hypothesis, however.
The dipolar fluctuations owe their existence to the fact that the pure Einstein lagrangian $`(1/8\pi G)\sqrt{g(x)}R(x)`$ has indefinite sign also for static fields. It is well known that the non-positivity of the Einstein action makes an Euclidean formulation of quantum gravity difficult; in that context, however, the “dangerous” field configurations have small-scale variations and could be eliminated, for instance, by some UV cut-off. This is not the case of the dipolar zero modes. They exist at any scale and do not make the Euclidean action unbounded from below, but have instead null (or $`\mathrm{}`$) action.
A static virtual source will generate a zero mode provided it satisfies the condition $`d^3xT_{00}(𝐱)=0`$ up to terms of order $`G^2`$. The cancellation of the terms of order $`G`$ (Section 2.2) is important from the practical point of view. In our earlier work on dipolar fluctuations we developed some general remarks based on the form of Einstein equations, and the result was that in order to generate a zero mode the positive and negative masses of the source should differ by a quantity of order $`G`$, namely $`Gm^2/rmr_{Schw.}/r`$; this is very small for weak fields, but sufficient to produce a “monopolar” component which complicates the situation. Explicit calculations in Feynman gauge now have shown that the terms of order $`G`$ cancel out exactly. This opens the way to an amusing “virtual source engineering” work, to find explicitly some zero modes and give quantitative estimates in specific cases.
When analysing the Wilson loops, we had already pointed out some differences in the behavior of gravity and ordinary gauge theories, essentially due to the different signs of the allowed physical sources. Here, again, these differences are apparent. In gauge theories the real sources can be both positive and negative; therefore one can close two Wilson lines at infinity and find the static potential. The virtual sources cannot give rise to strong static dipolar fluctuations, because the lagrangian is quadratic in the fields. On the contrary, in gravity there are no real negative sources, the potential is always attractive and Wilson lines cannot be closed; however, since the lagrangian on-shell is indefinite in sign and equal to $`\sqrt{g(x)}\mathrm{Tr}T(x)`$, we can construct static zero modes employing +/- virtual sources. Then, of course, we can Lorentz-boost these modes in all possible ways.
The paper is composed of two main Sections. Section 2 is devoted to the analysis of the dipolar fields and virtual sources. We start from some general features and then focus on two examples. Section 3 contains an extensive discussion. For a summary of the main contents see also the Conclusions Section.
### 1.1 Conventions. Sign of $`\mathrm{\Lambda }`$ vs. its classical effects
Let us define here our conventions. We consider a gravitational field in the standard metric formalism; the action includes possibly a cosmological term:
$`S`$ $`=`$ $`S_{Einstein}+S_\mathrm{\Lambda }`$ (1)
$`S_{Einstein}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle d^4x\sqrt{g(x)}R(x)}`$ (2)
$`S_\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{8\pi G}}{\displaystyle d^4x\sqrt{g(x)}}`$ (3)
with $`g_{\mu \nu }(x)=\eta _{\mu \nu }+h_{\mu \nu }(x)`$.
By varying this action with respect to $`\delta g_{\mu \nu }(x)`$ and using the relation
$$\frac{\delta \sqrt{g}}{\delta g_{\mu \nu }}=\frac{1}{2}\sqrt{g}g^{\mu \nu }$$
(4)
one finds the field equations
$$R_{\mu \nu }(x)\frac{1}{2}g_{\mu \nu }(x)R(x)+\mathrm{\Lambda }g_{\mu \nu }(x)=8\pi GT_{\mu \nu }(x)$$
(5)
The energy-momentum tensor of a perfect fluid has the form
$$[T_{\mu \nu }]=\mathrm{diag}(\rho ,p,p,p)$$
(6)
For a zero-pressure “dust” one has $`p=0`$.
Now let us introduce a signature for the metric. Articles in General Relativity or cosmology use most often the metric with signature $`(,+,+,+)`$, and the experimental estimates of $`\mathrm{\Lambda }`$ are mainly referred to this metric. It is important to fix the sign of the cosmological term with reference to the metric signature in a way which is clear both formally and intuitively.
If spacetime is nearly flat, we can take the cosmological term in (5) to the r.h.s., set $`g_{\mu \nu }(x)=\eta _{\mu \nu }`$ and regard it as a part of the source. We obtain, in matrix form
$$\left[R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R\right]=\left\{\mathrm{diag}(\mathrm{\Lambda },\mathrm{\Lambda },\mathrm{\Lambda },\mathrm{\Lambda })+8\pi G\mathrm{diag}(\rho ,p,p,p)\right\}[\mathrm{metric}(,+,+,+)]$$
(7)
Which sign for $`\mathrm{\Lambda }`$ allows to obtain a static solution? Even without finding explicitly this solution, we see that for $`\mathrm{\Lambda }>0`$ the “pressure” due to the cosmological term is positive and can sustain the system against gravitational collapse – especially in the case of a zero-pressure dust with $`p=0`$. At the same time, the mass-energy density due to the cosmological term is negative and subtracts from $`\rho `$, still opposing to the collapse.
In conclusion, with this convention on the metric signature a static solution of Einstein equations with a cosmological term can be obtained for $`\mathrm{\Lambda }>0`$. If we are not interested into a static solution, but into an expanding space à la Friedman-Walker, in that case the effect of a cosmological term with $`\mathrm{\Lambda }>0`$ will be that of accelerating the expansion. The most recent measurements of the Hubble constant from Type Ia supernovae suggest indeed that there is a cosmologically significant positive $`\mathrm{\Lambda }`$ in our universe.
In Quantum Field Theory, on the other hand, the signature $`(+,,,)`$ is more popular, such that the squared four-interval is $`x^2=t^2|𝐱|^\mathrm{𝟐}`$. Since we shall introduce some coupling of gravity to matter fields in the following, and make a correspondence to the Euclidean case, we prefer to use this latter convention. We then have, instead of eq. (7)
$$\left[R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R\right]=\left\{\mathrm{diag}(\mathrm{\Lambda },\mathrm{\Lambda },\mathrm{\Lambda },\mathrm{\Lambda })+8\pi G\mathrm{diag}(\rho ,p,p,p)\right\}[\mathrm{metric}(+,,,)]$$
(8)
In this case, a static solution – or an accelerated expansion – corresponds to $`\mathrm{\Lambda }<0`$.
## 2 The dipolar fluctuations
We consider the functional integral of pure quantum gravity, which represents a sum over all possible field configurations weighed with the factor $`\mathrm{exp}[i\mathrm{}S_{Einstein}]`$ and possibly with a factor due to the integration measure. The Minkowski space is a stationary point of the vacuum action and has maximum probability. “Off-shell” configurations, which are not solutions of the vacuum Einstein equations, are admitted in the functional integration but are strongly suppressed by the oscillations of the exponential factor.
Due to the appearance of the dimensional constant $`G`$ in the Einstein action, the most probable quantum fluctuations of the gravitational field “grow” at very short distances, of the order of $`L_{Planck}=\sqrt{G\mathrm{}/c^3}10^{33}cm`$. This led Hawking, Coleman and others to depict spacetime at the Planck scale as a “quantum foam” , with high curvature and variable topology. For a simple estimate (disregarding of course the possibility of topology changes, virtual black holes nucleation etc.), suppose we start with a flat configuration, and then a curvature fluctuation appears in a region of size $`d`$. How much can the fluctuation grow before it is suppressed by the oscillating factor $`\mathrm{exp}[iS]`$? (We set $`\mathrm{}=1`$ and $`c=1`$ in the following.) A naive dimensional estimate suggests that $`|R|`$ should not exceed $`G/d^4`$, but in fact only a non-perturbative calculation can give reliable results in the short-distance regime. The most accurate estimates of the critical exponents in lattice quantum gravity are those obtained by Hamber through the Euclidean Regge calculus , and show that the correct behavior in four dimensions is
$$|R|\frac{1}{L_{Planck}d}$$
(9)
This is a consequence of the fact that the critical exponent $`\nu `$, related to the derivative of the gravitational $`\beta `$-function in the vicinity of the UV fixed point, is very close to 1/3.
### 2.1 General features
There is another way, however, to obtain vacuum field configurations with action smaller than 1 in natural units. This is due to the fact that the Einstein action has indefinite sign. Consider the Einstein equations with a source $`T_{\mu \nu }(x)`$
$$R_{\mu \nu }(x)\frac{1}{2}g_{\mu \nu }(x)R(x)=8\pi GT_{\mu \nu }(x)$$
(10)
and their covariant trace
$$R(x)=8\pi G\mathrm{Tr}T(x)=8\pi Gg^{\mu \nu }(x)T_{\mu \nu }(x)$$
(11)
Let us consider a solution $`g_{\mu \nu }(x)`$ of equation (10) with a source $`T_{\mu \nu }(x)`$ obeying the additional integral condition
$$d^4x\sqrt{g(x)}\mathrm{Tr}T(x)=0$$
(12)
Taking into account eq. (11) we see that the pure Einstein action (2) computed for this solution is zero. Thus the tensor $`T_{\mu \nu }(x)`$ only serves as an auxiliary source in order to construct zero-modes for the action of pure gravity. Condition (12) can be satisfied by energy-momentum tensors that are not identically zero, provided they have a balance of negative and positive signs, such that their total integral is zero. Of course, they do not represent any acceptable physical source, but the corresponding solutions of (10) exist nonetheless, and are zero modes of the pure Einstein action.
We shall give two explicit examples of auxiliary sources (we shall call them “virtual sources” in the following, because they generate virtual field configurations): (i) a “mass dipole” consisting of two separated mass distributions with different signs; (ii) two concentric “+/- shells”. In both cases there are some parameters of the source which can be varied: the total positive and negative masses $`m_\pm `$, their distance, the spatial extension of the sources.
The procedure for the construction of the zero mode corresponding to the dipole is the following. One first considers Einstein equations with the virtual source without fixing the parameters yet. Then one solves them with a suitable method, for instance in the weak field approximation when appropriate. Finally, knowing $`g_{\mu \nu }(x)`$ one adjusts the parameters in such a way that condition (12) is satisfied.
### 2.2 Computation of $`\sqrt{g(x)}g^{00}(x)`$
Now suppose we have a suitable virtual source, with some free parameters, and we want to adjust them in such a way to generate a zero-mode $`g_{\mu \nu }(x)`$ for which $`S_{Einstein}[g]=0`$. We shall always consider static sources where only the component $`T_{00}`$ is non vanishing. The action of their field is
$$S_{zeromode}=\frac{1}{2}d^4x\sqrt{g(x)}g^{00}(x)T_{00}(x)$$
(13)
To first order in $`G`$, the field $`h_{\mu \nu }(x)`$ generated by a given mass-energy distribution $`T_{\mu \nu }(x)`$ is given by an integral of the field propagator $`P_{\mu \nu \rho \sigma }(x,y)`$ over the source:
$$h_{\mu \nu }(x)=d^4yP_{\mu \nu \rho \sigma }(x,y)T^{\rho \sigma }(y)$$
(14)
where in Feynman gauge $`P_{\mu \nu \rho \sigma }(x,y)`$ is given, with our conventions on the metric signature, by
$$P_{\mu \nu \rho \sigma }(x,y)=\frac{2G}{\pi }\frac{\eta _{\mu \rho }\eta _{\nu \sigma }+\eta _{\mu \sigma }\eta _{\nu \rho }\eta _{\mu \nu }\eta _{\rho \sigma }}{(xy)^2+i\epsilon }$$
(15)
Computing the integral over time in eq. (14) we obtain for our source
$`h_{\mu \nu }(𝐱)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑y_0{\displaystyle d^3yT^{00}(𝐲)P_{\mu \nu 00}(x,y)}`$ (16)
$`=`$ $`{\displaystyle \frac{2G}{\pi }}(2\eta _{\mu 0}\eta _{\nu 0}\eta _{\mu \nu }\eta _{00}){\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑y_0{\displaystyle d^3y\frac{T^{00}(𝐲)}{(x_0y_0)^2(𝐱𝐲)^2+i\epsilon }}`$
$`=`$ $`2G(2\eta _{\mu 0}\eta _{\nu 0}\eta _{\mu \nu }\eta _{00}){\displaystyle d^3y\frac{T^{00}(𝐲)}{|𝐱𝐲|}}`$
Thus we have
$`\sqrt{g(x)}g^{00}(x)`$ $`=`$ $`\left[1+{\displaystyle \frac{1}{2}}\mathrm{Tr}h(𝐱)+o(G^2)\right]\left[1+h^{00}(𝐱)+o(G^2)\right]`$ (17)
$`=`$ $`1+{\displaystyle \frac{1}{2}}\mathrm{Tr}h(𝐱)+h^{00}(𝐱)+o(G^2)`$
$`=`$ $`1+2G\left[{\displaystyle \frac{1}{2}}(2\eta _{\mu 0}\eta _{\nu 0}\eta _{\mu \nu }\eta _{00})\eta ^{\mu \nu }+(\eta _{00})^2\right]{\displaystyle d^3y\frac{T^{00}(𝐲)}{|𝐱𝐲|}}+o(G^2)`$
$`=`$ $`1+o(G^2)`$
and finally the action is
$$S_{zeromode}=\frac{1}{2}d^4xT_{00}(x)+o(G^2)$$
(18)
Therefore provided the integral of the mass-energy density vanishes, the action of our field configuration is of order $`G^2`$, i.e., practically negligible, as we check now with a numerical example. Let us choose the typical parameters of the source as follows:
$`r`$ $``$ $`1cm`$
$`m`$ $``$ $`10^kg10^{37+k}cm^1`$ (19)
(implying $`r_{Schw.}/r10^{29+k}`$). We assume in general an adiabatic switch-on/off of the source, thus the time integral contributes to the action a factor $`\tau `$. We shall keep $`\tau `$ (in natural units) very large, in order to preserve the static character of the field. Here, for instance, let us take $`\tau 1s310^{10}cm`$. With these parameters we have
$$S_{zeromode}^{orderG^2}\tau \frac{G^2m_\pm ^2}{r^3}10^{20+3k}$$
(20)
Thus the field generated by a virtual source with typical size (19), satisfying the condition $`d^3xT^{00}(𝐱)=0`$, has negligible action even with $`k=6`$ (corresponding to apparent matter fluctuations with a density of $`10^6g/cm^3`$ !) This should be compared to the huge action of the field of a single, unbalanced virtual mass $`m`$; with the same values we have
$$S_{singlem}=\frac{1}{16\pi G}d^4x\sqrt{g(x)}R(x)=\frac{1}{2}d^4x\sqrt{g(x)}\mathrm{Tr}T(x)\frac{1}{2}\tau m+o(G^2)10^{47+k}$$
(21)
This example shows that the cancellation of the first order term in (17) allows to obtain a simple lower bound on the strength of the fluctuations. In principle, however, one could always find all the terms in the classical weak field expansion, proportional to $`G`$, $`G^2`$, $`G^3`$, etc., and adjust $`T_{00}`$ as to have $`S_{zeromode}=0`$ exactly. They can be represented by those Feynman diagrams of perturbative quantum gravity which contain vertices with 3, 4 … gravitons but do not contain any loops. The ratio between each contribution to $`S`$ and that of lower order in $`G`$ has typical magnitude $`r_{Schw.}/r`$, where $`r_{Schw.}=2\pi Gm_\pm `$ is the Schwarzschild radius corresponding to one of the two masses and $`r`$ is the typical size of the source. For a wide range of parameters, this ratio is very small, so the expansion converges quickly. From now on we agree that the “$`o(G^2)`$” term in eq. (18) comprises all the terms quadratic in the field, like for instance that arising from the expansion of $`\sqrt{g(x)}`$.
### 2.3 Explicit examples of static virtual sources
(i) The mass dipole
As an example of unphysical source which satisfies (12) one can consider the static field produced by a “mass dipole”. Certainly negative masses do not exist in nature; here we are interested just in the formal solution of (10) with a suitable $`T_{\mu \nu }`$, because for this solution we have $`d^4x\sqrt{g}R=0`$. Let us take the following $`T_{\mu \nu }`$ of a static dipole centered at the origin ($`m_+,m_{}>0`$):
$$T_{\mu \nu }(𝐱)=\delta _{\mu 0}\delta _{\nu 0}\left[\frac{m_+}{r_+^3}f_+(𝐱)\frac{m_{}}{r_{}^3}f_{}(𝐱)\right]$$
(22)
where
$$f_\pm (𝐱)f\left(\frac{𝐱\pm 𝐚}{r_\pm }\right)$$
(23)
and $`f(𝐱)`$ is a smooth test function with range $`1`$ and normalized to 1, which represents the mass density. Thus we have a positive source of mass $`m_+`$ and radius $`r_+`$ (placed at $`𝐱=𝐚`$) and a negative source with mass $`m_{}`$ and radius $`r_{}`$ (placed at $`𝐱=𝐚`$). The radii of the two sources are such that $`ar_\pm r_{Schw.}`$, where $`r_{Schw.}`$ is the Schwarzschild radius corresponding to the mass $`m_+`$.
The mass $`m_{}`$ is in general slightly different from $`m_+`$ and chosen in such a way to compensate the small difference, due to the $`\sqrt{g}g^{00}`$ factor, between the integrals
$$I_+=d^4x\sqrt{g(x)}g^{00}(x)\frac{f_+(𝐱)}{r_+^3}\mathrm{and}I_{}=d^4x\sqrt{g(x)}g^{00}(x)\frac{f_{}(𝐱)}{r_{}^3}$$
(24)
The action of the dipole is
$$S_{Dipole}=\frac{1}{2}d^4xT_{00}(𝐱)=\frac{1}{2}\tau (m_+m_{})+o(G^2)$$
(25)
The condition for $`S_{Dipole}=0`$ is $`m_+=m_{}`$, apart from terms of order $`G^2`$ (i.e., our dipoles have in reality a tiny monopolar component).
Also note that the values of the masses and the radii $`r_\pm `$ (both of order $`r`$) can vary in a continuous way – under the only condition that $`m_+=m_{}`$. This implies that these (non singular) “dipolar” fields constitute a subset with nonzero volume in the functional integration. Actually, they are only a subset of the larger class of solutions of the Einstein equations with sources satisfying eq. (12).
(ii) The concentric +/- shells
Consider two concentric spherical shells in contact, the internal one with radii $`r_1`$, $`r_2`$, and the external one with radii $`r_2`$, $`r_3`$ ($`r_1<r_2<r_3`$). Let the internal shell have mass density $`\rho _1`$ and the external shell density $`\rho _2`$, with opposite sign. The condition for zero action requires, up to terms of order $`G^2`$, that the total positive mass equals the total negative mass, i.e.,
$$\rho _1(r_2^3r_1^3)+\rho _2(r_3^3r_2^3)=0$$
(26)
(more generally, if the densities $`\rho _1`$ and $`\rho _2`$ are not constant throughout the shells, one has a suitable integral condition).
The spherical symmetry of the corresponding field configuration offers some advantages when one computes the contributions to the cosmological term and the Newtonian self energy (compare Sect. 3.1).
### 2.4 Contribution of virtual dipoles to the cosmological and $`R^2`$ terms
In the previous Sections we have seen that the pure Einstein action admits zero-modes having the form of virtual dipole field configurations with a small monopole residual. These field configurations are characterized by the parameters $`r_\pm `$ (radii of the virtual +/- sources), $`a`$ (distance between the sources) and $`m_\pm `$ (masses of the sources). We worked out these configurations as solutions of the linearized Einstein equations. We also checked that the weak field approximation is appropriate in a whole “macroscopic” range of the parameters $`r_\pm `$, $`a`$ and $`m`$. This is possible because these configurations (unlike the spacetime foam at the Planck scale) yield $`d^4x\sqrt{g(x)}R(x)=0`$ thanks to a cancellation between the $`R`$ contributions in two distinct regions of space. Similar considerations can be done for the field of the concentric +/- shells.
It is natural to ask whether the dipolar fluctuations can be suppressed by other terms present in the gravitational action besides the pure Einstein term. Possible candidates are the $`R^2`$ terms (usually relevant, however, only at very small distance) and the cosmological term. Let us first look at the latter (see also our general remarks on the role of a cosmological constant in quantum gravity in Section 3.2).
When a static source is spherically symmetric, we can use outside it the exact Schwarzschild metric with invariant interval
$$ds^2=\left(1\frac{2GM}{r}\right)^1dt^2\left[\left(1\frac{2GM}{r}\right)dr^2+r^2d\theta ^2+r^2\mathrm{sin}^2\theta d\varphi ^2\right]$$
(27)
The determinant of this metric equals that of flat space, so the presence of one single spherically-symmetric source does not change the volume of the outer space and does not contribute to the cosmological term.
We shall therefore handle separately the cases of the mass dipole and the concentric +/- shells.
The mass dipole
In the linearized approximation the integral $`S_\mathrm{\Lambda }=(\mathrm{\Lambda }/8\pi G)d^4x\sqrt{g(x)}`$ for a dipolar fluctuation can be splitted into the sum of the integrals of the field $`h_+(𝐱)`$ generated by the positive mass and the field $`h_{}(𝐱)`$ generated by the negative mass. Both fields are spherically symmetric, thus there is no contribution of order $`G`$ to $`S_\mathrm{\Lambda }`$ outside the sources.
To order $`h^2(Gm)^2`$ the field outside the sources differs from the sum of their Schwarzschild fields, and we do have some contributions to the cosmological term, but they are very small. One finds, inserting the numerical values (19) and the current estimate for $`|\mathrm{\Lambda }|G`$, namely $`|\mathrm{\Lambda }|G10^{116}`$
$$\mathrm{\Delta }S_{\mathrm{\Lambda },outside}\tau ^2|\mathrm{\Lambda }|Gm^210^{22+2k}$$
(28)
On the other hand, the integrals of $`\sqrt{g(x)}`$ inside the sources contribute to the action already at first order in $`h_{\mu \nu }`$. Let us use the explicit solutions in Feynman gauge found in the previous section and disregard the effect of the positive source inside the negative one and viceversa. (This will give small corrections proportional to $`a/r_\pm `$, but does not change the magnitude orders.) We denote by $`\omega (𝐱)`$ the characteristic function of a 3-sphere with unit radius placed at the origin of the coordinates, and define
$$\omega _\pm (𝐱)\omega \left(\frac{𝐱\pm 𝐚}{r_\pm }\right)$$
(29)
We then have, to leading order
$`\mathrm{\Delta }S_{\mathrm{\Lambda },inside}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{8\pi G}}\left[{\displaystyle d^4x\frac{1}{2}\mathrm{Tr}h_+(𝐱)\omega _+(𝐱)}+{\displaystyle d^4x\frac{1}{2}\mathrm{Tr}h_{}(𝐱)\omega _{}(𝐱)}\right]`$ (30)
$`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{8\pi G}}{\displaystyle \frac{\tau }{2}}{\displaystyle d^3x\omega _+(𝐱)\left(\frac{4m_+G}{r_+^3}\right)d^3y\frac{f_+(𝐲)}{|𝐱𝐲|}}+`$
$`+{\displaystyle \frac{\mathrm{\Lambda }}{8\pi G}}{\displaystyle \frac{\tau }{2}}{\displaystyle d^3x\omega _+(𝐱)\left(\frac{4m_{}G}{r_{}^3}\right)d^3y\frac{f_{}(𝐲)}{|𝐱𝐲|}}`$
In the double integrals we can suitably shift the variables by $`\pm 𝐚`$ and re-scale them as $`𝐱𝐱^{}r_\pm `$, $`𝐲𝐲^{}r_\pm `$, obtaining a pure number $`\xi `$ of order 1 multiplied by $`r_+^5`$ and $`r_{}^5`$, respectively. Finally we obtain
$$\mathrm{\Delta }S_{\mathrm{\Lambda },inside}=\frac{\xi }{4\pi }\tau \mathrm{\Lambda }(m_+r_+^2m_{}r_{}^2)$$
(31)
with
$$\xi =d^3x^{}d^3y^{}\frac{\omega (𝐱^{})f(𝐲^{})}{|𝐱^{}𝐲^{}|}$$
(32)
With the usual values we find, apart from an adimensional constant of order 1
$$\mathrm{\Delta }S_{\mathrm{\Lambda },inside}10^{3+k}$$
(33)
This means that a relatively small increase in the value of $`|\mathrm{\Lambda }|`$ would be sufficient to suppress the strongest fluctuations (except for those with $`r_+=r_{}`$ exactly).
The concentric +/- shells
In this case $`\mathrm{\Delta }S_{\mathrm{\Lambda },outside}`$ vanishes exactly. Inside the source we have to leading order
$$\mathrm{\Delta }S_{\mathrm{\Lambda },inside}=\frac{\mathrm{\Lambda }\tau }{8\pi G}\frac{1}{2}d^3x\mathrm{Tr}h(𝐱)$$
(34)
Since $`\mathrm{Tr}h(𝐱)=4V_{Newt.}(𝐱)`$ (compare eq. (16)), the integral is a special case of one we shall compute in Section 3.1 The result is
$$\mathrm{\Delta }S_{\mathrm{\Lambda },inside}=\mathrm{\Lambda }\tau mr^2Q(\beta )$$
(35)
where $`r_2r`$, $`r_3\beta r`$ and $`Q(\beta )`$ is an adimensional polynomial which can be either positive or negative, depending on the ratio $`|\rho _1|/|\rho _2|`$. The magnitude order is the same as for the dipole.
Finally, a word about the $`R^2`$ term. It is typically of the form
$$S_{R^2}=\alpha d^4x\sqrt{g(x)}R^2(x)$$
(36)
where $`\alpha `$ is a (small) adimensional coupling and $`R^2`$ can be replaced by more complex scalars like $`R_{\mu \nu \rho \sigma }R^{\mu \nu \rho \sigma }`$ etc. For an order of magnitude estimate it suffices to multiply the square of the curvature in the sources, namely $`R^2(G\mathrm{Tr}T)^2=(Gm/r^3)^2`$ by their volume $`V^{(4)}\tau r^3`$. We find in this way, still with the same parameters,
$$S_{R^2}\alpha \tau G^2\frac{m^2}{r^3}\alpha 10^{48+2k}$$
(37)
Thus the allowed values for $`m`$ are very large, i.e., there is no significant suppression of the virtual dipoles by the $`R^2`$ terms at this scale.
## 3 Discussion
### 3.1 Why are these fluctuations paradoxical
The order of magnitude estimates given in the previous Section show that the dipolar vacuum fluctuations allowed in the functional integral formulation of pure Einstein quantum gravity (i.e., such to give $`S1`$ in natural units) are very intense also at macroscopic scale.
One may think that such large fluctuations, if real, would not remain unnoticed. Even though vacuum fluctuations are homogeneous, isotropic and Lorentz-invariant, they could manifest themselves as noise of some kind. Most authors are skeptic about the possibility of detecting the noise due to spacetime foam , but the virtual dipole fluctuations described in this paper are much closer to the laboratory scale. Observable quantities, like for instance the invariant intervals $`ds^2=g_{\mu \nu }dx^\mu dx^\nu `$ and the connection coefficients $`\mathrm{\Gamma }_{\mu \nu }^\rho `$ could then exhibit strong fluctuations.
The existence of these fluctuations would be paradoxical, however, already at the purely conceptual level. Common wisdom in particle physics states that the vacuum fluctuations in free space correspond to virtual particles or intermediate states which live very short, i.e., whose lifetime is close to the minimum allowed by the Heisenberg indetermination relation. Let us first give a brief formal justification of this rule, and then compare it to our dipole fluctuations.
It is often the case that a quantum field theory has an imaginary time formulation, where the (positive-definite) lagrangian density corresponds to the original hamiltonian density $`H`$. For a scalar field, for instance, one has $`H=(1/2)[(\varphi )^2+(\mathrm{grad}\varphi )^2+m^2\varphi ^2]`$ and the Euclidean functional integral is given by $`z_{Eucl}=d[\varphi ]\mathrm{exp}[𝑑td^3xH(t,𝐱)]`$. A field fluctuation localized to a region of size $`\tau V^{(3)}`$ is weighed in the functional integral by the factor $`\mathrm{exp}[\tau V^{(3)}H]=\mathrm{exp}[\tau E]`$ and is thus effectively suppressed unless approx. $`\tau E<1`$. Another notable example is the electromagnetic field. Also in this case the analytical continuation of the lagrangian $`L=(1/8\pi )[𝐄^2𝐁^2]`$ yields the energy density $`H=(1/8\pi )[𝐄^2+𝐁^2]`$; to check this, one just needs to impose the $`A_0=0`$ gauge and remember that only the electric field contains time derivatives of $`𝐀`$.
Now let us estimate the product $`E\tau `$ for the dipolar fluctuations. The total energy of a static gravitational field configuration vanishing at infinity is the ADM energy. Since the source of a dipolar fluctuation satisfies the condition $`d^3xT_{00}(𝐱)=0`$ up to terms of order $`G^2`$, the dominant contribution in the ADM energy is the Newtonian binding energy .
The binding energy of the field generated by a source of mass $`m`$ and size $`r`$ is of the order of $`EGm^2/r`$, where the exact proportionality factor depends on the details of the mass distribution. For a dipolar field configuration characterized by masses $`m_+`$ and $`m_{}`$ and radii of the sources $`r_+`$ and $`r_{}`$, the total gravitational energy is of the order of
$$E_{tot}Gm_\pm ^2\left(\frac{1}{r_{}}+\frac{1}{r_+}\right)$$
(38)
(disregarding the interaction energy between the two sources, proportional to $`1/a1/r`$). For an order of magnitude estimate with the parameters (19) we can suppose that $`r_+`$ and $`r_{}`$ are both of the order of $`1cm`$. We then have $`E_{tot}Gm_\pm ^210^{12+k}cm^1`$. Remembering that $`k`$ can take values up to $`k=6`$, we find for these dipolar fluctuations $`\tau E_{tot}10^{28}`$!
(For comparison, remember the case of a “monopole” fluctuation of virtual mass $`m`$ and duration $`\tau `$. The condition $`S<1`$ implies $`\tau m<1`$. The dominant contribution to the ADM energy is just $`m`$, thus the rule $`E\tau <1`$ is respected.)
The Newtonian binding energy of the concentric +/- shells is given, like in electrostatics, by the formula $`E=(1/2)d^3x\rho (𝐱)V_{Newt.}(𝐱)`$. For general values $`r_1`$, $`r_2`$, $`r_3`$ of the radii and $`\rho _1`$, $`\rho _2`$ of the densities (constrained by the zero total mass condition (26)), one obtains a complicated expression, namely
$`E`$ $`=`$ $`{\displaystyle \frac{\pi \rho _1}{90r_2^2}}\{\rho _2(r_2^2+r_2r_3+r_3^2)(6r_2^515r_2^4r_3+10r_2^3r_3^2r_3^5)\rho _1[9r_1^711r_1^6r_2r_1^5r_2^210r_1^4r_2^3+`$ (39)
$`+5r_1^3(r_2^4+2r_2^3r_3+2r_2r_3^32r_3^4)+r_1^2r_2^5+r_1r_2^6+2r_2^3(3r_2^45r_2^3r_35r_2r_3^3+5r_3^4)]\}`$
We can study the sign and magnitude of $`E`$ setting $`r_2=r`$, $`r_1=\alpha r`$ ($`0<\alpha <1`$) and $`r_3=\beta r`$ ($`\beta >1`$). We express $`\alpha `$ in terms of $`\beta `$ using (26) and finally obtain
$$E=\frac{Gm}{r}P(\beta )$$
(40)
where $`P(\beta )`$ is a polynomial which is positive if $`|\rho _1|>|\rho _2|`$ (the repulsion between the two shells predominates) and negative if $`|\rho _1|<|\rho _2|`$ (the attraction inside each shell predominates).
This result is quite interesting, because
(i) Unlike the formula for the energy of the dipolar field, it does not contain any approximation to order $`G`$.
(ii) From the physical point of view it is reasonable to admit – remembering that we are in a weak-field regime and forgetting general covariance for a minute – that the binding energy is localized within the surface of the outer shell (the field is $`o(G^2)`$ outside). The energy density is therefore of the order of $`\frac{|E|}{r^3}\frac{Gm}{r^4}10^{29+k}cm^4`$ (with the parameters (19)), and can take both signs. This value looks quite large, even though the Ford-Roman inequalities or similar bounds do not apply to quantum gravity, where the metric is not fixed but free to fluctuate, and there is in general no way to define a local energy density (except outside the sources – see ).
### 3.2 A scale-dependent $`\mathrm{\Lambda }`$ ?
We have seen that a vacuum energy or cosmological term in the gravitational action is able to cut-off part of the dipolar fluctuations. This works better at large scales, because the $`\mathrm{\Lambda }`$-term does not contain any field derivatives. We may also hypothesize that the effective value of $`\mathrm{\Lambda }`$ at scales of the order of $`1cm`$ is larger than the value observed at cosmological scale. In the following we summarize some theoretical arguments supporting this idea. One would have, in other words, a small, negative, scale-dependent $`\mathrm{\Lambda }_{eff}`$, a sort of residual of purely gravitational self-adjustment processes taking place at the Planck scale.
We already mentioned the role played by the cosmological constant at the classical level. In particular, looking for solutions of Einstein equations of the Friedman-Robertson-Walker type, i.e. with an expanding space, one finds well-defined relations between the Hubble constant, the density of various kinds of matter, and $`\mathrm{\Lambda }`$ . In the last years, most estimates have given a negative value $`\mathrm{\Lambda }`$ (in our conventions) of the order of $`10^{50}cm^2`$.
The effect of a cosmological term in the quantum field theory of gravity is less clear. On one hand, there are some “naive” expectations; on the other hand, formal results which are however difficult to interpretate.
The naive view consists in disregarding the effect of the cosmological term on the global geometry of spacetime, as compared to the effect of matter or pre-existing (null) curvature. Therefore one just expands the gravitational action around a flat background and studies quantum fluctuations. These are determined to leading order by the part of the action quadratic in $`h_{\mu \nu }`$. In spite of the different tensorial form of the Einstein term $`𝑑x\sqrt{g}R`$ and the cosmological term $`𝑑x\sqrt{g}`$, their quadratic parts are similar. In Feynman gauge they are both proportional to the quantity $`\left[2\mathrm{T}\mathrm{r}h^2(\mathrm{Tr}h)^2\right]`$, multiplied by $`^\mu _\mu `$ in the case of the Einstein term and by $`\mathrm{\Lambda }`$ in the case of the cosmological term. Thus in this approximation the cosmological term corresponds to a mass term for the graviton; the mass is real for $`\mathrm{\Lambda }<0`$ and imaginary for $`\mathrm{\Lambda }>0`$ (in our conventions – see Section 1.1). This implies respectively a finite range propagator, á la Yukawa, with range of the order of $`|\mathrm{\Lambda }|^{1/2}`$, or the existence of unstable modes growing in time like real exponentials . Intuitively, the reason for this behavior is clear (see also ), because a positive $`\mathrm{\Lambda }`$ corresponds to a positive mass-energy density, which is gravitationally unstable.
In pure quantum gravity the curvature of the classical background is solely determined by $`\mathrm{\Lambda }`$, and therefore the previous approach does not really make sense. For instance, if $`\mathrm{\Lambda }<0`$, then the solution of the classical Einstein equations is a spacetime with curvature radius of the order of $`\mathrm{\Lambda }^{1/2}`$; the Yukawa range predicted by the flat space expansion would then coincide with the size of the universe. There have thus been some attempts at quantizing the gravitational action with respect to a background with constant curvature (de Sitter or anti-de Sitter). The theory is mathematically very difficult ; there is some evidence, however, that the graviton stays massless, while novel strong infrared effects would arise (due to the dimensional self-coupling $`\mathrm{\Lambda }`$), which might force the renormalized value of $`\mathrm{\Lambda }`$ to “relax towards zero”.
The Euclidean theory of pure quantum gravity is obtained from the Lorentzian theory in our conventions with the standard analytical continuation $`t_{Lor}it_{Eucl}`$. In the lattice approach in 4D , $`G`$ and $`\mathrm{\Lambda }`$ are entered as bare couplings at the beginning, and then the discretized space evolves according to a Montecarlo algorithm. Unlike in perturbation theory, where a flat background is introduced by hand, here flat space appears dynamically; namely, the average value of the curvature is found to vanish on a transition line in the bare-couplings space. This line separates a “smooth-phase”, with small negative curvature, from a “rough”, collapsed, unphysical phase, with large positive curvature. The collapse can be understood observing that the cosmological action is of the form $`\mathrm{\Lambda }V^{(4)}`$, where $`V^{(4)}`$ is the volume of the lattice, thus when $`\mathrm{\Lambda }_{eff}=R`$ is positive, the volume tends to decrease.
It turns out that as the continuum limit is approached, the adimensional product $`|\mathrm{\Lambda }_{eff}|G_{eff}`$ behaves like
$$|\mathrm{\Lambda }_{eff}|G_{eff}(l_0/l)^\gamma $$
(41)
where $`l`$ is the scale, $`l_0`$ is the lattice spacing, $`\gamma `$ a critical exponent and the sign of $`\mathrm{\Lambda }_{eff}`$ is negative. Furthermore, one can reasonably assume that $`l_0L_{Planck}`$, and that the scale dependence of $`G_{eff}`$ is much weaker than that of $`\mathrm{\Lambda }_{eff}`$.
A scale dependence of $`\mathrm{\Lambda }_{eff}`$ like that in eq. (41) also implies that any bare value of $`\mathrm{\Lambda }`$, expressing the energy density associated to the vacuum fluctuations of the quantum fields including the gravitational field itself, approaches zero at long distances just by virtue of the gravitational dynamics, without any need of a fine tuning. One would have, in other words, a purely gravitational solution of the cosmological constant problem.
It is remarkable that the conclusions of Euclidean lattice theory concerning the instability with $`\mathrm{\Lambda }>0`$ agree qualitatively with those obtained in the naive approach; and this in spite of the fact that the above argument concernig the volume of spacetime does not hold in the Lorentzian theory because in this case both positive and negative volume variations are suppressed by the oscillating factor $`\mathrm{exp}[iS]`$ in the functional integral.
### 3.3 Local changes in $`\mathrm{\Lambda }`$
The ability of the $`\mathrm{\Lambda }`$-term to cut-off part of the dipole fluctuations has an inevitable consequence. Consider the coupling of gravity to a scalar field $`\varphi `$, with lagrangian density
$$L=\frac{1}{2}\left(_\alpha \varphi ^\alpha \varphi m^2\varphi ^2\right)=\frac{1}{2}\left[\left(\frac{\varphi }{t}\right)^2(\mathrm{grad}\varphi )^2m^2\varphi ^2\right]$$
(42)
and energy-momentum tensor
$$T_{\mu \nu }=\mathrm{\Pi }_\mu \varphi _\nu \varphi g_{\mu \nu }L=_\mu _\nu \varphi g_{\mu \nu }L$$
(43)
The interaction term in the gravitational action is
$$S_{matter}=\frac{1}{2}d^4x\sqrt{g(x)}T^{\mu \nu }(x)h_{\mu \nu }(x)$$
(44)
and to lowest order in $`h_{\mu \nu }`$ we have
$$S_{matter}=\frac{1}{2}d^4x\left(h_{\mu \nu }^\mu \varphi ^\nu \varphi \mathrm{Tr}hL\right)$$
(45)
On the other hand, the cosmological action is, still to lowest order in $`h_{\mu \nu }`$ and expanding $`\sqrt{g}=1+\frac{1}{2}\mathrm{Tr}h+\mathrm{}`$
$$S_\mathrm{\Lambda }=\frac{\mathrm{\Lambda }}{8\pi G}d^4x\left(1+\frac{1}{2}\mathrm{Tr}h\right)$$
(46)
We can say that to lowest order the coupling of gravity to the field $`\varphi `$ produces a typical source term for $`h_{\mu \nu }`$, of the form $`h_{\mu \nu }^\mu \varphi ^\nu \varphi `$, and subtracts from the cosmological constant $`\mathrm{\Lambda }`$ the local density $`8\pi GL(x)`$, because we can write, apart from an additive constant,
$$S_{matter}+S_\mathrm{\Lambda }=\frac{1}{2}d^4xh_{\mu \nu }^\mu \varphi ^\nu \varphi +\frac{1}{2}d^4x\mathrm{Tr}h\left(\frac{\mathrm{\Lambda }}{8\pi G}L\right)$$
(47)
This separation of the matter coupling in two parts looks in general quite arbitrary, but it can be useful if the lagrangian density is such to affect locally the “natural” cosmological term and set free gravitational fluctuations corresponding to virtual mass densities much larger than the real density of the field $`\varphi `$.
An example will clarify our point. Suppose that $`\varphi `$ represents some coherent fluid with the density of ordinary matter ($`1g/cm^3`$). We have seen that at the scale of 1 $`cm`$ the dipolar fluctuations are cut-off according to eq. (31). For $`\mathrm{\Lambda }`$ equal to the cosmologically observed value of $`10^{50}cm^2`$, the exponent $`k`$ can take values up to $`k=3`$, corresponding to fluctuations with virtual sources of density $`10^3g/cm^3`$. (This is a prudent estimate; for short-lasting fluctuations – less than $`1s`$ – and for those with $`r_{}=r_+`$, the virtual mass density can be even higher.)
If the value of $`L`$ in some region is comparable to $`\mathrm{\Lambda }/8\pi G`$, this can introduce an inhomogeneity in the cut-off mechanism. The result will be a local inhomogeneity of the dipolar fluctuations, which, given their strength, could dominate the effects of the coupling $`h_{\mu \nu }_\mu \varphi _\nu \varphi `$ to the real matter.
Note that the magnitude of $`L`$ depends on whether $`\varphi `$ is itself “on shell” or not. For a free, spatially homogeneous scalar field, for instance, the Klein-Gordon equation implies $`\varphi =const.e^{\pm imt}`$. Therefore on shell one has $`L=0`$ exactly, even though the single terms in $`L`$ can well be (for atomic-scale masses and gradients) of the order of $`10^{33}cm^4`$.
The mechanism sketched above also has an Euclidean analogue , but a better understanding of the dipolar fluctuations is necessary before any progress in this direction can be made.
## 4 Conclusions
In the first part of this work we have studied the general features of “dipolar” zero modes of the pure Einstein action, giving some explicit examples in the weak-field approximation. We used a method based upon the classical Einstein equations with suitable virtual sources. Our aim was to prove in a rigorous way the null-action property of these modes. For applications to the quantum case we made reference to the (Lorentzian) functional integral. This represents just one of the possible approaches to quantum gravity, but in fact also the Planck-scale fluctuations have been studied through integral functional techniques . It should be stressed that the numerical estimates presented in Section 2.2 are only lower limits based on specific examples. The strength of the fluctuations can be in general larger.
In the Discussion Section we have been less concerned with rigor. We have described some paradoxical features of the large dipole fluctuations, and possible suppression processes. The ADM energy of the dipolar fields can be both positive and negative, and turns out to be very large compared to their inverse duration $`\tau ^1`$. If we admit (as is quite reasonable for the +/- shells) that this energy is localized, the corresponding density appears large, too – even though the Ford-Roman inequalities or similar bounds do not apply to quantum gravity, where the metric is not fixed.
The hypothesis of a scale-dependent cosmological constant remains at present speculative, yet only the $`\mathrm{\Lambda }`$-term seems to be capable of suppression at large scales. From the purely phenomenological point of view, the existence of a negative (in our conventions) $`\mathrm{\Lambda }_{eff}`$, which reduces to the observed $`\mathrm{\Lambda }10^{50}cm^2`$ at cosmological scale but is some orders of magnitude larger at $`cm`$ scale, is probably less disturbing than the existence of large quantum fluctuations .
Independently from the effective-$`\mathrm{\Lambda }`$ hypothesis, the results of Sections 2.4 and 3.3 show that any local vacuum term of the form $`g_{\mu \nu }(x)L(x)`$ acts as a cutoff for the dipolar fluctuations, especially for those at large scale. This can cause local inhomogeneities, which are usually important when dealing with vacuum fluctuations in quantum field theory, and deserve further investigation.
Acknowledgments \- This work was supported in part by the California Institute for Physics and Astrophysics via grant CIPA-MG7099. The author is grateful to C. Van Den Broeck for useful discussions.
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# Phonon Dispersion Effects and the Thermal Conductivity Reduction in GaAs/AlAs Superlattices
## I Introduction
Superlattice structures have been proposed to be materials with a high thermoelectric figure of merit ZT, for both in-plane and cross-plane current flow. In the latter case, the improvement in thermoelectric performance is attributable to a reduced lattice thermal conductivity rather than to a higher electronic conductivity. Experimentally, a factor-of-ten reduction in the component of lattice thermal conductivity along the growth axis, $`\kappa _{\mathrm{},zz}`$, is observed in GaAs/AlAs and Bi<sub>2</sub>Te<sub>3</sub>/Sb<sub>2</sub>Te<sub>3</sub> superlattices (SLs). Theoretically, the thermal conductivity of Si/Ge SLs has been studied previously by Hyldgaard and Mahan and by Chen . Within the context of a very simplified model of the phonon dynamics,the calculations of Ref. 9 were able to reproduce the factor-of-ten reduction in the thermal conductivity along the growth direction. By contrast, Chen’s extensive work focussed on the role of thermal boundary resistance. The present paper investigates the reduction associated with SL induced changes of the phonon dispersion based on a realistic, computationally intensive treatment of the phonon spectra and dynamics.
The origin of the observed reduction in thermal conductivity may be explained qualitatively as follows. First, we note that, in the experimental work of Capinski and Maris and Capinski et al on (GaAs)<sub>n</sub>/(AlAs)<sub>n</sub> SLs for $`n`$ up to 40, the phonon mean free path inferred from the thermal conductivity, heat capacity and Debye velocity is greater than 370 Å at all termperatures for which measurements exist, always large compared to the size of the SL unit cell, 5.66 Å. Thus, the phonon transport lies in a regime where the SL phonon dispersion relation and lifetime, and not those of the bulk constituents, determine the thermal conductivity. According to the expression for the thermal conductivity derived from the phonon Boltzmann equation in the relaxation-time approximation (see Eq. (6)), the thermal conductivity depends on (1) a quantity representing the contribution of the SL phonon dispersion relation, and (2) the lifetime $`\tau `$, which contains phonon-phonon, interface and defect scattering effects. We shall focus only on item (1), whose effect can be computed within our realistic, albeit complicated, lattice-dynamical model. The effect of the SL geometry is to introduce anticrossings and new gaps in the phonon dispersion relation, when the magnitude of the phonon wavevector along the growth direction equals an integer multiple of $`\pi /d`$, where $`d`$ is the period of the SL along the growth direction. The consequent flattening of the phonon branches near the Brillouin zone edge leads to a lowering of phonon velocities in the growth direction, and hence a reduction in thermal conductivity.
We describe in Section II the lattice-dynamical model for the SL, which is a generalization of the 11-parameter rigid-ion model of Kunc . It incorporates short-range interactions to next nearest neighbors and the long-range Coulomb interaction. The construction of the SL dynamical matrix is outlined in Section II and in the Appendix.
The formalism is applied to a (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub> SL in Section III. The phonon dispersion relation will be seen to display the flattening expected for a SL. Critical points, especially at the high-symmetry points $`\mathrm{\Gamma }`$, X and Z, produce sharp peaks in the density of states in the SL. Miniband formation and anticrossings in the SL phonon dispersion relation lead to a three-fold reduction, relative to bulk, in the contribution of the phonon dispersion relation to the thermal conductivity along the growth direction. The present results contrast with those of Hyldgaard and Mahan , who found that, in their simplified model, the order-of-magnitude reduction of $`\kappa _{zz}/\tau `$ was attributable to effects related to the SL dispersion relation alone. In our more realistic treatment of the SL dynamics, a significant three-fold reduction in the lifetime is also needed to explain the experimental reduction in $`\kappa _{zz}`$ by a factor of ten. Finally, the sensitivity of the decrease in $`\kappa _{zz}`$ to differences in masses or force constants between the GaAs and AlAs layers is investigated; differences in the force constants are found to play a markedly greater role in the reduction of $`\kappa _{zz}`$ than differences in mass.
Section IV is devoted to discussion and conclusions which are broadened by examining the dependence of our results on the period $`n`$ of the (GaAs)<sub>n</sub>/(AlAs)<sub>n</sub> SL for $`n=2,3`$ and 6. These results permit some discussion of interface effects on the phonon lifetime and a more detailed comparison of the present work with that of Refs. 9 and 10. The reduction in the contribution of the phonon dispersion relation to the thermal conductivity of the SL relative to bulk is computed, and found to be approximately independent of $`n`$ for the small values considered here. The lifetime for both the SL and bulk was determined using the experimental thermal conductivities of the SL and of bulk GaAs of Capinski et al . The lifetimes are found to be smaller for the 2x2 SL and larger and roughly equal for the 3x3 and 6x6 SLs, consistent with the presence of greater interface scattering in the 2x2 SLs.
## II Formalism
The lattice dynamics can be treated realistically via the 11-parameter rigid-ion model of Kunc , which has been successfully applied to zinc-blende-structure compounds. Consider (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub>. In the SL unit cell, consisting of a layer of GaAs above a layer of AlAs in the growth direction, there will be 3 pairs of GaAs followed by 3 pairs of AlAs, or 3 Ga, 3 Al and 6 As atoms in all, which may be indexed by $`\kappa =1,\mathrm{},12`$. Letting $`u_\alpha (\mathrm{}\kappa )`$ denote the displacement in the direction $`\alpha =x,y,z`$ of the $`\kappa `$-th atom in the $`\mathrm{}`$-th unit cell, plane-wave solutions of the form
$$u_\alpha (\mathrm{}\kappa )=M_\kappa ^{1/2}e^{i(𝐤𝐱(\mathrm{}\kappa )\omega _j(𝐤)t)}w_\alpha (\kappa |𝐤j)$$
(1)
are assumed, where $`𝐱(\mathrm{}\kappa )`$ is the equilibrium position of the $`\kappa `$-th atom in the $`\mathrm{}`$-th unit cell, and $`w_\alpha (\kappa |𝐤j)`$ satisfies the secular equation
$$\omega _j(𝐤)^2w_\alpha (\kappa |𝐤j)=\underset{\beta \kappa ^{}}{}C_{\alpha \beta }(\kappa \kappa ^{}|𝐤)w_\beta (\kappa ^{}|𝐤j).$$
(2)
The dynamical matrix $`\underset{¯}{C}`$ reflects the interatomic force constants of the crystal. In the present rigid-ion model, the interatomic forces are divided into (1) short-range forces extending to second nearest neighbors, and (2) the long-range Coulomb interaction. Accordingly, the dynamical matrix may be written
$$\underset{¯}{C}=\underset{¯}{C}_{sr}+\underset{¯}{C}_{\mathrm{Coul}}.$$
(3)
As a result of symmetry of the zinc-blende structure, the short-range forces to second nearest neighbors may be described by ten parameters for each material . For nearest and next nearest neighbor interactions in the SL unit cell, we employ the force constants determined for the constituent bulk materials separately, rotated by the appropriate point-group operation. Bulk GaAs and AlAs parameters are taken from the literature . In the SL, bulk parameters are used within each layer. For the interface atoms and Ga-Al bonds crossing the interface, we employ the average of the bulk parameters following Ren et al . For the Coulomb interaction, the atoms are treated as point charges. The Madelung sum, and its derivatives, are computed using the usual Ewald transformation , which has been generalized here for SLs. This is accomplished by separating the sum over the spatial index $`\mathrm{}`$ into a sum over layers normal to the growth direction, $`\mathrm{}_{}`$, and a sum along the growth axis, $`\mathrm{}_z`$. A two-dimensional Ewald transformation is performed in each layer; these results are then summed over $`\mathrm{}_z`$. The resulting expressions for the function $`\varphi (𝐤,𝐫)`$ and its derivatives (see the Appendix) are similar to those that arise in the usual three-dimensional Ewald procedure; however, the definite integrals differ and must be performed numerically. Detailed expressions for the Coulomb term are given in the Appendix.
The phonon Boltzmann equation in the relaxation-time approximation leads to the following expression for the lattice thermal conductivity:
$$\kappa _{ij}=\frac{d^3q}{(2\pi )^3}\underset{\alpha }{}\mathrm{}\omega _𝐪^{(\alpha )}\frac{\omega _𝐪^{(\alpha )}}{q_i}\frac{\omega _𝐪^{(\alpha )}}{q_j}\frac{dn(\omega _𝐪^{(\alpha )})}{dT}\tau _{\mathrm{ph}}(\omega _𝐪^{(\alpha )},T)$$
(4)
where $`n(\omega _q^{(\alpha )})`$ is the distribution function of the phonons, the sum is over branches $`\alpha `$, and $`\tau _{\mathrm{ph}}(\omega _𝐪^{(\alpha )},T)`$ is the lifetime. Eq. (4) can be written in terms of
$$\mathrm{\Sigma }_{ij}(\omega )=\frac{d^3q}{(2\pi )^3}\underset{\alpha }{}\mathrm{}\omega _𝐪^{(\alpha )}\frac{\omega _𝐪^{(\alpha )}}{q_i}\frac{\omega _𝐪^{(\alpha )}}{q_j}\delta (\omega \omega _𝐪^{(\alpha )})$$
(5)
as
$$\kappa _{ij}=𝑑\omega (dn(\omega )/dT)\mathrm{\Sigma }_{ij}(\omega )\tau _{\mathrm{ph}}(\omega ,T).$$
(6)
This paper will focus on the SL effects on the dispersion relation contained in $`\mathrm{\Sigma }_{ij}(\omega )`$. Relaxation-time effects associated for example with scattering from interfaces, defects, umklapp processes, etc. are not considered explicitly. However, the results will be used to infer some of their properties.
## III (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub> Superlattices
We focus first on the (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub> SL studied experimentally by Capinski and Maris . Using a picosecond pump-and-probe technique, they observed an order-of-magnitude reduction in the thermal conductivity along the growth direction, $`\kappa _{zz}`$, relative to bulk GaAs. The dispersion relation along the $`\mathrm{\Gamma }`$X and $`\mathrm{\Gamma }`$Z directions, which was generated numerically according to the method described in Section II, is shown in Fig. 1. The significance of the labelled features will be explained in the discussion of Fig. 2. Because the SL unit cell contains three unit cells each of bulk GaAs and bulk AlAs, respectively, arranged along the growth axis, the edge of the SL Brillouin zone in the growth direction, Z, is one-sixth as far from the center as it is in the in-plane directions, X and Y. As a result each of the six branches in the bulk material is folded back six times along $`\mathrm{\Gamma }`$Z, which is most easily seen for the longitudinal acoustic mode in Fig. 1. Bulk GaAs and bulk AlAs optical modes have no frequencies in common, and are therefore localized. This leads to (1) flat SL dispersion in optical modes, and (2) localization of AlAs optical modes to the AlAs layer. The GaAs optical modes are not localized, since they overlap with the acoustic modes as shown in Fig. 1 along $`\mathrm{\Gamma }`$X. Due to the non-analyticity of $`\underset{¯}{C}_{\mathrm{Coul}}`$ as $`q0`$ in the SL, $`\omega (|q|0)`$ differ along $`\mathrm{\Gamma }`$X and $`\mathrm{\Gamma }`$Z (See Ref. 14).
The density of states $`\rho (\omega )`$ versus frequency, computed using the tetrahedral integration method as presented in MacDonald et al , is given in Fig. 2. Note the transverse-acoustic features around $`\omega =80`$ to 100 cm<sup>-1</sup>, the longitudinal-acoustic features around 150 to 200 cm<sup>-1</sup>, the GaAs optical feature at 220 to 260 cm<sup>-1</sup>, and, separated in frequency at higher frequencies, the AlAs optical feature at 330 to 400 cm<sup>-1</sup>. Band gaps at the zone edge and anticrossings in Fig. 1 yield critical points which, depending on the amount of $`𝐪`$-space at those frequencies, produce sharp structure in the density of states. As shown in Fig. 2, this structure in the density of states can be correlated with features in the dispersion relation, usually at the $`\mathrm{\Gamma }`$ (for folded-back bands), X and Z points. The structures at $`\mathrm{\Gamma }`$, X and Z labelled in Fig. 2 are identified with the corresponding features in Fig. 1. The strength of each feature depends on the integral of $`\delta (\omega \omega _𝐪^{(\alpha )})`$ (cf. Fig. 2) at that frequency, which will be large if $`|𝐯|`$ is small. Surprisingly, most of the peaks in the density of states appear to be associated with the critical points at $`\mathrm{\Gamma }`$, X and Z. No such fine structure in the density of states exists in bulk GaAs or bulk AlAs, as may be seen for instance in Patel et al for GaAs. (The density of states for AlAs is not available in the literature, but our calculations confirmed the absence of fine structure for AlAs as well.)
Fig. 3 shows the results for $`\mathrm{\Sigma }(\omega )`$, as defined in Eq. (5), for bulk GaAs and the SL. Note that optical modes do not contribute appreciably to $`\mathrm{\Sigma }_{zz}(\omega )`$ in the SL, an effect of localization in the AlAs layers: flat dispersion leads to a vanishing of $`\omega _𝐪^{(\alpha )}/q_z`$. The fine structure in the density of states is also correlated with that in $`\mathrm{\Sigma }(\omega )`$: peaks in the density of states DOS $`𝑑S/|𝐯|`$ become dips in $`\mathrm{\Sigma }v^2𝑑S/|𝐯|`$. (Here, $`S`$ denotes the surface of constant frequency $`\omega `$ in the Brillouin zone; it consists not only of the surface containing the critical point, but also possibly of other surfaces at the same $`\omega `$ elsewhere in the zone.) Acoustic modes in the SL contribute less than they would in bulk because of the band-gap and anticrossing-induced reduction in $`v^2`$. This leads to a three-fold reduction in $`\kappa _{\mathrm{}}/\tau `$ at 300 K, determined here either by integrating Eq. (4) directly on a $`60\times 60\times 20`$ grid covering an irreducible wedge of the Brillouin zone, or by integrating Eq. (6). Experimentally, Capinski and Maris found a ten-fold reduction factor for (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub>. The full reduction factor is a product of the reduction due to $`\mathrm{\Sigma }`$ and that due to $`\tau `$. Assuming $`\tau `$ to be constant at any given temperature, it can be found by requiring equality in Eq. (4) or (6) to the experimental value for the thermal conductivity in bulk or SL, respectively, if the appropriate dispersion relation is used in computing $`\kappa _{zz}/\tau `$. The lifetime in bulk versus lifetime in SL, as well as the layer-width dependence of the results, will be discussed below. Finally, we note that $`\mathrm{\Sigma }_{zz}<\mathrm{\Sigma }_{xx}`$ in the SL for the simple physical reason that band flattening along the $`q_z`$-direction affects $`\omega /q_z`$ more than $`\omega /q_x`$.
The reduction in thermal conductivity due to $`\mathrm{\Sigma }`$ can be understood by means of an easily visualized picture in q-space. In bulk, the quantity $`q_x_\alpha (dn(\omega _𝐪^{(\alpha )})/dT)\omega _𝐪^{(\alpha )}(v_{𝐪,z}^{(\alpha )})^2`$ represents the contribution of the phonon dispersion relation in Eq. (4). In a range $`\mathrm{\Delta }q_x`$ of the integrand it is weighted by $`q_x`$ because, the dispersion relation being rotationally symmetric in the $`(q_x,q_y)`$ plane, we may integrate around the circle of radius $`q_x`$ to yield a properly weighted function of $`q_x`$ and $`q_z`$ alone. The quantity is plotted in Fig. 4(a),(b),(c) (bold line) together with the corresponding values for GaAs (light solid line) as a function of $`q_z`$ for three values of $`q_x`$. The SL contribution is reasonably localized in $`q_z`$. To gain physical insight, we replace it by the dashed rectangles (of equal area). This approximate localization is related to the reduction at $`q_z0`$ and $`\pi /d`$ due to miniband formation (that is, flattening of $`\omega `$ vs. $`q`$) and has dips from anticrossings, as in Fig. 4(a), for instance. The dependence on $`q_x`$ can then be summarized by the density plot in Fig. 4(d), comparing the SL on the left to bulk GaAs on the right. For GaAs the equivalent rectangles extend over the entire range of $`q_z`$ because there is no localization due to band flattening at the zone edge as in the case of the SL. The shading indicates the weight of each increment $`\mathrm{\Delta }q_x`$. We find that points around $`q_x=\pi /d`$ and at the zone edge ($`q_x=6\pi /d`$) contribute most to heat transport. The latter is due to the effect of the weighting by $`q_x`$ in the annular integration. In addition, this weighting causes the contribution from points near the origin to be very small. A simple numerical estimate leads to a reduction in $`\mathrm{\Sigma }`$ to 34%, which is to be compared to 36% in the exact calculation.
Finally, we discuss the sensitivity of variation of the mass differences, force constants and effective ionic charge $`e^{}`$ between layers on the thermal conductivity. This manifests itself through variations of the zone edge gap and hence the group velocity. This effect has been studied by interpolating $`e^{}`$, the force constants $`K`$, and cation mass $`M`$ between their AlAs and their GaAs values in the GaAs layer. Here $`K`$ refers collectively to the 10 Kunc parameters describing interatomic forces. Explicitly, starting from (AlAs)<sub>6</sub> we (i) vary $`e^2`$ and $`M`$ linearly in three neighboring AlAs positions to make a (AlAs)<sub>3</sub>/(GaAs)<sub>3</sub> SL. In the case of the masses we put $`M=(1\alpha )M_{\mathrm{Al}}+\alpha M_{\mathrm{Ga}}`$ so that at $`\alpha =0`$ and 1 correspond respectively to Al and Ga; (ii) vary $`e^2`$ and force constants $`K`$ linearly in $`\alpha `$ in the same way; and (iii) vary $`e^2`$, and both $`M`$ and $`K`$ linearly in $`\alpha `$. The results are given in Fig. 5. The thermal conductivity is seen to be far more sensitive to the variation of force constants than the variation in mass. The variation of the force constants alone produces band flattening which reduces the thermal conductivity by about 60%. The additional flattening due to changing masses leads to only a few percent additional reduction.
## IV Discussion and Conclusions
Before turning to the final results, we emphasize again that this paper is primarily concerned with the effects of superlattice induced changes of the phonon dispersion on the lattice thermal conductivity. Since a SL is a perfect crystal, the ideal structures under discussion here automatically include the coherent effects associated with perfect interfaces. The importance of including such effects was first pointed out by Hyldgaard and Mahan in connection with a simple, highly idealized model for Si/Ge SLs. Extensive work by Chen focussed on diffuse interface scattering of phonons. His model assumes SL layers sufficiently thick that the phonon spectrum in each layer corresponds to that of the bulk. A mixture of spectral and diffuse interface processes is found to be sufficient to explain the observed experimental reduction of the SL thermal conductivity relative to bulk. As already pointed out, the more realistically modeled results of the present paper for $`\kappa _{\mathrm{}}/\tau `$ lead to a three fold reduction without any assumptions about the SL phonon scattering mechanisms. Our results must therefore be viewed as complementary to those of Chen . Taken together, they suggest that a combination of phonon spectral changes and imperfect interfaces can account adequately for the observed reduction.
The present calculations permit some statements concerning lifetime effects from the dependence of the SL results on layer width. This dependence was studied numerically for 2x2, 3x3 and 6x6 GaAs/AlAs SLs. The results for
$$\overline{\mathrm{\Sigma }}_{zz}\kappa _{zz}/\tau =\frac{d^3q}{(2\pi )^3}\underset{\alpha }{}\mathrm{}\omega _𝐪^{(\alpha )}\frac{\omega _𝐪^{(\alpha )}}{q_i}\frac{\omega _𝐪^{(\alpha )}}{q_j}\frac{dn(\omega _𝐪^{(\alpha )})}{dT}$$
(7)
are given in Table I. $`\tau `$ is assumed constant. Given an experimental value for $`\kappa _{zz}`$ and the calculated value of $`\overline{\mathrm{\Sigma }}_{zz}`$, the lifetimes listed in Table I are found as $`\tau =\kappa _{zz}/\overline{\mathrm{\Sigma }}_{zz}`$; only the lifetime itself is given in Table I and not $`\overline{\mathrm{\Sigma }}_{zz}`$, but the ratios $`\overline{\mathrm{\Sigma }}_{zz}`$(SL)/$`\overline{\mathrm{\Sigma }}_{zz}`$(bulk) are listed because we are interested in
$$\frac{\tau _{\mathrm{SL}}}{\tau _{\mathrm{bulk}}}=\frac{\overline{\mathrm{\Sigma }}_{zz}(\mathrm{bulk})}{\overline{\mathrm{\Sigma }}_{zz}(\mathrm{SL})}\frac{\kappa _{zz}^{\mathrm{SL}}(\mathrm{expt})}{\kappa _{zz}^{\mathrm{bulk}}(\mathrm{expt})}.$$
(8)
The SL $`\overline{\mathrm{\Sigma }}_{zz}`$ is found to have a value about 40% of the bulk, and to be relatively insensitive to the SL period and temperature. For larger SL periods, there is more folding back, but this is compensated by the smaller size of the gaps at the zone edge, which is found in the computed $`\mathrm{\Gamma }`$Z dispersion relations to scale inversely with the SL period. The two effects balance, resulting in an approximately constant reduction factor. The phonon lifetimes for the 2x2, 3x3 and 6x6 SLs are also given in Table I. The bulk lifetime is the same to within 2% for the different SL periods, as it must be. The ratio of SL to bulk lifetimes is significantly smaller for the 2x2 SLs at each temperature while it is larger and roughly the same for the 3x3 and 6x6 SLs. Thus, the reduction in thermal conductivity may be divided into a dispersive part which is insensitive to $`n\times n`$ and a scattering part which is sensitive to interface scattering. The results for the 2x2 SL are possibly associated with the experimental difficulties associated with achieving sufficiently perfect interfaces for small $`n`$.
The present calculations imply: (1) that a similar reduction in the contribution of the SL phonon disperion relation to transport in the growth direction, and perhaps in $`\kappa _{zz}`$ itself, may be expected in any SL with similar mass or force-constant differences between layers. (2) If the lifetime is reduced in SLs by increased umklapp scattering, as suggested by Ren and Dow<sup>13</sup>, then the present calculations give an upper bound on the SL $`\kappa _{zz}`$ and a lower bound on the increase in $`ZT`$ in a SL relative to bulk.
###### Acknowledgements.
We wish to thank Dr. E. Runge for stimulating discussions. This work was supported by DARPA through ONR Contract No. N00014-96-1-0887 and NSF Grant Che9610501.
## A
This Appendix presents the detailed formulae for the Coulomb part of the dynamical matrix, derived using the Ewald procedure as described in the text of the paper. Letting $`\kappa `$ label the atoms of the SL unit cell, $`M_\kappa `$ be the mass of the $`\kappa `$-th atom and $`𝐱(\kappa \kappa ^{})`$ be the separation vector from the $`\kappa `$-th atom to the $`\kappa ^{}`$-th atom, we have, for $`\kappa \kappa ^{}`$,
$$C_{\alpha \beta }^{\mathrm{Coul}}(\kappa \kappa ^{}|𝐤)=\frac{e_\kappa e_\kappa ^{}}{\sqrt{M_\kappa M_\kappa ^{}}}e^{i𝐤𝐱(\kappa ^{}\kappa )}\frac{^2\varphi (𝐤,𝐫)}{r_\alpha r_\beta }|_{𝐫=𝐱(\kappa ^{}\kappa )}$$
(A1)
where $`\varphi (𝐤,𝐫)`$ is by definition
$$\varphi (𝐤,𝐫)=\underset{\mathrm{}}{}\frac{e^{i𝐤𝐱(\mathrm{})}}{|𝐱(\mathrm{})+𝐫|},$$
(A2)
$`𝐱(\mathrm{})`$ being the position of the $`\mathrm{}`$-th unit cell. The result of the Ewald procedure adapted to the superlattice is that $`\varphi (𝐤,𝐫)`$ can be written in the form
$`\varphi (𝐤,𝐫)`$ $`=`$ $`R{\displaystyle \underset{\mathrm{}}{}}H(|𝐱(\mathrm{})+𝐫|R)e^{i𝐤𝐱(\mathrm{})}`$ (A3)
$`+`$ $`{\displaystyle \frac{2\sqrt{\pi }}{v_{}}}{\displaystyle \underset{h_{},\mathrm{}_z}{}}{\displaystyle \frac{2}{|\tau (h_{})+𝐤_{}|}}I(\mathrm{},{\displaystyle \frac{|\tau (h_{})+𝐤_{}|}{2R}},{\displaystyle \frac{|\tau (h_{})+𝐤_{}||𝐱(\mathrm{}_z)+z|}{2}})e^{i(\tau (h_{})+𝐤_{})𝐫_{}+ik_z\widehat{z}𝐱(\mathrm{}_z)}`$ (A4)
where $`R`$ is an arbitrary cutoff (we always take $`R=3/a_0`$),
$$H(x)=\frac{2/\sqrt{\pi }}{x}_x^{\mathrm{}}e^{x^2}𝑑x^{},$$
(A5)
$`v_{}`$ is the area of the unit cell in the superlattice plane, $`\mathrm{}_z`$ labels the layers perpendicular to the growth ($`z`$) axis, $`h_{}=(h_x,h_y)`$ labels the cells located at reciprocal lattice vectors $`\tau (h_{})`$ in each plane, $`𝐤=(𝐤_{},k_z)`$, and
$$I(\alpha ,\beta ,\gamma )=_\beta ^\alpha e^{v^2\gamma ^2/v^2}𝑑v.$$
(A6)
For $`\kappa =\kappa ^{}`$ we have
$$C_{\alpha \beta }^{\mathrm{Coul}}(\kappa \kappa |𝐤)=\frac{e_\kappa }{M_\kappa }\left[e_\kappa \underset{\mathrm{}0}{}e^{i𝐤𝐱(\mathrm{})}\left(\frac{^2r^1}{r_\alpha r_\beta }\right)_{𝐫=𝐱(\mathrm{})}+\underset{\mathrm{}^{}\kappa ^{}0\kappa }{}e_\kappa ^{}\left(\frac{^2r^1}{r_\alpha r_\beta }\right)_{𝐫=𝐱(\mathrm{}^{}\kappa ^{},0\kappa )}\right].$$
(A7)
The first term on the right is similar to the above expressions in Eqs. (A1)-(A3) but for the $`\mathrm{}=0`$ in Eq. (A3) term one substitutes $`(4/3\sqrt{\pi })\delta _{\alpha \beta }`$. The second term is given by
$$\frac{^2}{r_\alpha r_\beta }\left[I_0+I_1+I_2\right]_{𝐫=0}$$
(A8)
with
$$I_0=Re_\kappa H^0(|𝐫|R),$$
(A9)
$$I_1=R\underset{\mathrm{}^{}\kappa ^{}0\kappa }{}e_\kappa ^{}H(|𝐱(\mathrm{}^{})+𝐱(\kappa ^{}\kappa )+𝐫|R)$$
(A10)
and
$$I_2=\frac{2\sqrt{\pi }}{v_{}}\underset{h_{},\mathrm{}_z,\kappa ^{}}{}e_\kappa ^{}\frac{2}{|\tau (h_{})|}I(\mathrm{},\frac{|\tau (h_{})|}{2R},\frac{|\tau (h_{})||𝐱(\mathrm{}_z)+\widehat{z}𝐱(\kappa ^{}\kappa )|}{2})e^{i\tau (h_{})(𝐱(\kappa ^{}\kappa )+𝐫)}$$
(A11)
where
$$H^0(x)=\frac{2/\sqrt{\pi }}{x}_0^xe^{x^2}𝑑x^{}.$$
(A12)
A similar, but not identical, approach was used in Ref. 15.
FIGURES
FIG. 1. (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub> SL dispersion relation along the $`\mathrm{\Gamma }`$X=$`(2\pi /a_0,0,0)`$ and $`\mathrm{\Gamma }`$Z=$`(0,0,\pi /3a_0)`$ directions; $`a_0`$ is the conventional GaAs unit cell size. The labels $`aw`$ are defined in Fig. 2.
FIG. 2. Phonon density of states for the (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub> SL, with labelled critical points identified with features in the dispersion relation in Fig. 1.
FIG. 3. The transport quantities $`\mathrm{\Sigma }_{xx}(\omega )`$ and $`\mathrm{\Sigma }_{zz}(\omega )`$, defined by Eq. (5) in the text, for bulk GaAs (solid line) and the (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub> SL (dashed and bold lines).
FIG. 4. The transport quantity $`q_x(dn/dT)_\alpha \omega _𝐪^{(\alpha )}(v_{𝐪,z}^{(\alpha )})^2`$ for $`0q_z\pi /d`$ at fixed $`q_x`$, for (a) $`q_x=2\pi /d`$, (b) $`q_x=4\pi /d`$ and (c) $`q_x=6\pi /d`$; solid line: bulk GaAs; bold line: (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub> SL. (d) Density plot in the $`(q_x,q_z)`$ plane whose shading indicates the weight of each increment $`\mathrm{\Delta }q_x`$ along $`q_x`$ to the value of this transport quantity for the (GaAs)<sub>3</sub>/(AlAs)<sub>3</sub> SL and bulk GaAs.
FIG. 5. Dependence of $`\kappa _{zz}`$ as $`e^2`$ and the mass $`M`$ (solid line), the spring constants $`K`$ (long-dashed line) and both $`M`$ and $`K`$ (short-dashed line) are interpolated between their AlAs ($`\alpha =0`$) and GaAs values ($`\alpha =1`$) for the atoms in three contiguous layers in what is initially (AlAs)<sub>6</sub> at $`\alpha =0`$.
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# Dynamic splitting of a Bose-Einstein condensate
## I Introduction
Bose-Einstein condensation of an ideal gas is typically presented in introductory textbooks solely in terms of particle numbers. And quantum mechanically enhanced number densities were the ‘smoking gun’ observed in the first experiments on dilute gas condensates. But many phenomena of interacting condensates depend critically on the conjugate quantity to particle number, namely the quantum mechanical phase . One can have highly occupied states with or without phase coherence between them, and the presence or absence of phase coherence can make a dramatic difference in the physical properties of an ultra-cold gas. As in the case of a gas held in an optical lattice, which can be a superfluid or a Mott insulator depending on the strength of the lattice, the onset or loss of phase coherence can even be a phase transition .
Since the global phase of a condensate is unobservable, the simplest system in which phase coherence can be manifested consists of two states, occupied by a large number of bosons. As such a system can realistically be approximated by a condensate in a double well, it has recently attracted attention . In these works, Josephson oscillations and the phase coherence between two coupled condensate were studied considering time independent coupling parameters. In the disappearance of the phase coherence between the two wells due to a change in time of the tunneling coupling has been studied. The coupling parameters and the related phase coherence properties depend on the potential barrier between the two wells, since in the limit where the barrier is very low we have a single condensate and in the limit where the barrier is very high we have two completely separated condensates. Even if, at least in the high barrier limit, it has often been pointed out how to relate the coupling parameters (on-site energy and tunneling coupling) to the overlap of the wavefunctions localised in the two wells, this has never really been taken explicitly into account. Introducing the spatial degrees of freedom, as we did, allows us to relate all that to the potential barrier in a more than phenomenological way and becomes expecially important in the study of the dynamics of the process, because the wavefunctions change drastically in time and collective modes are excited.
In this paper we examine this problem and develop a method which can be extended to the case of many wells, in order to encompass the turning on of an optical lattice in a condensate, allowing to go beyond a Gross-Pitaevkii treatment, as the one done for example in . The physical situation is similar to the ones already realized experimentally, where a double well potential was created by shining a far-off resonant laser beam in the center of the magnetic trap (see e.g. ) or where an array of traps was created by on optical standing wave .
We are interested in the full dynamics of the process: in addition to the phase coherence properties, we want to study the excitation of the collective modes by taking the spatial dependence of the condensate wavefunction explicitly into account. The price of including all these effects without assuming mean field theory is that we must use a time-dependent variational approach, choosing variational ansatz to describe both the “internal” and “external” dynamics (that is, the distribution of particles between two motional states treated as given, and the evolution of the spatial wave functions of these states). This allows us to reduce the intractable full problem to a set of coupled differential equation for our few variational parameters. Although the time-dependent variational approach is not guaranteed to be quantitatively accurate, it allows qualitatively important processes to be investigated, and it has proven surprisingly reliable in previous applications to condensate physics . Here we use it to derive the coupling between the internal and external dynamics, investigate which are the conditions under which the two can be decoupled, and identify the typical time scales for both.
In Sec. II we briefly present our two-state model and discuss qualitatively the behaviour of the system that our model can explain. In Sec. III, we introduce explicitly the variational ansatz we choose to describe the phase dynamics and the collective modes. After introducing the time–dependent variational principle and deriving the Lagrangian, we write the equations of motion for the variational parameters. Then in Sec. III E we discuss the conditions under which it is an accurate approximation to neglect the coupling between internal and external degrees of freedom. This decoupling allows us to model both internal and external evolution in a still simpler way. The external dynamics is studied in Sec. IV, comparing simple analytic models with numerical results. Sec. V is devoted to the internal dynamics, where the statics is studied and analytic estimates obtained. With these decoupled studies to guide expectations, the full variational equations of motion, with coupled internal and external degrees of freedom, will be studied numerically in Sec. VI. The results are compared with those of the phase model describing the relative phase of two weakly coupled superfluids . We conclude with a general discussion of our results and their implications. Two appendices define the number difference and relative phase operators , derive the phase model Hamiltonian, and show that our variational ansatz adequately represents the evolution generated by it.
## II Model
Let us consider the situation in which we have $`N`$ bosons confined in a harmonic trap at zero temperature. We slowly deform the trap symmetrically around its center raising a potential barrier, until it becomes a double–well potential. We want to study the dynamics of the process as well as the final state of the bosons.
We will treat the problem in one dimension. This describes the situation in which we have a cigar-shaped condensate, and we deform it along the most elongated direction. Nevertheless, our model could be extendend in a straightforward way to $`2`$ and $`3`$ dimensions. This would just lead to more complicated equations, but would not affect the main results.
### A Two–mode model
The second quantization Hamiltonian describing the situation we have in mind and in which particles interact via a $`\delta `$-pseudopotential is
$`\widehat{H}(t)`$ $`=`$ $`{\displaystyle \widehat{\mathrm{\Psi }}^{}(z)\left(\frac{p^2}{2M}+V(z,t)\right)\widehat{\mathrm{\Psi }}(z)𝑑z}+`$ (1)
$`+`$ $`{\displaystyle \frac{1}{2}}g{\displaystyle \widehat{\mathrm{\Psi }}^{}(z)\widehat{\mathrm{\Psi }}^{}(z)\widehat{\mathrm{\Psi }}(z)\widehat{\mathrm{\Psi }}(z)𝑑z},`$ (2)
where $`\widehat{\mathrm{\Psi }}`$ is the bosonic field operator and $`V(z,t)`$ is the time–dependent potential which describes the deformed trap. Here, $`g`$ is an effective coupling constant which depends both on the $`s`$-wave scattering length and on the atomic distribution in the transverse directions.
We assume that at time $`t=0`$ the atoms are in the ground state of this Hamiltonian, and we want to determine the state after the potential is deformed. This problem cannot be solved even numerically (see, for example, and references therein). Therefore, we need to consider a simplified model which describes the main features of the process. If the process is completely adiabatic, the final state will be a ‘fragmented condensate’ with half of the particles in each of the potential wells: when the two condensates do not interact this is a much lower energy state than the one with phase coherence, because minimizing fluctuations in the relative particle number lowers the energy due to interparticle repulsion. Such a fragmented condensate can be regarded as having two entirely independent condensates. If we changed the potential very fast, then we would obtain a single condensate that oscillates in each of the potential wells. We would also expect to see collapses and revivals in the “condensate phase” , provided the losses are not important .
It is clear that the Gross–Pitaevskii Equation (GPE) will not give a good description of the splitting process, in principle. This equation describes the evolution of a single mode of the condensate $`\phi (z,t)`$, and, therefore, is not valid for fragmented condensates. In order to interpolate between the Gross-Pitaevskii limit of phase coherence and the limit of two independent condensates, one needs to consider at least two modes $`\phi _{1,2}(z,t)`$. Then, one can write the state of the atoms as (we assume $`N`$ to be even)
$`|\mathrm{\Phi }(t)`$ $`=`$ $`{\displaystyle \underset{m=N/2}{\overset{N/2}{}}}c_m(t)|m(t),`$ (3)
where
$`|m(t)`$ $`=`$ $`{\displaystyle \frac{a_1(t)^{\frac{N}{2}m}}{\sqrt{\left(\frac{N}{2}m\right)!}}}{\displaystyle \frac{a_2(t)^{\frac{N}{2}+m}}{\sqrt{\left(\frac{N}{2}+m\right)!}}}|vac,`$ (4)
and $`a_{1,2}`$ are the mode annihilation operators defined as
$`a_i^{}(t)={\displaystyle \phi _i(z,t)\mathrm{\Psi }^{}(z)𝑑z}.`$ (5)
We have to consider the evolution of the wavefunctions $`\phi _{1,2}`$ as well as of the coefficients $`c_m`$, which will be coupled and governed by the Hamiltonian
$`\widehat{H}(t)=`$ (6)
$`={\displaystyle \underset{ij=1,2}{}}a_i^{}a_j{\displaystyle \phi _i^{}(z,t)\left(\frac{p^2}{2M}+V(z,t)\right)\phi _j(z,t)𝑑z}`$ (7)
$`+{\displaystyle \frac{g}{2}}{\displaystyle \underset{ijlm=1,2}{}}a_i^{}a_j^{}a_la_m`$ (8)
$`{\displaystyle \phi _i^{}(z,t)\phi _j^{}(z,t)\phi _l(z,t)\phi _m(z,t)𝑑z}.`$ (9)
We will refer to the evolution of those wavefunctions as “external dynamics” and to the one of the coefficients as “internal dynamics”.
In order to find the mode–functions $`\phi _{1,2}`$ we can use the variational principle (in the same way as one derives the GPE from a Hartree-Fock ansatz). However, this also turns out to be very complicated. A way around that problem is to express $`\phi _{1,2}`$ in terms of some few variational parameters: we will use quasi–gaussian functions to describe those mode-functions. On the other hand, we could also use a variational principle to determine the evolution of the coefficients $`c_m`$. This again turns out to be very complicated, so that we will also use a Gaussian ansatz for them. Once they are known, in order to estimate when the condensate is fragmented, we can look at the eigenvalues of the single particle density operator corresponding to the internal dynamics only, that is, the matrix
$`\rho ={\displaystyle \frac{1}{N}}\left(\begin{array}{cc}a_1^{}a_1& a_1^{}a_2\\ a_2^{}a_1& a_2^{}a_2\end{array}\right),`$ (12)
where $`a_i^{}a_j`$ mean the expectation value on the state $`|\mathrm{\Phi }`$. The eigenvalues $`\lambda _\pm `$ indicate whether we have a single condensate ($`\lambda _+1`$, $`\lambda _{}0`$), or a fragmented one ($`\lambda _+\lambda _{}1/2`$).
As already extensively discussed in the literature , the two-mode model has a limited validity in the case of low barrier, when in principle one is not allowed to neglect higher excited modes. However, it becomes more and more accurate the higher the barrier gets, since in this case the two lower lying modes move closer together in energy compared to the higher ones. Hence it should allow a good description of the splitting process.
### B Qualitative behavior
Before presenting the numerical and analytical result coming from our analysis, we will briefly discuss the qualitative behavior that we expect from the model under study. We will show that the equations we will derive for the external and internal dynamics can be decoupled to a very good approximation. That is, we can first solve the equations for the external dynamics basically taking constant values for the variational parameters describing the coefficients $`c_m`$. Once these equations are solved, we can use the corresponding wavefunctions $`\phi _{1,2}`$ to calculate the time–dependent coefficients for the equations that describe the internal variational parameters. In summary, once we have solved the equations for the external parameters, we are left with a two–mode model with time–dependent coefficients which contain all information about the external dynamics: they will define the hopping and on-site interaction, whose competition determines the phase relation between the two modes.
Regarding the external dynamics, one can see two kinds of behaviors depending on the time scale $`\tau `$ at which the barrier is raised. The important time scale with which one has to compare is the oscillation period $`\tau _z`$ in the trapping potentials at each time. These periods change by roughly a factor of two between the initial harmonic potential and the final double well (with our specific choice of the trapping potential). Thus, if $`\tau \tau _z`$ the process will be adiabatic with respect to the external dynamics, which means that $`\phi _{1,2}(z,t)`$ will basically correspond to the two ground states of the right and left wells at the final time. If $`\tau \tau _z`$ we will have collective excitations, in which $`\phi _{1,2}(z,t)`$ oscillate strongly. In this case, we will have that the energy of the condensate $`E`$ (with respect to the ground state energy) has increased, so that it may be destroyed. Although we cannot account for the disappearance of the condensate within our model, we can estimate when this will happen just by considering the fact that under normal circumstances thermalization will occur, and, therefore, the condensate will be destroyed for $`E=K_BT_c>N\mathrm{}\omega _z`$, where $`\omega _z=2\pi /\tau _z`$, $`K_B`$ is the Boltzman constant, $`T_c`$ indicates the critical temperature and $`E`$ is the extra energy in the final state. We find that the condensate disappears for $`\tau \tau _z`$.
Regarding the internal dynamics, there is also a time scale in the problem which determines the dynamics; this is the revival time $`\tau _r`$. Given two condensates with an initial well defined relative phase, it is well–known that the relative phase first disappears (collapse) and then is restored at time $`\tau _r1/g`$ . If $`\tau \tau _r`$ the process will be adiabatic with respect to the internal dynamics, the phase coherence will be lost during the process and therefore we will end up with two independent condensates in each well, with no phase coherence at all (note that it makes no sense to talk about collapses or revivals in this situation). If $`\tau \tau _r`$ at the end of the process we will have two condensates with a good phase coherence. In that case, after the splitting, collapses and revivals could be observed provided the particle losses are practically absent.
In summary, we have two important time scales in the problem, namely $`\tau _z`$ and $`\tau _r`$. Typically, in experiments $`\tau _r\tau _z`$, so that it will be harder to be adiabatic with respect to the internal dynamics than to the external one. On the other hand, since $`\tau _r`$ is very long in practise, it will be hard to achieve $`\tau \tau _r`$ within the validity of our model, in which particle losses and other imperfections are not included.
Finally we notice that the external and internal dynamics depend in a very different way on the parameters $`g`$ and $`N`$. Very similar to what happens in the GPE, the equations of motion for the parameters describing the external dynamics depend almost only on the product $`gN`$. On the contrary, the on-site energy splitting scales like $`g`$ and the tunneling coupling like $`N`$ giving rise to very different internal dynamics. For increasing $`N`$ and decreasing $`g`$ the relative phase becomes better defined and the time $`\tau _r`$ required to destroy the phase coherence is longer. In particular in the limit $`N\mathrm{}`$, unless the tunneling coupling exactly vanishes, the GPE case of phase coherent condensate is recovered.
## III Variational ansatz
In this Section we introduce a variational ansatz to describe the internal and external dynamics. To describe the ground state of the system, which we will call equivalently static or equilibrium solution, two parameters are sufficient: $`z_0`$ which corresponds roughly to the center of the mode functions, and $`p`$ which is related to the width of the number distribution. To allow dynamic evolution we have to introduce also the variables $`\sigma _I`$ and $`x_I`$, which will vanish in the static case. For simplicity in the following we present the variational ansatz for the symmetric case, describing a spatially symmetric double–well potential and a symmetric atomic distribution among the two modes. Starting from symmetric initial conditions (no unbalance in the population of the two modes and symmetric mode functions $`\phi _1(z)=\phi _2(z)`$), the symmetry will be preserved under the time evolution. Nevertheless, it is possible to generalize our ansatz to describe asymmetric situations. We will briefly discuss the corresponding results in the conclusions.
### A Variational ansatz for the mode functions
For a single condensate in a harmonic trap, a Gaussian ansatz for the wavefunction has proved to be very useful and able to predict the excitation frequencies within a very high precision . In our case we make a similar choice. For a very high barrier we expect to find two separate condensates for each of which a Gaussian ansatz should be good. We then define the two functions for the left and right side:
$`\phi _{R,L}(z)=`$ (13)
$`=`$ $`\left({\displaystyle \frac{\sigma _R}{\pi }}\right)^{1/4}\mathrm{exp}\left[{\displaystyle \frac{\sigma _R(zz_0)^2}{2}}\right]\mathrm{exp}\left(i{\displaystyle \frac{\sigma _I}{2}}z^2\right).`$ (14)
In the case of low barrier (small $`z_0`$), these two functions have to be orthogonalized to satisfy the orthonormality property required for the two mode functions. Hence we define
$`\phi _\pm (z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}\sqrt{1\pm \mathrm{exp}(\sigma _Rz_0^2)}}}\left(\phi _R\pm \phi _L\right),`$ (16)
$`\phi _{1,2}(z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\phi _+\pm \phi _{}\right),`$ (17)
which are different from $`\phi _{R,L}`$ because of the two different normalization constants for $`\phi _+`$ and $`\phi _{}`$.
The variational parameters describing the mode functions $`\phi _{1,2}`$ are $`z_0`$ and $`\sigma _I`$. Physically, the excitation modes that these parameters can describe depend on whether the two Gaussians significantly overlap in space or not. If they do, as in the case of low barrier, then the excitation mode (changes in $`z_0`$) corresponds to a breathing mode. If they do not overlap, as for high barrier, then it corresponds to an oscillation mode in each of the potential wells. Allowing also $`\sigma _R`$ to vary (and adding an extra variable for the dynamics) one can describe more excitation modes. Instead we fix it to a constant value, since the curvature of the potential wells will be chosen to be always of the same order of magnitude. The value of $`\sigma _R`$ is not obvious since the overlap integrals depend strongly on it. For instance the hopping terms can be overestimated due to the long Gaussian tales. To fix $`\sigma _R`$, we identified a range of values for which the dynamic behaviour of the system was qualitatively the same and choose one of the values within this interval, which turned out to be lower than the one corresponding to the static solution.
Since the mode wavefunctions are linear combinations of Gaussians, if we choose a trapping potential of the form
$`V(z)={\displaystyle \frac{1}{2}}M\omega _T^2z^2+V_0(t)\mathrm{exp}[z^2/a(t)],`$ (18)
the integrals in Eq.(V A) can be performed analytically.
In the following, we will use dimensionless units: $`a_{ho}=\sqrt{\mathrm{}/M\omega _T}=1`$, $`\mathrm{}\omega _T=1`$ and $`\mathrm{}=1`$. So, all lengths will be measured in units of harmonic oscillator length, all energy in units of the trap frequency and all times in units of $`\omega _T^1`$.
### B Variational ansatz for the coefficients
We also take for the coefficients $`c_m(t)`$ a Gaussian distribution centered at $`m=0`$,
$`c_m`$ $`=`$ $`𝒩(p)\mathrm{exp}\left[\left({\displaystyle \frac{1}{4p}}+ix_I\right)m^2\right],`$ (19)
where $`𝒩(p)`$ is a normalization constant that depends on $`p`$ only. The variational parameters are $`p`$ and its conjugate one $`x_I`$. The parameter $`p`$ is directly related to the width of the distribution $`|c_m|^2`$, whereas $`x_I`$ contributes to the width of the Fourier transform of such a distribution (i.e., to the width of the phase distribution, see App. A).
### C Time-dependent variational principle
We study the dynamics using the time dependent variational principle. To derive the equations of motion one starts by writing the action $`S`$
$`S={\displaystyle 𝑑t\frac{\dot{\mathrm{\Phi }}|\mathrm{\Phi }\mathrm{\Phi }|\dot{\mathrm{\Phi }}}{2i}}\mathrm{\Phi }|\widehat{H}|\mathrm{\Phi }.`$ (20)
where $`\widehat{H}`$ and $`|\mathrm{\Phi }`$ have been defined in (6) and (3), respectively. In evaluating the term $`\mathrm{\Phi }|\dot{\mathrm{\Phi }}`$, one should remember that the state $`|\mathrm{\Phi }`$ depends on time both through the coefficients and the mode functions contained in $`|m`$:
$`\dot{\mathrm{\Phi }}|\mathrm{\Phi }={\displaystyle \underset{m}{}}\dot{c}_m^{}c_m+{\displaystyle \underset{mm^{}}{}}c_m^{}^{}c_m\dot{m}^{}|m.`$ (21)
The Lagrangian which follows takes the form
$`L`$ $`=`$ $`{\displaystyle \frac{1}{2i}}[{\displaystyle \underset{m}{}}\dot{c}_m^{}c_m+`$ (23)
$`+{\displaystyle \underset{ij}{}}a_i^{}a_j{\displaystyle }\dot{\phi }_i^{}\phi _jdzh.c.],`$
where $`\mathrm{\Phi }|\widehat{H}|\mathrm{\Phi }`$ is given by
$`=`$ (24)
$`={\displaystyle \underset{ij=1,2}{}}a_i^{}a_j{\displaystyle \phi _i^{}(z,t)\left(\frac{p^2}{2M}+V(z,t)\right)\phi _j(z,t)𝑑z}`$ (25)
$`+{\displaystyle \frac{1}{2}}g{\displaystyle \underset{ijlm=1,2}{}}a_i^{}a_j^{}a_la_m`$ (26)
$`{\displaystyle \phi _i^{}(z,t)\phi _j^{}(z,t)\phi _l(z,t)\phi _m(z,t)𝑑z}.`$ (27)
From the Lagrangian (23), in general one derives the equations of motion, carrying out the variation with respect to the discrete variables $`c_m`$ and the fields $`\phi _i`$. In our case, all the integrals and expectation values are functions of the variational parameters which can be calculate analytically. The only important overlap integrals containing time derivatives are
$`{\displaystyle \dot{\phi }_{1,2}^{}\phi _{1,2}𝑑z}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left({\displaystyle \frac{1}{\sigma _R}}+{\displaystyle \frac{z_0^2}{1e^{2\sigma _Rz_0^2}}}\right)\dot{\sigma }_I,`$ (29)
$`{\displaystyle \dot{\phi }_{1,2}^{}\phi _{2,1}𝑑z}`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \frac{z_0^2e^{2\sigma _Rz_0^2}}{1e^{2\sigma _Rz_0^2}}}\dot{\sigma }_I.`$ (30)
If one defines
$`N_\varphi =a_1^{}a_2+a_2^{}a_1`$ (31)
and
$`(p,x_I,z_0)`$ (32)
$`=`$ $`{\displaystyle \frac{N}{2}}\left({\displaystyle \frac{1}{\sigma _R}}+{\displaystyle \frac{z_0^2}{1e^{2\sigma _Rz_0^2}}}\right){\displaystyle \frac{N_\varphi }{2}}{\displaystyle \frac{z_0^2e^{2\sigma _Rz_0^2}}{1e^{2\sigma _Rz_0^2}}},`$ (33)
the Lagrangian becomes
$`L`$ $`=`$ $`p\dot{x}_I+\dot{\sigma }_I,`$ (34)
and the corresponding equations of motion are
$`\left(\begin{array}{cccc}0& 1& 0& \frac{}{p}\\ 1& 0& 0& \frac{}{x_I}\\ 0& 0& 0& \frac{}{z_0}\\ \frac{}{p}& \frac{}{x_I}& \frac{}{z_0}& 0\end{array}\right)\left(\begin{array}{c}\dot{p}\\ \dot{x}_I\\ \dot{z}_0\\ \dot{\sigma }_I\end{array}\right)=\left(\begin{array}{c}\frac{}{p}\\ \frac{}{x_I}\\ \frac{}{z_0}\\ \frac{}{\sigma _I}\end{array}\right)`$ (47)
$``$ $`\{\begin{array}{ccc}\hfill \dot{p}& =& \frac{}{x_I}+\frac{}{x_I}\dot{\sigma }_I\hfill \\ \hfill \dot{x}_I& =& \frac{}{p}+\frac{}{p}\dot{\sigma }_I\hfill \\ \hfill \frac{}{z_0}\dot{z}_0& =& \frac{}{\sigma _I}\left(\frac{}{p}\dot{p}+\frac{}{x_I}\dot{x_I}\right)\hfill \\ \hfill \frac{}{z_0}\dot{\sigma }_I& =& \frac{}{z_0}\hfill \end{array}`$ (52)
These equations describe the internal and external coupled dynamics of the splitting of the condensate. In what follows, we have solved them numerically in different regimes. Before presenting the results we will show that one can decouple the evolution of the external and internal variables, which helps to understand the dynamics. For that, we will introduce in the next subsection some analytical approximations to derive explicit formulae for the quantities $``$ and $``$ by replacing the discrete distribution $`c_m`$ by a continuous one, and treating the index $`m`$ as a continuous variable running from $`\mathrm{}`$ to $`\mathrm{}`$.
### D Continuous limit
In order to calculate $``$ and $``$ we have to evaluate expectation values of the form $`a_i^{}a_j`$ and $`a_i^{}a_j^{}a_la_m`$. We can do that if we replace the sums in $`m`$ by integrals extended from $`\mathrm{}`$ to $`\mathrm{}`$. When this replacement is valid, we can even calculate the width of the number distribution $`|c_m|^2`$, $`\sigma _m`$, as well as the one corresponding to the phase distribution, $`\sigma _\varphi `$ (see App.A). We obtain
$`\sigma _m`$ $`=`$ $`\sqrt{p},`$ (54)
$`\sigma _\varphi `$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{4p}}+4px_I^2}.`$ (55)
On the other hand, we have
$`a_{1,2}^{}a_{1,2}={\displaystyle \frac{N}{2}},`$ (57)
$`a_{1,2}^2a_{1,2}^2={\displaystyle \frac{N^2}{4}}+p,`$ (58)
$`a_1^{}a_2=a_2^{}a_1^{}=`$ (59)
$`={\displaystyle \frac{N}{2}}\mathrm{exp}\left({\displaystyle \frac{\sigma _\varphi ^2}{2}}\right)\left[1{\displaystyle \frac{2p}{N^2}}+{\displaystyle \frac{8p^2x_I^2}{N^2}}\right],`$ (60)
$`a_1^{}a_2^{}a_1a_2={\displaystyle \frac{N^2}{4}}p,`$ (61)
$`a_1^2a_2^2=`$ (62)
$`={\displaystyle \frac{N^2}{4}}\mathrm{exp}\left(2\sigma _\varphi ^2\right)\left[1{\displaystyle \frac{4p}{N^2}}+{\displaystyle \frac{64p^2x_I^2}{N^2}}\right].`$ (63)
Furthermore, in this limit we can determine the eigenvalues of the single particle density operator corresponding to the internal degrees of freedom, obtaining
$`\lambda _\pm {\displaystyle \frac{1\pm \mathrm{exp}\left(\sigma _\varphi ^2/2\right)}{2}}.`$ (64)
Notice that $`\lambda _+1`$ for $`pN/4,x_I0`$, i.e. $`\sigma _\varphi 0`$, which corresponds to the Gross-Pitaevskii limit; instead $`\lambda _\pm 1/2`$ for large $`\sigma _\varphi `$, giving a signature of the fragmentation of the condensate.
One can easily determine the limits of validity of this continuous approximation. On the one hand, the distribution in $`c_m`$ has to be sufficiently broad, which implies $`p1`$. On the other hand, $`x_I`$ has to be such that $`\sigma _\varphi \pi `$. For our numerical simulations, we corrected the expressions in Eqs.(III D) to make them valid $`p`$ and $`x_I`$: for $`p>1`$ we included the periodicity in $`x_I`$ and for $`p<1`$ we performed the exact sum over $`m`$, considering only the few number states which are populated.
### E Decoupling between external and internal dynamics
The coupling among the $`z_0`$ dynamics and the $`p`$ dynamics appears in the off-diagonal blocks in Eq.(47) and in the dependence of $``$ on all variational parameters. In the following we will analyze under which condition it is possible to decouple the internal ($`p,x_I`$) and external ($`z_0,\sigma _I`$) dynamics, so that one can study them independently one from the other.
#### 1 External dynamics
One can rewrite the equation of motion for $`z_0`$ is a more explicit form as
$`{\displaystyle \frac{}{z_0}}\dot{z}_0`$ $`=`$ $`{\displaystyle \frac{}{\sigma _I}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{z_0^2\mathrm{exp}(\sigma _Rz_0^2)}{1\mathrm{exp}(2\sigma _Rz_0^2)}}\dot{N}_\varphi .`$ (65)
Let us first see under which circumstances it is possible to neglect the off-diagonal blocks:
(i) in the low barrier limit, i.e. $`\mathrm{exp}(\sigma _Rz_0^2)1`$, the off-diagonal blocks can be neglected if $`\dot{N}_\varphi `$ plays no role: $`N_\varphi `$ evolves at the frequency which governs the internal dynamics $`\omega _p`$ (remember definition (31)), while the external dynamics is governed by the frequency $`\omega _z`$, usually of the order of the trapping frequency $`\omega _T`$. If $`\omega _p\omega _z`$, then the time average of $`\dot{N}_\varphi `$ vanishes and has no effect on the evolution of $`z_0`$. On the other hand $`\omega _pgNU_1\mathrm{exp}(\sigma _Rz_0^2/2)\omega _z`$ only for a chemical potential $`\mu gNU_1\omega _T`$. This happens either for very few particles, a case in which our model does not hold, since $`N1`$ is a fundamental assumption in our model, or for very weakly interacting particles, in which $`\dot{N}_\varphi 0`$, because the phase coherence is very difficult to be destroyed.
(ii) in the high barrier limit, normally $`\omega _p\omega _z`$ (see Fig. 5, small $`p`$ limit) and $`N_\varphi `$ might vary abruptly, since this is when the phase coherence is supposed to desappear. Anyway, since $`\mathrm{exp}(\sigma _Rz_0^2)0`$, the off-diagonal block can be neglected.
Now let us analyze the $`p`$$`x_I`$ dependence in $``$ (see Eqs.(24,III D)). In the on-site terms this dependence is of order $`p/N^2,x_I/N^21`$ and can be therefore safely neglected (Eqs.(57,58)). In the hopping terms, it scales like $`\mathrm{exp}(\sigma _\varphi ^2/\alpha )`$ ($`\alpha =2,1/2`$, see Eqs.(59,62)): it is strong only when the hopping terms are already small and negligible in comparison with the on-site terms, so it can also be neglected.
After these considerations, we conclude that the $`z_0`$-dynamics is, in a good approximation and in reasonable regimes, independent of the $`p`$-dynamics. This is confirmed by the numerical simulations.
#### 2 Internal dynamics
The off-diagonal blocks can be neglected in the $`p`$-dynamics if $`/z_00`$ or $`\mathrm{exp}(\sigma _Rz_0^2)0`$. If the barrier is raised starting from a condensate in the ground state, $`z_0`$ evolves almost adiabatically in the low barrier limit ($`/z_00`$); when it reaches the high barrier limit, then $`\mathrm{exp}(\sigma _Rz_0^2)0`$. Therefore, the off-diagonal terms can be neglected during the whole process.
Instead, the $`z_0\sigma _I`$ dependence in $``$ is strong. One can solve the complete coupled dynamics or substitute in $``$ the adiabatic solution for $`z_0`$ and compare the results. We will show some examples in the following and see that the difference is small.
## IV External dynamics: excitations
In the previous section, we have written the equations of motion describing the full coupled dynamics and demonstrated that if one splits a condensate starting from a ground state configuration, it is possible to decouple the internal and the external dynamics. In this section, we will use this result and study the external dynamics decoupling it from the internal one. Making use of the external static solution, in Sec. V we will discuss the static solution for the internal degrees of freedom. Finally in Sec. VI, we show the results for the internal dynamics, which are also useful as a check of the decoupling assumption.
The fact that the internal and external dynamics decouple under the condition discussed above, allows us to estimate the excitation of the collective modes using a very simple model. We point out that in this case, similarly to the Gross-Pitaesvkii equation, the external dynamics is governed by the product $`gN`$ and not by the two quantities separately.
We choose to raise the potential barrier with the following time dependence
$`V_0(t)={\displaystyle \frac{V_{0fin}}{2}}\left(\mathrm{tanh}{\displaystyle \frac{t}{\tau }}+1\right).`$ (66)
From the static solution, obtained by minimizing $`(p,x_I,z_0,\sigma _I,V_0)`$ with respect to all variational parameters simultaneously at fixed $`V_0`$, we know the equilibrium position of $`z_0`$ at any value of the barrier. The $`z_0`$ dynamics can be approximated very well by the dynamics in a harmonic potential whose center moves from $`z_{0in}`$ to $`z_{0fin}`$ following the adiabatic solution and whose frequency changes corresponding to the frequency of the small oscillations.
To get analytic estimations, we model this dynamics by fixing the frequency and shifting the center of the potential along a trajectory $`z_c(t)`$ for which an analytic solution exists. When the center follows an hyperbolic tangent trajectory with time constant $`\tau `$, the semi-amplitude of the oscillation is given by
$`\delta z_0=`$ (67)
$`\left(z_{0fin}z_{0in}\right){\displaystyle \frac{\pi \omega \tau }{4}}\mathrm{csch}\left({\displaystyle \frac{\pi \omega \tau }{4}}\right)\mathrm{sech}\left({\displaystyle \frac{\pi \omega \tau }{4}}\right).`$ (68)
Otherwise, if the center moves according to a linear ramp with time constant $`\tau `$, the average semi-amplitude of the oscillation is
$`\delta z_0`$ $`=`$ $`\sqrt{2}{\displaystyle \frac{z_{0fin}z_{0in}}{\omega \tau }}.`$ (69)
We compare in Fig.2 these expressions with the numerical results. For fast raising of the barrier the hyperbolic tangent shift gives a good estimate. For slower raising, instead, the amplitude of the oscillations is largely underestimated. In this case, the linear shift of the center is useful to give an upper bound. We have checked numerically (by changing the frequency of the harmonic potential in time according to the adiabatic solution) that the small discrepancy between the actual shift of the center and the hyperbolic tangent dependence is enough to produce such a big change in the amplitude of the final oscillations. The reason for this might be better understood comparing the time derivative of $`z_c(t)`$ in the two cases (see inset in Fig.2(a)). Our estimations are qualitative, but allow us to set some lower and upper bounds and deduce an adiabaticity condition $`\tau 1/\omega _z`$ for the external degrees of freedom. In the same way we can estimate the extra energy per particle due to the excitation of collective modes. Using the relation $`K_BT_cN\mathrm{}\omega _z`$ (the trapping frequency at the end of the process is $`\omega _z`$ and not $`\omega _T`$), we can estimate which is the fastest time scale which does not destroy the condensate. As expected, it is $`\tau 1/\omega _z`$.
## V Internal dynamics: static solution
In this section, we will study the complete splitting process, starting from a condensate trapped in a harmonic potential and ending with a potential barrier with height $`V_0`$. Apart from presenting the numerical results obtained by solving Eqs.(47), we will also introduce a simple two–mode model that allows a deeper understanding of the results obtained by the variational ansatz.
### A Two–mode model
Given the fact that the external dynamics is basically independent of the internal one, we can derive a simple model that accounts for most of the effects related to the internal dynamics. The Hamiltonian in Eq.(6) depends on the following overlap integrals
$`K_{ij}`$ $`=`$ $`{\displaystyle \phi _i^{}(z)\frac{p^2}{2M}\phi _j(z)𝑑z},`$ (71)
$`V_{ij}`$ $`=`$ $`{\displaystyle \phi _i^{}(z)V(z)\phi _j(z)𝑑z},`$ (72)
$`U_1`$ $`=`$ $`{\displaystyle |\phi _i(z)|^4𝑑z},`$ (73)
$`U_2`$ $`=`$ $`{\displaystyle |\phi _i(z)|^2|\phi _j(z)|^2𝑑z}=`$ (74)
$`=`$ $`{\displaystyle \phi _i^{}(z)\phi _i^{}(z)\phi _j(z)\phi _j(z)𝑑z},`$ (75)
$`U_3`$ $`=`$ $`{\displaystyle |\phi _i(z)|^2\phi _i^{}(z)\phi _j(z)𝑑z},`$ (76)
where $`i,j=1,2`$ . We use
$`a_1^{}a_1^{}a_1a_2+a_1^{}a_2^{}a_2a_2=`$ (77)
$`=\left(a_1^{}a_11+a_2^{}a_2\right)a_1^{}a_2Na_1^{}a_2,`$ (78)
to define two effective single-particle hopping terms $`J_{12}=K_{12}V_{12}gNU_3=K_{21}V_{21}gNU_3`$. The static solution for the external dynamics is known from the minimization of $`(p,x_I,z_0,\sigma _I,V_0)`$. Plugging the corresponding solution $`z_0(V_0)`$ in Eqs.(V A), we get Hamiltonian parameters depending only on the barrier height $`V_0`$ (in particular $`J_{12}=J_{21}=J`$). Neglecting constant terms, we write the simplified Hamiltonian
$`\widehat{H}`$ $`=`$ $`{\displaystyle \frac{1}{2}}gU_1\left[a_1^2a_1^2+a_2^2a_2^2\right]J\left[a_1^{}a_2+a_2^{}a_1\right]+`$ (79)
$`+`$ $`2gU_2a_1^{}a_2^{}a_1a_2+{\displaystyle \frac{1}{2}}gU_2\left[a_1^2a_2^2+a_2^2a_1^2\right].`$ (80)
As it is well–known (see App. A), under certain conditions we can replace this model Hamiltonian by a phase model of the form
$`\widehat{H}`$ $`=`$ $`g\left(U_12U_2\right){\displaystyle \frac{^2}{\varphi ^2}}JN\mathrm{cos}\varphi +`$ (82)
$`+g{\displaystyle \frac{N^2}{4}}U_2\mathrm{cos}2\varphi ,`$
where $`\varphi `$ represents the relative phase between the two modes. The overlap integral $`U_2`$ may be non-negligible at the beginning of the process when the two mode functions overlap in a sensible way. Instead at the end, when the two condensates are almost spatially separated it is very small. In this case $`U_20`$ and one recoveres the Josephson’s Hamiltonian . The ground state of such Hamiltonian is a localised wavefunction for $`JNgU_1`$ (corresponding to a broad number distribution) and a delocalised one for $`JNgU_1`$ (corresponding to a narrow number distribution).
### B Static solution and check of the Gaussian ansatz
For this simplified two–mode model, we write analytic approximated espressions for the static solution and check the validity of the Gaussian ansatz. As explained above, we let the parameters $`J`$, $`U_1`$, $`U_2`$ depend on the barrier height $`V_0`$, according to the static solution. The expectation value of the Hamiltonian can be now written as a function of the internal degrees of freedom only, $`(p,x_I,V_0)`$, and allows to study the coherence properties of the ground state at the different stages of the splitting process: for increasing barrier height, we expect the fluctuations in the number distribution to become smaller.
We solved the static problem numerically, finding the minimum of $``$ with respect to $`p`$ and $`x_I`$ for fixed $`V_0`$. Moreover, in the limits of large $`p`$ ($`\mathrm{exp}(\sigma _\varphi ^2/2)1`$) and small $`p`$ ($`p0`$) it is possible to get analytic estimation for the value of $`p`$ at equilibrium and the frequency $`\omega _p`$ of the small oscillations as a function of the overlap integrals . In the large $`p`$ limit, one finds
$`p_s={\displaystyle \frac{N}{4}}\sqrt{{\displaystyle \frac{2(JgNU_2)}{2JgNU_2+gN(U_12U_2)}}},`$ (84)
$`\omega _p=`$ (85)
$`2\sqrt{2(JgNU_2)[2JgNU_2+gN(U_12U_2)]};`$ (86)
in the small-$`p`$ limit ($`JNgU`$), where the continuum approximation for $`m`$ is not valid, $`p_s`$ and $`\omega _p`$ can be calculated with perturbation theory, considering the Hamiltonian in Eq.(79) with $`J=U_2=0`$ as unperturbed Hamiltonian and the number state $`|m=0`$ as unperturbed ground state. Then we get
$`p_s`$ $`=`$ $`\left(4\mathrm{log}{\displaystyle \frac{2gU_1}{JN}}\right)^1,`$ (88)
$`\omega _p`$ $`=`$ $`gU_1.`$ (89)
In Fig. 3 we plot the two analytic solutions for $`p_s`$, comparing them with the numeric solution for the minimum of $``$ using our ansatz and the corresponding value of $`p`$ obtained by minimizing the exact Hamiltonian for $`N=200`$ and $`g=5`$. The agreement between the variational solution and exact one is excellent and the analytic expressions interpolate correctly in the limits where they are expected to work.
Given the analytic espressions in Eqs.(V B,V B), we observe that the oscillation frequency for large $`p`$ coincides with the breathing frequency in the harmonic potential approximating the cosine potential, and the oscillation frequency for small $`p`$ corresponds to the revival time when the cosine potential is negligible (see Sec. VI A). Aware of the fact that it is not possible to define a transition point between the two regimes, being it a smooth transition, we still calculate the value of $`p`$ at which the two frequencies coincide, get $`p=0.125`$ and claim that the phase relation between the two condensates is smeared out for $`p<0.125`$.
## VI Results: different regimes
In this section, we show some numerical results obtained by integrating the equations of motion for the variational parameters. In those results the full coupled dynamics of the process were considered. It is possible to compare with the evolution of the phase distribution governed by the Hamiltonian (82) with time–dependent coefficients. Moreover, it is possible to get analytic estimates concerning the typical time scales of the process.
In the three cases presented below we will fix the product $`gN`$ in order to have similar external dynamics and better isolate the effect of different $`g`$ and $`N`$ on the phase properties of the system. At the end of the process, when the barrier has reached its final value, depending on $`N`$ and $`g`$, one can have $`gU_1JN`$ or $`JNgU_1`$ (we assume that $`U_2`$ is then negligible). The case $`gU_1JN`$ corresponds to the situation where the splitting process is completed and one expects a ground state with no well–defined relative phase. Instead in the case $`JNgU_1`$ the ground state is still characterized by a localized phase distribution, and even if the two condensates are almost spatially separated they cannot be considered as independent. In this sense, the splitting is not complete.
### A Complete splitting
We first analyse the case where in the final stage of the process, one has $`gU_1JN`$. Since in this subsection and in the following we fix the product $`gN`$, this case is obtained for relatively small $`N`$ and large $`g`$. Depending on the time scale of the process, it is possible to observe collapses and revivals or to reach a final fragmented condensate characterized by very small number fluctuations.
We assume that in the final stage of the process it is possible to neglect $`JN`$ and $`N^2U_2`$, since they depend on the overlap of the two mode functions. Then the eingenstates are the number states
$`\widehat{H}|m=gU_1m^2|m.`$ (90)
and the time evolution of the final state corresponds to
$`|\mathrm{\Phi }(t)`$ (91)
$`{\displaystyle \underset{m}{}}\mathrm{exp}\left({\displaystyle \frac{m^2}{4p_{fin}}}\right)\mathrm{exp}\left(igU_1m^2t\right)|m.`$ (92)
In the variational ansatz formulation, one can plot the constant energy trajectories for the final barrier height (see Fig. 4(a)). Then, the time evolution corresponds to $`p`$ and $`x_I`$ following one of these trajectories: $`p`$ keeps almost constant and $`x_I`$ evolves unbounded increasing linearly with time with “velocity” $`gU_1`$ (see Eq.(89)), which exactly reproduces the phase in Eq.(91).
The main features of this time evolution are the following: the width of the number distribution is constant (determined by $`p_{fin}`$) and do not evolve in time; instead the phase distribution collapses and revives: an initial distribution peaked around zero smears out, revives around $`\varphi =\pi `$ and so on. Defining the collapse time as the time $`\tau _c`$ when $`\sigma _\varphi \pi `$ and the revival time as the time $`\tau _r`$ when the original phase distribution is recovered shifted by $`\pi `$, one gets .
$`\tau _c`$ $``$ $`{\displaystyle \frac{2}{4p_{fin}gU_1}}`$ (94)
$`\tau _r`$ $`=`$ $`{\displaystyle \frac{\pi }{gU_1}}.`$ (95)
The collapse time is governed by $`p_{fin}`$, which depends in general on the barrier raising process and on the parameters $`g`$ and $`N`$.
The final width of the number distribution obtained after the raising of the barrier is completed depends on the time scale of the process. This is just given by the value $`p_{fin}`$ at which the number fluctuations are frozen. To evaluated it, we claim that as long as $`\omega _p>2\pi /\tau `$ the number fluctuations follow the static solution, and when $`\omega _p=2\pi /\tau `$, they are frozen out to a final value $`p_{fin}`$. Setting $`\omega _p=2\pi /\tau `$ in Eq.(LABEL:omegaplp) and substituting in Eq.(84), one gets
$`p_{fin}`$ $`=`$ $`{\displaystyle \frac{N}{4}}{\displaystyle \frac{1}{2[2JgNU_2+gN(U_12U_2)]}}{\displaystyle \frac{2\pi }{\tau }}`$ (96)
$``$ $`{\displaystyle \frac{1}{8gU_1}}{\displaystyle \frac{2\pi }{\tau }}.`$ (97)
Of course Eq.(96) is not valid for $`\tau \mathrm{}`$. For $`\tau >2\pi /gU_1=2\tau _r`$, we are in the adiabatic regime: $`\omega _p>2\pi /\tau `$ during all the process, and we expect to reach a completely delocalised relative phase.
We check this numerically and plot the results in Fig. 5 for several different splitting processes ($`N=2\times 10^2`$ and $`N=2\times 10^3`$). We actually find that for time scales $`\tau >2\tau _r`$ the final $`p`$ values lie in the region $`p<0.125`$. For faster time scales, the estimation in Eq.(96) was shown to be very good when the external degrees of freedom evolve adiabatically. Otherwise discrepancies can be observed (Fig. 5). It is not straighforward to estimate such discrepancies, since they depend on the exact dynamics and can be either positive or negative. Anyway, they are not striking and do not change from a situation in which the final relative phase is well defined to the opposite one.
#### 1 Collapses and revivals of the phase
Now we consider two of these splitting processes more in detail and compare quantitatively with the evolution of the phase distribution following the Hamiltonian in Eq.(82), where $`J`$, $`U_1`$ and $`U_2`$ vary in time with the barrier height according to the instantaneous static solution.
We have already mentioned that collapse and revival of the phase are predicted for two condensates with an initially good phase relation when the final tunneling coupling is negligible. Let us consider the case $`N=2\times 10^3`$, $`g=0.5`$ and $`\tau =4`$. In Fig. 6 we show the one-atom density matrix eigenvalues and the indetermination in the phase distribution $`\sigma _\varphi `$. The agreement between the results of the two simulations is perfect and our analytic estimations are also confirmed: we expect $`\tau _c8`$ and $`\tau _r50`$, which agree very well with the numerical results shown in Fig. 6.
The actual possibility of observing the revivals of the phase in an experiment is something outside our model. If they are actually destroyed by particles losses , one is left with two condensate with no phase relation, but higher number fluctations than they would have in the ground state.
#### 2 Final fragmented condensate
Another way to cut the initial condensate into two independent ones, is to raise the barrier much slower, so that the phase coherence is lost adiabatically all along the process. So, we now choose again $`N=2\times 10^3`$, $`g=0.5`$ but a longer time scale $`\tau =200`$.
The agreement between the results of the variational ansatz and the phase model is very good also in this case. The final state is characterised by much smaller number fluctations with respect to the previous case (see Fig. 7(a)) and by a complete delocalised relative phase as shown in Fig. 7(b).
From our analytic estimations, we actually expect to reach the static solution for $`\tau >2\tau _r100`$. Instead as it can be seen in Fig. 5 for $`N=2\times 10^3,\tau =200`$, this does not happen. Dealing with such small values of $`p_{fin}`$ it is very easy to get at the end very small excitations, which can be due both to the external degrees of freedom or to some excitations already present in the initial conditions. The important feature is that the final relative phase is anyway completely delocalised.
### B Incomplete splitting
We now analyse the case where in the final stage of the process, one has $`JNgU_1`$. This case is obtained for large $`N`$ and small $`g`$: in the limit of large $`N`$, even if $`J`$ might be very small, it can happen that $`JN>gU_1`$ and the cosine potential in the phase representation (Eq.(82)) is not negligible. This can be seen as a process in which the barrier is raised up to a level at which the splitting is not really completed.
In the case in which the cosine potential at the end of the raising process is still deep, so that the lowest levels can be approximated with harmonic oscillator levels spaced by $`\sqrt{2JNgU_1}`$, the time evolution follows
$`|\mathrm{\Phi }(t){\displaystyle \underset{n}{}}c_n\mathrm{exp}\left(in\sqrt{2JNgU_1}t\right)|n,`$ (98)
where $`|n`$ are the harmonic oscillator eigenstates and where the coefficients $`c_n`$ depend on the exact dynamics of the raising process. In particular, for symmetric initial conditions, the phase distribution is always symmetric and only the even eigenstates are populated. Then, the phase distribution breathes with a frequency $`2\sqrt{2JNgU_1}`$, remaining always centered at $`\varphi =0`$. Moreover we notice that in such a case, the width in the number distribution is not expected to be constant. In the variational ansatz treatment, we have to look again at the orbits in the phase space. The contour plot in Fig. 4(b) is just to show how the orbit modify from the limit of small $`N`$ to the limit of large $`N`$. So let us consider Fig. 4(c). The orbits around the minimum of $``$ represent a time evolution in which both the width of the number and phase distribution change in time. The frequency of the small oscillations around the equilibrium position can be calculted analytically for $`\mathrm{exp}(\sigma _\varphi ^2/2)1`$ (see Eq.(LABEL:omegaplp) for $`U_2=0`$) and coincides with the breathing frequency of the phase distribution in the case of the superposition of harmonic oscillator eigenstates if $`gNU_1J`$. If this condition is not satisifed, one is not in the weak coupling regime and the phase model is not valid.
The orbits in the $`p`$$`x_I`$ space are characterised by very large oscillations in $`p`$. Hence we cannot talk of frozen number fluctations. Nevertheless, with arguments similar to the ones used before, we can try to identify the orbit which describes the dynamics at the end of the process. We estimate the maximum $`p`$ value in such an orbit to be $`p_Mp_{fin}`$. The agreement with the numerical solution can be checked in Fig. 5 and it is within a factor of $`2`$. The adiabaticity condition in this case consists in requiring that the final state is superfluid, as the static solution would be. This means that the minimum value of $`p`$ corresponding to the same orbit as $`p_M`$ must be such that the phase coherence is still good. We found that a final phase coherence corresponding to a minimum of $`\lambda _+=1\beta `$ (with $`\beta 1/2`$) is reached in process with typical time scales $`\tau _\beta =2\pi /8JN\beta `$. Note that this condition is weaker than requiring $`p_{fin}=p_s`$. In fact $`\tau _\beta <\tau _s=2\pi /\sqrt{8NJgU_1}`$, since $`\beta <\sqrt{gU_1/8JN}`$, as it would correspond to the static solution. This means that we still allow even big oscillation of $`p`$ around the equilibrium value, as long as they do not destroy the phase coherence. This corresponds to a breathing of the phase distribution which never becomes completely smeared out. Moreover, while $`\tau _s`$ depends only on the product $`gN`$, the adiabatic time scale $`\tau _\beta `$ scales like 1/N, getting easier and easier to be met for large $`N`$.
#### 1 Final superfluid phase
As done before, we take now a particular case and check directly the results with the one obtained in the phase model. We choose $`N=2\times 10^5`$, $`g=0.005`$ and $`\tau =50`$ and find that the eigenvalues $`\lambda _\pm `$ oscillate with frequency $`\omega _p`$ (see Fig. 8(a)): $`\lambda _+`$ is always close to $`1`$ and $`\lambda _{}`$ close to $`0`$. This corresponds to the breathing of the phase distribution in the non negligible cosine potential or in the variational ansatz treatment to one of the orbits in Fig. 4(c). No complete smearing out of the phase is observed (see Fig. 8(b)). Those results are confirmed up to a good level by the phase model (see Fig. 8(a,b)), even if the oscillations of $`\sigma _\varphi `$ are damped, due to the anharmonicities of the cosine potential.
The conditions that preserve the superfluid state are that the final tunneling coupling $`JN`$ is comparable to the on-site interaction and that the time scale of the process is slow enough to allow final small oscillations around the equilibrium. To destroy the phase coherence in this parameter range, one should raise the barrier faster in order to get larger oscillations, or raise it higher (case described in Sec. VI A).
To summarise, in this Section we have determined analytic expressions for the adiabaticity conditions for the internal dynamics, i.e. we identified the time scale at which the barrier can be raised in order to obtain the same relative phase properties as expected for the ground state. Concerning a completed splitting process, a fragmented condensate without any relative phase (no collapses and revivals) and characterised by a very narrow number distribution is reached in raising process with time scale $`\tau >2\tau _r`$. This means that the process has to be slower for large $`N`$ and small $`g`$. Instead in an incomplete splitting process, the superfluid phase is preserved for $`\tau >2\pi /8JN\beta `$. This time scale becomes shorter for large $`N`$. Both conclusions agree with the fact that in the Gross-Pitaevskii limit ($`N\mathrm{}`$ and $`g0`$) the condensate is phase coherent.
## VII Discussion and conclusions
We have solved the two-mode model describing the splitting of a condensate by a potential barrier through a variational ansatz. We found coupling between the internal and the external dynamics of the mode functions. We have identified the regimes in which the two dynamics decouple, and have concluded that in the case of splitting starting from a condensate in the ground state, they do not influence each other in a dramatic way. Hence, the internal and external excitations created by raising the barrier can be estimated in a good approximation independently and have been characterised as a function of the interaction strength, the number of atoms and the time scale of the process.
From our analytic estimations, confirmed by numerical resuts, we were able to identify the time scales $`\tau _z`$, $`\tau _p`$ and $`\tau _\beta `$, which define the adiabatic regime for the external and internal dynamics (respectively in the case of final fragmented or phase coherent condensate)
$`\tau _z{\displaystyle \frac{1}{\omega _z}},\tau _p>2\tau _r=2{\displaystyle \frac{\pi }{gU_1}},\tau _\beta >{\displaystyle \frac{2\pi }{8JN\beta }}.`$ (99)
It is interesting to compare the adiabaticity condition for internal and external degress of freedom in the case of final fragmented condensate, when the splitting process can be consider to be completed. We normally found $`\tau _z<\tau _p`$; the case $`N=200`$, $`g=5`$ sets a boundary where internal and external degrees of freedom enter simultaneously the adiabatic regime for $`\tau 20`$ (see Fig. 5). To get $`\tau _z>\tau _p`$ one needs $`gU_1`$ to be large compared to the trapping frequency. The quantity $`gU_1`$ is similar to the derivative of the chemical potential with respect to the total number of atoms $`\mu /N`$. From a Thomas-Fermi solution one gets that either in $`1`$ or $`3`$ dimensions, it scales as a negative power of $`N`$ and a positive power of $`g`$. So one needs to have very large $`g`$ or very small $`N`$, and our model fails in both limits. Therefore, we claim that in the usual regime of many weakly interacting particles, the adiabaticity condition for the phase dynamics is more restrictive than the one for the external dynamics.
We carried out a comparison of our model with the phase model, finding a substancial very good agreement. The numerical solution of the phase model consists in the integration of a time dependent Schrödinger equation. In practice, the number of wavefunctions $`\phi _m(\varphi )=\mathrm{exp}(im\varphi )/\sqrt{2\pi }`$ that one has to use increases linearly with $`N`$, which leads to numerical problems. In this sense for large $`N`$, the variational ansatz is more convenient and has proved to give reliable results. Moreover the variational ansatz treatment allows to include the external degrees of freedom in a natural way.
In the previous results we did not take asymmetries into account. It is possible to include them in our ansatz, through the unbalance in population $`m_0`$ and through non symmetrically centered wave functions. In the case of complete splitting, one ends up with a final constant unbalance in population. The phase coherence shows the only new feature that the center of the phase distribution is now drifting with a velocity $`\mu _1\mu _2`$, where $`\mu _i`$ are the chemical potentials of the two separate condensates. A complete analysis of this case, which even consideres losses and fluctuations in the total number of particles, can be found in . Instead in the case of final phase coherent symmetric condensate, the asymmetry can destroy the phase coherence. The final unbalance in population might be so big, that the wavepacket describing the phase distribution flies above the cosine potential: depending on the “kinetic energy” $`gU_1m_0^2`$, the cosine potential may become negligible and the same features (collapses, revivals) as in the fragmented condensate case are observed.
Possible extensions of our model are the dynamic turning on of an optical lattice, where the initial harmonic trap is deformed into a many wells potential. The instantaneous version of such a process has recently been investigeted in . Another problem of great significance is the inverse process, i.e. the merging of two condensates. This could allow to refill a condensate and be an important step towards a continuous atom laser.
## ACKNOWLEDGMENTS
This work was supported by the European Union TMR network ERBFMRX-CT96-0002 and by the Austrian Science Foundation (Projekt Nr. Z30-TPH, Wittgenstein-Preis and SFB “Control and measurement of Coherent Quantum Systems”). C. M. is grateful to Y. Castin and A. Sinatra for useful discussions, and thanks E. Arimondo and G. La Rocca for comments.
## A Two–mode variational ansatz approach and phase model
Dealing with two coupled condensates, in this paper we have often talked about the number difference $`m`$ and relative phase $`\varphi `$. In this appendix, we will define number difference $`\widehat{m}`$ and relative phase $`\widehat{\varphi }`$ operators, and derive the phase model Hamiltonian (82), discussing the approximations involved.
We first define the operators
$`\widehat{L}_x`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(a_1^{}a_2+a_2^{}a_1\right),`$ (A2)
$`\widehat{L}_y`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left(a_1^{}a_2a_2^{}a_1\right),`$ (A3)
$`\widehat{L}_z`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(a_2^{}a_2a_1^{}a_1\right),`$ (A4)
such that the usual angular momentum commutation relations are fulfilled $`[\widehat{L}_i,\widehat{L}_j]=i\epsilon _{ijk}\widehat{L}_k`$. After a small amount of algebra, the exact two–mode Hamiltonian in Eq.(79) can be rewritten as
$`\widehat{H}=g(U_12U_2)\widehat{L}_z^22J\widehat{L}_x+gU_2\left(\widehat{L}_x^2\widehat{L}_y^2\right).`$ (A5)
In the subspace of fixed even total number of atoms $`N`$, the spectrum of $`\widehat{L}_z`$ is given by all integer number in the interval $`[N/2,N/2]`$. The operator $`\widehat{L}_z`$ coincides with the number difference operator $`\widehat{m}`$. We have often treated the eigenvalues $`m`$ as continuous, but in general care should be exercised. Given the phase operator $`\varphi `$ such that $`[\widehat{\varphi },\widehat{L}_z]=i`$, it is well–known that in the phase–representation $`\widehat{L}_z=\widehat{m}=i/\varphi `$ and the eigenstates of $`\widehat{L}_z`$ ($`\widehat{m}`$) with eigenvalue $`m`$ are $`\varphi |m=\mathrm{exp}(im\varphi )/\sqrt{2\pi }`$ .
The state of the system $`|\mathrm{\Phi }`$ can be in equivalent ways described as a superposition of eigenstates of $`\widehat{m}`$ ($`c_m`$ is the number distribution)
$`|\mathrm{\Phi }={\displaystyle \underset{m=N/2}{\overset{N/2}{}}}c_m|m.`$ (A6)
or by a wave function in the $`\varphi `$-representation given by
$`\mathrm{\Phi }(\varphi )=\varphi |\mathrm{\Phi }={\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle \underset{m=N/2}{\overset{N/2}{}}}c_m\mathrm{exp}(im\varphi ).`$ (A7)
Typical quantities characterizing $`|\mathrm{\Phi }`$ are
\- the width of the number distribution
$`\widehat{m}`$ $`=`$ $`{\displaystyle \underset{m}{}}m|c_m|^2,`$ (A8)
$`\sigma _m^2`$ $`=`$ $`{\displaystyle \underset{m}{}}m^2|c_m|^2\widehat{m}^2;`$ (A9)
\- the width of the phase distribution
$`\widehat{\varphi }`$ $`=`$ $`{\displaystyle _\pi ^\pi }𝑑\varphi \varphi \left|\mathrm{\Phi }(\varphi )\right|^2,`$ (A10)
$`\sigma _\varphi ^2`$ $`=`$ $`{\displaystyle _\pi ^\pi }𝑑\varphi \varphi ^2\left|\mathrm{\Phi }(\varphi )\right|^2\widehat{\varphi }^2`$ (A11)
The uncertainty relations are the same one as for angular momenta .
In order to write the Hamiltonian in the phase–representation, we make some approximations which will lead to the simple phase model Hamiltonian and discuss under which conditions such approximations are valid. We will describe the procedure only for the term $`\widehat{L}_x`$, since for the term $`\widehat{L}_x^2\widehat{L}_y^2`$ an analogous one applies. Using the raising and lowering operators $`\widehat{L}_\pm `$, we write
$`\varphi |2\widehat{L}_x|\mathrm{\Phi }={\displaystyle \underset{m}{}}c_m\varphi |\left(\widehat{L}_++\widehat{L}_{}\right)|m=`$ (A12)
$`={\displaystyle \frac{N+1}{\sqrt{2\pi }}}{\displaystyle \underset{m}{}}c_m{\displaystyle \frac{1}{2}}e^{im\varphi }`$ (A13)
$`\left[\sqrt{1\left({\displaystyle \frac{2m+1}{N+1}}\right)^2}e^{i\varphi }+\sqrt{1\left({\displaystyle \frac{2m1}{N+1}}\right)^2}e^{i\varphi }\right].`$ (A14)
We assume a narrow and centered number distribution (these assumptions will be quantified later on), so that we can expand the square roots at first order in $`(m/N)^2`$ and get
$`\varphi |2\widehat{L}_x|\mathrm{\Phi }=`$ (A15)
$`=(N+1)\mathrm{cos}\varphi \mathrm{\Phi }(\varphi )+`$ (A16)
$`+2(N+1)\left[{\displaystyle \frac{\stackrel{~}{m^2}+1/4}{(N+1)^2}}\mathrm{cos}\varphi +{\displaystyle \frac{\stackrel{~}{m}}{(N+1)^2}}\mathrm{sin}\varphi \right]+`$ (A17)
$`+o\left({\displaystyle \frac{\stackrel{~}{m^4}}{N^4}}\right)N\mathrm{cos}\varphi \mathrm{\Phi }(\varphi ),`$ (A18)
where we can roughly estimate
$`\left|\stackrel{~}{m^2}\right|^2\left|{\displaystyle \underset{m}{}}c_mm^2\varphi |m\right|^2<`$ (A19)
$`<|{\displaystyle \underset{m}{}}|c_m|m^2m^2|^2{\displaystyle \underset{m}{}}|c_m|^2m^2{\displaystyle \underset{m^{}}{}}m^2=`$ (A20)
$`=\left(\sigma _m^2+m^2\right){\displaystyle \underset{m=m\sigma _m}{\overset{m+\sigma _m}{}}}m^2\left(\sigma _m^2+m^2\right)^2\sigma _m<`$ (A21)
$`<\left(\sigma _m+|m|\right)^4\sigma _m,`$ (A22)
so that $`\left|\stackrel{~}{m^2}\right|<(\sigma _m+|m|)^2\sqrt{\sigma _m}`$ and where in a similar way $`|\stackrel{~}{m}|\left|_mc_mm\varphi |m\right|(\sigma _m+|m|)\sqrt{\sigma _m}`$. In an analogous way, the term $`\widehat{L}_x^2\widehat{L}_y^2`$ gives
$`4\left(\widehat{L}_x^2\widehat{L}_y^2\right)`$ (A23)
$`N^2\mathrm{cos}2\varphi 4\left[(\stackrel{~}{m^2}+1)\mathrm{cos}2\varphi +2\stackrel{~}{m}\mathrm{sin}2\varphi \right]`$ (A24)
$`N^2\mathrm{cos}2\varphi .`$ (A25)
All together the neglected terms have to be smaller than all the other terms in the Hamiltonian, which implies
$`\left|\stackrel{~}{m^2}\right|\left(\sigma _m+|m|\right)^2\sqrt{\sigma _m}N^2`$ (A26)
$`{\displaystyle \frac{2JgNU_2}{N}}g(U_12U_2).`$ (A27)
For $`|m|\sigma _m`$, which is the typical situation treated in this paper, from the first equation one gets $`\sigma _mN^{4/5}`$. In the opposite limit where $`|m|\sigma _m`$ instead one has
$`{\displaystyle \frac{|m|}{N}}{\displaystyle \frac{1}{\sigma _m^{1/4}}}.`$ (A28)
For $`\sigma _m0`$ this condition is not correct and becomes simply $`|m|N`$.
To summarise, in the specific case treated in this paper, where no unbalance of population between the wells was assumed, the important condition of validity for the phase model Hamiltonian is Eq.(A27), which written for $`U_2=0`$ takes the form
$`JgNU_1/2.`$ (A29)
We note that under such a condition, the ground state is characterised by $`\sigma _\varphi \sqrt{1/N}`$ or equivalently $`\sigma _m\sqrt{N/4}`$.
## B Comparison between variational ansatz and phase model
In this subsection we will compare the phase model with our variational ansatz. For simplicity we set $`U_2=0`$, but the following discussion could be repeated in the more general case. In particular for $`U_2=0`$ the phase model Hamiltonian reduces to the Josephson Hamiltonian
$`\widehat{H}=gU_1{\displaystyle \frac{^2}{\varphi ^2}}JN\mathrm{cos}\varphi .`$ (B1)
The quantities $`gU_1`$ and $`JN`$ are respectively the on-site energy splitting (charging energy) and the tunneling coupling (Josephson coupling). It describe accurately the case of high barrier, where the two condensates are almost spatially separated, leading to a negligible $`U_2`$, and are characterised by very small number fluctuations, making the Josephson Hamiltonian valid. The phase model Hamiltonian in Eq.(B1) describes the motion of a particle in a cosine potential and the classical limit is obtained for $`\sigma _m,\sigma _\varphi 0`$. In the following we discuss the classical and quantum limits comparing it directly with the variational ansatz results.
In order to allow a complete comparison with the Hamiltonian in Eq.(B1), we now describe the coefficients $`c_m`$ in Eq.(3) with four variational parameters $`p`$, $`x_I`$, $`m_0`$ and $`\phi `$
$`c_m=`$ (B2)
$`𝒩(p)\mathrm{exp}\left[\left({\displaystyle \frac{1}{4p}}+ix_I\right)(mm_0)^2\right]\mathrm{exp}[im\phi ],`$ (B3)
The variational parameters $`x_I`$ and $`\phi `$ are the variables conjugate to $`p`$ and $`m_0`$, allow the dynamic evolution and vanish in the ground state. In the limit of broad number distribution this ansatz describes a gaussian superposition of number states centered in $`m=m_0`$. The corresponding phase distribution $`\mathrm{\Phi }(\varphi )`$ is also a gaussian, centered in $`\varphi =\phi `$. Notice that in the case of symmetric splitting treated before $`m_0=0`$ and $`\phi =0`$ $`t`$.
Replacing again the sum over $`m`$ from $`N/2`$ to $`N/2`$ with an integral from $`\mathrm{}`$ to $`\mathrm{}`$ (for $`p1`$ and $`x_I`$ such that $`\sigma _\varphi \pi `$), the widths of the number and phase distribution are respectively given by Eqs.(III D) and the expectation values $`a_i^{}a_j`$ and $`a_i^2a_i^2`$ on the state $`|\mathrm{\Phi }`$ are now
$`\mathrm{\Phi }|a_{1,2}^{}a_{1,2}|\mathrm{\Phi }={\displaystyle \frac{N}{2}}m_0,`$ (B5)
$`\mathrm{\Phi }|a_{1,2}^2a_{1,2}^2|\mathrm{\Phi }={\displaystyle \frac{N^2}{4}}+pNm_0+m_0^2,`$ (B6)
$`\mathrm{\Phi }|a_1^{}a_2|\mathrm{\Phi }=\mathrm{\Phi }|a_2^{}a_1|\mathrm{\Phi }^{}=`$ (B7)
$`={\displaystyle \frac{N}{2}}\mathrm{exp}\left({\displaystyle \frac{\sigma _\varphi ^2}{2}}\right)\mathrm{exp}\left(i\phi \right)`$ (B8)
$`\left[1{\displaystyle \frac{2p}{N^2}}+{\displaystyle \frac{2}{N^2}}\left(m_0i2px_I\right)^2\right].`$ (B9)
#### a Classical limit
Let us first consider the classical limit and compare the results obtained with the variational ansatz with the phase model. We write the expectation value of the Hamiltonian (Eq.(79)) for $`U_2=0`$ in the classical limit by setting $`\sigma _m=0`$ and $`\sigma _\varphi =0`$
$`_{cl}=gU_1m_0^2JNcos\phi ;`$ (B10)
here we have assumed $`m_0N`$, neglected all terms $`o(1/N)1`$ and all terms independent of $`m_0`$ and $`\phi `$. One sees immediately that this expression coincide with the Hamiltonian in Eq.(B1), where the operators have been substituted with their expectation values ($`\widehat{m}=m_0`$ and $`\widehat{\varphi }=\phi `$). So, the analogy with the phase model in the classical limit is straighforward . In particular if $`J>0`$ the stable equilibrium position is $`\varphi =0`$; for a kinetic energy such that $`gU_1m_0^2<JN`$ the phase undergoes oscillations, otherwise if $`gU_1m_0^2>JN`$ the phase flies above the cosine potential.
#### b Quantum limit
In the quantum limit the width of the number and phase distribution start to play an important role. The comparison between the two models now is not so trivial, because we have on the one hand a wavefunction in the phase representation, $`\mathrm{\Phi }(\varphi )`$, and on the other hand $`4`$ variational parameters ($`p`$, $`x_I`$, $`m_0`$ and $`\phi `$) which follow a classical dynamics and should reproduce the features of $`\mathrm{\Phi }(\varphi )`$. We consider here the simplified case of a wavefunction $`\mathrm{\Phi }(\varphi )`$ symmetric around $`\varphi =0`$. In the variational ansatz picture, this corresponds to $`m_0=`$ and $`\phi =0`$ $`t`$. This last assumption is correct for symmetric initial conditions and preserved at all times, how can be checked esplicitly in the equations of motion.
Concerning the ground state of Hamiltonian (B1), two different regimes can be identified: for $`gU_1JN`$, the cosine potential can be neglected and the ground state is a flat wavefunction, which corresponds to a completely undefined relative phase between the two condensates. Instead for $`gU_1JN`$ the cosine potential is deep and the ground state is a localised wavefunction, which describe a state with very well defined relative phase. In the variational ansatz approach, to find the ground state we set $`x_I=0`$, $`m_0=0`$ and $`\phi =0`$, and get $`(p)=gU_1pJN\mathrm{exp}(1/8p)\left[12p/N^2\right]`$. For $`gU_1JN`$, the minimum of this Hamiltonian is $`p=N/4`$ corresponding to the Gross-Pitaevkii limit ($`\sigma _\varphi 1/\sqrt{N}`$) and if the hopping decreases $`p0`$ ($`\sigma _\varphi \pi `$), reproducing the same results as the phase model. A quantitative comparison shows perfect agreement.
The time evolution was discussed already in Sec. VI, where the analogy between the evolution of the phase distribution $`\mathrm{\Phi }(\varphi )`$ and the evolution of the variational parameters $`p`$ and $`x_I`$ was carried out. The variational ansatz is able to predict all the main features typical of the evolution of the phase distribution: collapses and revivals of the relative phase for $`gU_1JN`$ and breathing of the wavefunction in the cosine potential for $`JNgU_1`$. Moreover the frequency of the small oscillations around the equilibrium point calculated in the variational approach concides with the splitting of the energy levels of the phase model Hamiltonian both in the limit of negligible cosine potential and in the limit of deep cosine potential.
In this Appendix we have qualitatively compared the phase distribution described by the phase model with the Gaussian ansatz for the coefficients of the state $`|\varphi `$ in the number representation. We have shown that in the classical limit the expectation values $`m_0`$ and $`\phi `$ of number and phase (variational parameters indicating the centers of the respective Gaussian distributions) are governed by the classical phase model Hamiltonian. Moreover in the quantum description the parameters $`p`$ and $`x_I`$, related to the widths $`\sigma _m`$ and $`\sigma _\varphi `$, are able to reproduce the same features of the time evolution predicted by the phase model, whose limit of validity of the phase model (derived in App. A) is $`JgU_1N/2`$.
We discussed the limit of fragmented condensate ($`JNgU_1`$) and the limit of single phase coherent condensate ($`JNgU_1`$). As already mentioned, for fixed $`gN`$, the limit $`N\mathrm{}`$ corresponds to the Gross-Pitaevkii limit. Infact, unless $`J=0`$, if $`N\mathrm{}`$ the condensate is always phase coherent. For finite $`N`$, the above discussion allows to judge whether the Gross-Pitaevskii description is still valid or not.
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# Band and filling controlled transitions in exactly solved electronic models
## Abstract
We describe a general method to study the ground state phase diagram of electronic models on chains whose extended Hubbard hamiltonian is formed by a generalized permutator plus a band-controlling term. The method, based on the appropriate interpretation of Sutherland’s species, yields under described conditions a reduction of the effective Hilbert space. In particular, we derive the phase diagrams of two new models; the first one exhibits a band-controlled insulator-superconductor transition at half-filling for the unusually high value $`U_c=6t`$; the second one is characterized by a filling-controlled metal-insulator transition between two finite regions of the diagram.
1998 PACS number(s): 71.10.Pm; 71.27.+a; 05.30.-d
Metal-Insulator and Insulator-Superconductor Transitions in chain systems of interacting electrons have recently become a matter of great interest for the physics of new compounds and devices. Although many experimental data are nowadays at disposal, an important open question of this issue is to determine the Hamiltonian that could fairly describe these kinds of transitions. The task is particularly difficult just owing to the low dimensionality, which causes usual mean-field and perturbative approaches to often fail in providing reliable predictions. Fortunately, 1-dimensionality allows to exploit exact analysis techniques which can provide –although only for some particular cases– rigorous informations on the structure of the ground state and on low-energy excitations. For this reason a probative test for theoretical models is the comparison between experimental results and theoretical predictions on the ground state phase diagram.
The ground state is usually given as a function of the filling $`n`$, i.e. the number of effective carriers, and of a ‘band-parameter’, which indicates the intrinsic unit of energy of the system (its actual definition depends on the theoretical approach envisaged, see below). One can thus distinguish between filling-controlled (FC) and band-controlled (BC) transitions, according to which kind of parameter (number of carriers or energy scale respectively) is tuned to let the transition occur. Both kinds of transitions are very important in practical applications: BC transitions are relevant for instance in Vanadium oxides, where one can modify the band-width through hydrostatic pressure on the sample; FC transitions are frequent in perovskite-like materials such as $`\mathrm{R}_{1x}\mathrm{A}_x\mathrm{TiO}_3`$ (R=rare-earth ion, and A=alkaline-earth-ion), as well as in hole-doped compounds like $`\mathrm{La}_{1x}\mathrm{Sr}_x\mathrm{CuO}_{2.5}`$. In order to describe these materials, at least as far as their low energy excitations are concerned, where single band picture are often reliable, the class of extended Hubbard models provides an interesting starting point. For these models, which involve strong electronic correlations, the band-parameter is usually taken to be the on-site Coulomb repulsion ($`U`$), instead of $`w=4t`$ ($`t`$ = hopping amplitude), the latter being typical in mean-field approaches. A number of exact results have been obtained in terms of $`n`$ and $`U`$. For the ordinary Hubbard model, a FC metal-insulator-metal transition has been shown to occur at half filling ($`n=1`$) for any $`U>0`$ ($`U`$ being the on-site Coulomb repulsion); on the contrary, no BC transition takes place for $`U>0`$. More recently, some models ( and ) were solved in which a BC Insulator-Superconductor Transition occurs at half-filling at finite values of $`U>0`$, while the usual FC Metal-Insulator-Metal transition takes place for $`n=1`$ and $`U>U_c`$.
At the best of the authors’ knowledge, no detailed investigation has been devoted to either of the following issues: for the BC transitions it has not been pointed out yet what interaction terms are relevant to tune the critical value $`U_c`$ at which the transition occurs: this is quite important because $`U_c`$ can assume different values according to the chemical structure of the material. Secondly, for the FC transitions, all the above models provide an insulating state only at half filling; on the contrary, doped materials exhibit an insulating phase for a finite region of filling values. In this letter we examine the above subjects providing the exact ground state phase diagram of some 1-D electronic models. In particular, we obtain rigorous results which allow us to both discuss the dependence of $`U_c`$ of BC transitions on the Hamiltonian parameters, and to find a FC metal-insulator transition between two finite regions of the phase diagram.
We consider here the most general 1-band extended isotropic Hubbard model preserving the total spin and number $`N`$ of electrons, which reads
$``$ $`={\displaystyle \underset{j,k,\sigma }{}}[tX(\widehat{n}_{j,\sigma }+\widehat{n}_{k,\sigma })+\stackrel{~}{X}\widehat{n}_{j,\sigma }\widehat{n}_{k,\sigma }]c_{j,\sigma }^{}c_{k,\sigma }+U{\displaystyle \underset{j}{}}\widehat{n}_{j,}\widehat{n}_{j,}+{\displaystyle \frac{V}{2}}{\displaystyle \underset{j,k}{}}\widehat{n}_j\widehat{n}_k`$ (1)
$`+`$ $`{\displaystyle \frac{W}{2}}{\displaystyle \underset{j,k,\sigma ,\sigma ^{}}{}}c_{j,\sigma }^{}c_{k,\sigma ^{}}^{}c_{j,\sigma ^{}}c_{k,\sigma }+Y{\displaystyle \underset{j,k}{}}c_{j,}^{}c_{j,}^{}c_{k,}c_{k,}+P{\displaystyle \underset{j,k}{}}\widehat{n}_{j,}\widehat{n}_{j,}\widehat{n}_k+{\displaystyle \frac{Q}{2}}{\displaystyle \underset{j,k}{}}\widehat{n}_{j,}\widehat{n}_{j,}\widehat{n}_{k,}\widehat{n}_{k,},`$ (2)
In (1) $`c_{j,\sigma }^{},c_{j,\sigma }`$ are fermionic creation and annihilation operators on a 1-dimensional chain with $`L`$ sites, $`\sigma \{,\}`$, $`\widehat{n}_{j,\sigma }=c_{j,\sigma }^{}c_{j,\sigma }`$, $`\widehat{n}_j=_\sigma \widehat{n}_{j,\sigma }`$, and $`j,k`$ stands for neighboring sites. $`t`$ represents the hopping energy of the electrons (henceforth we set $`t=1`$), while the subsequent terms describe their Coulomb interaction energy in a narrow band approximation: $`U`$ parametrizes the on-site repulsion, $`V`$ the neighboring site charge interaction, $`X`$ the bond-charge interaction, $`W`$ the exchange term, and $`Y`$ the pair-hopping term. Moreover, additional many-body coupling terms have been included in agreement with : $`\stackrel{~}{X}`$ correlates hopping with on-site occupation number, and $`P`$ and $`Q`$ describe three- and four-electron interactions. In the following we shall identify the 4 physical states $`|`$, $`|`$, $`|0`$ and $`|`$ at each lattice site with the canonical basis $`e_\alpha `$ of $`^4`$, and denote $`n_{}=N_{}/L`$; $`n_{}=N_{}/L`$; $`n_o=N_o/L`$; $`n_{}=N_{}/L`$ the densities of these 4 species of physical states.
In it has been shown that, by fixing all the coupling constants of (1) to appropriate values, one can rewrite $``$ as a generalized permutator (GP) between neighboring sites (minus some constant terms). Here we add to the latter a further arbitrary term $`U_j\widehat{n}_j\widehat{n}_j`$, which is easily proved to commute with the GP. In matrix representation, the Hamiltonian (1) that we consider reads
$$H=\underset{\alpha \beta }{}\mathrm{\Pi }_{\alpha \beta }+UN_{}const.terms$$
(3)
where $`\mathrm{\Pi }_{\alpha \beta }`$ acts as a GP $`\mathrm{\Pi }`$ (see below) whenever 2 neighboring sites of the chain are occupied by $`e_\alpha `$ and $`e_\beta `$, otherwise it gives zero. The constant terms are of the form $`\overline{U}N_{}+\overline{\mu }N+\overline{c}𝕀`$, where $`\overline{U},\overline{\mu }`$ and $`\overline{c}`$ are fixed values.
The purpose of this letter is to show how to investigate the ground state phase diagram of (3) as a function of the band parameter $`U`$ and the filling of the carriers $`n`$.
Let us first recall some basic properties of the GP’s. With respect to an ordinary permutator, a generalized permutator can either permute or leave unchanged the states of the 2 neighboring sites (including a possible additional sign); explicitly:
$$\mathrm{\Pi }(e_\alpha e_\beta )=\theta _{\alpha \beta }^d(e_\alpha e_\beta )+\theta _{\alpha \beta }^o(e_\beta e_\alpha )$$
(4)
where $`\theta _{\alpha \beta }^d`$ and $`\theta _{\alpha \beta }^o`$ are two discrete valued (0, -1 or 1) functions determining on $`\mathrm{\Pi }`$ the positions and the signs of the diagonal and off-diagonal entries respectively. Also, $`\theta ^d`$ and $`\theta ^o`$ are ‘complementary’, i.e. $`|\theta _{\alpha \beta }^d|=1|\theta _{\alpha \beta }^o|`$, so that $`\mathrm{\Pi }`$ has only one non-vanishing entry for row or column. Moreover, $`\theta _{\alpha \beta }^o=\theta _{\beta \alpha }^o`$, hence $`\mathrm{\Pi }`$ is a symmetric matrix. The set of couples of subscripts $`(\alpha ,\beta )`$ for which $`\theta _{\alpha \beta }^d0`$ (resp.$`\theta _{\alpha \beta }^o0`$) is denoted by $`𝒜^d`$ (resp.$`𝒜^o`$). It is easily seen that $`𝒜^d`$ is always of the form $`𝒜^d=_i𝒮_i\times 𝒮_i`$, the $`𝒮_i`$’s being disjoint subsets of the set $`𝒮=\{1,2,3,4\}`$. By varying the functions $`\theta ^d`$ and $`\theta ^o`$ one obtains different kinds of GP’s.
In order to determine the ground state phase diagram of (3), the difficult task is to calculate the contribution to the energy of the first term. To solve this issue, one can reconduct it to a problem defined in a smaller Hilbert space. Indeed it can be seen that:
1) Under the conditions precised below, a generalized permutator between physical states is equivalent to an ordinary permutator between the so called Sutherland’s species (SS). The latter do not need to be identified with the physical species (PS) of the states. Indeed, the number of PS is determined by the nature of the problem (in our cases they are always the four $`e_\alpha `$), while the number of SS is determined by the structure of the GP entering the Hamiltonian. In particular, it may happen that different PS constitute a single SS, so that the number of the latter is $`4`$, leading to a reduction of the dimensionality of the effective Hilbert space. Through a suitable identification of the Sutherland’s species, the first term of (3) can be rewritten in the form:
$$H_0=\underset{A}{}p_AN_{AA}\underset{A>B}{}\sigma _{AB}\mathrm{\Phi }_{AB},$$
(5)
where $`p_A=\pm 1`$ determines the nature of the $`A`$-th species, even ($``$)/odd ($`𝒪`$) for $`+1/1`$; $`N_{AA}`$ is the number of neighboring sites occupied by the same species $`A`$, and $`\mathrm{\Phi }_{AB}`$ permutes objects of species $`A`$ and $`B`$ that occupy two neighboring sites, otherwise it gives zero. The $`\sigma _{AB}`$ are signs. For a given GP, the SS are to be identified through the subsets $`𝒮_i`$’s of $`𝒜^d`$. In practice, the reduction to Sutherland’s species is possible if: $`a)`$ $`\theta _{\alpha \beta }^d=p_i\alpha ,\beta 𝒮_i`$ and $`b)`$ $`\theta _{\alpha \beta }^o=\sigma _{ij},\alpha 𝒮_i,\beta 𝒮_j`$, where $`ij`$.
2) In the case where (5) has $`\sigma _{AB}`$=+1 $`A,B`$, and $`\sigma _{AB}`$ independent of $`B`$ for $`B𝒪`$ and $`A`$ (Sutherland’s Hamiltonian), it is possible to reduce the number of even species down to only one. Indeed in this case the ground state energy of hamiltonian (5) for a system with $`x`$ even species and $`y`$ odd species is equal to that of the same hamiltonian acting on a system with the same number of odd species but just one even species collecting all the previous ones (as implied by a simple extension of Sutherland’s theorem, see ).
Remark: For a given GP, the fulfillment of the conditions given at points 1) and 2) depends also on the normalization chosen to define the basis vectors. It is worth emphasizing that some GP, though apparently violating the above requirements, can be brought to fulfill them through a mere redefinition of the phase of a given physical species $`\overline{\alpha }`$, i.e. $`|e_{\overline{\alpha }}_j(1)^j|e_{\overline{\alpha }}_j`$. We shall make use of this remark in the following.
To illustrate how the above observations can be exploited, we start with a known case, the AAS model , which differs from the ordinary Hubbard model only for a correlated-hopping term ($`X=1;\stackrel{~}{X}=V=W=Y=P=Q=0`$). This model is of the form (3), with $`\overline{U}=4`$, $`\overline{\mu }=2`$, and $`\overline{c}=1`$. Its GP reduces to (5) by identifying the following two Sutherland’s species: $`A=\{|,|\}`$ (odd) and $`B=\{|0,|\}`$ (even). In this formalism the model is nothing but a free spinless fermion model, and its energy per site is given by
$$ϵ=2n_A1\frac{2}{\pi }\mathrm{sin}(\pi n_A)+(U\overline{U})n_{}\overline{\mu }n\overline{c},$$
(6)
where $`n_A=n_{}+n_{}`$. The phase diagram as a function of the filling $`n=N/L`$ and $`U`$ can be easily derived by exploiting the identity $`n_{}=(nn_A)/2`$ and minimizing $`ϵ`$ with respect to $`n_A`$ ($`n_A[0,n]`$ for $`0n1`$ and $`n_A[0,2n]`$ for $`1n2`$) and coincides with that derived in . We also notice that the model with $`X=1;\stackrel{~}{X}=2;V=W=Y=P=Q=0`$ has the same energy as AAS. In fact, using the above Remark and redefining the basis vector $`|_j(1)^j|_j`$, it can be cast in the form (3) with the same GP as AAS.
The method just outlined allows the solution of a wide class of models whose Hamiltonian has the form (3). In particular, as the AAS model displays both a BC and a FC transitions, here we apply it to the study of similar, but new, models in which further terms are included, yielding a change in the value of the parameters characterizing the transition. We consider the model with coupling constants $`X=1;\stackrel{~}{X}=(1\sigma );Y=\sigma ;P=1;Q=2`$ where $`\sigma =\pm 1`$. The resulting Hamiltonian has the form (3) with $`\overline{U}=2`$, $`\overline{\mu }=2`$, $`\overline{c}=1`$ in both cases $`\sigma =\pm 1`$. $`\mathrm{\Pi }`$ has diagonal entries characterized by the subsets $`𝒮_1=\{1,2\}`$ , $`𝒮_2=\{3\}`$ , $`𝒮_3=\{4\}`$, and off-diagonal entries $`\theta _{\alpha \beta }^o=+1`$ $`\alpha 𝒮_1,\beta 𝒮_2`$ and $`\theta _{\alpha \beta }^o=\sigma `$ $`\alpha 𝒮_1`$ or $`𝒮_2,\beta 𝒮_3`$. Both the conditions $`a)`$ and $`b)`$ to identify Sutherland’s species are thus fulfilled, and the species read: $`A=\{|,|\}`$ (which is ‘odd’ because $`\theta _{\alpha \beta }^d=1`$ if $`\alpha ,\beta 𝒮_1`$) ; $`B=|0`$ (‘even’ because $`\theta _{33}^d=+1`$) and $`C=|`$ (‘even’ because $`\theta _{44}^d=+1`$). For the case $`\sigma =+1`$ one can straightforwardly apply Sutherland’s theorem (see point 2)) to reduce the number of even species to 1, ending up with a free spinless fermion problem, where occupied sites are represented by $`A`$ and empty sites by $`B`$ and $`C`$. The ground state energy $`ϵ`$ per site has the same form as (6), in which again $`n_{}=(nn_A)/2`$, and it has to be minimized with respect to $`n_A`$. For the case $`\sigma =1`$, before using Sutherland’s theorem, one has to apply again the Remark, changing $`|_j(1)^j|_j`$. The expression of $`ϵ`$ is identical.
The phase diagram is given in fig.1; the lower region I is characterized by $`n_A=0`$, so that only doubly occupied or empty sites are present in the ground state; in this region the ground state $`|\mathrm{\Psi }_0`$ is made of the so-called (pure) eta-pairs, i.e. $`|\mathrm{\Psi }_0=(\eta _\phi ^{})^{N/2}|0`$, where $`\eta _\phi ^{}=_je^{i\phi j}c_j^{}c_j^{}=_kc_k^{}c_{\phi k}^{}`$ with pair momentum $`\phi =0,\pi `$. In the case $`\sigma =+1`$ we have 0-pairs, whereas if $`\sigma =1`$ the pairs have $`\pi `$-momentum. The latter case is particularly important because the $`\pi `$-pairs (and not $`0`$-pairs) are expected to survive as the constraint $`X=1`$ is relaxed (see ). In region II, delimited by $`U_{\text{II-III}}=24\mathrm{cos}(\pi n)`$, we have the simultaneous presence of empty ($``$), singly occupied ($``$ $``$) and doubly occupied ($``$) sites; this is called mixed region and the ground state is $`|\mathrm{\Psi }_0=(\eta _\phi ^{})^{N/2}|U=\mathrm{}`$, where $`|U=\mathrm{}`$ are the eigenstates of the $`U=\mathrm{}`$ Hubbard model. In both region I and II the ground state is superconducting, because the 2-particle reduced density matrix exhibits long range correlation, i.e. $`g(i,j)=\mathrm{\Psi }_0|c_i^{}c_i^{}c_jc_j|\mathrm{\Psi }_00`$ for $`|ij|+\mathrm{}`$. Finally, the region III-a ($`0n1`$) is made of singly occupied and empty sites; in this region the ground-state of the $`U=\mathrm{}`$ Hubbard model is eigenstate of the Hamiltonian and is metallic. The region III-b ($`1n2`$) is the particle-hole transformed of III-a, and the metallic carriers are holes. One can show that at half-filling the system is an insulator with gap $`\mathrm{\Delta }=U6`$.
With respect to the AAS model we observe that the pair-hopping term has two main effects: first it removes the degeneracy in $`\phi `$ in this region (only $`\phi =0`$ or $`\pi `$ survive, according to the sign $`\sigma `$ of $`Y`$); secondly it raises the borderline of such region upwards: in fact it can be generally shown that a pair hopping term acts as an effective attraction ($`|Y|`$) renormalizing the Coulomb repulsion $`U`$. The superconducting region II of our model is enhanced also with respect to that of the EKS model. Indeed, although the pair-hopping term is also present in the EKS model (the borderlines of region I coincide), its effect is strongly reduced near half-filling due to the Coulomb attraction term between neighboring sites ($`V=1`$), which is known to compete with the formation of on-site pairs.
As a consequence, the BC insulator-superconductor transition occurring at half-filling corresponds to the maximum critical value $`U_c^{max}=6`$, higher than for all other exactly solved models. This is important because higher values of $`U_c`$ reduce the probability that thermal fluctuations may destroy the superconducting phase.
Because of the particle-hole symmetry of the models we have considered so far, the insulating phase can exist just at half filling. In order to investigate FC transitions between finite metal-insulator regions of the phase diagram, we now discuss a simple model not particle-hole invariant, describing a competition between the $`U=\mathrm{}`$ Hubbard model (excluding doubly-occupancy), and the pair-hopping (favouring the formation of pairs), modulated by the band parameter $`U`$ (explicitly $`X=\stackrel{~}{X}=1`$; $`Y=\sigma `$; $`V=W=P=Q=0`$). It is easy to realize that (up to the application of the Remark) this model can be set in the form (3) ($`\overline{U}=2`$, $`\overline{\mu }=2`$, and $`\overline{c}=1`$). The GP is now equivalent to an ordinary permutation between the two Sutherland’s species: $`A=\{|,|,|\}`$ (which is ‘odd’ because $`\theta _{\alpha \beta }^d=1`$ if $`\alpha ,\beta 𝒮_1=\{1,2,4\}`$); $`B=|0`$ (‘even’ because $`\theta _{33}^d=+1`$). The ground state energy per site is still given by eq. (6), where now $`n_{}=nn_A`$, . The phase diagram –obtained by minimizing $`ϵ`$ at fixed $`n`$ with respect to $`n_A`$ in the range $`n/2n_A\mathrm{min}(n,1)`$– is presented in fig. 2, and exhibits again four regions. However, due to the absence of particle-hole invariance, the shape is not symmetric around half-filling.
In region I (just doubly occupied and empty sites) only the $`Y`$ and $`U`$ terms act: the model behaves like a spin isotropic XX model ($`\stackrel{~}{S}_i^+=c_i^{}c_i^{}`$, $`\stackrel{~}{S}_i^{}=c_ic_i`$) with the $`U`$ term acting as a magnetic field; it is well known that at $`U=0`$ the correlation function has a power law decay $`g(i,j)|ij|^{1/2}`$, whereas $`g`$ is not known for non-vanishing magnetic field. However, as far as $`U2`$, long-range order arises for any non-zero value of anisotropy. The borderline of this region is given by $`U_{\text{I-II}}=2\mathrm{cos}(\pi n/2)`$. Notice that this region raises up to positive values of $`U`$ for $`1n2`$. The mixed region II is entered as the double occupancy begins to decrease from its maximum value, yielding the increase of the local magnetic moment $`M_0=3/4L^1_j\mathrm{\Psi }_0|(\widehat{n}_{j,}\widehat{n}_{j,})^2|\mathrm{\Psi }_0=`$ $`3/4(2\pi ^1\mathrm{arccos}(U/2)n)`$. The value of $`n_{}`$ reaches its minimum for $`U_{\text{II-IIIa}}=2\mathrm{cos}(\pi n)`$ when $`n1`$, and for $`U_{\text{II-IIIb}}=2`$ when $`n1`$. Correspondingly, regions III-a and III-b are entered. The former is metallic, the ground-state is that of the $`U=\mathrm{}`$ Hubbard model, and the system behaves like a Tomonaga-Luttinger liquid. The most interesting feature is that region III-b is a finite insulating region. More precisely, at exactly half filling the gap is $`\mathrm{\Delta }=U2`$, while for $`1<n2`$ no empty site is present, and the model behaves like the Hubbard model in the atomic limit. Hence here the FC transition takes place between two finite regions, in analogy with experimental observations on chain hole-doped compounds. Interestingly, for the the special value $`U=2`$, our model and the $`U`$-supersymmetric model coincide. As a consequence, our ground state energy in this case is equal to that obtained in .
In this letter we have presented exact ground state phase diagrams of two electron models, and studied their BC and FC transitions. Our analysis supports the relevance of the pair-hopping term in raising the critical value of $`U`$ for BC superconducting-insulator transitions, as well as the importance of particle-hole not invariant terms in the apperance of a finite insulating region. The method we used can be implemented on all those models described by Hamiltonian (3) in which the GP verifies conditions a) and b). We stress that such GPs all correspond to integrable models , i.e. they are solutions of the Yang-Baxter Equation (consistency equation for factorizability). The Hamiltonians exhibit therefore a set of conserved quantities mutually commuting.
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# Compactness Theorems for Geometric Packings
## 1. Introduction
Given a collection $`𝒜=\{A_1,A_2,\mathrm{}\}`$ of subsets of $`^n`$, a packing of $`𝒜`$ into another set $`C^n`$ is a way of fitting each of the sets $`A_i`$ inside $`C`$ without overlap. By a positioning of a set $`A_i`$ we mean the image of $`A_i`$ under a rigid motion of $`^n`$, i.e., some combination of translations, rotations, and reflections. To avoid ambiguity about points on the boundaries of the $`A_i`$, we say more precisely that these positionings of the $`A_i`$ must be contained inside $`C`$ and that their interiors must be pairwise disjoint. One can also speak of oriented packings, where the sets $`A_i`$ may be translated and rotated but not reflected, and also translated packings, where the $`A_i`$ may be translated but neither rotated nor reflected. We also refer to a translated packing as a parallel packing, particularly when each set $`A_i`$ is a brick (a product $`[x_1,y_1]\times \mathrm{}\times [x_n,y_n]`$ of closed intervals). If the union of the repositioned sets $`A_i`$ is all of $`C`$, we call the packing a tiling of $`C`$.
It is often difficult to determine whether a particular collection $`𝒜`$ can be packed into some target set $`C`$. One representative example is the collection $`𝒜=\{A_1,A_2,\mathrm{}\}`$ where each $`A_i`$ is a rectangle of dimensions $`\frac{1}{i}\times \frac{1}{i+1}`$. Since the total area of these rectangles is 1, it is conceivable that $`𝒜`$ can tile a unit square (generally or even with a parallel tiling); but this problem, first posed by Moser (see and \[2, Section D5\]), is unsolved. One can instead ask the apparently weaker question of whether for every positive number $`\epsilon `$, the collection $`𝒜`$ can be packed inside a square of side length $`1+\epsilon `$ (see for example ). A similar situation holds with the collection $`𝒜=\{S_2,S_3,\mathrm{}\}`$ where each $`S_i`$ is a square of side length $`\frac{1}{i}`$. Conceivably this collection will tile a rectangle of area $`\frac{\pi ^2}{6}1`$ (and perhaps even one with dimensions $`(\frac{\pi ^2}{6}1)\times 1`$), but it is even unknown whether for every positive number $`\epsilon `$ the collection $`𝒜`$ can be packed into rectangles with area $`\frac{\pi ^2}{6}1+\epsilon `$. For both these problems, results of Paulhus shows that $`\epsilon `$ can at least be taken smaller than $`10^9`$.
The purpose of this paper is to show that the weaker “for every $`\epsilon `$” versions of these two packing problems are actually equivalent to the stronger tiling versions. Our methods apply in a somewhat more general setting, and we state the following two theorems as representative of what can be deduced. For the first theorem, we use the notation $`\lambda C=\{\lambda y:yC\}`$ for the homothetic expansion/dilation (or simply homothet) of $`C`$ by the constant factor $`\lambda >0`$.
###### Theorem 1.
Let $`𝒜`$ be a collection of subsets of $`^n`$, and let $`C`$ be a compact subset of $`^n`$. If for every $`\epsilon >0`$ there exists a packing of $`𝒜`$ into the homothet $`(1+\epsilon )C`$, then there exists a packing of $`𝒜`$ into $`C`$ itself. In particular, if there exist packings of $`𝒜`$ into closed balls of radius $`R+\epsilon `$ for every $`\epsilon >0`$, then there exists a packing of $`𝒜`$ into a closed ball of radius $`R`$. These statements remain true if “packing” is replaced by “oriented packing” or “translated packing”.
We remark that the collection $`𝒜`$ may have any cardinality. Of course, the hypothesis that the target set $`C`$ be compact is equivalent to $`C`$ being both closed and bounded; both of these conditions on $`C`$ are necessary. There are obvious counterexamples if $`C`$ is not required to be closed—for example, we can take $`C`$ to be the open unit disk in $`^2`$ and $`𝒜`$ to be the collection consisting solely of $`\overline{C}`$, the closure of $`C`$. The theorem also fails if $`C`$ is closed but not bounded: for example, we can again take $`𝒜`$ to consist solely of the closed unit disk in $`^2`$, and $`C`$ to be the the closed region $`\{(x,y):1x,|y|11/x\}`$.
###### Theorem 2.
Let $`𝒜`$ be a collection of subsets of $`^n`$. If there exist packings of $`𝒜`$ into bricks of volume $`V+\epsilon `$ for every $`\epsilon >0`$, then there exists a packing of $`𝒜`$ into a brick of volume $`V`$. In fact a stronger statement is true: let $`\{B_1,B_2,\mathrm{}\}`$ be a sequence of bricks in $`^n`$, with the dimensions of the $`j`$th brick $`B_j`$ being $`b_{j1}\times \mathrm{}\times b_{jn}`$. Set $`V=inf_j\{volB_j\}`$, and assume that $`volB_j>V`$ for every $`j`$. Suppose that there exists a packing of $`𝒜`$ into each brick $`B_j`$. Then there exists a packing of $`𝒜`$ into some brick $`B`$ with dimensions $`b_1\times \mathrm{}\times b_n`$, satisfying $`volB=V`$ and $`b_mlim\; sup_j\{b_{jm}\}`$ for each $`1mn`$. These statements remain true if “packing” is replaced by “oriented packing” or “translated packing”.
The equivalence of the weak and strong versions of the two packing problems mentioned in the introductory remarks follow as immediate corollaries of Theorem 2:
###### Corollary 1.
Let $`𝒜`$ be the collection of rectangles of dimensions $`1\times \frac{1}{2}`$, $`\frac{1}{2}\times \frac{1}{3}`$, $`\frac{1}{3}\times \frac{1}{4}`$, $`\frac{1}{4}\times \frac{1}{5}`$, …. Suppose that for every $`\epsilon >0`$, the collection $`𝒜`$ can be packed into a square of area $`1+\epsilon `$. Then $`𝒜`$ tiles a square of area 1. If the given packings are parallel packings, then $`𝒜`$ parallel-tiles a square of area 1.
###### Corollary 2.
Let $`𝒜`$ be the collection of squares of side lengths $`\frac{1}{2}`$, $`\frac{1}{3}`$, $`\frac{1}{4}`$, …. Suppose that for every $`\epsilon >0`$, the collection $`𝒜`$ can be packed into a rectangle of area $`\frac{\pi ^2}{6}1+\epsilon `$. Then $`𝒜`$ tiles a rectangle of area $`\frac{\pi ^2}{6}1`$. If the given packings are into rectangles of height 1, then $`𝒜`$ tiles a rectangle of dimensions $`1\times (\frac{\pi ^2}{6}1)`$. In either case, if the given packings are parallel packings, then $`𝒜`$ parallel-tiles the resulting rectangle of area $`\frac{\pi ^2}{6}1`$.
The aforementioned work of Paulhus makes a convincing argument that the “for every $`\epsilon `$” versions of these two packing questions have affirmative answers (since obstacles to finding rectangle tilings generally arise from the largest rectangles). In light of Corollaries 1 and 2, it therefore seems likely that tilings (indeed, parallel tilings) do exist in both cases.
As can be inferred from the title of this paper, the methods used to establish Theorems 1 and 2 are topological in nature. The intuitive idea is to convert a sequence of packings of the collection $`𝒜`$ in the hypothesized sets into a “limiting packing” of $`𝒜`$ into the desired target set. To this end, we will show how the set of packings of $`𝒜`$ can be naturally regarded as a topological space, and then use a compactness argument to show the existence of a “limiting packing” of some sort; it then remains to show that this packing is a valid packing into the type of set required by Theorem 1 or 2.
In Section 2 we set the notation to be used throughout this paper and exhibit simple properties of the defined objects that follow easily from elementary point-set topology. Section 3 contains the proofs of Theorems 1 and 2, modulo an important proposition whose proof will be deferred until Section 4 in order to clarify the issues involved in the proofs of the theorems themselves. In Section 5 we remark on some modified versions of Theorems 1 and 2 that can be proved using these methods, without going into the details of the proofs.
## 2. Notation and Basic Topological Facts
The methods that we use are valid for collections $`𝒜`$ of subsets of $`^n`$ of any cardinality, but for the sake of notational simplicity we work under the assumption that our collection $`𝒜=\{A_1,A_2,\mathrm{}\}`$ is countably infinite. In addition, we argue throughout with the understanding that we are allowing translations, rotations, and reflections and thus permitting the most general kinds of packings; at the beginning of Section 5 we will explain how our arguments extend to the more restrictive classes of oriented packings and parallel packings.
For any subset $`C`$ of $`^n`$, we denote by $`𝒫(𝒜,C)`$ the set (possibly empty a priori) of all packings of $`𝒜`$ into $`C`$. We mention at the outset that translated copies of the target space $`C`$ are equivalent to each other for the purposes of deciding whether there exists a packing of $`𝒜`$ into $`C`$—indeed, there is a natural bijection between the set of packings of $`𝒜`$ into $`C`$ and the set of packings of $`𝒜`$ into some translated copy of $`C`$. Similarly, we may modify the collection $`𝒜`$ by replacing each set $`A_i`$ by any translated copy of $`A_i`$, and still retain in essence the same set $`𝒫(𝒜,C)`$. For instance, it will often be convenient for us to assume that each set $`A_i`$ contains the origin in $`^n`$. We also note that if $`C`$ is a subset of $`D`$ then certainly $`𝒫(𝒜,C)𝒫(𝒜,D)`$.
Let $`O(n)`$ denote the $`n`$-dimensional orthogonal group, i.e., the set of all $`n\times n`$ matrices $`\theta `$ with real entries such that $`\theta ^1=\theta ^T`$. Every rigid motion of $`^n`$ can be identified with an element of the product space $`O(n)\times ^n`$ as follows: if $`\sigma =(\theta ,\xi )`$ is an element of $`O(n)\times ^n`$, then $`\sigma `$ acts on a point $`x`$ of $`^n`$ by the rule $`\sigma (x)=\xi +\theta x`$. (Throughout this paper we will maintain the notational conventions that elements of $`O(n)\times ^n`$ will be denoted by $`\sigma `$ or $`\tau `$, and that $`\theta `$ and $`\xi `$ will denote the $`O(n)`$\- and $`^n`$-components, respectively, when it is necessary to refer to these components separately.) Certainly these rigid motions $`\sigma `$ act on subsets $`A`$ of $`^n`$ as well, and we will write $`\sigma (A)=\{\xi +\theta x:xA\}`$ for the image. Any positioning of the set $`A`$ in $`^n`$, using translations, rotations, and/or reflections, can be realized as $`\sigma (A)`$ for some element $`\sigma `$ of $`O(n)\times ^n`$.
Define the topological space $`(^n)`$ to be the product space $`(O(n)\times ^n)^{\mathrm{}}`$, and for any subset $`D`$ of $`^n`$ define the subspace $`(D)=(O(n)\times D)^{\mathrm{}}`$ of $`(^n)`$. Since every positioning of a set $`A`$ in $`^n`$ corresponds uniquely to an element $`\sigma `$ of $`O(n)\times ^n`$, the space $`(^n)`$ parametrizes all possible positionings of the collection $`𝒜`$ in $`^n`$, and certain positionings among these will correspond to packings of $`𝒜`$ into a target set $`C`$. More precisely, if $`IntA`$ denotes the interior of $`A`$, we can write
(1)
$$\begin{array}{cc}\hfill 𝒫(𝒜,C)=\{S=\{\sigma _i\}(^n):& i,\sigma _i(A_i)C;\hfill \\ & ij,Int(\sigma _i(A_i))Int(\sigma _j(A_j))=\mathrm{}\}.\hfill \end{array}$$
(In general we will let $`S`$ and $`T`$ denote elements of $`(^n)`$ or of its subsets.) As a result, the set $`𝒫(𝒜,C)`$ can be given the subspace topology induced by the product topology on $`(^n)`$. The key to the proof of Theorem 1 is to exploit this topological structure on $`(^n)`$ to show that $`𝒫(𝒜,C)`$ is a nonempty subspace under the stated hypotheses, and the proof of Theorem 2 proceeds similarly after a suitable brick $`B`$ is chosen as the ultimate target set.
We now exhibit several facts, which follow from the definitions of the above notation together with elementary point-set topology, that will be useful to us later. As a final piece of notation, let
$$\mathrm{\Delta }_r(x)=\{y^n:|yx|<r\}$$
represent the open ball in $`^n`$ of radius $`r`$ and center $`x`$.
###### Fact 1.
For any element $`\sigma `$ of $`O(n)\times ^n`$, any point $`x`$ of $`^n`$, and any positive number $`r`$, we have $`\sigma (\mathrm{\Delta }_r(x))=\mathrm{\Delta }_r(\sigma (x))`$.
This follows directly from the fact that the elements $`\sigma `$ of $`O(n)\times ^n`$ correspond to rigid motions (isometries) of $`^n`$, i.e., $`|\sigma (y)\sigma (x)|=|yx|`$ for any points $`x,y^n`$.
###### Fact 2.
Each element $`\sigma `$ of $`O(n)\times ^n`$ is a homeomorphism of $`^n`$ onto itself; in particular, $`\sigma ^1`$ is well-defined.
Certainly $`\sigma `$, being an isometry, is continuous. Moreover, it is easy to see that if $`\sigma =(\theta ,\xi )`$, then $`\tau =(\theta ^1,\theta ^1\xi )`$ is an element of $`O(n)\times ^n`$ which inverts the action of $`\sigma `$ on $`^n`$. Therefore $`\sigma `$ is continuously invertible as well, hence a homeomorphism.
###### Fact 3.
For any element $`\sigma `$ of $`O(n)\times ^n`$ and any subset $`A`$ of $`^n`$, we have $`\sigma (Int(A))=Int(\sigma (A))`$.
This is an immediate consequence of the fact that $`\sigma `$ is a homeomorphism of $`^n`$.
###### Fact 4.
Let $`D`$ be a subset of $`^n`$, and let $`\{x_n\}`$ be a sequence of points of $`^n`$, all but finitely many of which belong to $`D`$. If $`\{x_n\}`$ converges to some point $`x`$, then $`x\overline{D}`$.
###### Fact 5.
Every closed subset of a compact space is itself compact.
###### Fact 6.
In a compact topological space, every sequence has a convergent subsequence.
These three statements are simple consequences of elementary point-set topology; see for instance Munkres , Sections 2.10, 3.5, and 3.7, respectively.
###### Fact 7.
If $`C`$ is a compact subset of $`^n`$, then the space $`(C)`$ is also compact.
The orthogonal group $`O(n)`$ is compact (it is clearly bounded, since each column is a unit vector in $`^n`$ and hence each entry is at most 1 in absolute value; and it is closed since it is the preimage of the identity matrix under the continuous map $`\theta \theta ^T\theta `$). Since $`(C)=(O(n)\times C)^{\mathrm{}}`$, Fact 7 therefore follows from Tychonov’s theorem that arbitrary products of compact spaces are compact (see \[4, Section 5.1\]). The compactness of these spaces $`(C)`$ is crucial to our proofs of Theorems 1 and 2.
###### Fact 8.
If $`𝒜=\{A_1,A_2,\mathrm{}\}`$ is a collection of subsets of $`^n`$, each containing the origin, then $`𝒫(𝒜,C)`$ is a subset of $`(C)`$.
We can justify this fact as follows: if $`0A`$ and $`\sigma =(\theta ,\xi )`$, then $`\xi =\xi +\theta (0)\sigma (A)`$. Thus if $`\sigma (A)C`$, we must have $`\xi C`$. Fact 8 then follows from the definition (1) of $`𝒫(𝒜,C)`$ by applying this reasoning to each image $`\sigma _i(A_i)`$.
###### Fact 9.
If $`𝒜=\{A_1,A_2,\mathrm{}\}`$ and $`𝒞=\{C_1,C_2,\mathrm{}\}`$ are collections of subsets of $`^n`$, then $`𝒫(𝒜,_{k=1}^{\mathrm{}}C_k)=_{k=1}^{\mathrm{}}𝒫(𝒜,C_k)`$.
This follows immediately from unfolding the definitions of $`𝒫(𝒜,_{k=1}^{\mathrm{}}C_k)`$ and $`_{k=1}^{\mathrm{}}𝒫(𝒜,C_k)`$ using equation (1). In words, Fact 9 states that any packing of $`𝒜`$ into the set $`_{k=1}^{\mathrm{}}C_k`$ is simultaneously a packing of $`𝒜`$ into each set $`C_k`$.
## 3. Proofs of Theorems 1 and 2
In this section we state the following crucial proposition from which we deduce Theorems 1 and 2:
###### Proposition 1.
Let $`C`$ be a closed subset of $`^n`$, and let $`𝒜`$ be any collection of subsets of $`^n`$. Then the space $`𝒫(𝒜,C)`$ is a closed subset of $`(^n)`$.
The proof of Proposition 1, while not tricky, is somewhat long-winded, and therefore we defer it to the next section. Assuming the validity of Proposition 1, we can establish Theorems 1 and 2 by means of the following lemma:
###### Lemma 2.
Let $`𝒜=\{A_1,A_2,\mathrm{}\}`$ and $`𝒞=\{C_1,C_2,\mathrm{}\}`$ be collections of subsets of $`^n`$. For each $`k1`$ define $`D_k=_{j=k}^{\mathrm{}}C_j`$, and suppose that $`D_1`$ is bounded. If there exist packings of $`𝒜`$ into $`C_j`$ for each $`j1`$, then there exists a packing of $`𝒜`$ into the set $`_{k=1}^{\mathrm{}}\overline{D}_k`$.
The set $`_{k=1}^{\mathrm{}}\overline{D}_k`$ can be compared to the related set $`_{k=1}^{\mathrm{}}D_k`$, which is simply the lim sup of the sets $`C_j`$ (the set of all points that are contained in infinitely many of the $`C_j`$). In fact, $`_{k=1}^{\mathrm{}}\overline{D}_k`$ is precisely the set of all points $`x^n`$ such that every neighborhood of $`x`$ intersects infinitely many of the $`C_j`$.
Proof: By translating the sets $`A_i`$ if necessary, we may assume that each $`A_i`$ contains the origin. By hypothesis, there exists a packing of $`𝒜`$ into each $`C_j`$, so we may choose
$$T_j𝒫(𝒜,C_j)𝒫(𝒜,\overline{D}_j)𝒫(𝒜,\overline{D}_1)$$
for each $`j1`$. The set $`\overline{D}_1`$ is closed and bounded, hence compact, and so by Fact 7 the space $`(\overline{D}_1)`$ is also compact. Since the sets $`A_i`$ all contain the origin, the space $`𝒫(𝒜,\overline{D}_1)`$ is contained in $`(\overline{D}_1)`$ by Fact 8; we know by Proposition 1 that $`𝒫(𝒜,\overline{D}_1)`$ is a closed set, and so it is itself compact by Fact 5. Therefore by Fact 6, the sequence $`\{T_j\}`$ of points in $`𝒫(𝒜,\overline{D}_1)`$ has a convergent subsequence. By replacing the sequence $`\{T_j\}`$ by this subsequence, we may assume that the $`T_j`$ converge to some element $`T𝒫(𝒜,\overline{D}_1)`$.
It remains to show that this element $`T`$ in fact represents a packing of $`𝒜`$ into $`_{k=1}^{\mathrm{}}\overline{D}_k`$. For each $`k1`$, the sequence $`T_j`$ is contained (except for at most the first $`k1`$ terms) in $`𝒫(𝒜,\overline{D}_k)`$. Since this set is closed by Proposition 1, we see by Fact 4 that the limit $`T`$ is itself an element of $`𝒫(𝒜,\overline{D}_k)`$. Because this is true for all $`k1`$, Fact 9 implies
$$T\underset{k=1}{\overset{\mathrm{}}{}}𝒫(𝒜,\overline{D}_k)=𝒫(𝒜,\underset{k=1}{\overset{\mathrm{}}{}}\overline{D}_k),$$
which establishes the lemma. ∎
Proof of Theorem 1: Since $`C`$ is compact, it is contained in some ball of radius $`R`$ centered at the origin, and therefore each set $`(1+\frac{1}{j})C`$ is contained in the ball of radius $`2R`$ around the origin. Therefore under the hypothesis that there exist packings of $`𝒜`$ into each set $`(1+\frac{1}{j})C`$, we may apply Lemma 2 to conclude that there exists a packing of $`𝒜`$ into the set $`_{k=1}^{\mathrm{}}\overline{D}_k`$, where we have put
(2)
$$D_k=\underset{j=k}{\overset{\mathrm{}}{}}(1+\frac{1}{j})C.$$
All that remains to establish the theorem is to show that $`_{k=1}^{\mathrm{}}\overline{D}_k`$ is contained in $`C`$; in other words, we need to show that for every $`xC`$, there exists some $`k1`$ such that $`x\overline{D}_k`$.
If $`xC`$ then, since $`C`$ is compact (hence closed), there exists a positive number $`\epsilon `$ such that $`\mathrm{\Delta }_\epsilon (x)C=\mathrm{}`$. We claim that
(3)
$$\text{for every }j>2|x|\epsilon ^1,\mathrm{\Delta }_{\epsilon /2}(x)(1+\frac{1}{j})C=\mathrm{}.$$
To see this, suppose that there did exist a point $`y`$ in $`\mathrm{\Delta }_{\epsilon /2}(x)(1+\frac{1}{j})C`$. Since $`y(1+\frac{1}{j})C`$, if we set $`z=(1+\frac{1}{j})^1y`$ then $`zC`$, and by our choice of $`\epsilon `$ we therefore have $`|xz|\epsilon `$. On the other hand, since $`y\mathrm{\Delta }_{\epsilon /2}(x)`$,
$$|xz||xy|+|yz|<\frac{\epsilon }{2}+|y(1+\frac{1}{j})^1y|=\frac{\epsilon }{2}+\frac{|y|}{j+1}.$$
The fact that $`y\mathrm{\Delta }_{\epsilon /2}(x)`$ forces $`|y|<|x|+\epsilon /2`$, and so
$$|xz|<\frac{\epsilon }{2}+\frac{|x|+\epsilon /2}{j+1}<\frac{\epsilon }{2}+\frac{|x|+\epsilon /2}{2|x|/\epsilon +1}=\epsilon $$
by our choice of $`j`$. This contradiction establishes equation (3).
If we set $`k=2|x|\epsilon ^1+1`$, we see from equation (3) and the definition (2) of $`D_k`$ that $`\mathrm{\Delta }_{\epsilon /2}(x)D_k=\mathrm{}`$, which implies that $`x\overline{D}_k`$ as desired. This establishes the theorem.∎
Proof of Theorem 2: First we make some reductions in the problem. By translating each set $`A_i`$ if necessary we may assume that each $`A_i`$ contains the origin. Similarly, by translating each brick $`B_j`$ if necessary, we may assume that each $`B_j`$ is contained in the positive orthant of $`^n`$ and has one vertex at the origin, that is, $`B_j=[0,b_{j1}]\times \mathrm{}\times [0,b_{jn}]`$. Next, by passing to a suitable subsequence of the $`B_j`$, we may also assume that $`volB_j`$ decreases monotonically to $`V`$. At this point we make the assumption that the dimensions $`b_{jm}`$ of the bricks $`B_j`$ are bounded uniformly in $`j`$ and $`m`$; at the end of the proof we will show why this assumption is legitimate. By passing once again to a suitable subsequence of the $`B_j`$, we may therefore assume that for each $`1mn`$ the sequence $`\{b_{jm}\}`$ converges to some number $`b_m`$, say.
Since the $`b_{jm}`$ are uniformly bounded, the sets $`B_j`$ are all contained in a single bounded region of $`^n`$, and thus we may apply Lemma 2 to conclude that there exists a packing of the set $`𝒜`$ into $`_{k=1}^{\mathrm{}}\overline{D}_k`$, where we have put $`D_k=_{j=k}^{\mathrm{}}B_j`$. The theorem will therefore be established if we can demonstrate that the intersection $`_{k=1}^{\mathrm{}}\overline{D}_k`$ is contained in the brick $`B=[0,b_1]\times \mathrm{}\times [0,b_n]`$. For any natural numbers $`k`$ and $`m`$ with $`1mn`$, define $`d_{km}=sup_{jk}\{b_{jm}\}`$. Then for $`jk`$ it is clear that $`B_j`$ is contained in the closed set $`[0,d_{k1}]\times \mathrm{}\times [0,d_{kn}]`$, and so $`\overline{D}_k`$ is contained in the same closed set. Consequently,
$$\begin{array}{cc}\hfill \underset{k=1}{\overset{\mathrm{}}{}}\overline{D}_k& \underset{k=1}{\overset{\mathrm{}}{}}\left([0,d_{k1}]\times \mathrm{}\times [0,d_{kn}]\right)\hfill \\ & =[0,inf_k\{d_{k1}\}]\times \mathrm{}\times [0,inf_k\{d_{kn}\}]\hfill \\ & =[0,lim\; sup_j\{b_{j1}\}]\times \mathrm{}\times [0,lim\; sup_j\{b_{jn}\}]\hfill \\ & =[0,b_1]\times \mathrm{}\times [0,b_n]=B.\hfill \end{array}$$
This establishes the theorem, modulo the assumption that the $`b_{jm}`$ are uniformly bounded. This assumption does not hold for a general collection of bricks of bounded volume, as the simple example $`[0,n]\times [0,1/n]`$ in $`^2`$ demonstrates. However, in the most natural case—where at least one of the sets $`A_i`$ has nonempty interior—we will be able to deduce from the existence of a packing of $`𝒜`$ into each brick $`B_j`$ that the $`b_{jm}`$ are uniformly bounded. In the contrary (less interesting) case, it will also be possible to reduce to the situation where the $`b_{jm}`$ are uniformly bounded by a somewhat different method.
Case 1. At least one of the sets $`A_i`$ has nonempty interior.
Choose an integer $`k`$ such that the set $`A_k`$ has nonempty interior, and then choose $`\eta >0`$ such that $`A_k`$ contains some open ball of radius $`\eta `$. Since there exists a packing of $`𝒜`$ into each brick $`B_j`$, we see in particular that each $`B_j`$ contains some open ball of radius $`\eta `$. Certainly then the dimensions $`b_{j1},\mathrm{},b_{jn}`$ of each brick $`B_j`$ must satisfy $`b_{jm}\eta `$ for each $`1mn`$, and so for each $`j1`$ and $`1mn`$,
$$0<b_{jm}=\frac{volB_j}{b_{j1}\mathrm{}b_{j,m1}b_{j,m+1}\mathrm{}b_{jn}}\frac{volB_1}{\eta ^{n1}},$$
since we have reduced to the case where the $`volB_j`$ are monotonically decreasing. This shows that the $`b_{jm}`$ are indeed uniformly bounded.
Case 2. All of the $`A_i`$ have empty interiors.
We claim that if there exists a packing of $`𝒜`$ into each brick $`B_j=[0,b_{j1}]\times \mathrm{}\times [0,b_{jn}]`$, then there also exists a packing of $`𝒜`$ into the smaller brick $`B_j^{}=[0,b_{j1}^{}]\times \mathrm{}\times [0,b_{jn}^{}]`$ where we have defined $`b_{jm}^{}=\mathrm{min}\{b_{jm},diamB_1\}`$. If we can justify this assertion, the theorem is established in this case as well since the $`b_{jm}^{}`$ are certainly uniformly bounded by $`diamB_1`$.
For a collection $`𝒜`$ of sets with empty interiors, the packing condition that the positionings of the sets $`A_i`$ must have disjoint interiors is no condition at all; in other words, there exists a packing of the entire collection $`𝒜`$ into $`C`$ if and only if there exists individual positionings of each set $`A_i`$ into $`C`$. Moreover, we can modify any positioning $`\sigma _i(A_i)`$ into the brick $`B_j`$ so that it becomes a positioning of $`A_i`$ into $`B_j^{}`$, by taking the rotated/reflected set $`\theta _i(A_i)`$ and translating it just enough to lie the positive orthant of $`^n`$. More precisely, if $`\sigma _i=(\theta _i,\xi _i)`$ is such that $`\sigma _i(A_i)B_j`$, then we define $`\sigma _i^{}=(\theta _i,\xi _i^{})`$ where the $`m`$th coordinate $`\xi _{im}^{}`$ of the vector $`\xi _i^{}^n`$ is given by
$$\xi _{im}^{}=\left|inf\{t\pi _i(\theta _i(A_i))\}\right|;$$
here $`\pi _i`$ denotes the projection map in the $`i`$th coordinate from $`^n`$ to $``$.
The fact that $`\sigma _i^{}(A_i)`$ is contained in the positive orthant of $`^n`$ follows immediately from the definition of the $`\xi _{im}^{}`$. Also, we are assuming that $`A_i`$ contains the origin, and so $`\xi _i`$ is an element of $`\sigma _i(A_i)`$; since $`\sigma _i(A_i)`$ is contained in the positive orthant, it follows that $`\xi _{im}^{}\xi _{im}`$, and consequently $`\sigma _i^{}(A_i)`$ is contained in the brick $`B_j`$. Finally, since $`A_i`$ contains the origin it is clear that $`\xi _{im}^{}diamA_i`$, and since there exists a packing of $`𝒜`$ into $`B_1`$ we certainly have $`diamA_idiamB_1`$. Therefore $`\sigma _i^{}(A_i)`$ is indeed contained in the brick $`B_j^{}`$.
Making this modification for each set $`A_i`$ results in a packing of the entire collection $`𝒜`$ into the smaller brick $`B_j^{}`$ (again, the assumption that the $`A_i`$ have empty interiors means that we do not need to worry about the relative positionings of the various $`A_i`$). As remarked earlier, this justifies the assumption that the dimensions of our bricks are uniformly bounded, since we may replace $`B_j`$ by $`B_j^{}`$ throughout.
This completes the proof of the theorem. ∎
In summary, we have established Theorems 1 and 2 modulo a proof of Proposition 1; this proof will be the subject of the following section.
## 4. Proof of Proposition 1
Proposition 1 is essentially a consequence of the fact that the action on $`^n`$ of the space of rigid motions $`O(n)\times ^n`$ is continuous. The following two lemmas, which give concrete statements of the continuity of this action, will enable us to establish Proposition 1. We note that the space $`O(n)\times ^n`$ can in fact be regarded as a metric space, inheriting as it does the standard metric from $`^{n^2}\times ^n`$: if $`\sigma =(\theta ,\xi )`$ and $`\sigma ^{}=(\theta ^{},\xi ^{})`$ are two elements of $`O(n)\times ^n`$, then the distance between them is
(4)
$$d(\sigma ^{},\sigma )=\left(|\theta ^{}\theta |^2+|\xi ^{}\xi |^2\right)^{1/2}=\left(\underset{l=1}{\overset{n}{}}\underset{m=1}{\overset{n}{}}(\theta _{lm}^{}\theta _{lm})^2+\underset{m=1}{\overset{n}{}}(\xi _m^{}\xi _m)^2\right)^{1/2},$$
considering $`\theta `$ and $`\theta ^{}`$ here simply as $`n^2`$-tuples of real numbers rather than elements of $`O(n)`$.
###### Lemma 3.
Let $`y`$ be a point in $`^n`$ and $`U`$ be an open subset of $`^n`$. Suppose that $`\sigma `$ is an element of $`O(n)\times ^n`$ such that $`\sigma (y)U`$. Then there exists a positive real number $`\delta `$ such that, for every $`\sigma ^{}O(n)\times ^n`$ satisfying $`d(\sigma ^{},\sigma )<\delta `$, we have $`\sigma ^{}(y)U`$.
Proof: For any $`y^n`$ and any pair $`\tau =(\theta ,\xi )`$, $`\tau ^{}=(\theta ^{},\xi ^{})`$ of elements of $`O(n)\times ^n`$, we have
(5)
$$|\tau (y)\tau ^{}(y)|=|\xi +\theta y\xi ^{}\theta ^{}y||\xi \xi ^{}|+|(\theta \theta ^{})y|,$$
We certainly have $`|\xi \xi ^{}|d(\tau ,\tau ^{})`$ by the definition (4) of the metric $`d`$. On the other hand, all entries of the matrix $`\theta \theta ^{}`$ are also at most $`d(\tau ,\tau ^{})`$ in absolute value, while the entries of the vector $`y`$ are at most $`|y|`$ in absolute value. Therefore each entry of $`(\theta \theta ^{})y`$ is bounded by $`n|y|d(\tau ,\tau ^{})`$ in absolute value, and so the inequality (5) becomes the upper bound
(6)
$$|\tau (y)\tau ^{}(y)|d(\tau ,\tau ^{})+\left(\underset{m=1}{\overset{n}{}}\left(n|y|d(\tau ,\tau ^{})\right)^2\right)^{1/2}=(n^{3/2}|y|+1)d(\tau ,\tau ^{})$$
(we have made no effort to obtain a strong constant in the inequality).
Now if $`\sigma `$ is an element of $`O(n)\times ^n`$ such that $`\sigma (y)`$ lies in the open set $`U`$, then there exists some positive number $`\epsilon `$ such that $`\mathrm{\Delta }_\epsilon (\sigma (y))U`$. If we set $`\delta =\epsilon (n^{3/2}|y|+1)^1`$, then for any $`\sigma ^{}O(n)\times ^n`$ such that $`d(\sigma ^{},\sigma )<\delta `$, the upper bound (6) tells us that
$$|\sigma ^{}(y)\sigma (y)|(n^{3/2}|y|+1)d(\sigma ^{},\sigma )<\epsilon ,$$
and therefore $`\sigma ^{}(y)\mathrm{\Delta }_\epsilon (\sigma (y))U`$ as desired. ∎
###### Lemma 4.
Let $`U_1`$ and $`U_2`$ be open subsets of $`^n`$. Suppose that $`\sigma _1`$ and $`\sigma _2`$ are elements of $`O(n)\times ^n`$ such that $`\sigma _1(U_1)\sigma _2(U_2)\mathrm{}`$. Then there exists a positive real number $`\delta `$ such that, for every $`\sigma _1^{},\sigma _2^{}O(n)\times ^n`$ satisfying $`d(\sigma _1^{},\sigma _1)<\delta `$ and $`d(\sigma _2^{},\sigma _2)<\delta `$, we have $`\sigma _1^{}(U_1)\sigma _2^{}(U_2)\mathrm{}`$.
Proof: Since $`\sigma _1(U_1)`$ and $`\sigma _2(U_2)`$ are open sets that are not disjoint, we can choose a point $`x^n`$ and a positive number $`\epsilon `$ such that $`\mathrm{\Delta }_\epsilon (x)\sigma _1(U_1)\sigma _2(U_2)`$. Using Fact 2 we may set $`y_1=\sigma _1^1(x)`$ and $`y_2=\sigma _2^1(x)`$, so that $`\mathrm{\Delta }_\epsilon (y_1)U_1`$ and $`\mathrm{\Delta }_\epsilon (y_2)U_2`$; we also set
$$\delta =\frac{\epsilon }{n^{3/2}\mathrm{max}\{|y_1|,|y_2|\}+1},$$
Then for $`i=1`$ or 2, for any $`\sigma _i^{}O(n)\times ^n`$ such that $`d(\sigma _i^{},\sigma _i)<\delta `$ the upper bound (6) tells us that
$$|\sigma _i^{}(y_i)x|=|\sigma _i^{}(y_i)\sigma _i(y_i)|(n^{3/2}|y_i|+1)d(\sigma _i^{},\sigma _i)<\epsilon ,$$
so that $`x\mathrm{\Delta }_\epsilon (\sigma _i^{}(y_i))=\sigma _i^{}(\mathrm{\Delta }_\epsilon (y_i))\sigma _i^{}(U_i)`$ by Fact 1. In particular, this shows that $`x`$ is an element of $`\sigma _1^{}(U_1)\sigma _2^{}(U_2)`$, which is therefore nonempty as desired. ∎
Proof of Proposition 1: Let $`T=\{\tau _i\}`$ be a point in $`(^n)𝒫(𝒜,C)`$. From the definition (1) of $`𝒫(𝒜,C)`$, one of the following two cases must hold.
Case 1. There exists a $`k1`$ such that $`\tau _k(A_k)C`$.
Choose a point $`x\tau _k(A_k)C`$, and set $`y=\tau _k^1(x)A_k`$ (using Fact 2). Applying Lemma 3 with $`\sigma =\tau _k`$ and $`U=^nC`$, we see that there exists a positive number $`\delta `$ such that, for every $`\sigma ^{}O(n)\times ^n`$ satisfying $`d(\sigma ^{},\tau _k)<\delta `$, we have $`\sigma ^{}(y)^nC`$, that is, $`\sigma ^{}(y)C`$.
Now define the open neighborhood $`𝒮`$ of $`T`$ in $`(^n)`$ by
$$𝒮=\{S=\{\sigma _i\}(^n):d(\sigma _k,\tau _k)<\delta \}.$$
For every $`S𝒮`$, we see that $`\sigma _k(y)C`$ by our choice of $`\delta `$. On the other hand, certainly $`\sigma _k(y)\sigma _k(A_k)`$, and so $`S`$ is not a packing of $`𝒜`$ into $`C`$. Since this is true for any $`S𝒮`$, we see that $`𝒮(^n)𝒫(𝒜,C)`$.
Case 2. There exist positive integers $`kl`$ such that $`Int(\tau _k(A_k))Int(\tau _l(A_l))\mathrm{}`$.
Applying Lemma 4 with $`\sigma _1=\tau _k`$, $`\sigma _2=\tau _l`$, $`U_1=Int(A_k)`$, and $`U_2=Int(A_l)`$, we see that there exists a positive real number $`\delta `$ such that, for every $`\sigma _1^{},\sigma _2^{}O(n)\times ^n`$ satisfying $`d(\sigma _1^{},\tau _k)<\delta `$ and $`d(\sigma _2^{},\tau _l)<\delta `$, we have
$$Int(\sigma _1^{}(A_k))Int(\sigma _2^{}(A_l))=\sigma _1^{}(Int(A_k))\sigma _2^{}(Int(A_l))\mathrm{}$$
(here we have used Fact 3). Now define the open neighborhood $`𝒮`$ of $`T`$ in $`(^n)`$ by
$$𝒮=\{S=\{\sigma _i\}(^n):d(\sigma _k,\tau _k)<\delta \text{ and }d(\sigma _l,\tau _l)<\delta \}.$$
For every $`S𝒮`$, we see that $`Int(\sigma _k(A_k))Int(\sigma _l(A_l))\mathrm{}`$ by our choice of $`\delta `$, and so $`S`$ is not a packing of $`𝒜`$ with disjoint interiors. Since this is true for any $`S𝒮`$, we see that $`𝒮(^n)𝒫(𝒜,C)`$.
In either case we see that $`(^n)𝒫(𝒜,C)`$ contains an open neighborhood $`𝒮`$ of $`T`$, which shows that $`(^n)𝒫(𝒜,C)`$ is an open set, i.e., $`𝒫(𝒜,C)`$ is a closed subset of $`(^n)`$. ∎
## 5. Generalizations of Theorems 1 and 2
We end by briefly discussing some extensions of Theorems 1 and 2 that can be established by the methods of this paper. First, in the statements of these two theorems we have claimed that “packings” may be replaced by “oriented packings”. This is true because the positionings allowed in oriented packings (translations and rotations, but not reflections) are parametrized by $`O(n)^+\times ^n`$, where $`O(n)^+`$ is the index-2 subgroup of $`O(n)`$ consisting of the orthogonal matrices of determinant 1. Because this subgroup $`O(n)^+`$ is a compact space in its own right, the analogous statement to Fact 7 for $`^+(C)=(O(n)^+\times C)^{\mathrm{}}`$ is also true, and thus all of the arguments of this paper go through for oriented packings upon simply replacing $`(C)`$ by $`^+(C)`$ at each occurrence. In the case of translated packings, where neither rotations nor reflections are allowed, we can similarly replace each occurrence of $`(C)`$ by $`C^{\mathrm{}}`$ and the arguments proceed unchanged (if we like, we can think of the space $`C^{\mathrm{}}`$ as $`(\{I_n\}\times C)^{\mathrm{}}`$, where $`\{I_n\}`$ is the compact subgroup of $`O(n)`$ consisting only of the identity matrix).
It is clear that many variations on Theorems 1 and 2 could be stated by changing the sequence of sets into which $`𝒜`$ can be packed. The important thing is for this sequence $`C_j`$ (which is a shrinking sequence of homothets in Theorem 1, and a sequence of bricks of varying dimensions in Theorem 2) to have enough structure for the limiting set $`_{k=1}^{\mathrm{}}\overline{D}_k`$ to be identified, where $`D_k=_{j=k}^{\mathrm{}}C_j`$ as defined in the statement of Lemma 2. This limiting set would be easy to determine if the $`C_j`$ were ellipsoids or simplices of varying dimensions, just to name two possible applications.
Finally we note two ways in which the hypotheses of Theorems 1 and 2 can be weakened. Instead of requiring that the collection $`𝒜`$ can be packed into each set $`C_j`$, we can require only that for each $`j1`$ the contracted collection $`(1\frac{1}{j})𝒜=\{(1\frac{1}{j})A_1,(1\frac{1}{j})A_2,\mathrm{}\}`$ can be packed into $`C_j`$. This is actually easily seen to be equivalent to the current statements of Theorems 1 and 2. However, we obtain genuinely stronger theorems by weakening the hypothesis in the following way: for every $`j1`$, we require only that the finite collection $`\{A_1,\mathrm{},A_j\}`$ can be packed into the set $`C_j`$. We leave the details of this variation to the reader.
Acknowledgements. The author acknowledges the support of Natural Sciences and Engineering Research Council grant number A5123. The author would also like to thank Mark Hamilton for his comments on a preliminary version of this paper.
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# de Broglie oscillation, rest mass and inertia
## 1 Introduction
Newton’s classical mechanics and gravitational theory underwent remarkable changes after the advent of Einstein’s special and general theories of relativity. In these theories a 4-dimensional picture of our physical world was presented providing deep insight of man about the nature of physical laws. In this respect, a 4-dimensional description appeared as the key feature in investigation of the fundamental physical principles. It turned out that the 4-dimensional analysis of some physical phenomena may give the precise 3-dimensional descriptions of our experiences about them. The basic element in this 4-dimensional analysis was the abstract object called world line element. The invariance property of this object under Lorentz transformations or general coordinate transformations played the key role in developing special and general theories of relativity.
Regarding the key role of the world line described above it is appealing to investigate more on this object in a rather abstract way. In this respect we are looking for some global features associated with the world line which may emerge due to an invariance symmetry called reparametrization invariance. We pay attention to the interesting role of this symmetry in investigation of the two mysteries in modern physics, namely de Broglie’s periodic phenomenon (oscillation) and the inertial properties of matter.
Louis de Broglie had presented the central point of his hypothesis on waves of matter in a remarkable periodic phenomenon by the famous Einstein-de Broglie formula $`\mathrm{}\omega _0=m_0c^2`$, where $`\omega _0`$ as the natural frequency of this phenomenon is related to the rest mass $`m_0`$ of the particle. Indeed, as is well-known, this phenomenon leads, through Lorentz transformation of the coordinate system, to de Broglie’s waves of matter. Hence the entire quantum mechanics stands on de Broglie oscillation and any attempt to find some insight into what stands behind this phenomenon is of particular interest.
On the other hand, the issue of inertia as an unresolved mystery in modern physics has recently been the subject of intense investigations. Vigier , used the Dirac vacuum to explain, alternative to Machian viewpoint , the origin of inertia as a necessary consequence of the real particle motions described by the Einstein-de Broglie-Bohm (E.d.B.B) formalism of quantum mechanics. Recently, Rueda $`etal.`$ in a series of papers have explained the inertia as scattering-like process of the zero point field (ZPF) radiation subject to the electromagnetic vacuum. They showed that the scattering of the incoming ZPF flux by the fundamental particles (quarks and leptons) within the object generates a reaction force which may account, at least in part, for inertia. Moreover, the inertial mass of an object is that fraction of the energy of the ZPF radiation enclosed within the object that interacts with it. More recently, they suggested that this interaction takes place at the Compton frequency of the particle and the inertial mass of the electron is then the reaction force due to resonance scattering of the ZPF at that frequency. Furthermore, they suggested that de Broglie oscillation is due to a resonant interaction with the ZPF, hence they were able to interrelate the two above mentioned mysteries.
In this paper we take an alternative approach to study and interrelate both of two mysteries. We seek a Hamiltonian constraint structure originating from the point particle’s action constructed by the world line element. As is well known, the action of a relativistic free massive point particle has a Hamiltonian constraint structure , leading to the mass-shell condition, where the space-time coordinates are assumed to be as dynamical variables. Indeed, this constraint structure is independent of the space-time dimensions, so one may be allowed to reconstruct an equivalent structure without the need to use the space-time coordinates as dynamical variables. The motivation to follow this approach is to derive some valuable information about a constraint dynamics which induces the local dynamics of the point particle in space-time.
We use the world line of a relativistic free massive point particle as one dynamical variable and parametrize its action by a world line parameter. We then derive the associated constraint structure. This provides a framework to justify the possible origins of de Broglie oscillation, rest mass and inertia.
de Broglie oscillation and the rest mass may naturally emerge due to the quantum constraint associated with one first class constraint (in the terminology of Dirac) appearing in the model. Indeed, using the Einstein-de Broglie formula, this quantum constraint becomes equivalent to an eigenvalue equation whose eigenstate $`e^{i\omega _0\tau }`$ corresponds to de Broglie’s localized oscillating field and its energy $`\mathrm{}\omega _0`$ corresponds to the rest mass of the point particle ($`\omega _0`$ is the Compton frequency). It is shown that what plays the role behind de Broglie oscillation and the rest mass may be this quantum constraint which is imposed on the rest point particle.
On the other hand, since the mass-shell condition as the 4-dimensional analogue of our desired first class constraint holds at all times it leads to derivation of the relativistic equation of motion and then a force in the form $`\frac{dP_i}{d\tau }`$ is interpreted physically as inertia. Although the physical mechanism behind this interpretation may be the ZPF scattering , but the causal origin of this reaction force seems to be in the reparametrization invariance of the point particle system which may be induced as a property resulting from the same invariance property in large universe scale. In this approach, the local inertial properties of the point particle are generated due to the “timeless” property of the universe and its causal structure. Inertia of the particle then becomes a direct evidence indicating that no absolute time exists in the universe.
## 2 Constraint System
In this section we will use Dirac’s formalism of Hamiltonian constraint systems . We start with the action of a relativistic free massive point particle
$$I=m_0c_{s_1}^{s_2}𝑑𝐬$$
(1)
where $`m_0`$ is the rest mass of the particle, $`c`$ is the velocity of light and $`d𝐬`$ is the world line element. We now parametrize the action as
$$I=m_0c_{\tau _1}^{\tau _2}\dot{𝐬}𝑑\tau $$
(2)
where $`\dot{𝐬}=\frac{d𝐬}{d\tau }`$ and $`\tau `$ is the world line parameter usually taken as proper time. As mentioned in the introduction, we have taken $`𝐬`$ as the only dynamical variable in order to find some global characteristics of the model. The linear action (2) admits a Hamiltonian constraint structure with the primary constraint
$$\varphi P_s+m_0c0$$
(3)
where $`P_s`$ is the momentum conjugate to $`𝐬`$. The original Hamiltonian $`H_0`$ is zero for the linear action, so the total Hamiltonian is given
$$H_T=\lambda (\tau )\varphi $$
(4)
where $`\lambda (\tau )`$ is an arbitrary function of time. Hence the dynamics is fully controlled by the constraint (3). The consistency condition
$$0\dot{\varphi }=\{\varphi ,H_T\}$$
(5)
is automatically satisfied since $`\{\varphi ,\varphi \}0`$, so that there is no secondary constraints. Now, the constraint (3) is by definition a first class constraint which generates the following infinitesimal local gauge transformations
$$\delta 𝐬=ϵ(\tau )\{𝐬,\varphi \}=ϵ(\tau )$$
(6)
$$\delta P_s=ϵ(\tau )\{P_s,\varphi \}=0$$
(7)
where $`ϵ(\tau )`$ is an infinitesimal arbitrary function of time. It follows from (6) that the variable $`𝐬`$ has a gauge orbit. On the other hand, from (7) we find that the momentum $`P_s`$ is a gauge invariant quantity on this orbit and should have a physical content. The equations of motion are
$$\dot{𝐬}=\{𝐬,H_T\}=\lambda (\tau )$$
(8)
$$\dot{P_s}=\{P_s,H_T\}=0$$
(9)
which integrate to
$$𝐬(\tau )=^\tau \lambda (\tau ^{})𝑑\tau ^{}$$
(10)
$$P_s=\text{Const.}$$
(11)
These are consistent with (6) and (7) provided that $`\delta \lambda =\dot{ϵ}(\tau )`$. With the identification
$$ϵ(\tau )=\lambda (\tau )\eta (\tau )$$
where $`\eta (\tau )`$ is defined for an infinitesimal world line parametrization as
$$\eta (\tau )=\tau \stackrel{~}{\tau }(\tau )$$
one can see that, using equations of motion (8) and (9), the infinitesimal local gauge transformations (6), (7) correspond to the world line reparametrizations. The freedom in the choice of Lagrange multiplier $`\lambda (\tau )`$ corresponds to the freedom in the choice of world line parametrization. Hence, choosing the parametrization means the gauge-fixing of the system through a choice for $`\lambda (\tau )`$. However, different functions $`\lambda (\tau )`$ leading to the same value of the action $`I`$ lie on the same gauge orbit and describe identically the same physical situation. <sup>1</sup><sup>1</sup>1 Note that as a result of equation (8) and $`\delta \lambda =\dot{ϵ}(\tau )`$ with $`ϵ(\tau _i)=0`$ we have
$$\delta I=m_0c\delta _{\tau _1}^{\tau _2}\lambda (\tau )𝑑\tau =0.$$
Therefore, the action $`I`$ is the Teichmüller parameter labelling the gauge orbit of the system. For positive energy with $`\lambda >0`$ we have forward propagation in time for $`I<0`$ (a particle) and backward propagation for $`I>0`$ (an antiparticle) .
## 3 de Broglie oscillation and rest mass
The constraint (3) as the heart of this model merits considerable attention. We know from special theory of relativity the well-known relation
$$E^2=P^2c^2+m_0^2c^4$$
(12)
which relates the energy content $`E`$ of a point particle to its spatial momentum $`P`$ and rest mass $`m_0`$. Using the constraint (3) $`|P_s|=m_0c`$ , we may rewrite (12) as <sup>2</sup><sup>2</sup>2The notation $``$ in Eq. (3) does not mean an approximate relation. It merely tells us that the constraint $`|P_s|=m_0c`$ should be imposed after derivation of the equations of motion .
$$E^2=(P^2+P_s^2)c^2$$
(13)
reflecting the fact that energy content of a point particle emerges due to two types of linear momentums, its usual spatial one $`P`$ and the intrinsic one $`P_s`$ conjugate to the world line coordinate. For a particle at rest in space-time $`P=0`$, equation (13) becomes
$$E=|P_s|c$$
(14)
which relates the energy content of the point particle to its intrinsic momentum $`P_s`$ on the world line. This is like the corresponding formula $`E=Pc`$ for massless particles (photons) for which we have
$$\mathrm{}\omega =Pc$$
(15)
where the angular frequency $`\omega `$ is related to the spatial momentum $`P`$ of these particles according to wave-particle duality. Insisting on the generality of wave-particle duality we may relate, as well, an intrinsic angular frequency $`\omega _0`$ (as the wave characteristic) to the intrinsic momentum $`P_s`$ (as the particle characteristic) of the rest point particle as
$$\mathrm{}\omega _0=|P_s|c.$$
(16)
Therefore, using $`|P_s|=m_0c`$, Einstein-de Broglie relation $`\mathrm{}\omega _0=m_0c^2`$ emerges naturally due to generalization of wave-particle duality to the rest particles <sup>3</sup><sup>3</sup>3de Broglie introduced the formula $`\mathrm{}\omega _0=m_0c^2`$ directly without mentioning to the wave-particle duality as its justification. Indeed, Einstein-de Broglie relation in the form $`\mathrm{}\omega _0=m_0c^2`$ for a rest mass does not indicate wave-particle duality in the sense of de Broglie’s waves of matter. This is because Einstein-de Broglie relation reads as the Compton wavelength $`\lambda =\frac{\mathrm{}}{m_0c}`$ and if one wants to describe it in the sense of de Broglie’s waves of matter then a rest particle should have the spatial momentum $`P=m_0c`$ which is a contradiction. In other words, although the wavelength $`\lambda `$, in the left hand side, manifests the wave aspect but the rest mass $`m_0`$, in the right hand side, by no means indicates the particle aspect. The particle aspect of the rest particles emerges whenever the intrinsic momentum $`P_s`$ is attributed to them. This gives rise to the wave-particle duality expressed as the formula $`\lambda =\frac{\mathrm{}}{P_s}`$. This formula relates the Compton wavelength $`\lambda `$ to the particle’s intrinsic momentum $`P_s`$ on the world line coordinate. . Einstein-de Broglie relation as realization of mass-energy equivalence relates the concept of rest mass to a type of energy which appears to be the energy of massless particles namely, $`\mathrm{}\omega _0`$. Therefore, intuitively, one may think that a massive rest particle as viewed by a rest inertial observer is nothing but a photon-like particle, moving with light velocity $`\lambda =c`$ on its world line, whose energy content appears, through mass-energy equivalence, as the particle’s rest mass<sup>4</sup><sup>4</sup>4Here, only the mass property of the particles are compared without referring to their other properties such as spin or charge.. This may shed light on the origin of de Broglie’s oscillation $`e^{i\omega _0\tau }`$ as it seems to be the time dependent part of a wave propagation along the world line associated with the photon-like particle. This is because a rest observer also moves with the light velocity $`\lambda =c`$ on its own world line so only the time evolving part of this wave propagation manifests.
The wave-particle duality for a rest particle becomes more apparent if we try to derive de Broglie oscillation from this constraint system. We remind that the entire quantum mechanics stands on wave-particle duality with wave functions and eigenvalues as indicating the wave and particle aspects respectively. So, if we can derive de Broglie oscillation itself from quantum mechanics then we may conclude that this periodic phenomenon follows, in principle, the wave-particle duality. To this end, we pay attention for quantization of the present model.
Quantization of this constraint system is done by operating on the Hilbert subspace $`\psi >`$ by the operator form of the constraint (3) as
$$(\widehat{P}_s+m_0c)\psi >=0.$$
(17)
Using Einstein-de Broglie relation and $`i\mathrm{}\frac{}{\tau }=c\frac{\mathrm{}}{i}\frac{}{𝐬}`$ for the rest particle obeying $`d𝐬=\lambda d\tau `$ with $`\lambda =c`$, the quantum constraint (17) becomes the eigenvalue equation
$$i\mathrm{}\frac{}{\tau }\psi >=E\psi >$$
(18)
with energy and wavefunction as
$$E=\mathrm{}\omega _0$$
$$\psi (\tau )>e^{i\omega _0\tau }.$$
(19)
Remarkably, the wave function (19) as the physical solution satisfying the quantum constraint (17) is exactly what we have known as de Broglie’s periodic phenomenon and the eigenvalue $`\mathrm{}\omega _0`$ is the corresponding energy. This wave function is the physical state describing the photon-like particle, as is observed by the rest observer. If we believe that quantum mechanics, as stands on de Broglie’s oscillation, provides us the proper observables with physical reality then by translating the quantum constraint (17) into an eigenvalue equation (18) we are able to derive the rest mass $`m_0`$, as a physical observable $`\frac{\mathrm{}\omega _0}{c^2}`$, from de Broglie’s periodic phenomenon. In other words, the rest mass $`m_0`$ itself does not qualify as an observable; it is $`\frac{\mathrm{}\omega _0}{c^2}`$ which physically we observe as the rest mass .
Moreover, the quantum constraint (17) or equivalently Eq (18) seems to describe the quantum motion well-known as Schrödinger Zitterbewegung.
## 4 Inertia
The issue of inertia is basically correlated with the phenomenon of acceleration. Suppose we have a rest point particle in an inertial reference frame. Actually, since the particle does not undergo any external force the action for a free point particle (1) is then applied and leads to the first class constraint
$$\varphi P_s+m_0c0.$$
In order to find what effects such a constraint structure generates in space-time we resort to the 4-dimensional analogue of the constraint (3) as the mass-shell condition
$$\psi P^\mu P_\mu +m_0^2c^40$$
(20)
subject to the total Hamiltonian
$$H_T=\lambda (\tau )\psi $$
where $`P_\mu `$ is the 4-momentum of the particle defined on the Minkowski space-time with metric (- + + +). Time independence of the mass-shell condition (20) $`\dot{\psi }=\{\psi ,H_T\}=0`$, means trivially
$$\frac{dP_\mu }{d\tau }=0$$
(21)
which is the equation of motion for a free particle where $`\tau `$ is the proper time. However, we find that the following non-trivial condition
$$P^\mu \frac{dP_\mu }{d\tau }=0$$
(22)
with $`\frac{dP_\mu }{d\tau }0`$ may also be satisfied in consistency with the first class constraint (20). To this end, we define the generalized spatial force $`Q_i`$ as the force of constraint
$$Q_iF_i\frac{dP_i}{d\tau }i=1,2,3$$
(23)
where $`F_i`$ indicates the spatial components of the four-force $`F_\mu `$ and $`\frac{dP_i}{d\tau }`$ is assumed to be a force with physical origin in the same foot as $`F_i`$. The physical nature of the force $`\frac{dP_i}{d\tau }`$ may be justified according to the recent results obtained by Rueda et.al. in that a reaction (physical) force $`\underset{r}{\overset{}{F}}=\frac{d\underset{r}{\overset{}{P}}}{d\tau }=\frac{d\stackrel{}{P}}{d\tau }`$ (generated by ZPF scattering) is interpreted as inertia so that one can rewrite equation (23) as $`Q\stackrel{}{F}+\underset{r}{\overset{}{F}}`$ . It is easy to show that $`Q_i=0`$ satisfies the condition (22). To see, we rewrite (22) as
$$E\frac{dE}{d\tau }P_i\frac{dP_i}{d\tau }c^2=0$$
(24)
which leads to
$$F_i=\frac{dP_i}{d\tau }$$
(25)
where
$$\frac{dE}{d\tau }=F_iV_i,E=\gamma m_0c^2andP_i=\gamma m_0V_i$$
have been used and $`\gamma `$ is the relativistic factor <sup>5</sup><sup>5</sup>5Here, the four-momentum denoted by $`K_\mu `$ in is exchanged symbolically by $`F_\mu `$ for convenience.. Therefore, the vanishing of the force of constraint $`Q_i`$ implies the mass-shell condition to be hold at all times. The interpretation of vanishing $`Q_i`$ is that the acceleration, as well as uniform motion, is also consistent with this first class constraint such that the active force $`F_i`$ is balanced, according to Newoton’s third law, by an equal and opposite directional (physical) force $`\frac{dP_i}{d\tau }`$ which we shall interpret it as inertia. As far as we concern with the present constraint system it is reasonable to say that during the acceleration of particle a reaction force against the agent of acceleration is generated as the price to be paid to preserve the constraint (20) which defines, on the other hand, the causal structure of space-time in the model.
As is usual in the study of non-inertial frames in classical mechanics, we may also transfer the term $`\frac{dP_i}{d\tau }`$ to the right hand side of equation $`Q_i=0`$ to obtain the equation (25) which is interpreted as the relativistic equation of motion. These interpretations will be reasonable if we apply the general law of motion in the form
$$\text{Physical Force}=\frac{d}{d\tau }(\text{Particle Momentum})$$
(26)
in both inertial and accelerated (co-moving) frames of reference. Then the equation
$$F_i\frac{dP_i}{d\tau }=0$$
(27)
compared to (26) implies that in the left hand side the agent of particle dynamics, namely the physical motive force $`F_i`$ co-moving with the particle experiences an equal and opposite directional physical force $`\frac{dP_i}{d\tau }`$ as the inertia which results in the right hand side a zero acceleration of particle in the co-moving frame. A realistic example for this situation is to suppose an accelerated rocket to which a particle has been attached by an string. In the co-moving frame, the rocket will then experience the inertia of particle as a reaction force against the active force imposed by string, but no acceleration is attributed to the particle simply because it moves in contact with the rocket. On the other hand, the equation written as
$$F_i=\frac{dP_i}{d\tau }$$
compared to (26) implies that the physical force $`F_i`$ (in the left hand side) produces a dynamics for the particle according to $`\frac{dP_i}{d\tau }`$ (in the right hand side) in the inertial frame and expresses the relativistic equation of motion.
Equations (25) and (27) may be also interpreted as expressing the circular motion of the point particle in both inertial and rotating frames of reference, respectively. In this respect, $`F_i`$ would be the centripetal force and the term $`\frac{dP_i}{d\tau }`$ would indicate the centrifugal force in the rotating frame.
We conclude that in the context of present paper inertia as the reaction force against acceleration may emerge due to resistance of the first class constraint to breakdown. The breakdown occurs either when an absolute time parameter exists or when the causal structure of space-time breaks down. Therefore, it may be said that the absence of an absolute time parameter or, equivalently, the causal structure described by the mass-shell condition in this system is the possible conceptual origin of the inertia. On the other hand, we know from cosmology that the universe, as a whole, may be described by a similar constraint structure as here with a first class constraint known as Wheeler-DeWitt equation. This equation also implies that there is no absolute time parameter external to the universe.
For an interested Machian, it would be appealing to relate the whole universe and the inertia of particle through the general property of the universe:“There is no absolute time”. In this regard, the inertia of a particle emerging from reparametrization invariance of its action may be induced locally by this global cosmological property. This is because it is not reasonable to have no absolute time in cosmological scale but have an absolute time in local scales.
The old question that “How a material object understands its motion to manifest the inertial effects? ” was answered by Newton as: “It feels the action of absolute space”, and by Mach as: “It feels the action of cosmic matter”. Another answer, based on this approach, may be given as: “It knows the absence of absolute time and feels the action of space-time causal structure”. Therefore, the universal property of reparametrization invariance seems to be a missed chain which links the universe to the local particle dynamics without the need to action at distance violating the causality principle. In simple words, the state of being at rest, in uniform motion and in accelerated motion of a particle are correlated with the state of universe through this universal property. In this regard, the first and second laws of Newton seem to be two different manifestations of this universal property ( see equations (21) and (25) ). Although this approach to the issue of inertia is not fully Machian <sup>6</sup><sup>6</sup>6Wheeler-DeWitt equation is also valid for an empty universe. Therefore, in the context of present paper the particle may have inertia even in an empty universe. (as there is no direct role playing by the matter content of the universe) but, unlike the Machian viewpoint, has the advantage of being compatible with general relativity as a realization of the reparametrization invariant theory. If correct, this approach would substitute for Mach’s principle and imply that the inertial effects are pure relativistic in nature preserving the causality in general agreement with .
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# Ground state laser cooling using electromagnetically induced transparency
## Abstract
A laser cooling method for trapped atoms is described which achieves ground state cooling by exploiting quantum interference in a driven $`\mathrm{\Lambda }`$-shaped arrangement of atomic levels. The scheme is technically simpler than existing methods of sideband cooling, yet it can be significantly more efficient, in particular when several motional modes are involved, and it does not impose restrictions on the transition linewidth. We study the full quantum mechanical model of the cooling process for one motional degree of freedom and show that a rate equation provides a good approximation.
Laser cooling has played a central role in the preparation of fundamental quantum mechanical atomic systems , for example in experiments which study the quantum statistical properties of atoms or which use trapped ions for processing information at the quantum level . In the latter context, laser cooling of trapped ions to the ground state of the confining potential is a fundamental step in the preparation of the ion trap quantum computer . Furthermore, the same techniques that allow ground state cooling are at the basis of coherent manipulation, i.e. gate operations, in quantum computation schemes with trapped ions. For both purposes, cooling and gate operations, the speed of the manipulation has become an important issue , because higher speed means that competing heating or decoherence due to coupling to the environment has less opportunity to perturb the desired processes.
Efficient ground state laser cooling of single trapped ions has been achieved using two-level sideband cooling and Raman sideband cooling , and recently these methods have been transferred to two ions and to atomic gases . These techniques involve laser excitation of an atom with two internal levels $`|g`$ and $`|e`$, which in the case of Raman sideband cooling is designed from a $`\mathrm{\Lambda }`$-shaped three-level atom by Raman coupling . Both techniques rely on several conditions : (i) the motional spectrum of the system has equidistant levels $`|n`$, which is true when the particle (or particles) are trapped in a harmonic potential; (ii) the amplitude of the oscillations of the trapped particles is much smaller than the wavelength of the cooling laser (”Lamb-Dicke regime”); (iii) The linewidth $`\gamma `$ of the internal transition is much smaller than the distance between any pair of motional energy levels (”strong confinement”). For the case of a single particle confined in a harmonic oscillator potential with frequency $`\nu `$, the strong confinement condition is $`\gamma \nu `$.
Under these conditions it is possible to selectively excite sidebands of the optical resonance, i.e. transitions corresponding to a fixed change of the vibrational quantum number $`n`$ to $`n^{}`$, by tuning the laser into resonance with that transition, while all other transitions are well off resonance and thus only negligibly excited. Specifically, for sideband cooling transitions $`|g,n|e,n1`$ are induced by tuning the laser to $`\omega _a\nu `$, i.e. to the ”red sideband” of the bare atomic resonance at frequency $`\omega _a`$. When a spontaneous decay from $`|e`$ to $`|g`$ takes place or, in the case of Raman sideband cooling, when the atom is optically pumped back from $`|e`$ to $`|g`$, this decay occurs with highest probability on the transition $`|e,n1|g,n1`$ due to the Lamb-Dicke condition. Thus in one fluorescence cycle the system is cooled, on average, by one vibrational quantum. The cooling limit is determined by the equilibrium between these cooling cycles and heating processes. Heating is induced by off-resonant excitation of the $`|g,n|e,n`$ ”carrier” transition followed by a $`|e,n`$ to $`|g,n+1`$ spontaneous emission event, or by excitation of a $`|g,n|e,n+1`$ ”blue sideband” transition. Since the selective excitation of the $`|g,n|e,n1`$ sideband is at the basis of this technique, this imposes a limitation on the intensity of the cooling laser and thus also on the cooling speed. In particular, high laser intensity leads to increased off-resonant excitation of carrier transitions which limits the final ground state occupation of the cooling process.
In this paper we describe a method for ground state cooling of atoms with a multi-level structure which eliminates the carrier excitation by electromagnetically induced transparency . The technique is based on continuous laser excitation and has several advantages over both 2-level and Raman sideband cooling. Unlike in 2-level sideband cooling, no strong confinement is required, instead two dipole-allowed transitions are used, neither of which has to fulfil the relation $`\gamma \nu `$. Unlike Raman sideband cooling which involves an additional repumping laser, only two lasers are needed in our method. Finally, as will be shown, by cancelling the carrier transition our scheme provides more efficient ground state cooling than sideband cooling methods, in particular for simultaneous cooling of several modes of vibration. This work extends previous analyses of laser-cooling in a three-level atomic system , which however focused on different cooling mechanisms, as will be discussed below.
Electromagnetically induced transparency (EIT) arises in three- (or multi-) level systems and consists in the cancellation of the absorption on one transition induced by simultaneous coherent driving of another transition. The phenomenon is also called ”coherent population trapping” or ”dark resonance” and has been demonstrated in many systems including single trapped ions . It belongs to a large class of quantum interference effects in multi-level systems and can be understood as a destructive interference of the two pathways to the excited level . It is also at the basis of velocity selective coherent population trapping (VSCPT), a laser cooling method for free atoms which achieves sub-recoil temperatures . Here, we use this situation to suppress absorption on the $`|g,n|e,n`$ transition, while enhancing the absorption on the $`|g,n|e,n1`$ sideband transition, thus decreasing the heating and increasing the cooling rate.
Let us for the moment neglect the motional degrees of freedom and consider a 3-level atom with ground state $`|g`$, stable or metastable state $`|r`$ and excited state $`|e`$ in $`\mathrm{\Lambda }`$-configuration as shown in Fig. 1. State $`|e`$ has linewidth $`\gamma `$ and is coupled to both $`|g`$ and $`|r`$ by dipole transitions. The transition $`|r|e`$ is excited by an intense ”coupling” laser field of frequency $`\omega _r`$, Rabi frequency $`\mathrm{\Omega }_r`$ and detuning $`\mathrm{\Delta }_r=\omega _r\omega _{re}`$, where $`\omega _{re}`$ is the frequency of the bare atomic transition $`|r|e`$. The absorption spectrum observed by exciting the transition $`|g|e`$ with another ”cooling” laser at frequency $`\omega _{ge}+\mathrm{\Delta }_g`$ and Rabi frequency $`\mathrm{\Omega }_g`$ is described by a Fano-like profile , whose zero corresponds to the case $`\mathrm{\Delta }_g=\mathrm{\Delta }_r`$ and which is asymmetric for $`\mathrm{\Delta }_r0`$, see Fig. 1. The same spectrum describes the rate at which photons are scattered from state $`|e`$, and one can infer from it the cooling effect of the laser excitation on the ion .
In the case $`\mathrm{\Delta }_r>0`$ which is displayed in Fig. 1, the two components of the absorption spectrum, i.e. the broad resonance at $`\mathrm{\Delta }_g0`$ with linewidth $`\gamma ^{\prime \prime }\gamma `$ and the narrow resonance at $`\mathrm{\Delta }_g\mathrm{\Delta }_r`$ with linewidth $`\gamma ^{}\gamma `$, correspond to the dressed states of the system atom + coupling laser , see Fig. 1c. These dressed states, and hence the maxima of the narrow and broad curve, are shifted from $`\mathrm{\Delta }_r`$ by $`+\delta `$ and $`\mathrm{\Delta }_r\delta `$, respectively, with
$$\delta =(\sqrt{\mathrm{\Delta }_r^2+\mathrm{\Omega }_r^2}|\mathrm{\Delta }_r|)/2$$
(1)
being the AC Stark shift induced by the coupling laser.
With the harmonic motion taken into account, the zero of the Fano-like profile at $`\mathrm{\Delta }_g=\mathrm{\Delta }_r`$ corresponds to the $`|g,n|e,n`$ transition which is therefore cancelled. Then, by choosing $`\mathrm{\Delta }_r>0`$ and a suitable Rabi frequency $`\mathrm{\Omega }_r`$, the spectrum can be designed such that the $`|g,n|e,n1`$ (red) sideband corresponds to the maximum of the narrow resonance, whereas the blue sideband falls into the region of the spectrum of small excitation probability, as shown in Fig. 1b. The condition on the laser parameters for enhancing the red-sideband absorption while eliminating the carrier is therefore:
$$\mathrm{\Delta }_g=\mathrm{\Delta }_r;\delta \nu $$
(2)
The laser parameters of Eq. (2) are easily achievable in single ion experiments where typically $`\gamma 2\pi \times `$ 20 MHz and $`\nu 2\pi \times `$ 1 MHz.
Note that the detunings are different from Raman sideband cooling (RSC) where $`\mathrm{\Delta }_g=\mathrm{\Delta }_r\nu `$. Furthermore, here both lasers must be blue-detuned from their respective atomic resonances, whereas in RSC they can be tuned either both below or both above resonance. Moreover, in RSC the bare states $`|g`$, $`|r`$, are coupled under saturation to $`|e`$ (this is the situation described in ) whereas in the new cooling scheme, both multiple scattering on the transition $`|r|e`$ and the quantum interference at $`\mathrm{\Delta }_g=\mathrm{\Delta }_r`$ are crucial for the cooling process, which is therefore adequately described by transitions between state $`|g`$ and the two dressed states, see Fig. 1c.
The mechanism is theoretically modelled as follows. We start with the master equation for the full three-level system and one motional degree of freedom. In the Lamb-Dicke regime, the master equation can be reduced to a rate equation projected on the internal state $`|g,n`$ provided that $`\mathrm{\Omega }_g\mathrm{\Omega }_r`$ and that the transition to the narrow dressed state of linewidth $`\gamma ^{}`$ is not saturated, i.e. $`\mathrm{\Omega }_g\sqrt{\gamma \gamma ^{}}`$. Then, in second order of the expansion in $`\mathrm{\Omega }_g/\sqrt{\gamma \gamma ^{}}`$ the dynamics is described by an equation for the populations $`P(n)`$ of the vibrational number states $`|n`$ :
$`{\displaystyle \frac{d}{dt}}P(n)`$ $`=`$ $`\eta ^2[A_{}((n+1)P(n+1)nP(n))+`$ (3)
$`+`$ $`A_+(nP(n1)(n+1)P(n))]`$ (4)
Here, $`\eta `$ is the Lamb-Dicke parameter, defined as $`\eta =|\stackrel{}{k_g}\stackrel{}{k_r}|a_0`$ with $`a_0`$ rms size of the ground state of the harmonic oscillator and $`\stackrel{}{k_g}`$ ($`\stackrel{}{k_r}`$) cooling (coupling) laser wave vector . The coefficients $`A_\pm `$ have the form
$$A_\pm =\frac{\mathrm{\Omega }_g^2}{\gamma }\frac{\gamma ^2\nu ^2}{\gamma ^2\nu ^2+4\left[\mathrm{\Omega }_r^2/4\nu \left(\nu \mathrm{\Delta }\right)\right]^2}$$
(5)
where we have set $`\mathrm{\Delta }_r=\mathrm{\Delta }_g=\mathrm{\Delta }`$. Note that while Eq. (3) has the same structure as the rate equation that describes cooling of a two-level atom in the Lamb-Dicke regime , the particular form of $`A_\pm `$ in Eq. (5) contains the full quantum interference around $`\mathrm{\Delta }_g=\mathrm{\Delta }_r`$. Solving Eq. (3) for the steady state value $`n_S`$ of the mean vibrational quantum number $`n=nP(n)`$ we get
$$n_S=\frac{A_+}{A_{}A_+}=\frac{\gamma ^2\nu ^2+4[\mathrm{\Omega }_r^2/4\nu (\nu +\mathrm{\Delta })]^2}{4\mathrm{\Delta }\nu (\mathrm{\Omega }_r^24\nu ^2)}$$
(6)
Eq. (6) has a pole ($`A_+=A_{}`$) at $`\mathrm{\Delta }=0`$, where the spectrum is symmetric, and at $`\mathrm{\Omega }_r=2\nu `$, where the value of the absorption spectrum at the two frequencies $`\mathrm{\Delta }_g\pm \nu `$ is the same. For properly chosen parameters, however, a value of $`n_S`$ close to zero can be reached, as shown in the example of Fig. 2. The optimum value $`n_S=(\gamma /4\mathrm{\Delta })^2`$ is found when the second term in the numerator of Eq. (6) vanishes, which corresponds precisely to the condition $`\delta =\nu `$ in Eq. (2).
From Eq. (3) the time dependence of $`n`$ follows,
$$\dot{n}=\eta ^2(A_{}A_+)n+\eta ^2A_+$$
(7)
where $`\eta ^2(A_{}A_+)`$ is the cooling rate. This rate together with $`n_S`$ determines the efficiency of the cooling technique which is compared to conventional sideband cooling below.
The dynamics of the full system for any set of parameters can be calculated with a quantum Monte-Carlo simulation . In Fig. 3 we plot the result of such a calculation and compare it to the rate equation solution of Eq. (7). We see that the rate equation provides a good description of the cooling. In this example, 99% occupation of the ground state is achieved.
In order to compare our EIT-cooling scheme with conventional 2-level sideband cooling, we use the results of the respective rate equations for both schemes and we allow for an angle $`\varphi `$ between the motional axis and the direction of the laser beam (of $`\stackrel{}{k_g}\stackrel{}{k_r}`$ in the case of EIT-cooling). We denote by $`A_\pm ^{\mathrm{SC}}`$ the coefficients for sideband cooling corresponding to Eq. (5), c.f. . Assuming a linewidth $`\gamma _{\mathrm{SC}}=\gamma ^{}`$ for the sideband cooling transition, and the same degree of saturation in both schemes, i.e. $`(\mathrm{\Omega }_{\mathrm{SC}}/\gamma _{\mathrm{SC}})^2=\mathrm{\Omega }_g^2/\gamma \gamma ^{}`$, the following relation holds:
$$A_\pm ^{\mathrm{SC}}=A_\pm +\frac{\mathrm{\Omega }_{\mathrm{SC}}^2}{\gamma _{\mathrm{SC}}}\left(\frac{\alpha }{\mathrm{cos}^2\varphi }\right)\frac{\gamma _{\mathrm{SC}}^2}{\gamma _{\mathrm{SC}}^2+4\nu ^2}$$
(8)
Here we have used that $`\gamma ^{}\gamma \frac{\nu }{\mathrm{\Delta }}`$ for our conditions $`\delta =\nu \mathrm{\Delta },\mathrm{\Omega }_r`$ , and $`\alpha =_1^1𝑑uN(u)u^2`$ with $`N(\mathrm{cos}\vartheta )`$ being the azimuthal dipole pattern of spontaneous emission on the sideband cooling transition. The additive term in Eq. (8) highlights that the difference between sideband cooling and EIT-cooling lies in the heating that accompanies carrier absorption, which in EIT-cooling is cancelled. This leads to the remarkable result that in contrast to any other cooling scheme, in EIT-cooling the theoretical limit $`n_S`$ does not depend on the angle between the direction of the laser beams (of $`\stackrel{}{k_g}\stackrel{}{k_r}`$) and the motional axis. For the typical condition of sideband cooling $`\gamma _{\mathrm{SC}}\nu `$ we find, using Eq. (6),
$$n_S_{\mathrm{EIT}}=\frac{n_S_{\mathrm{SC}}}{1+\frac{4\alpha }{\mathrm{cos}^2\varphi }}$$
(9)
Thus, given the same cooling rate, $`n_S_{\mathrm{EIT}}<n_S_{\mathrm{SC}}`$. The factor between the two steady state values is $`\frac{5}{29}`$ for all three motional degrees of freedom, assuming the ion is cooled in three dimensions by a single pair of laser beams, $`\alpha =\frac{2}{5}`$ , and $`\stackrel{}{k_g}\stackrel{}{k_r}`$ has the same angle with all axes of the trap. It becomes even smaller when $`\stackrel{}{k_g}\stackrel{}{k_r}`$ is at a large angle with the motional axis. This makes EIT-cooling a significantly more efficient technique in many typical experimental situations . Moreover, numerical studies show that the efficiency can be increased further if the cooling laser is tuned not exactly to the dark resonance but slightly above such that the combined heating of carrier and blue sideband transitions is minimised.
In conclusion, we have presented a laser cooling technique for trapped particles which exploits quantum interference, or electromagnetically induced transparency, in a 3-level atom. By appropriately designing the absorption profile with a strong coupling laser, the cooling transitions induced by a cooling laser are enhanced while heating by resonant absorption is suppressed. The method is not based on the strong confinement condition and requires only two continuous lasers. We have derived a simple model for describing the cooling process and shown that it is in good agreement with a full quantum Monte Carlo treatment. With the same cooling rate as in conventional sideband cooling, much higher ground state occupation is achieved by our method, in particular if three-dimensional cooling is considered. The technical requirements, two lasers with a well-controlled frequency difference, are met by most existing single ion experiments, and they are less stringent than for both Raman sideband cooling and ordinary 2-level sideband cooling. Furthermore, the method is insensitive to laser frequency fluctuations as long as the laser linewidth is small compared to the trap frequency.
Simultaneous cooling in three dimensions can be achieved with this method if the trap frequencies along the axes are similar, so that the red sideband for each oscillator falls into the neighbourhood of the maximum of the narrow resonance. Similarly, the method can be applied to simultaneously cool several axial modes of an N-ion crystal in a linear trap, or atoms in anharmonic traps, where the energies of the motional states are not equidistant. An extension would be to use an atom with more than three levels, where multiple dark resonances occur (see, e.g., ) and to design the absorption spectrum such that both the carrier and the blue sideband transition vanish.
Note: Following the ideas of this proposal, the method has meanwhile been experimentally demonstrated .
We thank J. I. Cirac, P. Lambropoulos, D. Leibfried, C. Roos, F. Schmidt-Kaler, H. Walther, and P. Zoller for many stimulating discussions, and acknowledge support by the European Commission (TMR networks ERB-FMRX-CT96-0077 and ERB-FMRX-CT96-0087; Marie Curie Program) and the German Science Foundation (SFB 276). CHK is thankful for hospitality at Innsbruck University.
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# Basic limitations for entanglement catalysis
## Abstract
In this paper we summarize the necessary condition for incomparable states which can be catalyzed under entanglement-assisted LQCC (ELQCC). When we apply an extended condition for entanglement transformation to entanglement-assisted local manipulation we obtain a fundamental limit for entanglement catalysts. Some relative questions are also discussed.
The existence of entanglement between spatially separated quantum systems is at the heart of quantum information theory. In the light of recent progress in quantum information theory entanglement is often viewed as the essential resource for processing and transmitting quantum information and forms the basis for many miraculous applications including quantum teleportation , quantum cryptography and quantum communication . For many practically applications of this nonlocal resource, one often restricts oneself to only preforming local quantum operations on respective subsystems of entangled state and exchanging classical information. In this paper transformations of this type will be referred to as “local transformation ”, or “ LQCC ” for short. LQCC forms a fundamental limit for quantification and manipulation of entanglement .
One can transform entangled quantum system from one situation to another under LQCC. However, on average, the entanglement degree of the initial system will be never enhanced. A quantitative way of expressing this fact is in terms of so-called entanglement monotones ( EMs ) . EMs provide necessary restrictions for entanglement transformation under LQCC. But, there arises naturally the question of what are the sufficient conditions for entanglement transformation. One very fruitful approach was achieved by Nielsen who provided a sufficient and necessary condition for pure state entanglement transformations . Inspired by Nielsen’s theory, Jonathan and Plenio investigated the behaviors of entanglement-assisted local manipulation of pure quantum states and presented a new concept called entanglement catalysis . Afterwards, Eisert and Wilkens made further studies for catalysis of entanglement manipulation for mixed states .
This remarkable phenomenon of entanglement catalysis can be described as follows. Let $`\rho _s`$ and $`\rho _t`$ be source state and target state respectively, which are taken from the state space S(H) over H, where H=$`H_AH_B`$ is the Hilbert space associated with a bipartite quantum system consisting of part A and B. The target state $`\rho _t`$ can not be reached by LQCC from the source state $`\rho _s`$ with certainty. But with the assistance of a particular known entangled state $`\rho _c`$ taken from the state space S($`\stackrel{~}{H}`$) over $`\stackrel{~}{H}`$ the transformation from $`\rho _s\rho _c`$ to $`\rho _t\rho _c`$ can be achieved by LQCC with 100% probability, where $`\stackrel{~}{H}`$ is a tensor product $`\stackrel{~}{H}=\stackrel{~}{H}_A\stackrel{~}{H}_B`$ of two Hilbert spaces belonging to systems A and B respectively. The auxiliary entangled state $`\rho _c`$, which plays an indeed catalyst role in this process, is left in exactly the same state and remain finally completely uncorrelated to the quantum system of interest. This counter-intuitive effect demonstrate that entanglement can be “ borrowed ”.
In this paper we first summarize the necessary condition for pure bipartite incomparable states which can be catalyzed under entanglement-assisted LQCC (ELQCC). Furthermore, starting from an extended condition for entanglement transformation, we find that catalytic transition processes will provide a fundamental limit for catalysts themselves.
Let us begin with Nielsen’s theorem.
Theorem ( Nielsen ): Let $`|\mathrm{\Psi }_1=_{i=1}^n\sqrt{\alpha _i}|i_A|i_B`$ and$`|\mathrm{\Psi }_2=_{i=1}^n\sqrt{\beta _i}|i_A|i_B`$ be pure bipartite states, where the Schmidt coefficients are ordered according to $`\alpha _1\alpha _2\mathrm{}\alpha _n>0`$ and $`\beta _1\beta _2\mathrm{}\beta _n>0`$ respectively. ( We can refer to such distributions as “ ordering Schmidt coefficients ”, or OSCs. ) Then the transformation from $`|\mathrm{\Psi }_1`$ to $`|\mathrm{\Psi }_2`$ with 100% probability can be realized using LQCC iff the OSCs{$`\alpha _i`$} are majorized by {$`\beta _i`$} , that is, iff for $`1ln`$
$$\underset{i=1}{\overset{l}{}}\alpha _i\underset{i=1}{\overset{l}{}}\beta _i.$$
(1)
Nielsen call the state $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_2`$ incomparable if neither state can not convert into the other with certainty. The current studies indicate only transitions between incomparable states may be catalyzed. Then, under ELQCC what conditions should be satisfied if the transformations between incomparable states are possible? One naturally desires to find analogous Nielsen’s criterion. Unfortunately, we at present do not know what are sufficient conditions for the existence of catalysts. To find appropriate catalysts one has to resort to numerical search. Nevertheless, we may provide two necessary conditions for entanglement catalysis.
Let $`|\mathrm{\Psi }_1=_{i=1}^n\sqrt{\alpha _i}|i_A|i_B`$ and$`|\mathrm{\Psi }_2=_{i=1}^m\sqrt{\beta _i}|i_A|i_B`$ ( m$``$n ) be pure bipartite states with OSCs {$`\alpha _i`$ ; i=1, …,n} {$`\beta _i`$ , i=1, … ,n} ( here $`\beta _i=0`$ if i=m+1, … , n ). Then $`|\mathrm{\Psi }_1`$ can be converted into $`|\mathrm{\Psi }_2`$ with certainty under ELQCC only if
$$(i)\text{ }\alpha _1\beta _1\text{ , }\alpha _n\beta _n$$
(2)
$$(ii)\text{ }S\left(\rho _1\right)S\left(\rho _2\right),$$
(3)
where $`S\left(\rho _i\right)`$ is the marginal Von Neumann entropy of state $`|\mathrm{\Psi }_i.`$ In ( ii ) saturation is reached iff $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_2`$ are locally unitarilly equivalent.
Condition ( i ) has been proved in . For condition ( ii ) we offer a brief proof in the following.
Proof: Suppose that $`|\mathrm{\Psi }_1`$ can be transformed into $`|\mathrm{\Psi }_2`$ under ELQCC. Then there exists a catalyst $`|\mathrm{\Psi }_c`$ such that $`|\mathrm{\Psi }_1|\mathrm{\Psi }_c`$ can be converted to $`|\mathrm{\Psi }_2|\mathrm{\Psi }_c`$ with 100% probability under LQCC. ( For simplicity, in the next context we indicate this process by $`|\mathrm{\Psi }_1|\mathrm{\Psi }_c|\mathrm{\Psi }_2|\mathrm{\Psi }_c.`$ ) In light of the non-increase of partial entropy under LQCC and additivity of entropy the inequality (3) can be obtained. In order to prove the saturation case we may refer to Theorem 1 in . If $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_2`$ are two marginally isentropic pure states, they are either locally unitarilly equivalent or else LQCC-incomparable. A manifest fact is that $`|\mathrm{\Psi }_1|\mathrm{\Psi }_c`$ and $`|\mathrm{\Psi }_2|\mathrm{\Psi }_c`$ is locally unitarilly equivalent if they are marginal isentropic. Furthermore, we may deduce that $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_2`$ are locally unitarilly equivalent.
The above conditions just put forward limits for the relationship between source and target. Suppose that two incomparable states to hold inequality (2) and (3) can be “catalyzed”. Are there any limitations for potential catalysts? To obtain our theorem the following lemma is necessary.
Lemma 1: Let $`|\mathrm{\Psi }_p`$ be p$`\times `$p-level maximal entangled state. Then the maximal probability of obtaining state $`|\mathrm{\Psi }_2`$ from $`|\mathrm{\Psi }_1`$ by means of LQCC is just equal to that of $`|\mathrm{\Psi }_1|\mathrm{\Psi }_p|\mathrm{\Psi }_2|\mathrm{\Psi }_p`$, i.e.
$$P_{\mathrm{max}}\left(|\mathrm{\Psi }_1|\mathrm{\Psi }_2\right)=P_{\mathrm{max}}\left(|\mathrm{\Psi }_1|\mathrm{\Psi }_p|\mathrm{\Psi }_2|\mathrm{\Psi }_p\right).$$
(4)
Proof: The maximal probability of $`|\mathrm{\Psi }_1|\mathrm{\Psi }_2`$ has been achieved by Vidal:
$$P_{\mathrm{max}}\left(|\mathrm{\Psi }_1|\mathrm{\Psi }_2\right)=\underset{l[1,n]}{\mathrm{min}}\frac{_{i=l}^n\alpha _i}{_{i=l}^n\beta _i}=\underset{l[1,n]}{\mathrm{min}}\frac{E_l\left(|\mathrm{\Psi }_1\right)}{E_l\left(|\mathrm{\Psi }_2\right)},$$
(5)
where $`E_l\left(|\mathrm{\Psi }_1\right)=_{i=l}^n\alpha _i`$ is called entanglement monotone of state $`|\mathrm{\Psi }_1`$. $`|\mathrm{\Psi }_1|\mathrm{\Psi }_p`$ and $`|\mathrm{\Psi }_2|\mathrm{\Psi }_p`$ have separately p-fold degenerate OSCs {$`\alpha _i^{}`$ ; i=1, … ,pn} and {$`\beta _i^{}`$ ; i=1, … ,pn}. if $`l`$ is a number between $`pk+1`$ and $`p\left(k+1\right)`$, $`k=0,1,\mathrm{},n1`$, we have:
$$\frac{E_l\left(|\mathrm{\Psi }_1|\mathrm{\Psi }_p\right)}{E_l\left(|\mathrm{\Psi }_2|\mathrm{\Psi }_p\right)}=\frac{\frac{lpk1}{p}E_{k+2}\left(|\mathrm{\Psi }_1\right)+\frac{p(k+1)l+1}{p}E_{k+1}\left(|\mathrm{\Psi }_1\right)}{\frac{lpk1}{p}E_{k+2}\left(|\mathrm{\Psi }_2\right)+\frac{p(k+1)l+1}{p}E_{k+1}\left(|\mathrm{\Psi }_2\right)},$$
(6)
where $`E_{n+1}\left(|\mathrm{\Psi }_1\right)=E_{n+1}\left(|\mathrm{\Psi }_2\right)=0.`$ Without loss of generality, let us set $`\frac{E_{k+2}\left(|\mathrm{\Psi }_1\right)}{E_{k+2}\left(|\mathrm{\Psi }_2\right)}>\frac{E_{k+1}\left(|\mathrm{\Psi }_1\right)}{E_{k+1}\left(|\mathrm{\Psi }_2\right)}`$. Taking advantage of the equivalence:
$$\frac{a}{b}<\frac{a+c}{b+d}\frac{a}{b}<\frac{c}{d},$$
(7)
we have the following relation:
$$\frac{E_l\left(|\mathrm{\Psi }_1|\mathrm{\Psi }_p\right)}{E_l\left(|\mathrm{\Psi }_2|\mathrm{\Psi }_p\right)}\frac{E_{k+1}\left(|\mathrm{\Psi }_1\right)}{E_{k+1}\left(|\mathrm{\Psi }_2\right)}.$$
(8)
Here, the saturation holds iff $`l=pk+1`$. Therefore,
$`P_{\mathrm{max}}\left(|\mathrm{\Psi }_1|\mathrm{\Psi }_p|\mathrm{\Psi }_2|\mathrm{\Psi }_p\right)=\underset{l[1,pn]}{\mathrm{min}}{\displaystyle \frac{_{i=l}^{pn}\alpha _i^{}}{_{i=l}^{pn}\beta _i^{}}}`$
$$=\underset{l[1,n]}{\mathrm{min}}\frac{_{i=l}^n\alpha _i}{_{i=l}^n\beta _i}=P_{\mathrm{max}}\left(|\mathrm{\Psi }_1|\mathrm{\Psi }_2\right).$$
(9)
This completes the proof of lemma 1.
The following theorem provide us with a fundamental limit for catalysts.
Theorem 1: Successful transformation from $`|\mathrm{\Psi }_1`$ to $`|\mathrm{\Psi }_2`$ can be reached under ELQCC only if any of a p$`\times `$p-level catalyst with OSCs {$`\gamma _i`$ ; i=1, … , p} meet the following relation:
$$p\gamma _pP_{\mathrm{max}}\left(|\mathrm{\Psi }_1|\mathrm{\Psi }_2\right).$$
(10)
Proof: Based on lemma 1 we have
$`P_{\mathrm{max}}(|\mathrm{\Psi }_1`$ $``$ $`|\mathrm{\Psi }_2)=P_{\mathrm{max}}(|\mathrm{\Psi }_1|\mathrm{\Psi }_p|\mathrm{\Psi }_2|\mathrm{\Psi }_p)`$ (11)
$``$ $`P_{\mathrm{max}}(|\mathrm{\Psi }_1|\mathrm{\Psi }_c|\mathrm{\Psi }_2|\mathrm{\Psi }_p)`$ (12)
$``$ $`P_{\mathrm{max}}(|\mathrm{\Psi }_2|\mathrm{\Psi }_c|\mathrm{\Psi }_2|\mathrm{\Psi }_p)`$ (13)
$``$ $`P_{\mathrm{max}}(|\mathrm{\Psi }_c|\mathrm{\Psi }_p),`$ (14)
where $`|\mathrm{\Psi }_p`$ and $`|\mathrm{\Psi }_c`$ are p$`\times `$p-level maximal entangled state and catalyst state respectively. The first inequality is satisfied due to the following fact: one can look on the process $`|\mathrm{\Psi }_1|\mathrm{\Psi }_p|\mathrm{\Psi }_1|\mathrm{\Psi }_c`$ as one of steps in the transformation from $`|\mathrm{\Psi }_1|\mathrm{\Psi }_p`$ to $`|\mathrm{\Psi }_2|\mathrm{\Psi }_p`$ while the transformation from $`|\mathrm{\Psi }_p`$ to $`|\mathrm{\Psi }_c`$ is deterministic under LQCC. If any an intermediate state can be deterministically arrived the probability from this intermediate state to target state can not be higher than the maximal probability from source to target. Otherwise, this will lead to contradiction. Depended on the similar options the second inequality can also be obtained. While, the last inequality is obvious. On the basis of (5) and (7), the maximal probability of $`|\mathrm{\Psi }_c|\mathrm{\Psi }_p`$ is as follows:
$$P_{\mathrm{max}}(|\mathrm{\Psi }_c|\mathrm{\Psi }_p)=p\gamma _p.$$
(15)
We thus complete the proof of our theorem.
As a direct consequence, this theorem can supply some concrete limits for entanglement catalysts under some special circumstances. For example, our choice of Schmidt coefficients for source state$`|\mathrm{\Psi }_1`$ and target state$`|\mathrm{\Psi }_2`$ is $`\alpha _1=\alpha _2=0.31`$, $`\alpha _3=0.3`$, $`\alpha _4=\alpha _5=0.04`$, $`\beta _1=0.48`$, $`\beta _2=0.24`$, $`\beta _3=\beta _4=0.14`$, $`\beta _5=0`$. Taking advantage of this theorem we can easily conclude that 2$`\times `$2-level catalysts do not exist for the process $`|\mathrm{\Psi }_1|\mathrm{\Psi }_2`$. A detailed analysis is as follows. Assume that there exists a 2$`\times `$2-level catalyst $`|\mathrm{\Psi }_c`$, with OSCs { $`x,1x`$ }. In view of the inequality relation of (10) we have $`2\times (1x)\frac{4}{7},x\frac{5}{7}`$. Hence, the first three OSCs of $`|\mathrm{\Psi }_1|\mathrm{\Psi }_c`$ must be $`0.31x`$, $`0.31x`$, $`0.3x`$. Similarly, the first three OSCs of $`|\mathrm{\Psi }_2|\mathrm{\Psi }_c`$ is either $`0.48x`$, $`0.24x`$, $`0.48(1x)`$ or $`0.48x`$, $`0.24x`$, $`0.14x`$. No matter which cases take place the relation $`_{i=1}^3\alpha _i^{}>_{i=1}^3\beta _i^{}`$ must be satisfied. ( Here, {$`\alpha _i^{}`$} and {$`\beta _i^{}`$} refer to OSCs of $`|\mathrm{\Psi }_1|\mathrm{\Psi }_c`$ and $`|\mathrm{\Psi }_2|\mathrm{\Psi }_c`$ respectively.) We thus have $`1_{i=1}^3\alpha _i^{}=E_4\left(|\mathrm{\Psi }_1|\mathrm{\Psi }_c\right)<1_{i=1}^3\beta _i^{}=E_4\left(|\mathrm{\Psi }_2|\mathrm{\Psi }_c\right)`$. In view of Nielsen’s theorem this implies there exist no $`2\times 2`$-level entangled states to hold that {$`\alpha _i^{}`$} can be majorized by {$`\beta _i^{}`$}. Seeing reference , we know $`|\mathrm{\Psi }_2`$ itself just a catalyst of the process $`|\mathrm{\Psi }_1|\mathrm{\Psi }_2`$. From this example we find that higher dimensional entangled states have exactly more powerful capability of catalysis than lower dimensional entangled states. The reasons rely on this fact: theorem 1 provides only quite slack bounds for the structure of OSCs of catalysts when marginal Hilbert space of entanglement catalysts have high dimensions. In other words, more degrees of freedom conceals in higher dimensional entanglement. When applying our theorem to generalized cases, we may acquire some interesting corollaries.
According to the inequality (10), entanglement catalysts can be divided into two sorts. We call those catalysts saturating inequality (10) “saturated catalysts”, or else “non-saturated catalysts”. At present, we do not know whether there are saturated catalysts for all pairs of convertible incomparable states under ELQCC. However, any of a pair of incomparable states with saturated catalysts must meet the following corollary.
Corollary 1: $`|\mathrm{\Psi }_1`$ can be deterministically transformed into $`|\mathrm{\Psi }_2`$ by using saturated catalysts-assisted LQCC only if
$$P_{\mathrm{max}}(|\mathrm{\Psi }_1^n|\mathrm{\Psi }_2^n)P_{\mathrm{max}}(|\mathrm{\Psi }_1|\mathrm{\Psi }_2).$$
(16)
Proof: Similar to the proof of theorem 1, we have
$`P_{\mathrm{max}}(|\mathrm{\Psi }_1^n|\mathrm{\Psi }_2^n)=P_{\mathrm{max}}(|\mathrm{\Psi }_1^n|\mathrm{\Psi }_p|\mathrm{\Psi }_2^n|\mathrm{\Psi }_p)`$
$`P_{\mathrm{max}}(|\mathrm{\Psi }_1^n|\mathrm{\Psi }_c|\mathrm{\Psi }_2^n|\mathrm{\Psi }_p)P_{\mathrm{max}}(|\mathrm{\Psi }_2^n|\mathrm{\Psi }_c|\mathrm{\Psi }_2^n|\mathrm{\Psi }_p)`$
$$P_{\mathrm{max}}(|\mathrm{\Psi }_c|\mathrm{\Psi }_p)=p\gamma _p=P_{\mathrm{max}}(|\mathrm{\Psi }_1|\mathrm{\Psi }_2).$$
(17)
This proves corollary 1.
One can make further extensions for the concept of entanglement-assisted transformations. For instance, there is a case of “quasi-catalysis” in which the transformation from $`|\mathrm{\Psi }_1`$ to $`|\mathrm{\Psi }_2`$ can not be preformed with certainty even under ELQCC, but the optimal probability of transformation may still be increased. Based on the analogous analysis, we have:
Corollary 2: The probability of transition from state $`|\mathrm{\Psi }_1`$ to $`|\mathrm{\Psi }_2`$ can be enhanced to $`P^{}`$ with the assistance of a $`p\times p`$-level entangled state $`|\mathrm{\Psi }_c=_{i=1}^p\sqrt{\gamma _i}|i|i`$ only if
$$p\gamma _pP_{\mathrm{max}}(|\mathrm{\Psi }_1|\mathrm{\Psi }_2)/P^{}.$$
(18)
Of late, one begins to consider practical applications for entanglement catalysis. Two schemes for quantum secure identification using catalysts have been presented by Barnum and Jensen et.al respectively. These proposals using entanglement catalysts have an attracting prospect relying on this fact: entanglement catalysts will not be depleted during quantum information processes, i.e. a protracted characteristic entanglement between quantum users may be employed repeatedly. However, Barnum’s protocol has been shown to be insecure . On the condition that all quantum operations are error-free and that the quantum channel is noiseless the quantum authentication protocol presented by Jensen and Schark appears to be secure even in the presence of an eavesdropper who has complete control over both classical and quantum communication channels at all times. How to develop secure and unjammable quantum authentication schemes is attracting more and more attention.
In conclusion, we have shown some necessary limitations for pure bipartite incomparable states which can be catalyzed under ELQCC and catalysts themselves. We find that the product of dimension and the final OSC of catalysts restricts the ability of catalysts. We believe that the results of the present paper can help in deeper understanding of entanglement-assisted local manipulations.
This project was supported by the National Nature Science Foundation of China and the Doctoral Education Fund of the State Education Commission of China.
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# Time Dependent Analysis of Decays Λ_𝑏→Λ+𝐷⁰ and Λ_𝑏→Λ+𝐷̄⁰
## Abstract
The time-dependent analysis of the decays $`\mathrm{\Lambda }_b\mathrm{\Lambda }+D^0(t)`$ and $`\mathrm{\Lambda }_b\mathrm{\Lambda }+\overline{D}^0(t)`$ is discussed. The effect of particle mixing due to time evolution of $`D^0`$ and $`\overline{D}^0`$ on the observables for these decays viz the branching ratio of decay widths and the asymmetry parameters $`\alpha ,\beta `$ and $`\gamma `$ are analysed. It is shown that it is possible to extract information about $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma }`$ from the experimental data for these observables. Here $`\mathrm{\Delta }m=m_{D_1^0}m_{D_2^0},\mathrm{\Gamma }`$ is the decay widths for $`D^{}s`$ and $`\widehat{\gamma }`$ is the weak phase.
In this paper, we discuss the time dependent analysis of the decays $`\mathrm{\Lambda }_b\mathrm{\Lambda }+D^0(t)`$ and $`\mathrm{\Lambda }_b\mathrm{\Lambda }+\overline{D}^0(t)`$. Due to $`D^0`$ and $`\overline{D}^0`$ mixing, a pure $`D^0(\overline{D}^0)`$ beam acquires a component of $`\overline{D}^0(D^0)`$ as it evolves. We analyse the effect of particle mixing on the diservables for these decays viz the ratio of decay widths and the asymmetry parameters, $`\alpha ,\beta ,`$ and $`\gamma .`$
These decays have been previously studied in references $`[1]`$ and $`[2]`$; especially the decays $`\mathrm{\Lambda }_b\mathrm{\Lambda }+D_{1,2}`$ where $`D_{1,2}`$ are CP-eigenstates. The decays are described by four amplitudes $`A_D(t),A_{\overline{D}}(t),B_D(t)`$ and $`B_{\overline{D}}(t)`$. Denoting these amplitudes as $`R_D(t)`$ and $`R_{\overline{D}}(t),`$ where $`R=A`$ or $`B`$, we get the time dependent amplitudes.
$`\left|R_D(t)\right|^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\mathrm{\Gamma }t}\left[(1+\mathrm{cos}\mathrm{\Delta }mt)R_D^22\mathrm{sin}\mathrm{\Delta }mt\mathrm{sin}\widehat{\gamma }R_DR_{\overline{D}}+(1\mathrm{cos}\mathrm{\Delta }mt)R_D^2\right]`$ (1)
$`ReA_D^{}(t)B_D(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\mathrm{\Gamma }t}\left[(1+\mathrm{cos}\mathrm{\Delta }mt)A_DB_D+\mathrm{sin}\widehat{\gamma }\mathrm{sin}\mathrm{\Delta }mt\text{ }(A_DB_{\overline{D}}+A_{\overline{D}}B_D)+(1\mathrm{cos}\mathrm{\Delta }mt)A_{\overline{D}}B_{\overline{D}}\right]`$ (2)
$`ImA_D^{}(t)B_D(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\mathrm{\Gamma }t}[(1+\mathrm{cos}\widehat{\gamma }\mathrm{sin}\mathrm{\Delta }mt)(A_{\overline{D}}B_DA_DB_{\overline{D}})`$ (3)
where we have explicitly exhibited the weak phase $`\widehat{\gamma }`$. After taking out the weak phase, these amplitudes are real, if we neglect the final state interactions. For $`\overline{D},`$ change $`D\overline{D}`$ and $`\mathrm{sin}\mathrm{\Delta }mt\mathrm{sin}\mathrm{\Delta }mt`$ in Eqs. $`(1),(2)`$ and $`(3).`$ We now take the time average of these amplitudes:
$$\overline{R}_D^2=\frac{_0^{\mathrm{}}R_D^2(t)𝑑t}{_0^{\mathrm{}}e^{\mathrm{\Gamma }t}𝑑t}$$
(4)
After taking the time average and neglecting terms of the order $`(\mathrm{\Delta }m/\mathrm{\Gamma })^2,`$ we obtain (similar results follow if instead of taking time average, we take $`t\frac{1}{\mathrm{\Gamma }})`$
$`\overline{R}_D^2`$ $`=`$ $`\left[R_D^2+(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma }R_DR_{\overline{D}}\right]`$ (5)
$`\overline{\alpha }_D`$ $`=`$ $`2\left|\stackrel{}{k}\right|[A_DB_D+{\displaystyle \frac{1}{2}}(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma }A_DB_{\overline{D}}+A_{\overline{D}}B_D)]/\overline{F}_D^2`$ (6)
$`\overline{\beta }_D`$ $`=`$ $`2\left|\stackrel{}{k}\right|\left(\mathrm{\Delta }m/\mathrm{\Gamma }\right)\mathrm{cos}\widehat{\gamma }{\displaystyle \frac{1}{2}}\left(A_{\overline{D}}B_{\overline{D}}A_DB_{\overline{D}}\right)/\overline{F}_D^2`$ (7)
$`\overline{\gamma }_D`$ $`=`$ $`{\displaystyle \frac{\left[(E_{}+m_{})A_D^2(E_{}m_{})B_D^2\right]+(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma }\left[(E_{}+m_{})A_DA_{\overline{D}}(E_{}m_{})B_DB_{\overline{D}}\right]}{\overline{F}_D^2}}`$ (8)
where
$`\overline{F}_D^2`$ $`=`$ $`(E_{}+m_{})\overline{A}_D^2+(E_{}m_{})\overline{B}_D^2`$ (9)
$`=`$ $`\left[(E_{}+m_{})A_D^2+(E_{}m_{})B_D^2\right]`$ (11)
$`+(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma }\left[(E_{}+m_{})A_DA_{\overline{D}}+(E_{}m_{})B_DB_{\overline{D}}\right]`$
For $`\overline{D},`$ change $`(\mathrm{\Delta }m/\mathrm{\Gamma })`$ to - $`(\mathrm{\Delta }m/\mathrm{\Gamma })`$ and $`A_D,B_DA_{\overline{D}},B_{\overline{D}}`$ in $`Eqs.(5),(6),(7),(8)`$ and $`(9).`$
It is convenient to put
$$A_D=\frac{a_D}{\sqrt{E_{}+m_{}}},B_D=\frac{b_D}{\sqrt{E_{}m_{}}}$$
(12)
In terms of these amplitudes, we have from Eqs. $`(69)`$,
$`\overline{F}_D^2`$ $`=`$ $`(a_D^2+b_D^2)\left[1+(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma }{\displaystyle \frac{a_Da_{\overline{D}}+b_Db_{\overline{D}}}{a_D^2a_{\overline{D}}^2}}\right]`$ (13)
$`\overline{\alpha }_D`$ $`=`$ $`{\displaystyle \frac{2a_Db_D}{a_D^2+b_{\overline{D}}^2}}\left[1+(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma }\left({\displaystyle \frac{1}{2}}({\displaystyle \frac{b_{\overline{D}}}{b_D}}+{\displaystyle \frac{a_{\overline{D}}}{a_D}}){\displaystyle \frac{a_Da_{\overline{D}}+b_Db_{\overline{D}}}{a_D^2+a_{\overline{D}}^2}}\right)\right]`$ (14)
$`\overline{\beta }_D`$ $`=`$ $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{cos}\widehat{\gamma }{\displaystyle \frac{a_{\overline{D}}b_Da_Db_{\overline{D}}}{a_D^2+b_D^2}}`$ (15)
$`\overline{\gamma }_D`$ $`=`$ $`{\displaystyle \frac{1}{a_D^2+b_D^2}}\left[\left(a_D^2b_D^2\right)+\left(\mathrm{\Delta }m/\mathrm{\Gamma }\right)\mathrm{sin}\widehat{\gamma }\left(\left(a_Da_{\overline{D}}+b_Db_{\overline{D}}\right)\left(a_D^2b_D^2\right){\displaystyle \frac{\left(a_D^2b_D^2\right)\left(a_Da_{\overline{D}}+b_Db_{\overline{D}}\right)}{a_D^2+b_D^2}}\right)\right]`$ (16)
For $`\overline{D}`$, change $`(\mathrm{\Delta }m/\mathrm{\Gamma })(\mathrm{\Delta }m/\mathrm{\Gamma }),`$ $`a_D,b_Da_{\overline{D}},b_{\overline{D}}`$ in Eqs. $`(11),(12),(13)`$ and $`\left(14\right).`$
To proceed further, we note that in the factorization ansatz $`[2]`$
$`a_{\overline{D}}`$ $`=`$ $`{\displaystyle \frac{|V_{ub}||V_{cs}|}{|V_{cb}||V_{us}|}}a_D\sqrt{\rho ^2+\eta ^2}a_D`$ (17)
$`b_{\overline{D}}`$ $`=`$ $`\sqrt{\rho ^2+\eta ^2}{\displaystyle \frac{b_D}{1+x}}.`$ (18)
$`a_D`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}|V_{cb}V_{us}|a_2F_D(m__bm_{})g_V\sqrt{E_{}+m_{}}`$ (19)
$`b_D`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}|V_{cb}V_{us}|a_2F_D(m_b+m_{})g_A\sqrt{E_{}m_{}}(1+x)`$ (20)
Here $`x=\frac{b_p}{b_f}`$ and $`b_p`$ is the baryon poles contribution which contributes only to $`b_D.b_f=a_2F_D(m_b+m_{})g_A.`$ Note that in Eq. $`(15),`$ we have used Wolfenstein parametrization $`[3]`$ of CKM matrix $`[4]`$. We will take $`g_V=g_A`$. Thus we can write
$$b_{\overline{D}}/a_{\overline{D}}=d,\text{ }b_D/a_D=d(1+x)$$
(21)
where
$$d=\frac{m__b+m_{}}{m__bm_{}}\sqrt{\frac{m_bm_{}}{m_b+m_{}}}0.946$$
(22)
on using $`m__b=5.624GeV`$ and $`m_{}=1.116GeV.`$
Using Eqs. $`(18)`$ and $`(15)`$, we obtain from Eqs. $`(11),(12),(13)`$ and $`(14).`$
$`\delta `$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }(_b+\overline{D}^0)}{(\rho ^2+\eta ^2)\mathrm{\Gamma }(_b+D^0)}}={\displaystyle \frac{1}{(\rho ^2+\eta ^2)}}{\displaystyle \frac{\overline{F}_{\overline{D}}^2}{\overline{F}_D^2}}`$ (23)
$`=`$ $`{\displaystyle \frac{1+d^2}{1+d^2(1+x)^2}}\left[1\left(\mathrm{\Delta }m/\mathrm{\Gamma }\right)\mathrm{sin}\widehat{\gamma }{\displaystyle \frac{1}{\sqrt{\rho ^2+\eta ^2}}}\left(1+{\displaystyle \frac{d^2}{1+d^2}}x\right)\right]`$ (24)
$`\overline{\alpha }_D`$ $`=`$ $`{\displaystyle \frac{2d(1+x)}{1+d^2(1+x)^2}}\left[1\sqrt{\rho ^2+\eta ^2}\left(\mathrm{\Delta }m/\mathrm{\Gamma }\right)\mathrm{sin}\widehat{\gamma }\left({\displaystyle \frac{x}{2(1+x)}}\right){\displaystyle \frac{1d^2(1+x)^2}{1+d^2(1+x)^2}}\right]`$ (25)
$`\overline{\beta }_D`$ $`=`$ $`\sqrt{\rho ^2+\eta ^2}\left(\mathrm{\Delta }m/\mathrm{\Gamma }\right)\mathrm{cos}\widehat{\gamma }{\displaystyle \frac{dx}{1+d^2(1+x)^2}}`$ (26)
$`\overline{\gamma }_D`$ $`=`$ $`{\displaystyle \frac{1}{1+d^2(1+x)^2}}\left[1d^2(1+x)^2+\sqrt{\rho ^2+\eta ^2}\left(\mathrm{\Delta }m/\mathrm{\Gamma }\right)\mathrm{sin}\widehat{\gamma }{\displaystyle \frac{2d^2x(1+x)}{1+d^2(1+x)^2}}\right]`$ (27)
$`\overline{\alpha }_{\overline{D}}`$ $`=`$ $`{\displaystyle \frac{2d}{1+d^2}}\left[1\left(\mathrm{\Delta }m/\mathrm{\Gamma }\right)\mathrm{sin}\widehat{\gamma }{\displaystyle \frac{1}{\rho ^2+\eta ^2}}{\displaystyle \frac{x(1d^2)}{2(1+d^2)}}\right]`$ (28)
$`\overline{\beta }_{\overline{D}}`$ $`=`$ $`{\displaystyle \frac{\left(\mathrm{\Delta }m/\mathrm{\Gamma }\right)\mathrm{cos}\widehat{\gamma }}{\sqrt{\rho ^2+\eta ^2}}}\text{ }{\displaystyle \frac{dx}{1+d^2}}`$ (29)
$`\overline{\gamma }_{\overline{D}}`$ $`=`$ $`{\displaystyle \frac{1d^2}{1+d^2}}+\left[1+{\displaystyle \frac{1}{\sqrt{\rho ^2+\eta ^2}}}\left(\mathrm{\Delta }m/\mathrm{\Gamma }\right)\mathrm{sin}\widehat{\gamma }{\displaystyle \frac{2d^2x}{(1d^4)}}\right]`$ (30)
We first note that the inteference effect is more pronounced in the observables for $`\overline{D}`$, since the admixture of D tends to enhance it by a factor$`\frac{1}{\sqrt{\rho ^2+\eta ^2}}`$. But since $`d^2=0.895,`$ the interference effect in $`\overline{\alpha }_{\overline{D}}`$ is neglegible. Also we note that this effect venishes in $`\alpha ,`$ $`\beta `$ and $`\gamma `$ for $`x=0`$. But there is no reason to believe that $`x=0`$ as shown in reference $`[2].`$ Just to give an estimate of the effect of $`D^0\overline{D}^0`$ mixing, using $`x=0.64`$ $`[2],`$ $`(\rho ,\eta )=(0.05,`$ $`0.36)`$ $`[5],`$ we get
$`1.38`$ $``$ $`\delta 2.03`$ (31)
$`0.006`$ $``$ $`\overline{\gamma }_{\overline{D}}0.104`$ (32)
Without interference effect $`\delta =1.70`$ and $`\overline{\gamma }_{\overline{D}}=0.055.`$ In deriving the inequality $`(27)`$, we have used the experimental $`[6]`$ upper limit $`|\mathrm{\Delta }m/\mathrm{\Gamma }|<0.10.`$
First we note that asymmetry parameter $`\beta `$ which characterizes CP-violation is a consequence of particle mixing. The experimental upper limit on $`|\mathrm{\Delta }m/\mathrm{\Gamma }|`$ gives
$`0.1(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\overline{\gamma }0.1`$
Thus if we plot $`\delta `$ and $`\overline{\gamma }_{\overline{D}}`$ as function of $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\overline{\gamma }`$ in the range $`0.1`$ to $`0.1`$, treating $`x`$ as a free parameter lying in the range $`1<x<1`$, we may be able to extract some information for $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\overline{\gamma }`$ from the experimental values of $`\delta `$ and $`\overline{\gamma }_{\overline{D}}`$. Experimentally, measurement of the ratio $`\delta `$ should not be very difficult. In Figs. $`1`$ and $`2`$ we have plotted $`\delta `$ and $`\overline{\gamma }_{\overline{D}}`$ vs $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\overline{\gamma }`$ for the four values of $`x`$ viz $`x=0.8,0.6,0.4,0.4,0.6`$ and $`0.8`$ taking $`\sqrt{\delta ^2+\eta ^2}=0.36.`$ If $`\mathrm{sin}\widehat{\gamma }`$ is too small, then $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{cos}\widehat{\gamma }`$ may be extracted form similar plot for $`\overline{\beta }_{\overline{D}}`$ as shown in $`Fig.3.`$
To conclude the mixing of $`D^0\overline{D}^0`$ has some observable effects in the decays $`_b+D^0`$ and $`_b+\overline{D}^0`$. The branching ratio $`\delta `$ for these decays can give some information to extract the value of $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma },`$ provided $`\mathrm{sin}\widehat{\gamma }`$ is not too small.
Acknowledgement
I am grateful to Prof. Riazuddin for helpfull discussions.
Figure Captions
1. Plot of $`\delta `$ vs $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma }`$ (cf. Eq. 20) for $`x=0.8,0.6,0.4,0.4,0.6`$ and $`0.8.`$
2. Plot of $`\overline{\gamma }_{\overline{D}}`$ vs $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{sin}\widehat{\gamma }`$ (cf. Eq. 26) for $`x=0.8,0.6,0.4,0.4,0.6`$ and $`0.8.`$
3. Plot of $`\overline{\beta }_{\overline{D}}`$ vs $`(\mathrm{\Delta }m/\mathrm{\Gamma })\mathrm{cos}\widehat{\gamma }`$ (cf. Eq. 25) for $`x=0.8,0.6,0.4,0.4,0.6`$ and $`0.8.`$
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# Fluctuating Filaments I: Statistical Mechanics of Helices
## I Introduction
Modern polymer physics is based on the notion that while real polymers can be arbitrarily complicated objects, their universal features are captured by a minimal model in which polymers are described as continuous random walks. While this approach has been enormously successful and led to numerous triumphs such as the understanding of rubber elasticity, the solution of the excluded volume problem and the theory of semi–dilute polymer solutions, it is ill–suited for the description of non–universal features of polymers which may depend on their chemical structure in a way that can not be captured by a simple redefinition of the effective monomer size or its second virial coefficient. For relatively simple synthetic polymers, such “local details” can be treated by polymer chemistry –type models (e.g., rotational isomer state model). However, chemically–detailed approaches become prohibitively difficult (at least as far as analytical modeling is concerned) in the case of complex biomolecules such as DNA, proteins and their assemblies and a new type of minimal models is needed to model recent mechanical experiments on such systems. Such an alternative approach is to model polymers in the way one usually thinks of them, i.e., as continuous elastic strings or filaments which can be arbitrarily deformed and twisted. However, while the theory of elasticity of such objects is well–developed, little is known about the statistical mechanics of fluctuating filaments with arbitrary natural shapes. The main difficulty is mathematical in origin: the description of three–dimensional filaments with non–circular cross–section and non–vanishing spontaneous curvature and twist, involves rather complicated differential geometry and most DNA–related theoretical studies of such models assumed circular cross–sections and focused on fluctuations around the straight rod configuration .
Recently, we reported a study of the effect of thermal fluctuations on the statistical properties of filaments with arbitrary spontaneous curvature and twist. In this work we present a detailed exposition of the theory and of its application to helical filaments. In Section II we introduce the description of the spatial configuration of the filament in terms of a triad of unit vectors oriented along the principal axes of the filament, and show that all the information about this configuration can be obtained from the knowledge of a set of generalized torsions. The elastic energy cost associated with any instantaneous configuration of the filament, is expressed in terms of the deviations of the generalized torsions that describe this configuration, from their spontaneous values in some given stress–free reference state. We use this energy to construct the statistical weights of the different configurations and show that the deviations of the generalized torsions behave as Gaussian random noises, whose amplitudes are inversely proportional to the bare persistence lengths that characterize the rigidity associated with the different deformation modes. We then derive the differential equations for the orientational correlation functions that can be expressed as averages of a rotation matrix which generates the rotation of the triad vectors as one moves along the contour of the filament. An expression for the persistence length in terms of one of the correlators is derived. In Section III we apply the general formalism to helical filaments and derive exact expressions for the correlators (see Appendix A) and for the effective persistence length of an untwisted helix. We show that the persistence length is, in general a non–monotonic function of the amplitudes of thermal fluctuations. We also show that in the weak fluctuation regime, our exact expressions for the correlators can be derived from a simplified long–wavelength description of the helix, which is equivalent to the incompressible rod–like chain model , and that the fluctuation spectrum is dominated by the Goldstone modes of this rod–like chain. Analytical expressions for the persistence length of a spontaneously twisted helix are derived (see Appendix B) and it is found that this length exhibits resonant–like dependence on the rate of twist. Finally, in Section IV we discuss our results and outline directions for future research.
## II General Theory of Fluctuating Filaments
A filament of small but finite and, in general, non–circular cross–section, is modeled as an inextensible but deformable physical curve parametrized by a contour length $`s`$ ($`0sL`$ where $`L`$ is the length of the filament). To each point $`s`$ one attaches a triad of unit vectors $`\left\{𝐭(s)\right\}`$ whose component $`𝐭_3`$ is the tangent vector to the curve at $`s`$, and the vectors $`𝐭_1(s)`$and $`𝐭_2(s)`$ are directed along the two axes of symmetry of the cross–section. The vectors $`\left\{𝐭(s)\right\}`$, together with the inextensibility condition $`d𝐱/ds=𝐭_3`$, give a complete description of the space curve $`𝐱(s)`$, as well as of the rotation of the cross–section (i.e., twist) about this curve.
The rotation of all the vectors $`𝐭_i`$ of the triad as one moves from point $`s`$ to point $`s^{}`$ along the line, is generated by the rotation matrix $`𝐑(s,s^{})`$
$$𝐭_i(s)=\underset{j}{}R_{ij}(s,s^{})𝐭_j(s^{})$$
(1)
The rotation matrix has the property
$$𝐑(s,s^{})=𝐑(s,s^{\prime \prime })𝐑(s^{\prime \prime },s^{})$$
(2)
where $`s^{\prime \prime }`$ is an arbitrary point on the contour of the filament. It satisfies the equation
$$\frac{R_{ij}(s,s^{})}{s}=\underset{k}{}\mathrm{\Omega }_{ik}(s)R_{kj}(s,s^{})$$
(3)
where
$$\mathrm{\Omega }_{ij}=\underset{k}{}\epsilon _{ijk}\omega _k$$
(4)
$`\epsilon _{ijk}`$ is the antisymmetric tensor and $`\left\{\omega _k\right\}`$ will be referred to as generalized torsions, for lack of a better term. The above equations are supplemented by the “initial” condition $`R_{ij}(s,s)=\delta _{ij}`$, where $`\delta _{ij}`$ is the Kronecker delta function. The formal solution of Eq. (3) is given by the “time–ordered” exponential
$$𝐑(s,s^{})=𝐓_s\mathrm{exp}\left(_s^{}^s𝑑s^{\prime \prime }𝛀(s^{\prime \prime })\right)=\underset{\mathrm{\Delta }s0^+}{lim}e^{𝛀(s_n)\mathrm{\Delta }s}\mathrm{}e^{𝛀(s_2)\mathrm{\Delta }s}e^{𝛀(s_1)\mathrm{\Delta }s}$$
(5)
The time–ordering operator with respect to $`s,`$ $`𝐓_s`$ is defined by the second equality in the above equation, where we broke the interval $`ss^{}`$ into $`n`$ parts of length $`\mathrm{\Delta }s`$ each, so that $`s_1=s^{}`$ and $`s_n=s^{}`$. The origin of the difficulty in calculating the above expression is that the matrices $`𝛀(s)`$ and $`𝛀(s^{})`$ do not commute for $`ss^{}`$ (this is related to the non–Abelian character of the rotation group in 3d).
Eq. (3) is equivalent to a set of generalized Frenet equations from which one can calculate the spatial configuration of the filament, given a set of generalized torsions $`\left\{\omega _k\right\}`$,
$$\frac{d𝐭_1}{ds}=\omega _2𝐭_3\omega _3𝐭_2,\frac{d𝐭_2}{ds}=\omega _1𝐭_3+\omega _3𝐭_1,\frac{d𝐭_3}{ds}=\omega _1𝐭_2\omega _2𝐭_1$$
(6)
Note that in the original Frenet description of space curves in terms of a unit tangent (which coincides with $`𝐭_3`$), normal ($`𝐧`$) and binormal ($`𝐛`$), one considers mathematical lines for which it would be meaningless to define twist about the centerline. The Frenet equations contain only two parameters: the curvature $`\kappa `$ and torsion $`\tau `$:
$$\frac{d𝐛}{ds}=\tau 𝐧,\frac{d𝐧}{ds}=\kappa 𝐭_3+\tau 𝐛,\frac{d𝐭_3}{ds}=\kappa 𝐧$$
(7)
The two frames are related through rotation by an angle $`\alpha `$ about the common tangent direction (see Fig. 1),
$$𝐭_1=𝐛\mathrm{cos}\alpha +𝐧\mathrm{sin}\alpha ,𝐭_2=𝐛\mathrm{sin}\alpha +𝐧\mathrm{cos}\alpha $$
(8)
Substituting this relation into Eqs. (6) and using Eqs. (7), we relate the generalized torsions $`\left\{\omega _k\right\}`$ to the curvature $`\kappa `$, torsion $`\tau `$ and twist angle $`\alpha `$,
$$\omega _1=\kappa \mathrm{cos}\alpha ,\omega _2=\kappa \mathrm{sin}\alpha ,\omega _3=\tau +d\alpha /ds$$
(9)
The theory of elasticity of thin rods is based on the notion that there exists a stress–free reference configuration defined by the set of spontaneous (intrinsic) torsions $`\left\{\omega _{0k}\right\}`$. The set $`\left\{\omega _{0k}\right\}`$ together with Eqs. (3) and (4) (with $`\omega _k\omega _{0k}`$) completely determines the equilibrium shape of the filament, in the absence of thermal fluctuations. Neglecting excluded–volume effects and other non–elastic interactions, it can be shown that the elastic energy associated with some actual configuration $`\left\{\omega _k\right\}`$ of the filament is a quadratic form in the deviations $`\delta \omega _k=\omega _k\omega _{0k}`$
$$U_{el}\left(\left\{\delta \omega _k\right\}\right)=\frac{kT}{2}_0^L𝑑s\underset{k}{}a_k\delta \omega _k^2$$
(10)
where $`T`$ is the temperature,$`k`$ is the Boltzmann constant, and $`a_i`$ are bare persistence lengths that depend on the elastic constants and on the principal moments of inertia with respect to the symmetry axes of the cross–section, in a model–dependent way. Thus, assuming anisotropic elasticity (with elastic moduli $`E_i`$) and a particular form of the deformation, one obtains $`a_1=E_1I_1/kT`$, $`a_2=E_1I_2/kT`$ and $`a_3=E_2(I_1+I_2)/kT`$ where $`I_i`$ are the principal moments of inertia. In general, the theory of elasticity of incompressible isotropic rods with shear modulus $`\mu `$ yields $`a_1=3\mu I_1/kT`$, $`a_2=3\mu I_2/kT`$, and $`a_3=C/kT`$ where the torsional rigidity $`C`$ is also proportional to $`\mu `$ and depends on the geometry of the cross–section (for an elliptical cross–section with semi–axes $`b_1`$ and $`b_2`$, $`C=\pi \mu b_1^3b_2^3/(b_1^2+b_2^2)`$ ).
The elastic energy $`U_{el}\left(\left\{\delta \omega _k\right\}\right)`$ determines the statistical weight of the configuration $`\left\{\omega _k\right\}`$. The statistical average of any functional of the configuration $`B(\left\{\omega _k\right\})`$ is defined as the functional integral
$$B\left(\left\{\omega _k\right\}\right)=\frac{D\left\{\delta \omega _k\right\}B\left(\left\{\omega _k\right\}\right)e^{U_{el}\left\{\delta \omega _k\right\}/kT}}{D\left\{\delta \omega _k\right\}e^{U_{el}\left\{\delta \omega _k\right\}/kT}}$$
(11)
Calculating the corresponding Gaussian path integrals we obtain
$$\delta \omega _i(s)=0,\delta \omega _i(s)\delta \omega _j(s^{})=a_i^1\delta _{ij}\delta (ss^{})$$
(12)
We conclude that fluctuations of generalized torsions at two different points along the filament contour are uncorrelated, and that the amplitude of fluctuations is inversely proportional to the corresponding bare persistence length.
The statistical properties of fluctuating filaments are determined by the orientational correlation functions, which can be expressed as averages of the elements of the rotation matrix,
$$𝐭_i(s)𝐭_j(s^{})=R_{ij}(s,s^{})=\underset{k}{}R_{ik}(s,s^{\prime \prime })R_{kj}(s^{\prime \prime },s^{})$$
(13)
The last equality was written using Eq. (2), with $`s>s^{\prime \prime }>s^{}`$. Inspection of Eqs. (5) and (4), shows that $`𝐑(s,s^{\prime \prime })`$ depends only the torsions $`\omega _k(s_1)`$ with $`s>s_1>s^{\prime \prime },`$ and that $`𝐑(s^{\prime \prime },s^{})`$ depends only on $`\omega _k(s_2)`$ with $`s^{\prime \prime }>s_2>s^{}.`$ Since fluctuations of the torsion in two non–overlapping intervals are uncorrelated (see Eq. (12)), the average of the product of rotation matrices splits into the product of their averages:
$$R_{ij}(s,s^{})=\underset{k}{}R_{ik}(s,s^{\prime \prime })R_{kj}(s^{\prime \prime },s^{})$$
(14)
In order to derive a differential equation for the averaged rotation matrix, we consider the limit $`\mathrm{\Delta }s=ss^{\prime \prime }0`$. Keeping terms to first order in $`\mathrm{\Delta }s`$ we find
$$\frac{R_{ij}(s,s^{})}{s}=\underset{k}{}\mathrm{\Lambda }_{ik}(s)R_{kj}(s,s^{})$$
(15)
where the matrix $`𝚲`$ is defined as
$$\mathrm{\Lambda }_{ik}\left(s\right)=\underset{\mathrm{\Delta }s0^+}{lim}\frac{\delta _{ik}R_{ik}(s,s\mathrm{\Delta }s)}{\mathrm{\Delta }s}$$
(16)
Analogously to Eq. (5), the formal solution of Eq. (15) can be written as a time–ordered exponential,
$$𝐑(s,s^{})=𝐓_s\mathrm{exp}\left(_s^{}^s𝑑s^{\prime \prime }𝚲(s^{\prime \prime })\right)$$
(17)
In order to calculate the matrix $`𝚲`$ we expand the exponential in Eq. (5) to second order in $`\mathrm{\Delta }s=ss^{}`$ and use the property of time–ordering operator
$$𝐓_s_{s\mathrm{\Delta }s}^s𝑑s_1_{s\mathrm{\Delta }s}^s𝑑s_2𝛀(s_1)𝛀(s_2)=_{s\mathrm{\Delta }s}^s𝑑s_1\left[_{s\mathrm{\Delta }s}^{s_1}𝑑s_2𝛀(s_1)𝛀(s_2)+_{s_1}^s𝑑s_2𝛀(s_2)𝛀(s_1)\right]$$
(18)
In order to average this equation, we first calculate the average of the product $`𝛀(s_1)𝛀(s_2),`$ using Eqs. (4) and (12)
$$𝛀(s_1)𝛀(s_2)=𝛀(s_1)𝛀(s_2)+\mathrm{𝐝𝐢𝐚𝐠}(\gamma _i)\delta (s_1s_2)$$
(19)
where $`\mathrm{𝐝𝐢𝐚𝐠}(\gamma _i)`$ is a diagonal matrix with elements
$$\gamma _i=\underset{k}{}\frac{1}{2a_k}\frac{1}{2a_i}$$
(20)
Using Eqs. (18) and (19), and keeping terms up to first order in $`\mathrm{\Delta }s`$ (upon integration, the contribution of $`𝛀(s_1)𝛀(s_2)`$ is of order $`\left(\mathrm{\Delta }s\right)^2`$), yields
$$\mathrm{\Lambda }_{ik}=\gamma _i\delta _{ik}+\underset{l}{}\epsilon _{ikl}\omega _{0l}\text{ }$$
(21)
The elements of the averaged rotation matrix are simply the correlators of the triad vectors (see Eq. (13)). From the knowledge of the above correlators one can calculate other statistical properties of fluctuating filaments, the most familiar of which is the persistence length $`l_p`$, that can be interpreted as an effective statistical segment length of a coarse–grained model, in which one replaces the filament by a random walk with the same contour length $`L`$ and rms end–to–end separation $`r^2`$:
$$l_p=\underset{L\mathrm{}}{lim}\frac{1}{L}r^2$$
The end–to–end vector is defined as $`𝐫=_0^L𝐭_3(s)𝑑s`$ and thus
$$l_p=\underset{L\mathrm{}}{lim}\frac{2}{L}_0^L𝑑s_0^s𝑑s^{}𝐭_3(s)𝐭_3(s^{})$$
(22)
The above equations describe the fluctuations of filaments of arbitrary shape and elastic properties, and in the following this general formalism is applied to helical filaments.
## III Fluctuating Helices
### III.1 Untwisted Helix: Correlation Functions and Persistence Length
Consider a helical filament without spontaneous twist, such that the generalized spontaneous torsions $`\left\{\omega _{0k}\right\}`$ are independent of position $`s`$ along the contour. In order to describe the stress–free configuration of such a filament, it is convenient to introduce the conventional Frenet triad which consists of the tangent, normal and binormal to the space curve spanned by the centerline, supplemented by a constant angle of twist $`\alpha _0`$ which describes the orientation of the cross–section in the plane normal to the centerline. According to the general relation between the two frames, Eq. (9), $`\omega _{01}=\kappa _0\mathrm{cos}\alpha _0`$, $`\omega _{02}=\kappa _0\mathrm{sin}\alpha _0`$ and $`\omega _{03}=\tau _0`$, where $`\kappa _0`$ and $`\tau _0`$ are the constant curvature and torsion of the space curve in terms of which the total spontaneous curvature that defines rate of rotation of the helix about its long axis, is given by $`\omega _0=(\kappa _0^2+\tau _0^2)^{1/2}`$. The corresponding helical pitch is $`2\pi \tau _0/\omega _0^2`$ and the radius of the helical turn is $`2\pi \kappa _0/\omega _0^2`$. We proceed to calculate the orientational correlation functions.
Since $`𝚲`$ is a constant matrix, Eq. (17) yields (for $`s_1>s_2`$)
$$𝐭_i\left(s_1\right)𝐭_j\left(s_2\right)=\left[e^{𝚲\left(s_1s_2\right)}\right]_{ij}$$
(23)
In order to calculate the matrix $`e^{𝚲\left(s_1s_2\right)}`$ we first find the eigenvalues $`\lambda _i`$ of the matrix $`𝚲`$, which are determined by the characteristic polynomial
$$\lambda ^3\gamma \lambda ^2+\mu \lambda \nu =0$$
(24)
where we introduced the notations
$`\gamma `$ $`=`$ $`\gamma _1+\gamma _2+\gamma _3=a_1^1+a_2^1+a_3^1,`$ (25)
$`\mu `$ $`=`$ $`\omega _0^2+\gamma _1\gamma _2+\gamma _2\gamma _3+\gamma _1\gamma _3,`$ (26)
$`\nu `$ $`=`$ $`\kappa _0^2\left(\gamma _1\mathrm{cos}^2\alpha _0+\gamma _2\mathrm{sin}^2\alpha _0\right)+\tau _0^2\gamma _3+\gamma _1\gamma _2\gamma _3`$ (27)
The solution of this cubic equation depends on the sign of the expression
$$\mathrm{\Delta }=27\left(\nu \nu _1\right)^2+4\left(\mu \gamma ^2/3\right)^3,\nu _1=\frac{1}{3}\gamma \mu \frac{2}{27}\gamma ^3$$
(28)
For $`\mathrm{\Delta }<0`$ all the roots $`\lambda _i`$ are real. In this parameter range, fluctuations are strong enough to destroy the helical structure on all length scales. In the limit of very strong fluctuations when the bare persistence lengths are much smaller than the radii of curvature $`\gamma \omega _0,`$ we have $`\lambda _i\gamma _i`$ and correlation functions become
$$𝐭_i\left(s_1\right)𝐭_j\left(s_2\right)=e^{\gamma _i\left(s_1s_2\right)}\delta _{ij}$$
(29)
with $`s_1s_2>0`$. Eq. (29) shows that although angular correlations remain on length scales smaller than $`1/\lambda _i`$, they are identical to those of a persistent rod and do not carry any memory of the original helix.
In the case $`\mathrm{\Delta }>0`$, there is one real eigenvalue, $`\lambda _1`$, and two complex ones, $`\lambda _{2,3}=\lambda _R\pm i\omega `$, where
$`\lambda _1`$ $`=`$ $`{\displaystyle \frac{K}{6}}2{\displaystyle \frac{\mu \gamma ^2/3}{K}}+{\displaystyle \frac{\gamma }{3}},\lambda _R={\displaystyle \frac{\gamma \lambda _1}{2}}`$ (30)
$`\omega `$ $`=`$ $`\sqrt{3}\left({\displaystyle \frac{K}{12}}+{\displaystyle \frac{\mu \gamma ^2/3}{K}}\right),K=12^{1/3}\left[9\left(\nu \nu _1\right)+\sqrt{3\mathrm{\Delta }}\right]^{1/3}`$ (31)
It is shown in Appendix A that the diagonal orientational correlation functions take the form
$$𝐭_i\left(s_1\right)𝐭_i\left(s_2\right)=\left(1c_ic_i^{}\right)e^{\lambda _1s}+\left(c_ie^{i\omega s}+c_i^{}e^{i\omega s}\right)e^{\lambda _Rs}$$
(32)
where $`s=s_1s_2>0`$. The complex coefficients $`c_i`$ are calculated in Appendix A.
In the limit of small fluctuations, $`\gamma \omega _0`$, we have
$$\lambda _1=\underset{i}{}\left(12c_i\right)\gamma _i,\lambda _R=\underset{i}{}c_i\gamma _i,2c_i=1\frac{\omega _{0i}^2}{\omega _0^2},\omega ^2=\omega _0^2$$
(33)
In this limit, it is easy to generalize our results for the diagonal correlators and write down expressions for all the orientational correlation functions:
$$𝐭_i(s_1)𝐭_j(s_2)=\frac{\omega _{0i}\omega _{0j}}{\omega _0^2}e^{\lambda _1s}+\left(\delta _{ij}\frac{\omega _{0i}\omega _{0j}}{\omega _0^2}\right)\mathrm{cos}(\omega _0s)e^{\lambda _Rs}\underset{k}{}\epsilon _{ijk}\frac{\omega _{0k}}{\omega _0}\mathrm{sin}(\omega _0s)e^{\lambda _Rs}$$
(34)
where $`s=s_1s_2>0`$. As expected, Eq. (34) satisfies the condition of orthonormality of triad vectors $`𝐭_i(s_1)𝐭_j(s_1)=\delta _{ij}`$ (this geometric condition must be satisfied for the instantaneous triad vectors, not only on the average). Note that in the limit of weak fluctuations the local helical structure is preserved on contour distances $`s<\lambda _R^1`$ and the period of rotation of the helix about its axis is given by its spontaneous value, $`2\pi \omega _0^1`$.
Using Eqs. (25)–(28) it can be shown that when $`\mathrm{\Delta }0`$, the total curvature of the helix vanishes as $`\omega \mathrm{\Delta }^{1/2}`$. Since $`\omega `$is positive for $`\mathrm{\Delta }>0`$ and vanishes for $`\mathrm{\Delta }0`$, in a loose sense it plays the role of an order parameter associated with helical order, and the point $`\mathrm{\Delta }=0`$ can be interpreted as the critical point at which a continuous helix to random coil transition takes place. However, although the dependence of $`\omega `$ on the various parameters exhibits surprisingly rich behavior, the investigation of the transition region is of limited physical significance. The change of the helical period from $`2\pi \omega _0^1`$ to infinity takes place in the “overdamped” regime where this period is larger than the persistence length ($`\omega \gamma `$), and local helical structure can no longer be defined in a statistically significant sense. An approximate but more physically meaningful criterion for the “melting” transition is that a helix of period $`2\pi \omega ^1`$ melts when the persistence length becomes of order of this period.
We now return to Eq. (22) for the persistence length. Using the matrix equation $`_0^{\mathrm{}}𝑑s\mathrm{exp}(𝚲s)=𝚲^1`$ and taking the appropriate matrix element we find:
$$l_p=2\frac{\tau _0^2+\gamma _1\gamma _2}{\kappa _0^2\left(\gamma _1\mathrm{cos}^2\alpha _0+\gamma _2\mathrm{sin}^2\alpha _0\right)+\left(\tau _0^2+\gamma _1\gamma _2\right)\gamma _3}$$
(35)
The above expression diverges in the limit of a rigid helix $`\gamma _i0`$ in which fluctuations have a negligible effect on the helix. Non–monotonic behavior is observed for “plate–like” helices, with large radius to pitch ratio, $`\kappa _0/\tau _0`$. When no thermal fluctuations are present ($`\gamma _i0`$), the effective persistence length approaches zero. Weak thermal fluctuations “inflate” the helix by releasing stored length (by a mechanism similar to the stretching of the “slinky” toy spring) and increase the persistence length. Eventually, in the limit of strong fluctuations, the persistence length vanishes again (as $`\gamma _3^1)`$ because of the complete randomization of the filament. Note that the sensitivity to the (constant) angle of twist increases with radius to pitch ratio.
In the opposite limit of “rod–like” helices $`\kappa _00`$, the effective persistence length approaches $`2/\gamma _3`$ and therefore depends on $`a_1`$ and $`a_2`$ only and not on $`\tau _0`$ and $`a_3`$ which describe the twist of the cross–section about the centerline. This agrees with the expectation that since straight inextensible rods do not have stored length, their end–to–end distance and persistence length are determined by random bending and torsion (writhe) fluctuations only and are independent of twist.
### III.2 Weak Fluctuations: The Rod–Like Chain Model
From the discussion in the preceding section we expect that in the presence of weak thermal fluctuations, the filament will maintain its helical structure locally and that fluctuations will only affect its large scale conformation by introducing random bending and torsion of the helical axis, as well as random rotation of the filament about this axis. We now rederive the expressions for the correlators, Eq. (34), using a different approach that relates our work to that of previous investigators and, in the process, leads to important insights about the nature of the long wavelength fluctuations that dominate the spectrum of fluctuations in this regime.
Note that in the absence of thermal fluctuations, $`\gamma _i=0`$, the triad vectors $`𝐭_i`$ attached to the helix can be expressed in terms of the space–fixed orthonormal triad $`\left\{𝐞\right\}`$ of vectors $`𝐞_i`$, where $`𝐞_3`$ is oriented along the long axis of the helix and $`𝐞_1`$ and $`𝐞_2`$ lie in the plane normal to it (Fig. 2). It is convenient to introduce the Euler angles $`\varphi _0(s)=\omega _0s,`$ $`\theta _0=\mathrm{arctan}(\kappa _0/\tau _0)`$ and $`\alpha _0`$ in terms of which the relation between the two frames is given by
$$𝐭^R(s)=𝐑_3(\alpha _0)𝐑_2(\theta _0)𝐑_3\left[\varphi _0(s)\right]𝐞,$$
(36)
where the rotation matrix
$$𝐑_3\left(\varphi _0\right)=\left(\begin{array}{ccc}\mathrm{cos}\varphi _0& \mathrm{sin}\varphi _0& 0\\ \mathrm{sin}\varphi _0& \mathrm{cos}\varphi _0& 0\\ 0& 0& 1\end{array}\right)$$
(37)
describes rotation by angle $`\varphi _0(s)`$ with respect to the $`𝐞_3`$ axis. The matrix
$$𝐑_2(\theta _0)=\left(\begin{array}{ccc}\mathrm{cos}\theta _0& 0& \mathrm{sin}\theta _0\\ 0& 1& 0\\ \mathrm{sin}\theta _0& 0& \mathrm{cos}\theta _0\end{array}\right)$$
(38)
gives the rotation by angle $`\theta _0`$ with respect to the $`𝐞_2^{}`$ axis ($`𝐞_2^{}=`$ $`𝐑_3\left[\varphi _0(s)\right]𝐞_2`$), and $`𝐑_3(\alpha _0)`$ is a rotation by angle $`\alpha _0`$ about the $`𝐞_3^{}`$ axis ($`𝐞_3^{}=𝐑_2(\theta _0)𝐞_3`$). Note that while the space–fixed $`𝐞`$ was taken as a conventional right–handed triad, we chose the helix–fixed $`𝐭`$ as a left–handed triad. Although this choice does not affect our previous results, it does affect the geometric relation between the two coordinate systems and, for consistency, we replaced the left–handed $`𝐭`$ by the right–handed one, $`𝐭^R=(𝐭_1,𝐭_2,𝐭_3)`$, in Eq. (36).
In the presence of weak thermal fluctuations, the axis of the helix slowly bends and rotates in space, resulting in rotation of the triad $`\left\{𝐞\right\}`$. Since with each point $`s`$ on the helix we can associate its projection
$$\sigma =\tau _0s/\omega _0$$
(39)
on the long axis of the helix (see Fig. 2), the rotation of the triad $`\left\{𝐞\right\}`$ as one moves along this axis is given by the generalized Frenet equations,
$$\frac{d𝐞_1}{d\sigma }=\varpi _2𝐞_3\varpi _3𝐞_2,\frac{d𝐞_2}{d\sigma }=\varpi _1𝐞_3+\varpi _3𝐞_1,\frac{d𝐞_3}{d\sigma }=\varpi _1𝐞_2\varpi _2𝐞_1$$
(40)
The generalized torsions, $`\varpi _i\left(s\right)`$, are Gaussian random variables determined by the conditions
$$\varpi _i\left(\sigma \right)=0,\varpi _i\left(\sigma \right)\varpi _j\left(\sigma ^{}\right)=\overline{a}_i^1\delta _{ij}\delta \left(\sigma \sigma ^{}\right)$$
(41)
where the constants $`\overline{a}_i`$ should be determined by the requirement that the resulting expressions for the correlators (the averages of the elements of the rotation matrix) coincide with these in Eq. (34). A calculation similar to that in the previous section yields the correlators
$$𝐞_i\left(\sigma \right)𝐞_j\left(\sigma ^{}\right)=\delta _{ij}\mathrm{exp}\left(\overline{\gamma }_i\left|\sigma \sigma ^{}\right|\right),$$
(42)
where, analogously to Eq. (20), we have
$$\overline{\gamma }_i=\underset{k}{}\frac{1}{2\overline{a}_k}\frac{1}{2\overline{a}_i}$$
(43)
Using Eqs. (36), the correlators of the original triad $`\left\{𝐭\right\}`$ can be expressed in terms of the correlators of the $`\left\{𝐞\right\}`$ triad. Comparing the results with Eq. (34), gives
$$\overline{a}_1^1=\overline{a}_2^1=\underset{i}{}\gamma _i\frac{\omega _{0i}^2}{\omega _0\tau _0},\overline{a}_3^1=\underset{i}{}\frac{1}{a_i}\frac{\omega _{0i}^2}{\omega _0\tau _0}$$
(44)
where the equality $`\overline{a}_1=\overline{a}_2`$ is the consequence of symmetry under rotation in the $`(𝐞_1,𝐞_2)`$ plane.
The correlators (41) can be derived from an effective free energy which describes the long wavelength fluctuations of the helical filament, on length scales larger than the period of the helix $`\omega _0^1`$.
$$U_{el}^{LW}=\frac{kT}{2}𝑑\sigma \left[\overline{a}_1\left(\varpi _1^2+\varpi _2^2\right)+\overline{a}_3\varpi _3^2\right]$$
(45)
This expression coincides with the elastic energy of a rod–like chain (RLC) introduced by Bouchiat and Mezard. The persistence length $`\overline{a}_1`$ describes the elastic response to bending and torsion of the effective rod–like filament. The persistence length $`\overline{a}_3`$ controls the elastic response of the RLC to twist about its axis. As a consequence of the fluctuation–dissipation theorem, it also determines the amplitude of fluctuations $`\mathrm{\Delta }\varphi `$ of the angle $`\varphi (\sigma )=\omega _0^2\sigma /\tau _0+\mathrm{\Delta }\varphi (\sigma )`$, where the correlator of the random angle of rotation about the axis of the RLC is given by
$$\left[\mathrm{\Delta }\varphi (\sigma )\mathrm{\Delta }\varphi (\sigma ^{})\right]^2=\overline{a}_3^1\left|\sigma \sigma ^{}\right|$$
(46)
In Eq. (44) we calculated the effective persistence lengths of this model ($`\overline{a}_i`$) in terms of the bare parameters of the underlying helical filament. In reference where the analysis begins with the RLC model, these corresponding persistence lengths were introduced by hand. The difference between the two models becomes important if one considers the combined application of extension and twist: while such a coupling appears trivially in models of stretched helical filaments, twist has no effect on the extension in the RLC model, in contradiction with experimental observations . Our analysis underscores the fact that the RLC model does not give a complete description of the fluctuating helix. Rather, it describes long wavelength fluctuations of the “phantom” axes $`\left\{𝐞_i\right\}`$ which, by themselves, contain no information about the local helical structure of the filament. In order to recover this information and construct the correlators of the original helix $`𝐭_i\left(s_1\right)𝐭_i\left(s_2\right)`$, one has to go beyond the RLC model and reconstruct the local helical geometry using the relation between $`𝐞_i`$ and the helix–fixed axes $`𝐭_i`$, Eqs. (36).
In deriving the expressions for the correlators $`𝐭_i(s)𝐭_j(0)`$ in terms of the correlators of the RLC model, we did not take into account the possibility of fluctuations of the twist angle of the cross–section of the helix about its centerline, $`\alpha _0\alpha (s)=\alpha _0+\mathrm{\Delta }\alpha (s)`$. From the fact that the resulting correlators coincide with the exact expressions, Eq. (34), we conclude that such fluctuations do not contribute to the correlators. This surprising result follows from the fact that in the weak fluctuation regime, the statistical properties of the helix are completely determined by the low energy part of the fluctuation spectrum. Such long–wavelength fluctuation modes (Goldstone modes) lead to the loss of helical correlations on length scales larger than all the natural length scales of the helix ($`s\gamma ^1\omega _0^1`$). These Goldstone modes are associated with spontaneously broken continuous symmetries and correspond to bending ($`\varpi _1`$ and $`\varpi _2`$) and twist ($`\varpi _3`$) modes of the RLC. It is important to emphasize that these modes correspond to different deformations of the centerline of the helix and not to twist of its cross–section about this centerline. Since the elastic energy, Eq. (10), depends on the spontaneous angle of twist of the helix about its centerline through the combinations $`\delta \omega _1=\kappa \mathrm{cos}\alpha \kappa _0\mathrm{cos}\alpha _0`$ and $`\delta \omega _2=\kappa \mathrm{sin}\alpha \kappa _0\mathrm{sin}\alpha _0,`$ we conclude that the energy is not invariant under global rotation of the cross–section about the centerline and that such a rotation is not a continuous symmetry of the helix. Therefore, twist fluctuations of the helical cross–section are not Goldstone modes and do not contribute to the correlators in the weak fluctuation limit.
Another interesting observation is that there is no contribution from compressional modes to the long–wavelength energy, Eq. (10). This is surprising since the RLC is a coarse–grained representation of the helix and the latter may be expected to behave as a compressible object, with accordion–like compressional modes. In order to check this point, we write down the spatial position of a point $`s`$ on the helix as
$$𝐱(s)=\overline{𝐱}(\sigma )+\delta 𝐱(s)$$
(47)
where $`\overline{𝐱}(\sigma )`$ describes the curve spanned by the long axis of the helix and, therefore, defines the spatial position of the point $`\sigma `$, Eq. (39), on the RLC contour. The deviation $`\delta 𝐱(s)`$ describes the rotation of the locally helical filament about this axis. Since the original filament is incompressible, it satisfies $`d𝐱/ds=𝐭_3`$. From Eq. (36) we obtain an expression for $`𝐭_3`$ which, upon substitution into the incompressibility condition and averaging over length scales $`\left\{|\varpi _i|^1\right\}s`$ $`\omega _0^1`$ (much larger than the inverse total curvature of the helix but much smaller than the radii of curvature of the RLC), yields
$$\frac{d\overline{𝐱}(\sigma )}{d\sigma }=𝐞_3(\sigma )$$
(48)
The fact that the long–wavelength fluctuations of the helix satisfy the above incompressibility conditions, implies that compressional fluctuations do not contribute to the long–wavelength correlators. The origin of this observation becomes clear if we recall that the energy of the helix depends on the spontaneous curvature $`\kappa _0`$ and torsion $`\tau _0`$ and, since compressional modes change the local curvature and torsion, they have a gap in the energy spectrum and their energy does not vanish even in the long–wavelength limit. We conclude that similarly to twist fluctuations of the helical cross–section, compressional modes are not Goldstone modes.
The above deliberations have profound consequences for the elastic response of the filament to long–wavelength perturbations, such as tensile forces and moments applied to its ends. Using the fluctuation–dissipation theorem, we conclude that as long as the deformation of the filament remains small (on scale $`\omega _0^1`$), these forces and moments do not induce twist of the cross–section of the helix about its centerline, and that the deformation can be completely described by the incompressible RLC model.
### III.3 Effect of Spontaneous Twist
We proceed to calculate the persistence length of a helix whose cross–section is twisted by an angle $`\alpha _0(s)=\dot{\alpha }_0s`$ about the centerline ($`\dot{\alpha }_0`$ is a constant rate of twist). It is convenient to rewrite Eq. (22) as:
$$l_p=\underset{L\mathrm{}}{lim}\frac{2}{L}_0^L𝑑s^{}_0^{Ls^{}}𝑑s𝐭_3(s+s^{})𝐭_3(s^{})$$
(49)
Recall that the correlator in the integrand of Eq. (49) is simply the $`33`$ element of the averaged rotation matrix, and is therefore the solution of equation Eq. (15), the coefficients of which are the elements of the matrix $`𝚲(s+s^{})`$ defined in Eq. (21). The diagonal elements of this matrix are constants ($`\gamma _i`$), while the non–diagonal elements are given by the expressions
$`\mathrm{\Lambda }_{12}(s+s^{})`$ $`=`$ $`\mathrm{\Lambda }_{21}(s+s^{})=\tau _0+\dot{\alpha }_0,`$
$`\mathrm{\Lambda }_{31}(s+s^{})`$ $`=`$ $`\mathrm{\Lambda }_{13}(s+s^{})=\kappa _0\mathrm{sin}\left(\dot{\alpha }_0s+\alpha _0\right),`$
$`\text{ }\mathrm{\Lambda }_{23}(s+s^{})`$ $`=`$ $`\mathrm{\Lambda }_{32}(s+s^{})=\kappa _0\mathrm{cos}\left(\dot{\alpha }_0s+\alpha _0\right),`$ (50)
where all the dependence on $`s^{}`$ is contained in $`\alpha _0=\alpha _0(s^{}).`$
The correlator in Eq. (49) decays exponentially fast with $`s`$, and thus the upper limit on the integral over $`s`$ can be extended to infinity. Since the correlator is a periodic function of $`\alpha _0,`$ the integration over $`s^{}`$ can be replaced by that over $`\alpha _0`$ and we obtain
$$l_p=_0^{2\pi }\frac{d\alpha _0}{\pi }_0^{\mathrm{}}𝑑s𝐭_3(s)𝐭_3(ss_1)$$
(51)
In deriving the above expression we assumed that the limit $`L\mathrm{}`$ is taken and that the total angle of twist is always large, $`L\dot{\alpha }_0`$ $`2\pi `$ (i.e., the product $`L\dot{\alpha }_0`$ remains finite for arbitrarily small $`\dot{\alpha }_0`$). This assumption will be used in the following analysis.
We first consider some limiting cases in which analytical results can be derived. In the limit of vanishing twist rates, $`\dot{\alpha }_00,`$ the persistence length is obtained by averaging Eq. (35) with respect to $`\alpha _0`$. This yields:
$$l_p=\frac{2(\tau _0^2+\gamma _+^2\gamma _{}^2)}{\sqrt{[\kappa _0^2\gamma _++(\tau _0^2+\gamma _+^2\gamma _{}^2)\gamma _3]^2\kappa _0^4\gamma _{}^2}}$$
(52)
where
$$\gamma _\pm (\gamma _1\pm \gamma _2)/2$$
(53)
with $`\gamma _{1\text{ }}`$and $`\gamma _2`$ defined in (20).
In the limit of large twist rates, $`\dot{\alpha }_0\mathrm{},`$ we can replace the denominator of Eq. (35) by its average with respect to $`\alpha _0`$. This yields
$$l_p=\frac{2\left(\tau _0^2+\gamma _+^2\gamma _{}^2\right)}{\kappa _0^2\gamma _++\left(\tau _0^2+\gamma _+^2\gamma _{}^2\right)\gamma _3}$$
(54)
Finally, when $`\gamma _1=\gamma _2`$ ($`a_1=a_2`$), the persistence length becomes independent of twist and can be derived from either of Eqs. (52) and (54), by substituting $`\gamma _{}=0.`$
We now consider the case of arbitrary twist rates and fluctuation amplitudes. The calculation involves the solution of linear differential equations with periodic coefficients and details are given in Appendix B. We obtain:
$$l_p=\frac{2\gamma _3^1}{1+\left(\mathrm{\Xi }1\right)^1+\left(\mathrm{\Xi }^{}1\right)^1}$$
(55)
An analytical expression for the complex function $`\mathrm{\Xi }(\dot{\alpha }_0)`$ is given in Appendix B.
In Fig. 3 we present a three–dimensional plot of the persistence length given in units of the helical pitch $`l^{}=l\omega _0^2/2\pi \tau _0,`$ as a function of the dimensionless rate of twist $`w=2\omega _0^1\dot{\alpha }_0`$ and of the logarithm of the bare persistence length $`a_1`$, for a “plate–like” helix with large radius to pitch ratio $`\kappa _0/\tau _0`$. Inspection of Fig. 3 shows that in the case of a circular cross–section with $`a_1=a_2=1000`$, the persistence length becomes independent of twist. With increasing asymmetry, $`a_1<a_2`$, a maximum appears at vanishing twist rates, accompanied by two minima at $`\dot{\alpha }_0=\pm \omega _0/2`$. The geometrical significance of the locations ($`\dot{\alpha }_0=0,\pm \omega _0/2`$) of these resonances is underscored by the observation that in the limit of vanishing pitch, a ribbon–like untwisted ($`\dot{\alpha }_0=0`$) helix degenerates into a ring. For $`\dot{\alpha }_0=\pm \omega _0/2`$, the cross–section of a twisted helix rotates by $`\pm \pi `$ with each period, and in the above limit the helix degenerates into a Mőbius ring. As asymmetry increases ($`a_1a_2`$), each extremum splits into a minimum and a maximum and eventually one obtains a dip at $`\dot{\alpha }_0=0`$, accompanied by two symmetrical peaks at $`\dot{\alpha }_0\pm \omega _0/2`$. Note that the persistence length is a non–monotonic function of the amplitude of thermal fluctuations (i.e., of $`1/a_1`$): it first slowly increases and eventually decreases rapidly with decreasing $`a_1.`$ Several two–dimensional plots of the persistence length as a function of the rate of twist, for different combinations of the bare persistence lengths $`a_i`$ are shown in Fig. 4. The detailed behavior of the persistence length depends sensitively on the choice of the parameters: for example, in the limit of weak fluctuations three maxima are observed in Fig. 4, instead of a maximum accompanied by two minima in Fig. 3. In all cases, the locations of the extrema are determined by geometry only: $`\dot{\alpha }_0=0`$, $`\pm \omega _0/2`$.
In order to demonstrate how the initial choice of the handedness of the helix breaks the symmetry between the effects of under and over–twist on the persistence length, in Fig. 5 we present a three–dimensional plot of the persistence length as a function of the dimensionless rate of twist $`w`$ and of the inverse radius of curvature $`\kappa _0`$, for helices with radius to pitch ratios of order unity and large asymmetry of the cross–section, $`a_1a_2`$. Note that for $`\kappa _0/\tau _0<1`$ (rod–like helices), there is a single broad maximum at $`\dot{\alpha }_0=\omega _0/2`$. Then, at $`\kappa _0/\tau _01,`$ a central peak appears at $`\dot{\alpha }_0=0`$. This peak grows much faster than the $`\dot{\alpha }_0=\omega _0/2`$ peak, with increasing $`\kappa _0/\tau _0.`$ At yet higher values of $`\kappa _0/\tau _0`$ another peak appears at $`\dot{\alpha }_0`$ $`=\omega _0/2`$ and eventually the amplitudes of the two Mőbius side–peaks become equal (and much smaller than the amplitude of the $`\dot{\alpha }_0`$=$`0`$ peak) in the limit of plate–like helices, $`\kappa _0/\tau _01`$ (see curve 1 in Fig. 4).
What is the origin of the Mőbius resonances observed in Figs. 3–5? Recall that the calculation of the persistence length of a twisted helix involves the solution of linear differential equations with periodic coefficients (Eqs. (65) in Appendix B). These equations were derived from linear differential equations with periodic coefficients and multiplicative random noise, Eqs. (3) and Eqs. (6), which are known to lead to stochastic resonances. Some physical intuition can be derived from the following argument. While the persistence length is a property of the space curve described by the Frenet triad, the microscopic Brownian motion of the filament arises as the result of random forces that act on its cross–section and therefore are given in the frame associated with the principal axes of the filament. Since the two frames are related by a rotation of the cross–section by an angle $`\alpha _0(s)`$, the random force in the Frenet frame is modulated by linear combinations of $`\mathrm{sin}\alpha _0(s)`$ and $`\mathrm{cos}\alpha _0(s)`$. This gives a deterministic contribution to the persistence length which, to lowest order in the force, is proportional to the mean square amplitude of the random force and therefore varies sinusoidally with $`\pm 2\alpha _0(s)`$. The Mőbius resonances occur whenever the total curvature of the helix $`\omega _0`$ coincides with the rate of variation of this deterministic contribution of the random force, $`\pm 2\dot{\alpha }_0`$.
## IV Discussion
In this work we studied the statistical mechanics of thermally fluctuating elastic filaments with arbitrary spontaneous curvature and twist. We constructed the equations for the orientational correlation functions and for the persistence length of such filaments. We would like to stress that our theory describes arbitrarily large deviations of a long filament from its equilibrium shape; the only limitation is that fluctuations are small on microscopic length scales, of the order of the thickness of the filament. Furthermore, since the equilibrium shape and the fluctuations of the filaments are completely described by the set of spontaneous torsions $`\left\{\omega _{0k}\right\}`$ and its fluctuations $`\left\{\delta \omega _k\right\}`$ respectively, our theory is set up in the language of intrinsic geometry of the space curves. All the interesting statistical information is contained in the correlators of the triad vectors $`\left\{𝐭\right\}`$ which can be expressed in terms of the known correlators of the fluctuations $`\left\{\delta \omega _k\right\}`$, using the Frenet equations. Since these equations describe pure rotation of the triad vectors, this has the advantage that fluctuations of the torsions introduce only random rotations of the vectors of the triad, and preserve their unit norm. The use of intrinsic geometry automatically ensures that the inextensibility constraint is not violated in the process of thermal fluctuations and therefore does not even have to be considered explicitly in our approach. We would like to remind the readers that the formidable mathematical difficulties associated with attempts to introduce this constraint, have hindered the development of persistent chain type models in the past and led to the introduction of the mean spherical approximation in which the constraint is enforced only on the average, and to perturbative expansions about the straight rod limit.
The general formalism was then applied to helical filaments both with and without twist of the cross–section about the centerline. In the latter case we found that weak thermal fluctuations are dominated by long wavelength Goldstone modes that correspond to bending and twist of the coarse–grained filament (the rod–like chain). Such fluctuations distort the helix on length scales much larger than its natural period but do not affect its local structure and, in particular, do not change the angle of twist of the cross–section about the centerline. Strong thermal fluctuations lead to melting of the helix, accompanied by complete loss of local helical structure. Depending on the parameters of the helix, the persistence length is a non–monotonic function of the strength of thermal fluctuations, and may first increase and then decrease as the amplitude of fluctuations is increased. Resonant peaks and dips in plots of the persistence length versus the spontaneous rate of twist are observed both for small twist rates and for rates equal to half the total curvature of the helix, phenomena which bear some formal similarity to stochastic resonances.
There are several possible directions in which the present work can be extended. We did not consider here the effects of excluded volume and other non–elastic interactions, on the statistical properties of fluctuating filaments. Such an analysis requires the introduction of a field theoretical description of the filaments. While this approach is interesting in its own right, we expect that the excluded volume exponent for the scaling of the end–to–end distance of a single filament will be identical to that of a Gaussian polymer chain (self–avoiding random walk). However, new effects related to liquid crystalline ordering are expected in dense phases of such filaments. Another possible extension of the model relates to the elasticity of random heteropolymers, with quenched distribution of elastic constants and/or spontaneous torsions.
A natural application of our theory involves the modeling of mechanical properties and conformational statistics of chiral biomolecules such as DNA and RNA. The advantage of our theory is that it allows us to take into account, in an exact manner, the effects of thermal fluctuations on the persistence length and other elastic parameters of the filament. Thus, the generalization of the theory to include the effect of tensile forces and torques applied to the ends of the filament, is expected to lead to new predictions for mechanical stretching experiments in the intermediate deformation regime, for tensile forces that affect the global but not the local (on length scales $``$ $`l_p`$) conformation of the filament. Measurements of the effect of elongation on thermal fluctuations of the molecule, can give information about its elastic constants, and help resolve long–standing questions regarding the natural curvature of DNA. It is interesting to compare our expression for the persistence length to that introduced by Trifonov et al. who proposed that the apparent persistence length $`l_a`$ of DNA depends not only on the rigidity (dynamic persistence length $`l_d`$), but also on the intrinsic curvature of the molecule (static persistence length $`l_s`$). The apparent persistence length is given in terms of the two others as
$$\frac{1}{l_a}=\frac{1}{l_d}+\frac{1}{l_s}$$
(56)
Note that the philosophy of the above approach is very similar to ours – we begin with filaments which have some given intrinsic length (spontaneous radius of curvature/torsion), and find that the interplay between this length and thermal fluctuations gives rise to a persistence length $`l_p`$. In fact, taking for simplicity the case of a circular cross–section, $`a_1=a_2`$, our expression Eq. (35), can be recast into the form of Eq. (56), with
$$l_a=l_p,\text{ }l_d=2a_1,l_s=\kappa _0^2\left(\gamma _1+\tau _0^2/\gamma _1\right)$$
(57)
Indeed, in our model, $`a_1`$ is the bare persistence length that determines the length scale on which the filament is deformed by thermal bending and torsion fluctuations. Our analog of the static persistence length $`l_s`$ depends on the spontaneous bending rate $`\kappa _0`$ and diverges in the case of a straight filament ($`\kappa _00`$), in which case $`l_al_d`$. If we make the further assumption that twist rigidity is much smaller than the bending rigidity, $`a_3a_1`$, the static persistence length becomes independent of the bending rigidity and depends on both the spontaneous curvature and the twist rigidity. Note, however, that the resulting $`\kappa _0^2`$ dependence of $`l_s`$ differs from the originally proposed one ($`\kappa _0^1`$).
Another possible application of our theory involves a new way of looking into the protein folding problem. Usually, one assumes that the folded conformation of proteins is determined by the interactions between the constituent amino–acids. A different approach, more closely related to the present work, would be to reverse the common logic: instead of trying to understand what kind of spatial structure will result for a given primary sequence of amino–acids, one can begin with a known equilibrium shape (native state) and attempt to identify the parameters of an effective filament (distributions of spontaneous torsions $`\left\{\omega _{0i}(s)\right\}`$ ) which will give rise to this three–dimensional structure. Knowledge about the fluctuations and the melting of proteins can then be used to determine the distribution of the bare persistence lengths $`\left\{a_i(s)\right\}`$. While the question of whether such an approach can be successfully implemented in order to determine the relation between primary sequence and ternary structure remains open, our insights about the statistical properties of fluctuating filaments are clearly applicable to modeling of $`\alpha `$helices and other elements (e.g., $`\beta `$sheets) of secondary structure of proteins.
> Appendix A: Calculation of Correlation Functions
We begin with the construction of the eigenvectors of the matrix $`𝚲`$, defined by Eq. (21), in the case $`\mathrm{\Delta }>0`$ (see Eq. (28)), when there is one real eigenvalue $`\lambda _1`$ and two complex ones, $`\lambda _R\pm i\omega `$. Expanding this matrix over its eigenvectors, we get
$$\mathrm{\Lambda }_{ij}=\lambda _1\overline{u}_iu_j+\left(\lambda _R+i\omega \right)\overline{v}_iv_j^{}+\left(\lambda _Ri\omega \right)\overline{v}_i^{}v_j$$
(58)
where the eigenvectors $`𝐮`$, $`\overline{𝐮}`$, $`𝐯`$, $`\overline{𝐯}`$ (and the complex conjugates of the latter two, $`𝐯^{}`$ and $`\overline{𝐯}^{}`$) obey the orthonormality conditions
$$\underset{i=1}{\overset{3}{}}\overline{u}_iu_i=\underset{i=1}{\overset{3}{}}\overline{v}_iv_i^{}=1,\underset{i=1}{\overset{3}{}}\overline{u}_iv_i=\underset{i=1}{\overset{3}{}}\overline{v}_iu_i=\underset{i=1}{\overset{3}{}}\overline{v}_iv_i=0$$
(59)
Using these conditions we can exponentiate the matrix $`𝚲`$
$$\left[e^{𝚲s}\right]_{ij}=\overline{u}_iu_je^{\lambda _1s}+\overline{v}_iv_j^{}e^{\left(\lambda _R+i\omega \right)s}+\overline{v}_i^{}v_je^{\left(\lambda _Ri\omega \right)s}$$
(60)
Since we are interested only in the diagonal elements of this matrix, it is convenient to introduce the notations
$$c_i=\overline{v}_iv_i^{},\underset{i=1}{\overset{3}{}}c_i=1$$
(61)
In addition, substituting $`s=0`$ in Eq. (60) we get
$$\overline{u}_iu_i=1c_ic_i^{}$$
(62)
In order to find the complex coefficients $`c_i`$ we write down expressions for diagonal elements of the matrices $`𝚲`$ and $`𝚲^2`$
$$\begin{array}{c}\gamma _i=\left(1c_ic_i^{}\right)\lambda _1+c_i\left(\lambda _R+i\omega \right)+c_i^{}\left(\lambda _Ri\omega \right)\\ \left(\gamma _i\lambda _1\right)^2\omega _0^2+\omega _{0i}^2=\left(1c_ic_i^{}\right)\lambda _1^2+c_i\left(\lambda _R+i\omega \right)^2+c_i^{}\left(\lambda _Ri\omega \right)^2\end{array}$$
(63)
Looking for the solution of these equations in the form $`c_i=Rec_i+iImc_i`$ we get expressions for real and imaginary parts of complex parameters $`c_i`$
$$\begin{array}{c}2Rec_i=\frac{\gamma _i^2+2\epsilon _i\left(\lambda _1+\lambda _R\right)+2\lambda _R\lambda _1+\omega _0^2\omega _{0i}^2}{\omega ^2+\left(\lambda _1\lambda _R\right)^2},\\ 2\omega Imc_i=\lambda _1\gamma _i+2\left(\lambda _R\lambda _1\right)Rec_i\end{array}$$
(64)
Appendix B: Persistence Length of Twisted Helix
Since the persistence length is defined by the $`33`$ element of the averaged rotation matrix, we will consider the $`i3`$ component of equation Eq. (15) which, using Eq. (13), can be expressed as an equation for the corresponding correlator:
$$\frac{dg_i}{ds}=\underset{l}{}\mathrm{\Lambda }_{il}(s+s^{})g_l,g_i(s,s^{})𝐭_i(s+s^{})𝐭_3(s^{})$$
(65)
with initial conditions $`g_1(0,s^{})=g_2(0,s^{})=0`$ and $`g_3(0,s^{})=1`$. The matrix $`𝚲\left(s+s^{}\right)`$ was defined in Eq. (50). Note that since the only $`s`$–dependent parameter of the helix is the angle of twist, the correlators $`g_i(s,s^{})`$ depend on $`s^{}`$ only through the parameter $`\alpha _0(s^{})=\alpha _0`$ and, in order to simplify the notation, we will omit the second argument of these functions in the following.
It is convenient to introduce the complex function
$$f\left(s\right)=\left[g_1(s)+ig_2(s)\right]e^{i\left(\dot{\alpha }_0s+\alpha _0\right)}$$
(66)
such that $`f`$ and $`g_3`$ obey the coupled equations
$`{\displaystyle \frac{df}{ds}}+\gamma _+f+\gamma _{}f^{}e^{2i\left(\dot{\alpha }_0s+\alpha _0\right)}`$ $`=`$ $`i\kappa _0g_3+i\tau _0f,`$
$`{\displaystyle \frac{dg_3}{ds}}+\gamma _3g_3`$ $`=`$ $`i\kappa _0{\displaystyle \frac{1}{2}}\left(ff^{}\right)`$ (67)
Taking a Laplace transform of these equations,
$$\stackrel{~}{f}\left(p\right)_0^{\mathrm{}}f\left(s\right)e^{ps}𝑑s,\stackrel{~}{g}_3\left(p\right)_0^{\mathrm{}}g_3\left(s\right)e^{ps}𝑑s$$
(68)
where $`p`$ is, in general, a complex parameter, we get
$`\left(p+\gamma _+i\tau _0\right)\stackrel{~}{f}\left(p\right)+i\kappa _0\stackrel{~}{g}_3\left(p\right)`$ $`=`$ $`\gamma _{}e^{2i\alpha _0}\stackrel{~}{f}^{}\left(p+2i\dot{\alpha }_0\right),`$ (69)
$`\left(p+\gamma _3\right)\stackrel{~}{g}_3\left(p\right)+i\kappa _0{\displaystyle \frac{1}{2}}\left[\stackrel{~}{f}\left(p\right)\stackrel{~}{f}^{}\left(p\right)\right]`$ $`=`$ $`1`$ (70)
In deriving these equations, we used the initial conditions, $`f\left(0\right)=0`$ and $`g_3\left(0\right)=1.`$ Substituting $`\stackrel{~}{g}_3`$ from Eq. (70) into (69), we get a closed equation for the complex function $`\stackrel{~}{f}:`$
$$\left[\left(p+\gamma _+i\tau _0\right)\left(p+\gamma _3\right)+\frac{\kappa _0^2}{2}\right]\stackrel{~}{f}\left(p\right)+i\kappa _0\frac{\kappa _0^2}{2}\stackrel{~}{f}^{}\left(p\right)+\gamma _{}\left(p+\gamma _3\right)e^{2i\alpha _0}\stackrel{~}{f}^{}\left(p+2i\dot{\alpha }_0\right)=0$$
(71)
Note that the persistence length is determined by $`\stackrel{~}{g}_3\left(0\right)`$ which can be expressed through $`\stackrel{~}{f}\left(0\right)\stackrel{~}{f}^{}\left(0\right)`$, Eq. (70). The latter functions can be calculated from Eq. (71), which upon substituting $`p=2in\dot{\alpha }_0`$ ($`n`$ integer), is recast in the standard form of difference equations,
$$a_n\kappa _0\stackrel{~}{f}\left(2in\dot{\alpha }_0\right)+2i\kappa _0\stackrel{~}{f}^{}\left(2in\dot{\alpha }_0\right)+2\gamma _{}b_ne^{2i\alpha _0}\stackrel{~}{f}^{}\left[2i\left(n1\right)\dot{\alpha }_0\right]=0$$
(72)
where we defined
$$a_n=1+2\left[\gamma _+i\left(\tau _0+2n\dot{\alpha }_0\right)\right]\left(\gamma _32in\dot{\alpha }_0\right)/\kappa _0^2,b_n=\left(\gamma _32in\dot{\alpha }_0\right)/\kappa _0$$
(73)
Since the persistence length is defined as the average of $`\stackrel{~}{g}_3\left(0\right)`$ with respect to $`\alpha _0`$, it is convenient to introduce dimensionless functions $`h_n`$ as:
$$h_n=\kappa _0_0^{2\pi }\frac{d\alpha _0}{2\pi }e^{2in\alpha _0}\stackrel{~}{f}\left(2in\dot{\alpha }_0\right)$$
(74)
We multiply Eq. (72) by $`\mathrm{exp}\left(2in\alpha _0\right)`$ and average it with respect to $`\alpha _0`$. Defining the parameter $`\epsilon =2\gamma _{}/\kappa _0`$ we rewrite Eq. (72) in the form
$$a_nh_n+2i\delta _{n0}h_n^{}+\epsilon b_nh_{1n}^{}=0$$
(75)
in which both $`h_n`$ and $`h_m^{}`$ enter. In order to derive closed equations for the set of $`\left\{h_n\right\}`$ only, we apply complex conjugation to the above equation and change $`nn`$. This yields
$$a_n^{}h_n^{}2i\delta _{n0}h_n+\epsilon b_nh_{n+1}=0$$
(76)
Substituting the equations for $`h_n^{}`$ and $`h_{1n}^{}`$ into (75) we find
$$\begin{array}{c}\left(a_n1/a_n^{}\epsilon ^2b_nb_{n1}/a_{1n}^{}\right)h_n+2i\left(11/a_n^{}\right)\delta _{n0}+\\ 2i\epsilon \delta _{n1}b_n/a_{1n}^{}+\epsilon h_{n+1}b_n/a_n^{}+\epsilon h_{n1}b_n/a_{1n}^{}=0\end{array}$$
(77)
Let us first consider the case $`n0,1`$. Introducing new variables $`y_n`$ by the equality $`h_{n+1}=\epsilon y_nh_n`$ we find
$$a_n1/a_n^{}\epsilon ^2b_nb_{n1}/a_{1n}^{}+\epsilon ^2y_nb_n/a_n^{}+y_{n1}^1b_n/a_{1n}^{}=0$$
(78)
We now define
$$A_n=\left(a_n1/a_n^{}\right)a_{1n}^{}/b_n\epsilon ^2b_{n1},B_n=a_{1n}^{}/a_n^{}$$
(79)
and get the following recurrence relation, valid for $`n=2,3,`$
$$A_n+1/y_{n1}+\epsilon ^2B_ny_n=0$$
(80)
We now take $`n=2`$ in the above equation, and solve for $`y_1`$in terms of $`y_2.`$ Repeating this procedure (expressing $`y_2`$ in terms of $`y_3`$, etc.) we can write the solution as a continued fraction
$$y_1=1/\left(A_2\epsilon ^2B_2/\left(A_3\epsilon ^2B_3/\left(A_4\mathrm{}\right)\right)\right)$$
(81)
Now consider the case $`n=1`$ in Eq. (77). Using the definitions of $`A_1`$ and $`B_1`$, Eq. (79), it can be recast into the form:
$$\left(A_1+\epsilon ^2B_1y_1\right)h_1+2i\epsilon +\epsilon h_0=0$$
(82)
In order to obtain a closed equation for $`h_0`$, we return to Eq. (76) with $`n=0,`$
$$a_0^{}h_0^{}2ih_0+\epsilon b_0h_1=0$$
(83)
Eliminating $`h_1`$ from the above two equations we find
$$h_0=2i+\mathrm{\Xi }h_0^{}$$
(84)
where, using Eq. (81), $`\mathrm{\Xi }`$ can be represented as a continued fraction:
$`\mathrm{\Xi }`$ $`=`$ $`a_0^{}/\left(1+\epsilon ^2b_0/\left(A_1+\epsilon ^2B_1y_1\right)\right)`$ (85)
$`=`$ $`a_0^{}/\left(1+\epsilon ^2b_0/\left(A_1\epsilon ^2B_1/\left(A_2\epsilon ^2B_2/\left(A_3\mathrm{}\right)\right)\right)\right)`$ (86)
The solution of Eq. (84) is:
$$h_0=2i\frac{1\mathrm{\Xi }}{1\left|\mathrm{\Xi }\right|^2}$$
(87)
Recall that $`h_0`$ was defined as the integral over $`\alpha _0`$ of the function $`\stackrel{~}{f}\left(0\right)`$ (Eq. (74)) which, in turn, determines the Laplace transform at $`p=0`$ of the correlator $`\stackrel{~}{g}_3`$ that appears in the definition of the persistence length, Eq. (51). Collecting the above expressions we find:
$$l_p=_0^{2\pi }\frac{d\alpha _0}{\pi }\stackrel{~}{g}_3\left(0\right)=\frac{2}{\gamma _3}\left[1i\frac{1}{2}\left(h_0h_0^{}\right)\right]=\frac{2\gamma _3^1}{1+\left(\mathrm{\Xi }1\right)^1+\left(\mathrm{\Xi }^{}1\right)^1}$$
(88)
Acknowledgment
We would like to thank A. Drozdov for illuminating discussions and D. Kessler for helpful comments on the manuscript. YR acknowledges support by a grant from the Israel Science Foundation. SP thanks M. Elbaum for hospitality during his stay at the Weizmann Institute.
Figure captions
Figure 1: Schematic drawing of a twisted ribbon–like filament. The vectors of the physical ($`𝐭_1,𝐭_2`$) and the Frenet ($`𝐛,𝐧`$) triad can be brought into coincidence through rotation by angle $`\alpha `$, about the common tangent ($`𝐭_3`$).
Figure 2: Schematic plot of section of a ribbon–like helix. The helix–fixed coordinate system $`𝐭`$ at contour point $`s^{}`$ is shown. The solid line describes the associated “rod–like chain” to which the coordinate system $`𝐞`$ is attached at point $`\sigma `$ on its contour. The points $`\sigma `$ and $`\sigma ^{}`$ on the rod–like chain are the projections of the points $`s`$ and $`s^{}`$ respectively.
Figure 3: Three–dimensional plot of the persistence length $`l^{}`$ as a function of the dimensionless rate of twist $`w`$ and of the bare persistence length $`a_1`$ (logarithmic scale), for a helical filament with spontaneous curvature $`\kappa _0=1`$, and torsion $`\tau _0=0.01`$ (in arbitrary units). The bare persistence lengths are $`a_2=1000`$, and $`a_3=5000`$.
Figure 4: Plot of the persistence length $`l^{}`$ as a function of the dimensionless rate of twist $`w`$ for a helical filament with spontaneous curvature $`\kappa _0=1`$, and torsion $`\tau _0=0.01`$ (in arbitrary units). The different curves correspond to different bare persistence lengths: (1) $`a_1=100`$, $`a_2=a_3=5000`$, (2) $`a_1=1`$, $`a_2=a_3=100`$, (3) $`a_1=0.1`$, $`a_2=a_3=10`$, (4) $`a_1=0.01`$, $`a_2=a_3=10`$. A magnified view of the region of small twist rates is shown in the insert.
Figure 5: Three–dimensional plot of the persistence length $`l^{}`$ as a function of the dimensionless rate of twist $`w`$ and of the spontaneous curvature $`\kappa _0`$, for a helical filament with spontaneous torsion $`\tau _0=1`$ (in arbitrary units). The bare persistence lengths are $`a_1=500`$, $`a_2=1`$ and $`a_3=500`$.
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# Untitled Document
CALT-68-2258 CITUSC/00-005 PUPT-1904 UMDEPP 00-046
Two Two-Dimensional Supergravity Theories from Calabi-Yau Four-Folds
S. James Gates, Jr., <sup>1</sup> E-Mail: gatess@wam.umd.edu Sergei Gukov,<sup>♣♢</sup><sup>2</sup> E-Mail: gukov@feynman.princeton.edu and Edward Witten<sup>♣♡</sup>
Department of Physics, University of Maryland at College Park,
College Park, MD 20742-4111, USA
California Institute of Technology, Pasadena, CA 91125, USA,
CIT-USC Center For Theoretical Physics
Joseph Henry Laboratories, Princeton University,
Princeton, NJ 08544, USA
School of Natural Sciences, Institute for Advanced Study,
Olden Lane, Princeton, NJ 08540, USA
We consider two-dimensional supergravity theories with four supercharges constructed from compactification of Type II string theory on a generic Calabi-Yau four-fold. In Type IIA and Type IIB cases, respectively, new superspace formulations of $`𝒩=(2,2)`$ and $`𝒩=(0,4)`$ dilaton supergravities are found and their coupling to matter multiplets is discussed.
May 2000
1. Introduction
For a long time, compactification of heterotic string theory on Calabi-Yau manifolds was the primary candidate for constructing realistic models in four dimensions with $`𝒩=1`$ supersymmetry. This was also a strong motivation to study Type II superstrings on Calabi-Yau three-folds which share many common properties with the corresponding heterotic compactifications. At the same time substantial progress has been made in understanding the mathematical aspects of Calabi-Yau three-folds, such as quantum cohomology and mirror symmetry . It was not until the discovery of F-theory that it was realized that $`𝒩=1`$ four-dimensional heterotic string vacua can be equivalently described as F-theory compactifications on elliptically fibered Calabi-Yau four-folds. Since then, the study of Calabi-Yau four-fold compactifications has become of particular importance for physical applications. Compactification of F-theory on an elliptically fibered Calabi-Yau four-fold is closely related to the corresponding compactifications of Type IIA and M-theory. Namely, when the area of the elliptic fiber shrinks to zero, M-theory compactification on a Calabi-Yau four-fold is well described by F-theory compactification on the same Calabi-Yau manifold. On the other hand, Type IIA string theory is related to M-theory via compactification on an extra circle (of small radius).
One of the most striking outcomes in the study of compactifications on Calabi-Yau three-folds is the great success in understanding non-perturbative phenomena in $`𝒩=2`$ field theories in four dimensions, see e.g. for introduction and references. On the other hand, understanding of Calabi-Yau four-fold compactifications still is quite far from that stage, so we shall not discuss non-perturbative phenomena in this paper. Instead, we consider classical supergravity theories interacting with two-dimensional non-linear sigma-models that can be constructed from Calabi-Yau four-folds. Surprisingly, it turns out that manifestly supersymmetric formulations of such theories has not been given previously.
On general grounds, a compactification of Type IIA (IIB) string theory on a Calabi-Yau four-fold leads to a $`𝒩=(2,2)`$ (resp. $`𝒩=(0,4)`$) effective field theory in two dimensions. In the low-energy limit the theory is described by supergravity coupled to matter. For example, from the Kaluza-Klein reduction of Type IIA string theory on a Calabi-Yau four-fold in section 3 we find that a suitable low-energy theory is $`𝒩=(2,2)`$ dilaton supergravity interacting with some number of chiral and twisted chiral multiplets. It is invariant under the “mirror transformation” which, acting on the matter fields, exchanges chiral multiplets and twisted chiral multiplets. We thus generalize the proposal of where it was suggested that the “kinematic structure” of the mirror transformation has its origin in a mapping between chiral and twisted chiral multiplets when these superfields are regarded as the fundamental degrees of $`𝒩=(2,2)`$ superstring theories. The generalization posits that this mapping also applies to the effective action. Another characteristic feature of this supergravity theory is that the supergravity multiplet contains a real dilaton field. In a special case, when all matter multiplets are chiral and massless, a component action of this $`𝒩=2`$ dilaton supergravity was constructed in . However, to describe Type IIA compactifications on Calabi-Yau four-folds we need a generalization of this theory that includes interaction with twisted chiral multiplets and the possibility to turn on the superpotential and as well the twisted superpotential. Thus, in sections 4 and 5 we present superspace construction of general $`𝒩=(2,2)`$ dilaton supergravity coupled to matter. The construction in section 4 is based on the Goldstone mechanism in the superspace formulation of non-minimal gauged $`𝒩=(2,2)`$ supergravity. Coupling of the new $`𝒩=(2,2)`$ dilaton supergravity to matter multiplets is the subject of section 5, where we discuss local integration in superspace. In section 6 we perform the Kaluza-Klein reduction of Type IIB string theory on a Calabi-Yau four-fold and describe the component action of the resulting $`𝒩=(0,4)`$ dilaton supergravity. Most of this section, as well as section 3, is not new and presented for the sake of completeness. The superspace formulation of the new $`𝒩=(0,4)`$ dilaton supergravity is presented in section 7. In the appendix A we present a straightforward but technical world-sheet calculation of string amplitudes corresponding to the target space metric of the effective two-dimensional theory, and in the appendix B we list extra derivative constraints arising from de-gauging $`𝒩=(2,2)`$ non-minimal supergravity. Appendix C contains components of the covariant derivative in $`𝒩=(2,2)`$ dilaton supergravity needed in section 5. Finally, in appendix D we repeat the derivation of the chiral density projection formula in $`𝒩=(2,2)`$ dilaton supergravity. We begin in the next section with a summary of notations and definitions used throughout the paper.
2. Calabi-Yau Four-folds: Some Conventions and Definitions
We study compactification of Type II string theory on $`M_{(2)}\times X`$ where $`M_{(2)}`$ is a maximally symmetric homogeneous two-dimensional space-time and $`X`$ is a Calabi-Yau four-fold. We use the following notations for the space-time indices. Capital letters $`M`$, $`N`$, $`\mathrm{}`$ run from 0 to 9 and denote ten-dimensional Lorentz indices. Latin letters $`m`$, $`n`$, $`\mathrm{}`$ and $`a`$, $`b`$, $`\mathrm{}`$ represent, respectively, real and holomorphic indices tangent to $`X`$. Greek letters $`\alpha `$, $`\beta `$, $`\mathrm{}`$ and $`\mu `$, $`\nu `$, $`\mathrm{}`$ are used for the two-dimensional spinor indices “$`+`$” and “$``$”, and light-cone indices “ ” and “ ”, correspondingly. Sometimes we also use capital latin letters $`A`$, $`B`$, $`\mathrm{}`$ to denote both spinor and vector indices.
A Calabi-Yau space $`X`$ is a compact Kähler manifold with complex dimension four and $`SU(4)`$ holonomy group. It follows that $`X`$ is a Ricci-flat manifold and, therefore, it can be used as a background for Type II string compactification. As a topological space, $`X`$ is classified by the Hodge numbers $`h^{p,q}`$ which count the number of harmonic $`(p,q)`$-forms $`\omega _i^{(p,q)}H^{p,q}(X)`$, $`i=1,\mathrm{}h^{p,q}`$. The non-vanishing cohomology groups have the following dimensions :
$$h^{1,1}=h^{3,3},h^{3,1}=h^{1,3},$$
$$h^{2,1}=h^{1,2}=h^{3,2}=h^{2,3},$$
$$h^{0,0}=h^{4,4}=h^{4,0}=h^{0,4}=1,$$
$$h^{2,2}=2(22+2h^{1,1}+2h^{3,1}h^{2,1}).$$
For the Euler number of $`X`$ we have:
$$\frac{\chi }{6}=8+h^{1,1}+h^{3,1}h^{2,1}$$
We denote by $`\mathrm{\Omega }`$ a covariantly constant $`(4,0)`$-form. The $`(1,1)`$\- and $`(3,1)`$-forms are related to the deformation parameters of the Kähler form and the complex structure of $`X`$, respectively. Namely, an arbitrary variation of the metric of the Calabi-Yau four-fold $`X`$ that respects $`SU(4)`$ holonomy looks like:
$$\delta g_{a\overline{b}}dz^ad\overline{z}^{\overline{b}}+\delta g_{ab}dz^adz^b+\mathrm{c}.\mathrm{c}.$$
where
$$\delta g_{ab}=\underset{j=1}{\overset{h^{3,1}}{}}\varphi ^jw_{ab}^j,i\delta g_{a\overline{b}}=\underset{i=1}{\overset{h^{1,1}}{}}s^i\omega _{a\overline{b}}^i.$$
By the appropriate contraction with $`\overline{\mathrm{\Omega }}`$, from the forms $`w_{ab}^j`$ we can construct elements in $`H^{1,3}(X)`$:
$$\omega _j^{(1,3)}=\overline{\mathrm{\Omega }}_{\overline{a}\overline{b}\overline{c}\overline{d}}g^{\overline{d}d}w_{df}^jd\overline{z}^{\overline{a}}d\overline{z}^{\overline{b}}d\overline{z}^{\overline{c}}dz^f$$
In what follows we will use some integrals over the Calabi-Yau space $`X`$ \[8,,9\]:
$$𝒱=_Xd^8z\sqrt{g}=\frac{1}{4!}_X𝒦𝒦𝒦𝒦$$
$$𝒢_{\varphi _i\overline{\varphi }_{\overline{j}}}=\frac{1}{4𝒱}_Xd^8z\sqrt{g}w_{iab}\overline{w}_{\overline{j}}^{ab}$$
$$𝒢_{\sigma _k\overline{\sigma }_l}=\frac{1}{2𝒱}_Xd^8z\sqrt{g}\omega _{a\overline{b}}^k\omega ^{la\overline{b}}$$
$$Y_{im\overline{n}}=_X\omega _i^{(1,1)}\omega _m^{(2,1)}\omega _{\overline{n}}^{(1,2)}$$
$$d_{ijkl}=_X\omega _i^{(1,1)}\omega _j^{(1,1)}\omega _k^{(1,1)}\omega _l^{(1,1)}$$
The notations (2.1) and (2.1) will become clear in the next section where we identify expectation values of the fields $`s_i`$ and $`\varphi _j`$ with the Kähler and complex structure moduli, respectively. In particular, we write:
$$𝒦=ig_{a\overline{b}}dz^ad\overline{z}^{\overline{b}}=\underset{i=1}{\overset{h^{1,1}}{}}s_i\omega _i^{(1,1)}$$
for the Kähler form on $`X`$. The moduli space of a Calabi-Yau space is locally a product of the moduli space of complex deformations, $`_c(X)`$, and the (complexified) moduli space of Kähler structure, $`_𝒦(X)`$. Notice, the metric $`𝒢_{\varphi _i\overline{\varphi }_{\overline{j}}}`$ defined above is the Weil-Petersson metric on the moduli space of complex structure of $`X`$, with the Kähler potential, cf. :
$$K(\varphi _i,\overline{\varphi }_{\overline{i}})=\mathrm{ln}\left(_X\mathrm{\Omega }\overline{\mathrm{\Omega }}\right).$$
3. Compactification of Type IIA String Theory on Calabi-Yau Four-folds
In this section we describe the effective two-dimensional theory constructed from compactification of Type IIA string theory on a Calabi-Yau four-fold $`X`$. When the volume of $`X`$ is large compared to the string scale, Type IIA supergravity is a good low-energy approximation to Type IIA string theory. Therefore, in the ‘large volume limit’ we may describe the low-energy effective theory studying compactification of Type IIA supergravity on the Calabi-Yau space $`X`$. With this motivation, let us start this section recalling some facts about Type IIA supergravity itself.
The bosonic field content of Type IIA supergravity contains the metric $`g_{MN}`$, the dilaton $`\phi `$, a vector field $`A_M`$, and tensor fields $`B_{MN}`$ and $`C_{MNP}`$. The bosonic part of the Lagrangian (in string frame) looks like:
$$L_{(10)}=\sqrt{g}\left[\frac{1}{2}e^{2\phi }(R^{(10)}+4(\phi )^2\frac{1}{12}H^2)\frac{1}{4}F^2\frac{1}{48}G^2\right]+\mathrm{}$$
where we introduced the gauge-invariant field strengths:
$$F=dA,H=dB,$$
$$G=dC,G^{}=G+AH.$$
With this choice of normalization the fields $`A`$, $`B`$, and $`C`$ transform in a natural, dilaton-independent way under gauge transformations.
The dots in the Lagrangian (3.1) stand for higher order terms among which we find the Chern-Simons term $`BGG`$ and the anomaly term $`BI_8`$, where the eight-form $`I_8`$ is proportional to the Euler density of $`X`$. After integration over a compact eight-manifold $`X`$ these topological terms produce a global anomaly \[10,,11\]:
$$N=\frac{\chi }{24}\frac{1}{2(2\pi )^2}_XGG$$
To cancel the tadpole for the $`B`$-field one has to introduce $`N`$ fundamental strings filling two-dimensional non-compact space.
The action of Type IIA supergravity is invariant under 16 left and 16 right supersymmetry transformations, such that the left supersymmetries are chiral while the right supersymmetries are anti-chiral with respect to the ten-dimensional chirality operator $`\mathrm{\Gamma }_{11}`$. Since $`X`$ admits a nowhere vanishing complex spinor of definite chirality, compactification of Type IIA string theory on $`X`$ is described by $`𝒩=(2,2)`$ supergravity theory coupled to matter. With the appropriate choice of orientation, the fundamental strings filling two-dimensional space-time do not break supersymmetry further.
To find the spectrum of the effective low-energy theory we perform Kaluza-Klein reduction of Type IIA supergravity to two dimensions. Below we describe the decomposition of Type IIA bosonic fields in harmonics of $`X`$. By supersymmetry, incorporation of fermionic zero-modes completes the resulting spectrum into appropriate $`𝒩=(2,2)`$ supermultiplets. Via dimensional reduction Type IIA dilaton $`\phi `$ becomes a real scalar field in the two-dimensional theory. The ten-dimensional metric $`g_{MN}`$ decomposes into the two-dimensional metric $`g_{\mu \nu }`$, $`h^{3,1}`$ complex scalars $`\varphi _i`$ and $`h^{1,1}`$ real scalars $`s_j`$ defined in (2.1). The antisymmetric tensor fields $`B_{MN}`$ and $`C_{MNP}`$ can be expanded into harmonic modes as follows:
$$B=\underset{i=1}{\overset{h^{1,1}}{}}r^i\omega _i^{(1,1)},$$
$$C=\underset{j=1}{\overset{h^{1,1}}{}}A_\mu ^j\omega _j^{(1,1)}+\underset{k=1}{\overset{h^{2,1}}{}}z^k\omega _k^{(2,1)}+\mathrm{c}.\mathrm{c}.$$
It is convenient to combine real fields $`s^i`$ and $`r^i`$ into complex scalars $`\sigma ^i`$. Taking into account the vector field $`A_\mu `$ from the Ramond-Ramond sector of Type IIA theory, we end up with the following list of $`𝒩=(2,2)`$ supermultiplets:
$$\mathrm{a}\mathrm{gravitational}\mathrm{multiplet}:g_{\mu \nu },A_\mu ,\phi $$
$$h^{3,1}\mathrm{chiral}\mathrm{multiplets}:\varphi _i,\overline{\varphi }_{\overline{i}}$$
$$h^{1,1}\mathrm{twisted}\mathrm{chiral}\mathrm{multiplets}:\sigma ^j,\overline{\sigma }^j,A_\mu ^j$$
$$h^{2,1}(\mathrm{twisted})\mathrm{chiral}\mathrm{multiplets}:z^k,\overline{z}^{\overline{k}}$$
Vector fields $`A_\mu `$ and $`A_\mu ^i`$ do not have propagating degrees of freedom in two dimensions and play the role of auxiliary fields in the supergravity multiplet and twisted chiral multiplets, respectively.
The complex scalar fields $`z^k`$ which come from $`(2,1)`$-modes take value in a torus. When background fluxes satisfy $`G\omega _k^{(2,1)}=0`$ and $`H^{}\omega _k^{(2,1)}=0`$ there is no superpotential for the corresponding harmonics $`z_k`$, so these fields are massless. A $`T`$-duality transformation on the torus then converts them from ordinary chiral superfields to twisted chiral superfields. It is natural to choose $`z_k`$ to be scalar components of chiral superfields. Indeed, if we start in eleven dimensions, the reduction of the $`C`$-field (3.1) yields $`h^{2,1}`$ complex scalar modes in three dimensions. These modes are scalar components of chiral superfields since there is no notion of “twisted chiral superfields” in three dimensions. After a further compactification on a circle they naturally remain as chiral superfields in two dimensions, but now a $`T`$-duality becomes possible and $`z_k`$ can be alternatively described if one prefers as twisted chiral superfields.
To find the effective action for the light fields we have to substitute (2.1), (3.1) and (3.1) in the Lagrangian (3.1) and integrate over the internal space $`X`$. Using the formulas (2.1), (2.1) and (2.1) we obtain the following effective action for the bosonic modes, cf. \[9,,5\]:
$$L_{(2)}=e^{2\phi }𝒱[R^{(2)}+4(\phi )^2𝒢_{\varphi _i\overline{\varphi }_{\overline{j}}}(_\mu \varphi ^i)(^\mu \overline{\varphi }^{\overline{j}})$$
$$\frac{1}{2}𝒢_{\sigma _i\overline{\sigma }_j}(_\mu \sigma ^i)(_\mu \overline{\sigma }^j)]\frac{1}{4}Y_{im\overline{n}}\sigma ^i(D_\mu z^m)(D^\mu \overline{z}^{\overline{n}})+\mathrm{}$$
where the covariant derivative $`D_\mu `$ acting on $`z^m`$ contains a connection corresponding to the holomorphic dependence of the basis of $`(2,1)`$-forms on the complex structure .
To summarize, we find that in the large volume limit compactification of Type IIA string theory on a Calabi-Yau four-fold $`X`$ leads to $`𝒩=(2,2)`$ dilaton supergravity coupled to a non-linear sigma-model. The target space of this sigma-model is parametrized by some number of chiral and twisted chiral multiplets. This agrees with the result of , where it was found that the most general $`𝒩=(2,2)`$ non-linear sigma-model is based on a target space with two non-commuting complex structures $`J_\pm `$, so that the space $`\mathrm{ker}(J_+J_{})`$ is parametrized by chiral superfields, while $`\mathrm{ker}(J_++J_{})`$ is parametrized by twisted chiral superfields.
This kind of dilaton supergravity coupled to $`𝒩=(2,2)`$ chiral matter was studied some time ago . However, for our purposes we need to generalize the component construction of to include twisted chiral multiplets. Furthermore, background fluxes of Ramond-Ramond field strengths induce effective superpotential \[13,,14\] and/or twisted chiral superpotential \[13,,14,,15\] in the two-dimensional theory. Hence, we have to incorporate these terms in the construction as well. The most elegant and convenient way to do this is in $`𝒩=(2,2)`$ superspace where the supersymmetry becomes manifest \[16,,17\]. In addition to the usual space-time coordinates $`x^\mu `$, $`𝒩=(2,2)`$ superspace is parametrized by anti-commuting coordinates $`\theta ^\alpha =(\theta ^+,\theta ^{})`$ and their complex conjugates $`\overline{\theta }^{\dot{\alpha }}=(\overline{\theta }^{\dot{+}},\overline{\theta }^\dot{})`$. Then, we expect that the action of the matter fields (3.1) can be written in a compact form, similar to the action of matter coupled $`𝒩=1`$ supergravity in four dimensions:
$$S=d^2xd^2\theta d^2\overline{\theta }E^1\mathrm{exp}(K).$$
We postpone the discussion of the superspace measure $`E`$ till the next sections where superspace formulation will be discussed in detail. Now we simply assume that the suitable measure exists. The main advantage of the superspace formulation is that due to the extended supersymmetry, all the term in the action (3.1) with up to two derivatives or four fermions are determined by a single real function $`K(\varphi _i,\overline{\varphi }_{\overline{i}},\sigma _j,\overline{\sigma }_j,z_k,\overline{z}_{\overline{k}},\phi )`$, the Kähler potential . It is invariant under the generalized Kähler transformation:
$$KK+\mathrm{\Lambda }_1(\varphi _i,\sigma _j,z_k)+\overline{\mathrm{\Lambda }}_1(\overline{\varphi }_{\overline{i}},\overline{\sigma }_j,\overline{z}_{\overline{k}})+\mathrm{\Lambda }_2(\varphi _i,\overline{\sigma }_j,\overline{z}_{\overline{k}})+\overline{\mathrm{\Lambda }}_2(\overline{\varphi }_{\overline{i}},\sigma _j,z_k)$$
The target space metric is given by the second derivative of the Kähler potential.
For the sake of simplicity, let us assume for a moment that $`h^{2,1}=0`$. Then the metric is block diagonal:
$$𝒢_{\varphi _i\overline{\sigma }_j}=\frac{^2K}{\varphi _i\overline{\sigma }_j}=0$$
From the condition (3.1) it follows that locally we can write the Kähler potential that gives the effective action (3.1) as:
$$K=K_c(\varphi _i,\overline{\varphi }_{\overline{i}})+K_𝒦(\sigma _j,\overline{\sigma }_j)$$
where $`K_c`$ is the Kähler potential (2.1) on the moduli space of the complex structure:
$$K_c(\varphi _i,\overline{\varphi }_{\overline{i}})=\mathrm{ln}\left(_X\mathrm{\Omega }\overline{\mathrm{\Omega }}\right).$$
Similar to the case of Calabi-Yau three-folds , one can verify that the metric $`𝒢_{\sigma _i\overline{\sigma }_j}`$ can be obtained from the Kähler potential:
$$K_𝒦(\sigma _j,\overline{\sigma }_j)=\mathrm{ln}\left(_X𝒦𝒦𝒦𝒦\right).$$
Indeed, if $`\omega ^{(2)}`$ is a harmonic 2-form on a Calabi-Yau four-fold $`X`$, its Hodge dual is given by the following neat formula:
$${}_{}{}^{}\omega _{}^{(2)}=\frac{1}{2}\omega ^{(2)}𝒦𝒦+\frac{2}{3}\frac{\left(_X\omega ^{(2)}𝒦𝒦𝒦\right)}{\left(_X𝒦𝒦𝒦𝒦\right)}𝒦𝒦𝒦$$
Therefore, we can write (2.1) in the following form:
$$𝒢_{\sigma _k\overline{\sigma }_l}=\frac{1}{2𝒱}_X\omega _k^{(1,1)}^{}\omega _l^{(1,1)}=$$
$$=\frac{1}{4𝒱}_X\omega _k^{(1,1)}\omega _l^{(1,1)}𝒦𝒦\frac{1}{72𝒱^2}\left(_X\omega _k^{(1,1)}𝒦𝒦𝒦\right)\left(_X\omega _l^{(1,1)}𝒦𝒦𝒦\right)$$
Using the explicit expression (2.1) for the Kähler form $`𝒦`$, it is easy to see that the above metric indeed follows from the Kähler potential (3.1):
$$𝒢_{\sigma _k\overline{\sigma }_l}=\frac{1}{2}\frac{^2K_𝒦(\sigma _j,\overline{\sigma }_j)}{\sigma _k\overline{\sigma }_l}$$
Hence, to the leading order the metric on the target space is Kähler, torsionless, and equal to the metric on the moduli space of the Calabi-Yau space $`X`$, $`_c(X)\times _𝒦(X)`$.
The classical action (3.1) is invariant under two $`U(1)`$ $`R`$-symmetries. We will denote their linear combinations as $`U(1)_A`$ and $`U(1)_V`$. The action of these symmetries on the supercharges can be represented as:
$$Q_{}\overline{Q}_{\dot{+}}$$
$$Q_+\overline{Q}_\dot{}$$
where the upper (lower) row is assigned a $`U(1)_A`$ charge $`+1`$ ($`1`$) while the right (left) column is assigned a $`U(1)_V`$ charge $`+1`$ ($`1`$). These R-symmetries are not symmetries of the string theory – there is no way to assign $`R`$-transformations to massive string modes to preserve them. Even though we will not explicitly include massive string modes in the present paper, we will include a superpotential and twisted chiral superpotential that violate the $`R`$-symmetries. Even in the absence of the superpotentials, higher derivative interactions among the massless fields obtained by integrating out massive string states would be expected to violate the $`R`$-symmetry.
The explicit expression for the chiral superpotentials generated by the most general Ramond-Ramond flux $`=(RR\mathrm{fields})`$ in terms of the Calabi-Yau moduli was derived in :
$$W(\varphi _i)=\frac{1}{2\pi }_X\mathrm{\Omega }G$$
and for the twisted chiral superpotential:
$$\stackrel{~}{W}(\sigma _j)=\frac{1}{2\pi }_Xe^𝒦$$
The superpotential $`W(\varphi _i)`$ and the twisted superpotential $`\stackrel{~}{W}(\sigma _j)`$ are holomorphic functions of the fields $`\varphi _i`$ and $`\sigma _j`$, respectively. Taking into account the superpotential terms, the action of the matter fields reads as:
$$S_{(2)}=d^2xd^2\theta d^2\overline{\theta }E^1e^K+d^2xd^2\theta ^1W(\varphi _i)+$$
$$+d^2x𝑑\theta ^+𝑑\theta ^\dot{}\stackrel{~}{}^1\stackrel{~}{W}(\sigma _j)+\mathrm{c}.\mathrm{c}.$$
Generic values of Ramond-Ramond fluxes completely break the $`𝒩=(2,2)`$ supersymmetry<sup>3</sup> Investigating the supersymmetry conditions as in \[13,,14,,19\], one can also show that any $`H`$-field flux breaks all the supersymmetry. A simple way to see this is to assume, on the contrary, that there exists a supersymmetric vacuum corresponding to a non-zero $`H`$-flux and consider a BPS soliton connecting such a vacuum to the vacuum with zero $`H`$-flux. In Type IIA string theory this soliton would correspond to an $`NS5`$-brane wrapped over a Poincaré dual supersymmetric 5-cycle. However, there is a contradiction since Calabi-Yau 4-folds do not have supersymmetric 5-cycles, see e.g. . Therefore, a non-zero $`H`$-flux lifts all the supersymmetric vacua. It is natural to interpret this in terms of the effective superpotential $`WC^{}H`$ for the scalar fields $`z_k`$.. In the two-dimensional theory this effect corresponds to generation of a superpotential that lifts (part of) supersymmetric vacua. However, if the vacuum values of the fields $`\varphi _i`$ and $`\sigma _j`$ satisfy:
$$\frac{DW}{D\varphi _i}=0\mathrm{and}\frac{D\stackrel{~}{W}}{D\sigma _i}=0,$$
then Type IIA compactification on the corresponding Calabi-Yau manifold is supersymmetric \[13,,14\]. From the formulas (2.1), (3.1) and the quantization condition of the $`G`$-flux it follows that there is a finite number of choices for $`[G]H^4(X)`$ corresponding to supersymmetric vacua. In particular, if $`h^{2,1}>8+h^{1,1}+h^{3,1}`$, then there are no such vacua at all.
From the superspace construction in section 5 it follows that in the equations (3.1) we should use the appropriate covariant derivatives:
$$\frac{DW}{D\varphi _i}=\frac{W}{\varphi _i}+\frac{K_c(\varphi _i,\overline{\varphi }_{\overline{i}})}{\varphi _i}W,\frac{D\stackrel{~}{W}}{D\sigma _i}=\frac{\stackrel{~}{W}}{\sigma _i}+\frac{K_𝒦(\sigma _j,\overline{\sigma }_j)}{\sigma _i}\stackrel{~}{W}$$
where $`K_c`$ and $`K_𝒦`$ are given by the tree-level formulas (2.1) and (3.1), respectively. A simple way to see that one has to use the covariant derivatives instead of ordinary ones is to consider first compactification of F-theory on the same Calabi-Yau space<sup>4</sup> Of course, here we assume that $`X`$ is elliptically fibered. The result, however, is independent of this assumption. $`X`$. In the component action of the effective $`𝒩=1`$ four-dimensional theory there is a scalar potential:
$$e^K\left(𝒢^{\varphi _i\overline{\varphi }_{\overline{j}}}(D_{\varphi _i}W)(D_{\overline{\varphi }_{\overline{j}}}\overline{W})3|W|^2\right)$$
where the covariant derivative $`D_{\varphi _i}W=\frac{DW}{D\varphi _i}`$ is defined in (3.1). After a further compactification on a torus $`T^2`$, this theory is dual to compactification of Type IIA string theory on $`X`$. It is clear that after the dimensional reduction of the four-dimensional component action the covariant derivatives (3.1) also appear in the component action of the two-dimensional theory in question.
These models can have a variety of $`T`$-duality symmetries. Of particular interest are mirror symmetries \[20,,21\]. A mirror symmetry is, of course, a symmetry that maps Type IIA string theory on a four-fold $`X`$ to Type IIA string theory on the mirror variety $`\stackrel{~}{X}`$, such that:
$$h^{p,q}(X)=h^{4p,q}(\stackrel{~}{X})$$
and the conformal field theories associated with $`X`$ and $`\stackrel{~}{X}`$ are equivalent. This operation corresponds to a transformation which exchanges chiral multiplets and twisted chiral multiplets. It can be interpreted in terms of the supergeometrical coordinate transformation $`\theta ^{}\theta ^\dot{}`$ that also exchanges chiral multiplets and twisted chiral multiplets and changes the superspace measure in a way consistent with other definitions of mirror symmetry. In particular, the latter implies that under the mirror symmetry we have $`\sigma _i\varphi _j`$ which is consistent with our interpretation of (vevs of) these fields as the Kähler and the complex structure moduli of the Calabi-Yau space $`X`$. Therefore, the mirror symmetry relates different quantum $`𝒩=(2,2)`$ theories also interchanging:
$$U(1)_AU(1)_V$$
$$\varphi _i\sigma _j$$
$$W(\varphi _i)\stackrel{~}{W}(\sigma _j)$$
The mirror map has no effect on the $`𝒩=(2,2)`$ dilaton supergravity itself, so that in the absence of matter fields it must be mirror-symmetric. It may seem that twisted chiral fields $`z_k`$ violate the invariance under (3.1). Recall that via a spacetime T-duality transformation, those fields can be described by either chiral or twisted chiral superfields. Hence mirror symmetry just exchanges these two descriptions.
By definition, the low-energy effective action (3.1) describes dynamics of the light modes in Type IIA string theory on $`X`$ in the large volume limit. In other words, tree-level amplitudes in Type IIA string theory must agree with the corresponding amplitudes in the effective two-dimensional theory. In the appendix A we illustrate this by a world-sheet calculation which independently proves that the target space metric is block-diagonal, cf. (3.1).
A superspace formulation of $`𝒩=(2,2)`$ dilaton supergravity that includes chiral and twisted chiral multiplets on equal footing does not seem to exist in the literature. Although a superspace model of $`𝒩=(2,2)`$ supergravity where the dilaton is a complex field was constructed in , we are interested in a theory where the supergravity multiplet contains a real dilaton field. A superspace formulation of such a supergravity theory is presented in the next section.
4. Superspace Formulation of $`𝒩=(2,2)`$ Dilaton Supergravity
In this section we present a superspace construction of $`𝒩=(2,2)`$ dilaton supergravity without gauged symmetry. This last property is a distinguishing feature of the new formulation since all the known $`𝒩=(2,2)`$ gravity theories have at least one gauged $`U(1)`$ $`R`$-symmetry (see for a general presentation). Theories where the entire $`U(1)_AU(1)_V`$ symmetry group is gauged are called non-minimal (or reducible), as opposed to minimal theories where only $`U(1)_A`$ or $`U(1)_V`$ factor is gauged. It is very well known how to obtain one supergravity theory with a smaller holonomy group from a supergravity theory with a larger holonomy group. This process has been for a long time called “de-gauging” (see , section 5.3.b.7).
The basic idea of de-gauging is to break the gauge symmetry introducing extra matter field in a Goldstone-like mechanism. For example, consider an abelian vector multiplet in a four-dimensional $`𝒩=1`$ gauge theory. It contains a $`U(1)`$ gauge vector field, gaugino and an auxiliary field. All the other fields can be set to zero by a supersymmetric choice of gauge, the so-called Wess-Zumino gauge. On the other hand, a massive gauge multiplet contains some extra component fields which could be eliminated in the massless multiplet. The reason is that the massive system no longer possess the $`U(1)`$ gauge invariance. This toy model teaches us that when a symmetry is broken in superspace, extra component fields not present in the gauge symmetric phase begin to appear. In other words, Goldstone supermultiplets must appear. And their component fields come from that part of the vector multiplet that was ignored in the symmetric phase. Following these steps, we construct $`𝒩=(2,2)`$ dilaton supergravity via de-gauging $`U(1)_AU(1)_V`$ non-minimal gauged supergravity. An advantage of this approach is that both the original and the resulting theories are manifestly invariant under the mirror symmetry (3.1). We also find that the new $`𝒩=(2,2)`$ supergravity multiplet contains a real dilaton field $`\phi `$, in accordance with the results of section 3 where we studied compactification of Type IIA string theory on Calabi-Yau four-folds. We hope that apart from this obvious application there may also be many other aspects of the new supergravity to explore. For example, it would be interesting to study black hole solutions in this dilaton supergravity, cf. .
The relation between different supergravity theories can be schematically represented in the form of the following diagram:
$$\begin{array}{ccccc}& & U(1)_AU(1)_V& & \\ & & & & \\ U(1)_A& & & & U(1)_V\\ & & & & \\ & & \mathrm{𝟏}& & \end{array}$$
where the theory with the trivial holonomy group is the new $`𝒩=(2,2)`$ dilaton supergravity we are going to construct. Notice, however, that various arrows in this diagram have different meaning. For example, minimal theories with either of the $`U(1)`$ R-symmeties gauged can be obtained by truncation of the non-minimal $`U(1)_AU(1)_V`$ gauged supergravity . On the other hand, the vertical arrow corresponds to de-gauging $`U(1)_AU(1)_V`$ symmetry, so that the total number of degrees of freedom increases. More precisely, it has to be a combination of consistent truncation of the non-minimal $`𝒩=(2,2)`$ supergravity to a minimal one plus a de-gauging of the latter. To see this, let us count the number of real Goldstone scalars. Since the broken $`U(1)_AU(1)_V`$ phase of the theory has exactly the same field content as the gauge symmetric phase plus the field content of the Goldstone multiplets minus the parts that go into the longitudinal components of the $`U(1)_AU(1)_V`$ gauge fields, we find that in total Goldstone multiplets should have three real scalars. However, there are no mirror-symmetric $`𝒩=(2,2)`$ multiplets with such a field content. Therefore, we conclude that the vertical arrow should be a more economical de-gauging. Indeed, if we were following another route via a minimal gauged supergravity, at the first step we would have to make a consistent truncation that would eliminate one of the gauge fields. In order to de-gauge the resulting minimal gauged supergravity we would have to introduce extra chiral (or twisted chiral) Goldstone superfield. In any case, one real scalar from this multiplet would become the longitudinal component of the gauge vector field, and the other would become a dilaton, in agreement with what we expect. Assuming that the latter route (which is not, unfortunately, mirror-symmetric) is equivalent to the vertical arrow on the above diagram, we expect that there is an economical de-gauging of the non-minimal $`𝒩=(2,2)`$ supergravity that leads to only one massless scalar (the dilaton).
To start the construction, let us arrange component fields found in the previous section into superfields. First we give the definitions in flat superspace and then extend them to curved superspace. Left-right symmetric $`𝒩=2`$ superspace is parametrized by the bosonic coordinates $`x^\mu =(x^{\text{ }\text{ }\text{ }\text{ }},x^{\text{ }})`$, two anti-commuting complex spinor coordinates $`\theta ^\alpha =(\theta ^+,\theta ^{})`$ and their complex conjugates $`\overline{\theta }^{\dot{\alpha }}=(\overline{\theta }^{\dot{+}},\overline{\theta }^\dot{})`$. The spinor derivatives $`D_+`$, $`D_{}`$, $`D_{\dot{+}}`$ and $`D_\dot{}`$ satisfy $`\{D_+,D_{\dot{+}}\}=_{\text{ }\text{ }\text{ }\text{ }}`$ and $`\{D_{},D_\dot{}\}=_{\text{ }}`$, with all other (anti-)commutators vanishing.
Irreducible matter superfields are defined by imposing some constraints on general complex superfields. The simplest constraints are linear in derivatives and look like:
$$D_{\dot{+}}\mathrm{\Phi }=D_\dot{}\mathrm{\Phi }=D_+\overline{\mathrm{\Phi }}=D_{}\overline{\mathrm{\Phi }}=0$$
for a chiral superfield $`\mathrm{\Phi }`$, and:
$$D_{\dot{+}}\mathrm{\Sigma }=D_{}\mathrm{\Sigma }=D_+\overline{\mathrm{\Sigma }}=D_\dot{}\overline{\mathrm{\Sigma }}=0$$
for a twisted chiral superfield $`\mathrm{\Sigma }`$.
In what follows we promote the complex scalar fields $`\varphi _i`$ and $`\sigma _j`$ defined in the previous section to the chiral and twisted chiral superfields $`\mathrm{\Phi }_i`$ and $`\mathrm{\Sigma }_j`$, respectively. Similarly, we will regard the compact fields $`z_k`$ as the scalar components of (twisted) chiral superfields $`𝒵_k`$.
Components of the chiral superfield $`\mathrm{\Phi }_i`$ can be obtained using the projection method:
$$\mathrm{\Phi }_i|=\varphi _i,\overline{\mathrm{\Phi }}_i|=\overline{\varphi }_i$$
$$D_+\mathrm{\Phi }_i|=\psi _+^i,D_{\dot{+}}\overline{\mathrm{\Phi }}_i|=\psi _{\dot{+}}^i$$
$$D_{}\mathrm{\Phi }_i|=\psi _{}^i,D_\dot{}\overline{\mathrm{\Phi }}_i|=\psi _\dot{}^i$$
$$\frac{i}{2}[D_+,D_{}]\mathrm{\Phi }_i|=A_i,\frac{i}{2}[D_{\dot{+}},D_\dot{}]\overline{\mathrm{\Phi }}_i|=\overline{A}_i$$
where, for example, $`\mathrm{\Phi }_i|`$ denotes the leading component of the superfield $`\mathrm{\Phi }_i`$, with all the $`\theta `$-coordinates put to zero. Similarly, we find the components of a twisted chiral multiplet $`\mathrm{\Sigma }_j`$:
$$\mathrm{\Sigma }_j|=\sigma _j,\overline{\mathrm{\Sigma }}_j|=\overline{\sigma }_j$$
$$D_+\mathrm{\Sigma }_j|=\zeta _+^j,D_{\dot{+}}\overline{\mathrm{\Sigma }}_j|=\zeta _{\dot{+}}^j$$
$$D_\dot{}\mathrm{\Sigma }_j|=\zeta _\dot{}^j,D_{}\overline{\mathrm{\Sigma }}_j|=\zeta _{}^j$$
$$\frac{i}{2}[D_+,D_\dot{}]\mathrm{\Sigma }_j|=B_j,\frac{i}{2}[D_{\dot{+}},D_{}]\overline{\mathrm{\Sigma }}_j|=\overline{B}_j$$
To define chiral and twisted chiral superfields in curved $`𝒩=(2,2)`$ superspace, the spinor derivatives $`D_\alpha `$ must be appropriately replaced by the covariant derivatives $`_\alpha `$.
In our construction of $`𝒩=(2,2)`$ dilaton supergravity we start with non-minimal gauged supergravity and then de-gauge $`U(1)_AU(1)_V`$ symmetry. If we introduce superfields $`𝒜_\alpha `$ and $`𝒜_\alpha ^{}`$ representing $`U(1)_V`$ and $`U(1)_A`$ gauge connections, and denote by $`\mathrm{\Lambda }_\alpha `$ the Lorentz spin-connection, then the covariant derivative in this theory has the form:
$$_\alpha =E_{\alpha }^{}{}_{}{}^{B}D_B+\mathrm{\Lambda }_\alpha 𝒳+𝒜_\alpha 𝒴+𝒜_\alpha ^{}𝒴^{}$$
where $`E_{\alpha }^{}{}_{}{}^{B}`$ is the supervielbein.
The Lorentz generators, $`𝒳`$, $`U(1)_V`$ symmetry generators, $`𝒴`$, and $`U(1)_A`$ symmetry generators, $`𝒴^{}`$, act on the covariant derivative $`_\alpha `$ in the following way :
$$[𝒳,_\pm ]=\pm \frac{1}{2}_\pm ,[𝒳,_{\dot{\pm }}]=\pm \frac{1}{2}_{\dot{\pm }}$$
$$[𝒴,_\pm ]=\frac{i}{2}_\pm ,[𝒴,_{\dot{\pm }}]=+\frac{i}{2}_{\dot{\pm }}$$
$$[𝒴^{},_\pm ]=\frac{i}{2}_\pm ,[𝒴^{},_{\dot{\pm }}]=\pm \frac{i}{2}_{\dot{\pm }}$$
The constraints which define non-minimal $`𝒩=2`$ $`U(1)_AU(1)_V`$ supergravity are given by:
$$\{_+,_+\}=0,\{_{},_{}\}=0$$
$$\{_+,_{\dot{+}}\}=i_{\text{ }\text{ }\text{ }\text{ }},\{_{},_\dot{}\}=i_{\text{ }}$$
$$\{_+,_{}\}=\frac{1}{2}\overline{R}(\overline{𝒳}i\overline{𝒴}^{}),\{_+,_\dot{}\}=\frac{1}{2}\overline{F}(\overline{𝒳}i\overline{𝒴})$$
where the chiral superfield $`R`$ and the twisted chiral superfield $`F`$ are related to the two-dimensional curvature $`R^{(2)}`$ and the abelian field strengths of the graviphoton gauge fields. The easiest way to see this is to compute the commutator \[22,,25\]:
$$[_{\text{ }\text{ }\text{ }\text{ }},_{\text{ }}]=\frac{1}{2}\left((^2R)\frac{1}{2}R\overline{R}+(_+_\dot{}F)\frac{1}{2}F\overline{F}\right)𝒳+$$
$$+\frac{i}{2}(_+_\dot{}F)𝒴+\frac{i}{2}(^2R)𝒴^{}+\mathrm{}+\mathrm{c}.\mathrm{c}.$$
where the dots stand for the covariant derivative terms like $`(_{}F)_{\dot{+}}`$, etc. If we set $`F=0`$ in the constraints (4.1), we obtain $`U(1)_A`$ minimal gauged supergravity theory. On the other hand, if we set $`R=0`$, we end up with $`U(1)_V`$ minimal theory. We denote the leading components of the superfields $`F`$ and $`R`$ as follows:
$$F|=G,R|=H$$
In order to obtain $`𝒩=(2,2)`$ supergravity theory without gauged symmetry, we replace the covariant derivative $`_\alpha `$ by $`\widehat{}_\alpha `$ which includes only derivatives and Lorentz generator, so that:
$$_\alpha =\widehat{}_\alpha +𝒜_\alpha 𝒴+𝒜_\alpha ^{}𝒴^{}$$
and the remaining derivatives are real. The new covariant derivative $`\widehat{}_\alpha `$ contains neither the $`U(1)_V`$ gauge connection nor the $`U(1)_A`$ gauge connection, so it can describe two-dimensional $`𝒩=(2,2)`$ supergravity without gauged symmetry.
Substituting (4.1) into (4.1), and using (4.1), we derive the form of the commutator algebra for the $`\widehat{}_\alpha `$ operators:
$$\{\widehat{}_+,\widehat{}_+\}=i(\lambda _++\stackrel{~}{\lambda }_+)\widehat{}_+,\{\widehat{}_{},\widehat{}_{}\}=i(\lambda _{}\stackrel{~}{\lambda }_{})\widehat{}_{}$$
$$\{\widehat{}_+,\widehat{}_{}\}=\frac{1}{2}\overline{R}𝒳+\frac{i}{2}(\lambda _{}+\stackrel{~}{\lambda }_{})\widehat{}_++\frac{i}{2}(\lambda _+\stackrel{~}{\lambda }_+)\widehat{}_{}$$
$$\{\widehat{}_+,\widehat{}_\dot{}\}=\frac{1}{2}\overline{F}𝒳+\frac{i}{2}(\lambda _\dot{}+\stackrel{~}{\lambda }_\dot{})\widehat{}_+\frac{i}{2}(\lambda _+\stackrel{~}{\lambda }_+)\widehat{}_\dot{}$$
$$\{\widehat{}_+,\widehat{}_{\dot{+}}\}=i\widehat{}_{\text{ }\text{ }\text{ }\text{ }}+\frac{i}{2}(\lambda _{\dot{+}}+\stackrel{~}{\lambda }_{\dot{+}})\widehat{}_+\frac{i}{2}(\lambda _++\stackrel{~}{\lambda }_+)\widehat{}_{\dot{+}}$$
$$\{\widehat{}_{},\widehat{}_\dot{}\}=i\widehat{}_{\text{ }}+\frac{i}{2}(\lambda _\dot{}\stackrel{~}{\lambda }_\dot{})\widehat{}_{}\frac{i}{2}(\lambda _{}\stackrel{~}{\lambda }_{})\widehat{}_\dot{}$$
where the new fields $`\lambda _\alpha `$ and $`\stackrel{~}{\lambda }_\alpha `$ have appeared. They are components of the non-minimal gauged supergravity multiplet that could be eliminated by a gauge transformation in the $`U(1)_AU(1)_V`$ gauge-symmetric phase. In the theory we are constructing the $`U(1)_AU(1)_V`$ gauge symmetry is broken, so that the fields $`\lambda _\alpha `$ and $`\stackrel{~}{\lambda }_\alpha `$ are dynamical. They have to be identified with the spinorial derivatives of a new matter Goldstone multiplet. To the leading order in $`\theta ^\alpha `$, the fields $`\lambda _\alpha `$ and $`\stackrel{~}{\lambda }_\alpha `$ can also be identified with the dilatino field of the new supergravity multiplet.
If we introduce four linear independent spinors (along with their conjugates):
$$\eta _+\lambda _++\stackrel{~}{\lambda }_+,\eta _{}\lambda _{}\stackrel{~}{\lambda }_{},$$
$$\stackrel{~}{\eta }_+\lambda _+\stackrel{~}{\lambda }_+,\stackrel{~}{\eta }_{}\lambda _{}+\stackrel{~}{\lambda }_{}$$
then the conditions (4.1) which define $`𝒩=(2,2)`$ dilaton supergravity can be written in the following simple form:
$$\{\widehat{}_+,\widehat{}_+\}=i\eta _+\widehat{}_+,\{\widehat{}_{},\widehat{}_{}\}=i\eta _{}\widehat{}_{}$$
$$\{\widehat{}_+,\widehat{}_{}\}=\frac{1}{2}\overline{R}𝒳+\frac{i}{2}\stackrel{~}{\eta }_{}\widehat{}_++\frac{i}{2}\stackrel{~}{\eta }_+\widehat{}_{}$$
$$\{\widehat{}_+,\widehat{}_\dot{}\}=\frac{1}{2}\overline{F}𝒳+\frac{i}{2}\stackrel{~}{\eta }_\dot{}\widehat{}_+\frac{i}{2}\stackrel{~}{\eta }_+\widehat{}_\dot{}$$
$$\{\widehat{}_+,\widehat{}_{\dot{+}}\}=i\widehat{}_{\text{ }\text{ }\text{ }\text{ }}+\frac{i}{2}\eta _{\dot{+}}\widehat{}_+\frac{i}{2}\eta _+\widehat{}_{\dot{+}}$$
$$\{\widehat{}_{},\widehat{}_\dot{}\}=i\widehat{}_{\text{ }}+\frac{i}{2}\eta _\dot{}\widehat{}_{}\frac{i}{2}\eta _{}\widehat{}_\dot{}$$
Note that these equations, as well as the conditions (4.1), are manifestly invariant under the mirror symmetry transformation (3.1).
To summarize, there exists a unique mirror-symmetric two-dimensional $`𝒩=(2,2)`$ dilaton supergravity defined by the set of constraints (4.1) on the covariant derivative $`\widehat{}_\alpha `$. The new supergravity theory does not have gauged symmetry and contains a real dilaton field $`\phi `$, as we will show in a moment.
Since we define new $`𝒩=(2,2)`$ dilaton supergravity theory imposing constraints (4.1) on the covariant derivatives $`\widehat{}_\alpha `$, the Bianchi identities in this theory may lead to further constraints on some fields<sup>5</sup> In ordinary field theories, the fields satisfy Bianchi identities because they are expressed in terms of the potentials; they are identities and impose no extra constraints.. For theories described by covariant derivatives $`\widehat{}_\alpha `$, the Bianchi identities are simply Jacobi identities:
$$[\widehat{}_{[\alpha },[\widehat{}_\beta ,\widehat{}_{\gamma )}\}\}=0$$
where $`[,\}`$ is the graded commutator, and $`[,)`$ stands for the graded antisymmetrization symbol. Instead of deriving derivative constraints on the spinor fields $`\lambda _\alpha `$ and $`\stackrel{~}{\lambda }_\alpha `$ directly from the Jacobi identities (4.1) we use an equivalent approach which is much easier. While substituting (4.1) into (4.1) one also finds terms proportional to gauge symmetry generators $`𝒴`$ and $`𝒴^{}`$. Vanishing of these terms leads to a set of constraints which is equivalent to the set of constraints obtained from the Jacobi identities (4.1). We outline the result in appendix B.
There are two simple solutions to the Jacobi identities corresponding to either $`\eta _\alpha `$ or $`\stackrel{~}{\eta }_\alpha `$ put to zero. An advantage of the first solution is that the covariant derivative $`\widehat{}_\alpha `$ anti-commutes with itself, $`\widehat{}_\alpha ^2=0`$, like in the usual gauged supergravity theories . On the other hand, in the second case we find especially simple form of the anti-commutators $`\{\widehat{}_+,\widehat{}_{}\}`$ and $`\{\widehat{}_+,\widehat{}_\dot{}\}`$. In both cases the remaining spinor superfields can be expressed in terms of an unconstraint real superfield $`V`$:
$$V=\overline{V}$$
so that the Jacobi identities (4.1) are satisfied. This means that (4.1) impose no further constraints on $`V`$, and only define the other superfields (like $`F`$ and $`R`$) in terms of the derivatives of $`V`$. It is natural to identify the dilaton field with the leading scalar component of $`V`$:
$$\phi =V|$$
Below we present more evidence for this identification. One might notice that a real superfield $`V`$ contains one massless vector field, in agreement with the result of the previous section<sup>6</sup> However, massless vector fields in two dimensions do not have propagating degrees of freedom. For the same reason two-dimensional superfield $`V`$ does not have an irreducible transverse component, unlike a similar four-dimensional superfield.. It is also worthwhile to stress here that massless superfield $`V`$ is not a Goldstone multiplet itself, but rather what remains after the Goldstone mechanism takes place. A nice property of this solution is that $`V`$ is manifestly mirror-symmetric.
Since local integration measures of the new $`𝒩=(2,2)`$ dilaton supergravity can be nicely derived from the corresponding expressions of the $`U(1)_AU(1)_V`$ theory only for the solution corresponding to $`\stackrel{~}{\eta }_\alpha =0`$, in what follows we discuss in detail only this case. Namely, we take the following ansatz for the spinors $`\lambda _\alpha `$:
$$\lambda _+=\stackrel{~}{\lambda }_+=i(\widehat{}_+V),\lambda _{}=\stackrel{~}{\lambda }_{}=i(\widehat{}_{}V)$$
$$\lambda _{\dot{+}}=\stackrel{~}{\lambda }_{\dot{+}}=i(\widehat{}_{\dot{+}}V),\lambda _\dot{}=\stackrel{~}{\lambda }_\dot{}=i(\widehat{}_\dot{}V)$$
which implies $`\eta _\alpha =2i\widehat{}_\alpha V`$ and $`\stackrel{~}{\eta }_\alpha =0`$.
Substituting (4.1) into (4.1), we find the following supergravity algebra:
$$\{\widehat{}_+,\widehat{}_+\}=2(\widehat{}_+V)\widehat{}_+,\{\widehat{}_{},\widehat{}_{}\}=2(\widehat{}_{}V)\widehat{}_{}$$
$$\{\widehat{}_+,\widehat{}_{}\}=\frac{1}{2}\overline{R}𝒳,\{\widehat{}_+,\widehat{}_\dot{}\}=\frac{1}{2}\overline{F}𝒳$$
$$\{\widehat{}_+,\widehat{}_{\dot{+}}\}=i\widehat{}_{\text{ }\text{ }\text{ }\text{ }}+(\widehat{}_{\dot{+}}V)\widehat{}_++(\widehat{}_+V)\widehat{}_{\dot{+}}$$
$$\{\widehat{}_{},\widehat{}_\dot{}\}=i\widehat{}_{\text{ }}+(\widehat{}_\dot{}V)\widehat{}_{}+(\widehat{}_{}V)\widehat{}_\dot{}$$
where the superfields $`\overline{R}`$ and $`\overline{F}`$ can be obtained from the Bianchi identities. Solving the set of constraints in appendix B we get:
$$\overline{R}=4\widehat{}_{}\widehat{}_+V,\overline{F}=4\widehat{}_\dot{}\widehat{}_+V$$
We also find the following expressions for the gauge connection:
$$\lambda _{\text{ }\text{ }\text{ }\text{ }}=\stackrel{~}{\lambda }_{\text{ }\text{ }\text{ }\text{ }}=i\widehat{}_{\text{ }\text{ }\text{ }\text{ }}V2\widehat{}_+\widehat{}_{\dot{+}}V+4(\widehat{}_+V)(\widehat{}_{\dot{+}}V)$$
$$\lambda _{\text{ }}=\stackrel{~}{\lambda }_{\text{ }}=i\widehat{}_{\text{ }}V2\widehat{}_{}\widehat{}_\dot{}V+4(\widehat{}_{}V)(\widehat{}_\dot{}V)$$
Further commutators of the covariant derivatives with vector indices follow from the consistency of the Bianchi identities (4.1):
$$[\widehat{}_+,\widehat{}_{\text{ }\text{ }\text{ }\text{ }}]=\left(2i\lambda _+\lambda _{\dot{+}}+(\widehat{}_{\dot{+}}\lambda _+)+(\widehat{}_+\lambda _{\dot{+}})\right)\widehat{}_+$$
$$[\widehat{}_{\dot{+}},\widehat{}_{\text{ }\text{ }\text{ }\text{ }}]=\left(2i\lambda _+\lambda _{\dot{+}}+(\widehat{}_{\dot{+}}\lambda _+)+(\widehat{}_+\lambda _{\dot{+}})\right)\widehat{}_{\dot{+}}$$
$$[\widehat{}_{},\widehat{}_{\text{ }}]=\left(2i\lambda _{}\lambda _\dot{}+(\widehat{}_\dot{}\lambda _{})+(\widehat{}_{}\lambda _\dot{})\right)\widehat{}_{}$$
$$[\widehat{}_\dot{},\widehat{}_{\text{ }}]=\left(2i\lambda _{}\lambda _\dot{}+(\widehat{}_\dot{}\lambda _{})+(\widehat{}_{}\lambda _\dot{})\right)\widehat{}_\dot{}$$
$$[\widehat{}_+,\widehat{}_{\text{ }}]=\left(\frac{i}{2}(\widehat{}_\dot{}\overline{R})+\frac{i}{2}(\widehat{}_{}\overline{F})+\frac{1}{2}\lambda _\dot{}\overline{R}\frac{1}{2}\lambda _{}\overline{F}\right)𝒳\frac{i}{2}\overline{R}\widehat{}_\dot{}\frac{i}{2}\overline{F}\widehat{}_{}$$
$$[\widehat{}_{\dot{+}},\widehat{}_{\text{ }}]=\left(\frac{i}{2}(\widehat{}_{}R)+\frac{i}{2}(\widehat{}_\dot{}F)\frac{1}{2}\lambda _{}R+\frac{1}{2}\lambda _\dot{}F\right)𝒳+\frac{i}{2}R\widehat{}_{}+\frac{i}{2}F\widehat{}_\dot{}$$
$$[\widehat{}_{},\widehat{}_{\text{ }\text{ }\text{ }\text{ }}]=\left(\frac{i}{2}(\widehat{}_{\dot{+}}\overline{R})+\frac{i}{2}(\widehat{}_+F)+\frac{1}{2}\lambda _{\dot{+}}\overline{R}\frac{1}{2}\lambda _+F\right)𝒳+\frac{i}{2}\overline{R}\widehat{}_{\dot{+}}+\frac{i}{2}F\widehat{}_+$$
$$[\widehat{}_\dot{},\widehat{}_{\text{ }\text{ }\text{ }\text{ }}]=\left(\frac{i}{2}(\widehat{}_+R)+\frac{i}{2}(\widehat{}_{\dot{+}}\overline{F})\frac{1}{2}\lambda _+R+\frac{1}{2}\lambda _{\dot{+}}\overline{F}\right)𝒳\frac{i}{2}R\widehat{}_+\frac{i}{2}\overline{F}\widehat{}_{\dot{+}}$$
where we used (4.1). It is worthwhile to stress here that one would obtain a different result de-gauging the corresponding commutators in the $`U(1)_AU(1)_V`$ non-minimal supergravity .
5. Lagrangians for Matter Multiplets Coupled to $`𝒩=2`$ Dilaton Supergravity
In order to couple matter fields to new $`𝒩=(2,2)`$ dilaton supergravity we have to repeat the analysis of . Up to terms with two derivatives or four fermions, the most general action of $`𝒩=(2,2)`$ supergravity coupled to chiral superfields $`\mathrm{\Phi }_i`$ and $`𝒵_k`$, and twisted chiral superfields $`\mathrm{\Sigma }_j`$ looks like (3.1):
$$S=d^2xd^2\theta d^2\overline{\theta }E^1(\mathrm{\Phi }_i,\overline{\mathrm{\Phi }}_{\overline{i}},\mathrm{\Sigma }_j,\overline{\mathrm{\Sigma }}_j,𝒵_k,\overline{𝒵}_{\overline{k}})+$$
$$+d^2xd^2\theta ^1W(\mathrm{\Phi }_i)+d^2x𝑑\theta ^+𝑑\theta ^\dot{}\stackrel{~}{}^1\stackrel{~}{W}(\mathrm{\Sigma }_j)+\mathrm{c}.\mathrm{c}.$$
In order to obtain the component action corresponding to (5.1), one needs the appropriate projection formulas. For gauged $`𝒩=(2,2)`$ supergravity theories such formulas were derived by Grisaru and Wehlau . In the case of the minimal $`U(1)_A`$ theory the local density projection formula has the following form:
$$d^2xd^4\theta E^1=d^2xe^1[^2+i\psi _{\text{ }}^\dot{}_+i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}_{}+$$
$$+(\frac{1}{2}\overline{H}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}})]\overline{}^2|$$
Here $`\psi _\mu ^\alpha `$ is the gravitino field and $``$ is an arbitrary scalar function of superfields. In fact, the same projection formula is also valid in the non-minimal $`U(1)_AU(1)_V`$ supergravity theory . Although (5.1) is a D-type superinvariant, sometimes it is called a chiral density projector because in the non-minimal $`𝒩=(2,2)`$ supergravity $`\overline{}^2`$ is a chiral superfield (for a general $``$). Therefore, replacing $`\overline{}^2`$ by an arbitrary covariantly chiral Lagrangian $`_c`$, we can obtain the component projection formula for any chiral superspace integral:
$$d^2xd^2\theta ^1_c=d^2xe^1[^2+i\psi _{\text{ }}^\dot{}_+i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}_{}+$$
$$+(\frac{1}{2}\overline{H}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}})]_c|$$
In particular, the superspace measures $`E^1`$ and $`^1`$ are related as follows:
$$d^2xd^4\theta E^1=d^2xd^2\theta ^1\overline{}^2|$$
By mirror symmetry the twisted chiral density projection formula in the $`U(1)_V`$ gauged supergravity theory has the following form:
$$d^2xd^4\theta E^1=d^2xe^1[_\dot{}_+i\psi _{\text{ }}^{}_++i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}_\dot{}+$$
$$+(\frac{1}{2}\overline{G}+\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }}^{})]_{\dot{+}}_{}|$$
As explained in , the derivation of this formula goes through as in for the case of $`U(1)_A`$ theory. In the case of the $`U(1)_VU(1)_A`$ gauged supergravity the symmetry between chiral and twisted chiral fields is restored by the contribution of the anticommutator term $`\{_+,_\dot{}\}\overline{F}|`$. More explicitly, the twisted chiral density projection formula (5.1) can be derived using the methods of or .
The projection formulas in the two-dimensional $`𝒩=(2,2)`$ dilaton supergravity can be obtained from (5.1) and (5.1) replacing $`_\alpha `$ by a new covariant derivative $`\widehat{}_\alpha `$. Thus, substituting (4.1) and (4.1) into (5.1) and using the commutation relations (4.1) – (4.1) we obtain the following density projection formula:
$$d^2xd^4\theta E^1=d^2xe^1[(\widehat{}_+(\widehat{}_+V))(\widehat{}_{}(\widehat{}_{}V))+i\psi _{\text{ }}^\dot{}(\widehat{}_+(\widehat{}_+V))$$
$$i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\widehat{}_{}(\widehat{}_{}V))+(\frac{1}{2}\overline{H}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}})]\widehat{\overline{}}^2|=$$
$$=d^2xe^1[\widehat{}_+\widehat{}_{}+i(\psi _{\text{ }}^\dot{}\lambda _{})\widehat{}_+i(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)\widehat{}_{}+$$
$$+(\frac{1}{4}\overline{H}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+(\psi _{\text{ }}^\dot{}\lambda _{})(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+))]\widehat{}_{\dot{+}}\widehat{}_\dot{}|$$
In order to convince even hard boiled sceptics that (5.1) is the right projector, in appendix D we repeat the calculation of Grisaru and Wehlau in the new $`𝒩=(2,2)`$ dilaton supergravity. As expected, the result is equivalent to (5.1).
Similarly, the twisted chiral density projection formula (5.1) yields:
$$d^2xd^4\theta E^1=d^2xe^1[(\widehat{}_\dot{}(\widehat{}_\dot{}V))(\widehat{}_+(\widehat{}_+V))i\psi _{\text{ }}^{}(\widehat{}_+(\widehat{}_+V))+$$
$$+i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\widehat{}_\dot{}(\widehat{}_\dot{}V))+(\frac{1}{2}\overline{G}+\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }}^{})]\widehat{}_{\dot{+}}\widehat{}_{}|$$
We note that for a given superspace Lagrangian $``$, both projection formulas (5.1) and (5.1) lead to the same result:
$$d^2xd^4\theta E^1=d^2xe^1[(\widehat{}_+(\widehat{}_+V))(\widehat{}_{}(\widehat{}_{}V))+i\psi _{\text{ }}^\dot{}(\widehat{}_+(\widehat{}_+V))$$
$$i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\widehat{}_{}(\widehat{}_{}V))+(\frac{1}{2}\overline{H}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}})]\widehat{\overline{}}^2|=$$
$$=d^2xe^1[(\widehat{}_\dot{}(\widehat{}_\dot{}V))(\widehat{}_+(\widehat{}_+V))i\psi _{\text{ }}^{}(\widehat{}_+(\widehat{}_+V))+$$
$$+i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\widehat{}_\dot{}(\widehat{}_\dot{}V))+(\frac{1}{2}\overline{G}+\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }}^{})]\widehat{}_{\dot{+}}\widehat{}_{}|$$
This follows from the corresponding property of the local density projectors in gauged $`𝒩=(2,2)`$ supergravity theory , and also can be verified explicitly using the commutation relations (4.1) in the new $`𝒩=(2,2)`$ dilaton supergravity.
Now we are ready to derive component actions for various superspace Lagrangians $``$. Let us start with a simple example corresponding to pure dilaton supergravity. Obviously, in order to reproduce the right exponential dependence on the dilaton field in (3.1), we have to take the function $``$ in the form:
$$_{\mathrm{grav}}=\mathrm{exp}(2V).$$
Substituting this in the projection formula (5.1) (or (5.1)) and seting all the fermions to zero, we obtain the following action for bosonic fields:
$$S_{\mathrm{grav}}=d^2xe^1\mathrm{exp}(2\phi )\left[R^{(2)}+4(_\mu \phi )(^\mu \phi )\right]$$
Here we used (4.1) and the formulas for the other components of the real superfield $`V`$ derived in appendix C. Clearly, the action (5.1) for dilaton and graviton fields agrees<sup>7</sup> We will account for the extra volume factor $`𝒱`$ in a moment. with the first two terms in the effective action (3.1) of Type IIA theory on a Calabi-Yau four-fold. Moreover, we note that bosonic action (5.1) has exactly the same form as the action of $`𝒩=0`$ dilaton gravity studied long time ago, see e.g. .
Now we consider coupling of $`𝒩=(2,2)`$ dilaton supergravity to matter fields. In particular, we are interested in superspace form of the effective action (3.1) describing compactification of Type IIA string theory on a Calabi-Yau four-fold $`X`$. Once again, to reproduce the exponential dependence on $`\phi =V|`$ we take the superspace action in the form:
$$d^2xd^4\theta E^1\mathrm{exp}(2V)$$
where $``$ is a function of all matter superfields but $`V`$. It is convenient to absorb $`\mathrm{exp}(2V)`$ in the definition of the supervielbein determinant $`E_0^1=E^1\mathrm{exp}(2V)`$, so that $`=1`$ corresponds to pure supergravity action (5.1), like in $`𝒩=1`$ four-dimensional theory. Commuting $`\mathrm{exp}(2V)`$ to the left in (5.1), we find the modified projection formula:
$$d^2xd^4\theta E_0^1=$$
$$=d^2xe^1\mathrm{exp}(2\phi )[(\widehat{}_+3(\widehat{}_+V))(\widehat{}_{}3(\widehat{}_{}V))+i\psi _{\text{ }}^\dot{}(\widehat{}_+3(\widehat{}_+V))$$
$$i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\widehat{}_{}3(\widehat{}_{}V))+(\frac{1}{2}\overline{H}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}})](\widehat{}_{\dot{+}}2(\widehat{}_{\dot{+}}V))(\widehat{}_\dot{}2(\widehat{}_\dot{}V))|$$
Applying this projection formula to an arbitrary function $`(\mathrm{\Phi }_i,\overline{\mathrm{\Phi }}_{\overline{i}},\mathrm{\Sigma }_j,\overline{\mathrm{\Sigma }}_j)`$ of chiral superfields $`\mathrm{\Phi }_i`$ and twisted chiral superfields $`\mathrm{\Sigma }_j`$ we get the action of the bosonic fields (the vielbein determinant $`e^1`$ is suppressed):
$$L=e^{2\phi }[R^{(2)}+4_\mu (\phi \frac{1}{2}\mathrm{log})^\mu (\phi \frac{1}{2}\mathrm{log})+$$
$$+\frac{1}{2}(\mathrm{log})_{\varphi _i\overline{\varphi }_j}(_\mu \varphi _i)(^\mu \overline{\varphi }_j)\frac{1}{2}(\mathrm{log})_{\sigma _i\overline{\sigma }_j}(_\mu \sigma _i)(^\mu \overline{\sigma }_j)+$$
$$+\frac{1}{2}ϵ^{\mu \nu }((\mathrm{log})_{\varphi _i\overline{\sigma }_j}(_\mu \varphi _i)(_\nu \overline{\sigma }_j)+(\mathrm{log})_{\sigma _j\overline{\varphi }_i}(_\mu \overline{\varphi }_i)(_\nu \sigma _j))]$$
The subscripts on $``$ denote derivatives with respect to the scalar fields, e.g. $`(\mathrm{log})_{\varphi _i\overline{\varphi }_j}=(^2/\varphi _i\overline{\varphi }_j)\mathrm{log}`$. Deriving (5.1) one may find helpful some formulas from appendix D where we discuss in detail the component action of a free chiral superfield. A careful reader may notice that (5.1) has the structure reminiscent of $`𝒩=1`$ supergravity in four dimensions. In particular, it is convenient to introduce the Kähler potential $`K`$:
$$=\mathrm{exp}(K)$$
so that the superspace action (5.1) takes the form:
$$d^2xd^4\theta E_0^1e^K=d^2xd^4\theta E^1e^{2V}e^K.$$
Performing the superspace integration, one finds Lagrangian for the bosonic fields:
$$L=e^{2\stackrel{~}{\phi }}[R^{(2)}+4(_\mu \stackrel{~}{\phi })(^\mu \stackrel{~}{\phi })\frac{1}{2}K_{\varphi _i\overline{\varphi }_j}(_\mu \varphi _i)(^\mu \overline{\varphi }_j)+$$
$$+\frac{1}{2}K_{\sigma _i\overline{\sigma }_j}(_\mu \sigma _i)(^\mu \overline{\sigma }_j)\frac{1}{2}ϵ^{\mu \nu }(K_{\varphi _i\overline{\sigma }_j}(_\mu \varphi _i)(_\nu \overline{\sigma }_j)+K_{\sigma _j\overline{\varphi }_i}(_\mu \overline{\varphi }_i)(_\nu \sigma _j))]$$
where we introduced a new dilaton field $`\stackrel{~}{\phi }=\phi +\frac{1}{2}K`$ invariant under generalized Kähler transformations (3.1):
$$KK+\mathrm{\Lambda }_1(\varphi _i,\sigma _j)+\overline{\mathrm{\Lambda }}_1(\overline{\varphi }_{\overline{i}},\overline{\sigma }_j)+\mathrm{\Lambda }_2(\varphi _i,\overline{\sigma }_j)+\overline{\mathrm{\Lambda }}_2(\overline{\varphi }_{\overline{i}},\sigma _j)$$
Under this Kähler transformation the original dilaton field $`\phi `$ is shifted in the opposite way, so that the linear combination $`\stackrel{~}{\phi }=\phi +\frac{1}{2}K`$ remains invariant. Since the Kähler metric is invariant under (5.1) as well, both the superspace action (5.1) and the corresponding component action (5.1) are manifestly invariant under the generalized Kähler transformations (5.1).
Now we are in position to identify the function $`K`$ that would reproduce the effective action (3.1) of Type IIA theory on a Calabi-Yau four-fold. Namely, the superspace action (5.1) gives the effective action (3.1) if $`K`$ is the total Kähler potential (3.1). This form of the superspace action might be expected for a number of reasons. First of all, it is similar to the superspace action of $`𝒩=1`$ supergravity in four dimensions. Moreover, Type IIA supergravity on a Calabi-Yau four-fold has a breathing mode corresponding to rescaling of the volume $`𝒱c^2𝒱`$ and simultaneous shift of the dilaton $`\phi \phi +\mathrm{log}c`$, cf. (3.1). Therefore, the superspace action is expected to be a function of $`2V+K`$, where $`K`$ is the total Kähler potential given by (3.1) - (3.1).
To summarize, we constructed superspace Lagrangians describing $`𝒩=(2,2)`$ dilaton supergravity coupled to matter superfields and found projection formulas that allow one to rewrite integrals over the entire superspace in terms of component fields. Therefore, we provide a superspace formulation of the effective field theories constructed from compactification of Type IIA string theory on Calabi-Yau four-folds, as well as more general $`𝒩=(2,2)`$ sigma-models with torsion coupled to dilaton supergravity, cf. .
Incorporation of superpotential terms is more subtle. These terms are superinvariants obtained by integration only over a half of the superspace, cf. (5.1). Unfortunately, unlike (5.1) and (5.1), the chiral and twisted chiral density projectors in $`𝒩=(2,2)`$ dilaton supergravity do not simply follow from the full superspace projector (5.1). However, by dimensional arguments and from an examination of the index structure of the possible terms, the chiral density projection formula must look like:
$$d^2xd^2\theta ^1W=d^2xe^1[\widehat{}^2+\mathrm{}\frac{1}{2}\overline{H}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}]W|$$
where the dots stand for term containing $`\lambda _\alpha `$ or terms linear in covariant derivatives<sup>8</sup> These terms will not affect the action of bosonic fields.. By the similar reasoning, the twsited chiral density projection formula must look like (5.1):
$$d^2xd^2\theta \stackrel{~}{}^1\stackrel{~}{W}=d^2xe^1[\widehat{}_\dot{}\widehat{}_++\mathrm{}\frac{1}{2}\overline{G}+\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^{\dot{+}}+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }}^{}]\stackrel{~}{W}|$$
Although normalization and coefficients in (5.1) and (5.1) may not be correct, the terms quadratic in the gravitino are rather general and, in particular, are independent on de-gauging. So, we infer that one effect of the superpotential is to produce a mass term for the gravitino fields :
$$m_{\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}}W,m_{\psi _{\text{ }\text{ }\text{ }\text{ }}^{}}\stackrel{~}{W}$$
It is this property of $`𝒩=(2,2)`$ supergravity that was needed in \[13,,14\] in order to find the superpotentials (3.1) and (3.1) induced by Ramond-Ramond fluxes. Moreover, let us demonstrate that (5.1) and (5.1) lead to the expected structure of the scalar potential (3.1). Extending the computation of (5.1), from (5.1) we get the action of the auxiliary fields:
$$L_{\mathrm{aux}}[|\frac{1}{2}Hi\overline{A}_i(\mathrm{log})_{\varphi _i}|^2+|\frac{1}{2}Gi\overline{B}_j(\mathrm{log})_{\sigma _j}|^2+(\mathrm{log})_{\varphi _i\overline{\varphi }_j}A_i\overline{A}_j$$
$$(\mathrm{log})_{\sigma _i\overline{\sigma }_j}B_i\overline{B}_j]i(W_{\varphi _i}(\mathrm{log})_{\varphi _i})A_iW(\frac{1}{2}\overline{H}+iA_i(\mathrm{log})_{\varphi _i})$$
$$i(\stackrel{~}{W}_{\sigma _i}(\mathrm{log})_{\sigma _i})B_i\stackrel{~}{W}(\frac{1}{2}\overline{G}+iB_i(\mathrm{log})_{\sigma _i})+\mathrm{c}.\mathrm{c}.$$
Integrating out the auxiliary fields and using (5.1), we find the expected scalar potential (3.1):
$$L_{\mathrm{aux}}e^K\left(K_{\varphi _i\overline{\varphi }_j}^1(D_{\varphi _i}W)(D_{\overline{\varphi }_j}\overline{W})K_{\sigma _i\overline{\sigma }_j}^1(D_{\sigma _i}\stackrel{~}{W})(D_{\overline{\sigma }_j}\overline{\stackrel{~}{W}})|W|^2|\stackrel{~}{W}|^2\right)$$
with the covariant derivatives (3.1).
6. Compactification of Type IIB String Theory on Calabi-Yau Four-folds
Compactification of Type IIB string theory on a Calabi-Yau four-fold $`X`$ leads to a chiral $`𝒩=(0,4)`$ supersymmetric effective field theory in two non-compact dimensions. In the low-energy limit this theory is described by $`𝒩=(0,4)`$ supergravity coupled to scalar superfields. In this section we perform a Kaluza-Klein reduction on a Calabi-Yau four-fold $`X`$ and, in particular, find that the supergravity multiplet includes a real scalar dilaton field instead of an $`SU(2)`$ gauge field \[28,,29\]. A manifestly supersymmetric formulation of such $`𝒩=(0,4)`$ supergravity theory will be presented in the next section. In this section we show that Kaluza-Klein harmonics combine into $`𝒩=(0,4)`$ scalar superfields \[29,,30\]. Furthermore, since we deal with $`𝒩=(0,4)`$ supersymmetry, there is a significant difference between right-movers and left-movers. Namely, all left-moving modes are supersymmetry singlets. We will also see that the difference between the zero-point energy of the left-movers and the right-movers is proportional to the Euler number of $`X`$ .
In the large volume limit the bosonic spectrum of light modes in Type IIB string theory includes the metric $`g_{MN}`$, the dilaton $`\phi `$, the axion $`l`$, the 4-form tensor $`D_{MNPQ}`$ and two tensor fields $`B_{MN}^{RR}`$ and $`B_{MN}^{NS}`$. Together with the fermionic superpartners all these fields fit into Type IIB supergravity multiplet. Non-perturbative Type IIB string theory is invariant under $`SL(2,\text{ZZ})`$ duality group. In order to see the action of this group on the supergravity fields, it is convenient to define the following quantities:
$$\lambda =l+ie^\phi $$
$$𝐌=\frac{1}{\mathrm{Im}\lambda }\left(\begin{array}{cc}|\lambda |^2& \mathrm{Re}\lambda \\ \mathrm{Re}\lambda & 1\end{array}\right)$$
$$𝐇_{MNP}=\left(\begin{array}{c}_{[M}B_{NP]}^{NS}\\ \\ _{[M}B_{NP]}^{RR}\end{array}\right)$$
and
$$F_{MNPQR}=_{[M}D_{NPQR]}+\frac{3}{4}B_{[MN}^{NS}_PB_{QR]}^{RR}$$
Then, the field strength $`F_{MNPQR}`$ is a singlet under the $`SL(2,\text{ZZ})`$ duality group, while $`𝐇`$ transforms as a “vector”. Finally, $`SL(2,\text{ZZ})`$ acts on a complex scalar $`\lambda `$ in the usual way:
$$\lambda \frac{a\lambda +b}{c\lambda +d}$$
where the integer numbers $`a`$, $`b`$, $`c`$ and $`d`$ satisfy $`adbc=1`$.
The five-form field strength $`F`$ is self-dual:
$$F=F$$
Although this equation can not be derived from any action, for a moment we ignore this subtlety and write bosonic “Type IIB supergravity Lagrangian” simply as:
$$L_{(10)}=\sqrt{g}\left[\frac{1}{4}R^{(10)}+\frac{1}{16}\mathrm{Tr}\left(𝐌𝐌\right)+\frac{3}{16}𝐇^T𝐌𝐇+\frac{5}{24}F^2+\mathrm{}\right]$$
where the dots stand for higher derivative terms.
As in Type IIA theory, in order to find the zero-mode spectrum we have to expand Type IIB supergravity fields in harmonic $`(p,q)`$-forms on the space $`X`$. The metric modes are exactly the same as in (2.1). Namely, from the reduction of $`g_{MN}`$ we find $`h^{1,1}`$ real scalars $`s_i`$, $`h^{3,1}`$ complex scalars $`\varphi _j`$ and the two-dimensional metric $`g_{\mu \nu }`$. It turns out that one of the scalars $`s_i`$ comes into the two-dimensional supergravity multiplet. Namely, it is the Kaluza-Klein mode corresponding to the volume of the Calabi-Yau four-fold:
$$𝒱=d_{ijkl}s^is^js^ks^l$$
where $`d_{ijkl}`$ are the intersection numbers of $`X`$ given by (2.1). With this mode excluded, the Kähler deformations of the metric yield $`h^{1,1}1`$ scalars $`\stackrel{ˇ}{s}_i=𝒱^{\frac{1}{4}}s_i`$ satisfying the condition:
$$d_{ijkl}\stackrel{ˇ}{s}^i\stackrel{ˇ}{s}^j\stackrel{ˇ}{s}^k\stackrel{ˇ}{s}^l=1$$
Expanding the doublet of tensor fields as:
$$B^{NS}=\underset{i=1}{\overset{h^{1,1}}{}}r_i\omega _i^{(1,1)},B^{RR}=\underset{i=1}{\overset{h^{1,1}}{}}t_i\omega _i^{(1,1)},$$
we get pairs of real scalars $`r_i`$ and $`t_i`$, $`h^{1,1}`$ in number. All these modes are both right-moving and left-moving.
Expansion of the self-dual field $`D`$ is a bit subtle. Namely, instead of $`D`$ one has to expand the field strength $`F_i_\mu u_i\omega _i^{(4)}`$ and impose the self-duality condition (6.1). Depending on whether the form $`\omega _i^{(4)}H^4(X,\mathrm{IR})`$ is self-dual or anti-self-dual the scalar field $`u_i`$ is left-moving or right-moving, respectively. Therefore, we have to distinguish carefully self-dual and anti-self-dual harmonics of $`F`$. To this end we recall some topological properties of Calabi-Yau four-folds. There is a decomposition of the space of the middle dimensional forms on $`X`$:
$$H^4(X,\mathrm{IR})=B_+(X)B_{}(X)$$
where we denote by $`B_+(X)`$ (resp. $`B_{}(X)`$) the space of (anti-)self-dual 4-forms on $`X`$. Let us call the corresponding dimensions $`b_\pm =\mathrm{dim}B_\pm (X)`$. Then $`b_+`$ and $`b_{}`$ are related by the Hirzebruch signature <sup>9</sup> The explicit form of the Pontryagin classes is given by:
$$p_1=\frac{1}{2}\mathrm{tr}R^2,p_2=\frac{1}{4}\mathrm{tr}R^4+\frac{1}{8}(\mathrm{tr}R^2)^2$$
(see e.g. ):
$$\tau (Q)=b_+b_{}=\frac{1}{45}_X(7p_2p_1^2)=\frac{\chi }{3}+32$$
of the quadratic form $`Q(\omega _1,\omega _2)=_X\omega _1\omega _2`$. On the other hand, we also have $`b_4=b_++b_{}=2+2h^{3,1}+h^{2,2}`$. Hence, using (2.1) and (2.1) we find:
$$b_+=47+3h^{1,1}+4h^{3,1}2h^{2,1}$$
and
$$b_{}=1+h^{1,1}+2h^{3,1}$$
Actually, we can be a little bit more precise. All the forms of Hodge type $`(3,1)`$ or $`(1,3)`$ are anti-self-dual, while the $`(4,0)`$\- and $`(0,4)`$-forms on a Calabi-Yau four-fold are self-dual . Therefore, from (6.1) and (6.1) we find that:
$$b_+^{(2,2)}=45+3h^{1,1}+4h^{3,1}2h^{2,1}$$
and
$$b_{}^{(2,2)}=h^{1,1}1$$
Now we expand the self-dual field strength $`F`$ as:
$$F=\underset{i=1}{\overset{b_+^{2,2}}{}}(_\mu u_i)\omega _i^{(+)}+\underset{j=1}{\overset{b_{}^{2,2}}{}}(_\mu v_j)\omega _j^{()}+\underset{k=1}{\overset{h^{3,1}}{}}(_\mu p_k)\omega _k^{(3,1)}+(_\mu q)\mathrm{\Omega }+\mathrm{c}.\mathrm{c}.$$
where the scalar fields $`u_i`$ and $`v_j`$ are real, while $`p_k`$ and $`q`$ are complex. As we explained above, $`u_i`$ and $`q`$ must be left-moving, $`v_j`$ and $`p_k`$ must be right-moving. In particular, the former are singlets with respect to four left supercharges satisfying:
$$\{Q_+^i,Q_+^j\}=\delta ^{ij}P_+$$
Note that the Kaluza-Klein modes of the self-dual field $`F`$ associated with $`(2,1)`$-harmonic forms give 2-form field strengths of massless vector fields in two dimensions and, therefore, do not lead to new propagating degrees of freedom.
Unlike $`𝒩=(2,2)`$ theory constructed from compactification of Type IIA string theory, in Type IIB compactification on on a Calabi-Yau four-fold $`X`$ not all the fermionic modes can be determined by $`𝒩=(0,4)`$ supersymmetry. In the right sector we still can use supersymmetry arguments to conclude that two-dimensional supergravity multiplet contains $`2\mathrm{ind}(/D)`$ fermions and $`2\mathrm{ind}(/D)`$ Rarita-Schwinger fields that come from the corresponding spin-$`\frac{1}{2}`$ and spin-$`\frac{3}{2}`$ fields in Type IIB supergravity. On a Calabi-Yau four-fold the Dirac index is given by:
$$\mathrm{ind}(/D)=\frac{1}{1440}_X(\frac{7}{4}p_1^2p_2)=2$$
in accordance with $`𝒩=(0,4)`$ supersymmetry. Furthermore, all the right-moving scalars found above ($`\varphi _i`$, $`\overline{\varphi }_{\overline{i}}`$, $`p_i`$, $`\overline{p}_{\overline{i}}`$, $`\stackrel{ˇ}{s}_j`$, $`r_j`$, $`t_j`$, $`v_j`$, $`\phi `$ and $`l`$) are accompanied by right-moving fermions. Simple counting gives:
$$n_+=4h^{3,1}+4h^{1,1}$$
for the total number of the right-moving fermions <sup>10</sup> Note that $`n_+`$ is divisible by 4.. There are also left-moving fermions which are supersymmetry singlets. The number of left-moving fermions, however, is not determined by supersymmetry. So, it has to be computed separately. Since the fermions in question come from the Type IIB gravitinos, their number (minus the number of right-moving fermions) is given by the Rarita-Schwinger index:
$$n_{}n_+=2\mathrm{ind}(/D_{3/2})$$
Using (6.1) and the explicit expression for the Rarita-Schwinger index on a Calabi-Yau four-fold $`X`$:
$$\mathrm{ind}(/D_{3/2})=\frac{1}{180}_X(\frac{37}{4}p_1^231p_2)=4h^{1,1}4h^{3,1}+4h^{2,1}$$
we obtain:
$$n_{}=4h^{2,1}$$
Now we are ready to assemble the supermultiplets. Combining the left-moving bosonic modes with the fermion fields we get the following supermultiplets:
$$\mathrm{a}\mathrm{gravitational}\mathrm{multiplet}:g_{\mu \nu },𝒱$$
$$h^{3,1}\mathrm{scalar}\mathrm{multiplets}\mathrm{\Phi }_i:\varphi _i,\overline{\varphi }_i,p_i,\overline{p}_i$$
$$h^{1,1}\mathrm{scalar}\mathrm{multiplets}\mathrm{\Sigma }_j:\stackrel{ˇ}{s}_j,r_j,t_j,v_j,\phi ,l$$
Performing a reduction of the ten-dimensional supersymmetry conditions one can easily check that these fields indeed represent bosonic components of the supermultiplets as stated. Note, all the matter multiplets include four real scalar fields in accordance with the general classification of scalar superfields \[28,,32,,29,,30\]. However, the content of the gravitational multiplet is different from what was usually studied in $`𝒩=(0,4)`$ supergravity theories \[28,,29\]. In compactification of Type IIB string theory on a Calabi-Yau four-fold we find that the gravitational multiplet includes a real scalar $`𝒱`$ instead of $`SU(2)`$ gauge field. In the next section we describe the superspace formulation of this $`𝒩=(0,4)`$ supergravity using the Goldstone approach.
In order to find the low-energy effective action one has to substitute (2.1), (6.1) and (6.1) into (6.1). Integrating over the internal space by means of the formulas (2.1), (2.1) we get the following two-dimensional action for bosonic fields:
$$L_{(2)}=\sqrt{g}𝒱[\frac{1}{4}R^{(2)}+\frac{1}{2}(_+\lambda )(_{}\lambda )+\frac{1}{2}𝒢_{\varphi _i\overline{\varphi }_{\overline{j}}}(_+\varphi ^i)(_{}\overline{\varphi }^{\overline{j}})+\frac{1}{2}𝒢_{\varphi _i\overline{\varphi }_{\overline{j}}}(_0p^i)(_+\overline{p}^{\overline{j}})+$$
$$+\frac{1}{2}𝒢_{\sigma _i\overline{\sigma }_j}\left((_+s^i)(_{}s^j)+(_+r^i)(_{}r^j)+(_+t^i)(_{}t^j)+(_0v^i)(_+v^j)\right)+$$
$$+\frac{1}{2}Q_{ij}(_0u^i)(_{}u^j)+\frac{1}{2}(_0q)(_{}q)+\mathrm{c}.\mathrm{c}.]$$
This Lagrangian describes non-linear sigma-model interacting with $`𝒩=(0,4)`$ supergravity. The target space of the left-moving fields is the cotangent bundle to the moduli space of the Calabi-Yau manifold $`X`$, $`T^{}_c(X)\times T^{}_𝒦(X)`$, cf. . Since the moduli space itself is a Kähler space, this result agrees with the general analysis of $`𝒩=(0,4)`$ supersymmetric sigma-models. According to , $`𝒩=(0,4)`$ supersymmetric sigma-model is based on a target space which has three covariantly constant (with respect to $`_+`$) complex structures which obey the quaternionic algebra:
$$J_rJ_s=\delta _{rs}+f_{rs}^{}{}_{}{}^{t}J_t$$
Another interesting feature that we expect to see in this $`𝒩=(0,4)`$ theory is $`SL(2,\mathrm{IR})`$ symmetry of classical Type IIB supergravity. Apart from $`\phi `$, $`l`$, $`r_i`$ and $`t_i`$, all the scalar fields listed above are singlets with respect to this symmetry. The complex scalar $`\lambda =l+ie^\phi `$ transforms as (6.1) under $`SL(2,\mathrm{IR})`$ duality transformation, with $`a`$, $`b`$, $`c`$ and $`d`$ real numbers obeying $`adbc=1`$. The doublet of real scalar fields $`(r_i,t_i)`$ transforms as a vector under the general $`SL(2,\mathrm{IR})`$ transformation. In other words, only the scalar multiplets $`\mathrm{\Sigma }_i`$ transform non-trivially under this symmetry, while all the other fields, including the supergravity itself, are $`SL(2,\mathrm{IR})`$-singlets.
Finally, we remark that from T-duality with Type IIA string theory on a Calabi-Yau four-fold we expect an anomaly similar to (3.1) in Type IIB compactification on a Calabi-Yau four-fold. Recall that due to the global anomaly (3.1), in Type IIA vacuum we had to include $`N=\frac{\chi }{24}`$ fundamental strings filling two-dimensional space-time to cancel the tadpole. Under a T-duality in one of the space-time directions the winding modes of these strings transform into $`\frac{\chi }{24}`$ momentum modes:
$$P_+P_{}=\frac{\chi }{24}$$
in the Type IIB vacuum corresponding to compactification on the same Calabi-Yau four-fold $`X`$. Here, for the sake of simplicity, we assumed that there are no background fluxes. One can interpret (6.1) as the difference in the zero-point energy of the left-moving and the right-moving Kaluza-Klein modes . In order to see this, we note that a free boson on a circle has vacuum energy $`\frac{1}{24}`$ and a periodic fermion has vacuum energy $`+\frac{1}{24}`$. Therefore, due to $`𝒩=(0,4)`$ supersymmetry, in the right sector bosonic and fermionic contributions cancel each other, i.e. $`P_+=0`$. In the left sector we have $`48+6h^{1,1}+6h^{3,1}2h^{2,1}`$ bosonic modes corresponding to the fields $`\varphi _i`$, $`\overline{\varphi }_{\overline{i}}`$, $`\stackrel{ˇ}{s}_j`$, $`r_j`$, $`t_j`$, $`u_i`$, $`q`$, $`\overline{q}`$, $`\phi `$ and $`l`$ along with $`n_{}=4h^{1,1}`$ fermionic modes, cf. (6.1). Hence, the total vacuum momentum in the left sector is non-zero and is given by the following formula:
$$P_{}=\frac{1}{24}(48+6h^{1,1}+6h^{3,1}6h^{2,1})$$
Using $`P_+=0`$ and the explicit expression (2.1) for the Euler number, one can easily obtain the formula (6.1).
7. Superspace Formulation of $`𝒩=(0,4)`$ Dilaton Supergravity
In this section we construct $`𝒩=(0,4)`$ dilaton supergravity that arise, for example, in Type IIB superstring compactification on Calabi-Yau four-folds. As we demonstrated in the previous section such a theory has a number of distinct features which are absent in the existing superspace formulations of two-dimensional $`𝒩=(0,4)`$ supergravities. Namely, unlike the standard formulations with gauged $`SU(2)`$ $`R`$-symmetry \[28,,29\], new supergravity does not have a gauged symmetry and the supergravity multiplet contains a real dilaton field $`𝒱`$. Below we present a superspace construction of this theory obtained via de-gauging $`SU(2)`$ symmetry.
First let us remind that gauged $`𝒩=(0,4)`$ supergravity is defined in superspace by the following set of constraints :
$$[_{+i},_{+j}\}=0,[_{+i},_{\dot{+}}{}_{}{}^{j}\}=i2\delta _i^j_{\text{ }\text{ }\text{ }\text{ }},$$
$$[_{+i},_{\text{ }\text{ }\text{ }\text{ }}\}=0,[_{+i},_{\text{ }}\}=i[\overline{\mathrm{\Sigma }}^+{}_{i}{}^{}𝒳\overline{\mathrm{\Sigma }}^+{}_{j}{}^{}𝒴_{i}^{}{}_{}{}^{j}],$$
$$[_{\text{ }\text{ }\text{ }\text{ }},_{\text{ }}\}=\frac{1}{2}[\mathrm{\Sigma }^{+i}_{+i}+\overline{\mathrm{\Sigma }}^+{}_{i}{}^{}_{\dot{+}}^{}{}_{}{}^{i}+𝒳+i_i{}_{}{}^{j}𝒴_{j}^{}{}_{}{}^{i}]$$
on the covariant derivatives $`_A(_{+i},_{\dot{+}}{}_{}{}^{i},_{\text{ }\text{ }\text{ }\text{ }},_{\text{ }})`$:
$$_A=E_A{}_{}{}^{B}D_{B}^{}+\mathrm{\Lambda }_A𝒳+i𝒜_{Ai}{}_{}{}^{j}𝒴_{j}^{}{}_{}{}^{i}.$$
Here $`E_{A}^{}{}_{}{}^{B}`$ is the supervielbein, and $`𝒳`$ and $`𝒴_j^i`$ are the Lorentz and $`SU(2)`$ symmetry generators, respectively. The superfield $`𝒜_{Ai}^j`$ is $`SU(2)`$ gauge connection, while $`\mathrm{\Lambda }_A`$ stands for the Lorentz spin-connection. We write $`[,\}`$ for the graded (anti-)commutator. Finally, $`D_A`$ denotes the flat space fermi and bose derivatives $`D_A(\overline{D}_{+i},D_+{}_{}{}^{i},_{\text{ }\text{ }\text{ }\text{ }},_{\text{ }}`$).
The Lorentz generators act on $`_A`$ as follows:
$$[𝒳,_{+i}\}=\frac{1}{2}_{+i},[𝒳,_{\dot{+}}{}_{}{}^{i}\}=\frac{1}{2}_{\dot{+}}{}_{}{}^{i},$$
$$[𝒳,_{\text{ }\text{ }\text{ }\text{ }}\}=_{\text{ }\text{ }\text{ }\text{ }},[𝒳,_{\text{ }}\}=_{\text{ }}.$$
Similarly, for the action of $`SU(2)`$ gauge symmetry generators we have:
$$[𝒴_j{}_{}{}^{k},_{+i}\}=\delta _i^k_{+j}\frac{1}{2}\delta _j^k_{+i},$$
$$[𝒴_j{}_{}{}^{k},_{\dot{+}}{}_{}{}^{i}\}=\delta _j^i_{\dot{+}}{}_{}{}^{k}+\frac{1}{2}\delta _j^k_{\dot{+}}{}_{}{}^{i},$$
$$[𝒴_j{}_{}{}^{k},_{\text{ }\text{ }\text{ }\text{ }}\}=0,[𝒴_j{}_{}{}^{k},_{\text{ }}\}=0.$$
The constraints (7.1) lead to a set of Bianchi identities that are solved if:
$$_{\dot{+}}{}_{}{}^{i}\overline{\mathrm{\Sigma }}_{}^{+j}=0,_{+i}\mathrm{\Sigma }^{+j}=\frac{1}{2}\delta _i^j+i_i{}_{}{}^{j},$$
$$_{+i}=i2_{\text{ }\text{ }\text{ }\text{ }}\overline{\mathrm{\Sigma }}^+{}_{i}{}^{},_{+i}_j{}_{}{}^{k}=2\delta _i^k_{\text{ }\text{ }\text{ }\text{ }}\overline{\mathrm{\Sigma }}^+{}_{j}{}^{}+\delta _j^k_{\text{ }\text{ }\text{ }\text{ }}\overline{\mathrm{\Sigma }}^+{}_{i}{}^{}.$$
The first step in obtaining two-dimensional $`𝒩=(0,4)`$ supergravity theory that does not contain a gauged $`SU(2)`$ is to note that the covariant derivative in (7.1) can be split as:
$$_A=\widehat{}_A+i𝒜_{Ak}{}_{}{}^{l}𝒴_{l}^{}{}_{}{}^{k}.$$
Since the superspace covariant derivative $`\widehat{}_A`$ does not contain the $`SU(2)`$ connection nor the generator, it can not describe two-dimensional $`𝒩=(0,4)`$ supergravity with gauged $`SU(2)`$ symmetry. We next use (7.1) to derive the form of the commutator algebra for the $`\widehat{}_A`$ operators.
A straightforward set of calculations leads to:
$$[\widehat{}_{+i},\widehat{}_{+j}\}=i[𝒜_{+ij}{}_{}{}^{k}+𝒜_{+ji}{}_{}{}^{k}]\widehat{}_{+k},$$
$$[\widehat{}_{+i},\widehat{}_{\dot{+}}{}_{}{}^{j}\}=i2\delta _i^j\widehat{}_{\text{ }\text{ }\text{ }\text{ }}+i\overline{𝒜}{}_{+}{}^{j}{}_{i}{}^{}{}_{}{}^{k}\widehat{}_{+k}^{}+i𝒜_{+ik}{}_{}{}^{j}\widehat{}_{\dot{+}}^{}{}_{}{}^{k},$$
$$[\widehat{}_{+i},\widehat{}_{\text{ }\text{ }\text{ }\text{ }}\}=i𝒜_{\text{ }\text{ }\text{ }\text{ }i}{}_{}{}^{k}\widehat{}_{+k}^{},[\widehat{}_{+i},\widehat{}_{\text{ }}\}=i𝒜_{\text{ }i}{}_{}{}^{k}\widehat{}_{+k}^{}i\mathrm{\Sigma }^+{}_{i}{}^{}𝒳,$$
$$[\widehat{}_{\text{ }\text{ }\text{ }\text{ }},\widehat{}_{\text{ }}\}=\frac{1}{2}[\mathrm{\Sigma }^{+i}\widehat{}_{+i}+\overline{\mathrm{\Sigma }}^+{}_{i}{}^{}\widehat{}_{\dot{+}}^{}{}_{}{}^{i}+𝒳]$$
where the connection superfields now explicitly appear on the right-hand side of the equations. The leading component of $`𝒜_{+ij}^k`$ is a component of the gauged supergravity multiplet that could be eliminated by a gauge transformation in the $`SU(2)`$ gauge-symmetric phase.
At this stage, we have completed half of the de-gauging process. The second half consists of specifying the spinorial $`SU(2)`$ connections in terms of some components of another (matter) multiplet that is consistent with the two-dimensional $`𝒩=(0,4)`$ supergravity theory. For this purpose, we introduce the second of the four distinct $`𝒩=(0,4)`$ scalar multiplets (SM-II) that were discussed in \[29,,30\]. In the case of rigid supersymmetry this multiplet is described by:
$$D_{+i}𝒱=i\lambda ^{}{}_{i}{}^{},\overline{𝒱}=𝒱,\phi ^i{}_{i}{}^{}=0,$$
$$D_{+i}\phi _j{}_{}{}^{k}=2\delta _i^k\lambda ^{}{}_{j}{}^{}\delta _j^k\lambda ^{}{}_{i}{}^{},\overline{\phi }_j{}_{}{}^{i}=\phi _i{}_{}{}^{j},$$
$$D_{\dot{+}}{}_{}{}^{i}\lambda _{}^{}{}_{j}{}^{}=\delta _j^i_{\text{ }\text{ }\text{ }\text{ }}𝒱+i_{\text{ }\text{ }\text{ }\text{ }}\phi _j{}_{}{}^{i},D_{+i}\lambda ^{}{}_{j}{}^{}=0.$$
In the locally supersymmetric theory one has to replace $`D_\alpha `$ by $`_\alpha `$.
The triplet of spin-$`0`$ fields $`\phi _j^i`$ can be “eaten” by the triplet of spin-$`1`$ fields in the minimal $`𝒩=(0,4)`$ supergravity multiplet and thus become their longitudinal components via the usual Goldstone mechanism. This eliminates the local $`SU(2)`$ symmetry. The scalar and spinor field of the matter multiplet become the dilaton and dilatino.
We are thus led to conjecture that the new form of two-dimensional $`𝒩=(0,4)`$ supergravity with a component spectrum given by:
$`e_\mu ^\nu `$ graviton,
$`\psi _\mu ^{+i}`$$`SU(2)`$-doublet, gravitino,
$`A_{\mu i}^j`$ $`SU(2)`$-triplet, vector auxiliary fields,
$`\lambda ^{}_i`$ $`SU(2)`$-doublet, dilatino field and
$`𝒱`$ real dilaton field
may be constructed with (7.1) as its starting point. We note that the chirality of the dilatino is opposite to that of the gravitino.
In order to gain a control over the component field content of the theory, we must impose the following constraints:
$$𝒜_{+ij}{}_{}{}^{k}=[2\lambda ^{}{}_{j}{}^{}\delta _{i}^{k}\lambda ^{}{}_{i}{}^{}\delta _{j}^{k}],$$
$$\overline{𝒜}_+{}_{j}{}^{i}{}_{}{}^{k}=[2\overline{\lambda }^{}{}_{j}{}^{}\delta _{i}^{k}\overline{\lambda }^{}{}_{i}{}^{}\delta _{j}^{k}].$$
At lowest order in $`\theta ^\alpha `$, the field $`𝒜_{+ij}^k`$ is a component field that is absent in the $`SU(2)`$ gauge-symmetric phase of the theory; it can be set to zero in the Wess-Zumino gauge. On the other hand, when the $`SU(2)`$ symmetry is broken, part of this field becomes dynamical. In general, $`𝒜_{+ij}^k`$ contains $`SU(2)`$ representations of spin-$`\frac{1}{2}`$ and spin-$`\frac{3}{2}`$. However, the above constraints eliminate the pure spin-$`\frac{3}{2}`$ representation of $`SU(2)`$.
With this result substituted into (7.1),
$$[\widehat{}_{+i},\widehat{}_{+j}\}=i[\lambda ^{}{}_{i}{}^{}\widehat{}_{+j}^{}+\lambda ^{}{}_{j}{}^{}\widehat{}_{+i}^{}],$$
$$[\widehat{}_{+i},\widehat{}_{\dot{+}}{}_{}{}^{j}\}=i2\delta _i^j\widehat{}_{\text{ }\text{ }\text{ }\text{ }}+i[2\overline{\lambda }^k\delta _i^j\overline{\lambda }^j\delta _i^k]\widehat{}_{+k}+i[2\lambda ^{}{}_{k}{}^{}\delta _{i}^{j}\lambda ^{}{}_{i}{}^{}\delta _{k}^{j}]\widehat{}_{\dot{+}}{}_{}{}^{k},$$
$$[\widehat{}_{+i},\widehat{}_{\text{ }\text{ }\text{ }\text{ }}\}=i𝒜_{\text{ }\text{ }\text{ }\text{ }i}{}_{}{}^{k}\widehat{}_{+k}^{},[\widehat{}_{+i},\widehat{}_{\text{ }}\}=i𝒜_{\text{ }i}{}_{}{}^{k}\widehat{}_{+k}^{}i\mathrm{\Sigma }^+{}_{i}{}^{}𝒳,$$
$$[\widehat{}_{\text{ }\text{ }\text{ }\text{ }},\widehat{}_{\text{ }}\}=\frac{1}{2}[\mathrm{\Sigma }^{+i}\widehat{}_{+i}+\overline{\mathrm{\Sigma }}^+{}_{i}{}^{}\widehat{}_{\dot{+}}^{}{}_{}{}^{i}+𝒳]$$
we can calculate the Bianchi identities:
$$[[\widehat{}_{+i},\widehat{}_{+j}\},\widehat{}_{+k}\}+[[\widehat{}_{+k},\widehat{}_{+i}\},\widehat{}_{+j}\}+[[\widehat{}_{+j},\widehat{}_{+k}\},\widehat{}_{+i}\}=0$$
$$[[\widehat{}_{+i},\widehat{}_{+j}\},\widehat{}_{\dot{+}}{}_{}{}^{k}\}+[[\widehat{}_{+i},\widehat{}_{\dot{+}}{}_{}{}^{k}\},\widehat{}_{+j}\}+[[\widehat{}_{+j},\widehat{}_{\dot{+}}{}_{}{}^{k}\},\widehat{}_{+i}\}=0$$
These will be satisfied if:
$$\widehat{}_{+i}\lambda ^{}{}_{j}{}^{}=\lambda ^{}{}_{i}{}^{}\lambda _{}^{}{}_{j}{}^{},$$
$$\widehat{}_{+i}\overline{\lambda }^j=i\lambda ^{}{}_{i}{}^{}\overline{\lambda }_{}^{j}i𝒜_{\text{ }\text{ }\text{ }\text{ }i}{}_{}{}^{j}+\delta _j^i\widehat{}_{\text{ }\text{ }\text{ }\text{ }}𝒱,etc.$$
Let us now briefly comment on density projectors in the new $`𝒩=(0,4)`$ dilaton supergravity theory. For a general superspace Lagrangian $``$, the component action of $`𝒩=(0,4)`$ gauged supergravity can be obtained by means of the following projection formula:
$$d^2xd^2\theta ^{\text{ }\text{ }\text{ }\text{ }}d^2\overline{\theta }^{\text{ }\text{ }\text{ }\text{ }}E^1=\frac{1}{2}d^2xd^2\theta ^{\text{ }\text{ }\text{ }\text{ }}^1[\frac{1}{2}C_{ij}\widehat{}_{\dot{+}}{}_{}{}^{i}\widehat{}_{\dot{+}}^{}{}_{}{}^{j}]|+$$
$$+\frac{1}{2}d^2xd^2\overline{\theta }^{\text{ }\text{ }\text{ }\text{ }}\overline{}^1\left[\frac{1}{2}C^{ij}\widehat{}_{+i}\widehat{}_{+j}\right]|$$
where the corresponding chiral and anti-chiral density projector formulas look like:
$$d^2xd^2\theta ^{\text{ }\text{ }\text{ }\text{ }}^1|=id^2x\left[\frac{1}{2}e^1C^{ij}(\widehat{}_{+i}+i4e\overline{\psi }_{\text{ }\text{ }\text{ }\text{ }}{}_{}{}^{+}{}_{i}{}^{})\right]\widehat{}_{+j}|$$
$$d^2xd^2\overline{\theta }^{\text{ }\text{ }\text{ }\text{ }}\overline{}^1|=id^2x\left[\frac{1}{2}e^1C_{ij}(\widehat{}_{\dot{+}}{}_{}{}^{i}+i4e\psi _{\text{ }\text{ }\text{ }\text{ }}{}_{}{}^{+i})\right]\widehat{}_{\dot{+}}{}_{}{}^{j}|$$
Similar formulas also hold in the new two-dimensional $`𝒩=(0,4)`$ dilaton supergravity. The explicit expressions for the density projectors can be obtained by a straightforward but tedious computation substituting (7.1) into (7.1).
‘‘I found a way to make it work."
Stanislaw Ulam
Acknowledgments
We are grateful to Marc Grisaru, Martin Ro$`\stackrel{ˇ}{\mathrm{c}}`$ek and John H. Schwarz for useful discussions. The research of S.J.G. is supported by the NSF grant No PHY-98-02551; S.G. is supported in part by the Caltech Discovery Fund, grant RFBR No 98-02-16575 and Russian President’s grant No 96-15-96939. The work of E.W. is supported in part by NSF Grant PHY-9513835 and the Caltech Discovery Fund.
Appendix A. World-Sheet Calculation of Type IIA String Amplitudes
Consider compactification of Type IIA string theory on a Calabi-Yau four-fold $`X`$. Let us further assume that there are no background Ramond-Ramond fluxes, so that two-dimensional space-time is flat. From the world-sheet viewpoint, this compactification corresponds to adjoining $`c=(12,12)`$ $`𝒩=(2,2)`$ superconformal field theory (SCFT) to free conformal theory with central charge $`c=(3,3)`$ that is responsible for the two-dimensional space-time. It is the first part that will be interesting to us. Namely, we are going to show that two-point correlation function of vertex operators corresponding to chiral and twisted chiral superfields is zero, i.e. that the Zamolodchikov metric on the moduli space of Calabi-Yau four-folds is block diagonal (3.1).
$`𝒩=(2,2)`$ superconformal algebra consists of two $`𝒩=2`$ supervirasoro algebras — one left-moving and one right-moving — each generated by an energy-momentum tensor $`T`$, a current $`J`$ and two weight $`3/2`$ supercurrents $`G^\pm `$ with $`J`$-charge $`Q=\pm 1`$. Recall that in a Kaluza-Klein reduction two-dimensional chiral superfields come from harmonic $`(3,1)`$-forms on $`X`$, while twisted chiral superfields correspond to harmonic $`(1,1)`$-forms. Similar to the three-fold case, we identify these fields with marginal operators in $`(c,c)`$ and $`(a,c)`$ multiplets, respectively. Let us call this operators $`\mathrm{\Phi }_{(1,1)}`$ and $`\mathrm{\Phi }_{(1,1)}`$. They are neutral, $`(Q,\overline{Q})=(0,0)`$, and have conformal weight $`1/2`$. The lowest components of (anti-)chiral multiplets must satisfy $`2h=Q`$ and $`2\overline{h}=\overline{Q}`$, so $`\mathrm{\Phi }_{(1,1)}`$ and $`\mathrm{\Phi }_{(1,1)}`$ are not the lowest components in the corresponding multiplets. They can be obtained in the operator product expansion of the supercurrents with operators $`\mathrm{\Psi }`$:
$$2G^{}(w,\overline{w})\mathrm{\Psi }_{(\pm 1,1)}(z,\overline{z})=\frac{1}{wz}\mathrm{\Phi }_{(\pm 1,1)}(z,\overline{z})+\mathrm{reg}$$
where “reg” stands for the regular part. Another operator product expansion that will be useful to us is the following:
$$2G^{}(w,\overline{w})\mathrm{\Phi }_{(\pm 1,1)}(z,\overline{z})=\mathrm{reg}$$
Now we are ready to demonstrate (3.1). Consider a matrix element of the target space metric $`𝒢_{\varphi _i\overline{\sigma }_j}`$ that mixes $`(1,1)`$ and $`(3,1)`$ moduli:
$$\frac{𝒢_{\varphi _i\overline{\sigma }_j}}{|zz^{}|^4}=\mathrm{\Phi }_{(1,1)}^{\varphi _i}(z,\overline{z})\mathrm{\Phi }_{(1,1)}^{\overline{\sigma }_j}(z^{},\overline{z}^{})=$$
$$=\frac{dw}{2\pi i}2G^{}(w,\overline{w})\mathrm{\Psi }_{(1,1)}^{\varphi _i}(z,\overline{z})\mathrm{\Phi }_{(1,1)}^{\overline{\sigma }_j}(z^{},\overline{z}^{})=0$$
In the last equality we used the fact (A.1) that the operator product of $`G^{}(w,\overline{w})`$ and $`\mathrm{\Phi }_{(1,1)}^{\overline{\sigma }_j}(z^{},\overline{z}^{})`$ has no singularity as $`wz^{}`$.
One might think that $`𝒢_{\varphi _i\sigma _j}`$ would be non-zero since the corresponding OPE has a singular part:
$$2G^\pm (w,\overline{w})\mathrm{\Phi }_{(\pm 1,1)}(z,\overline{z})=\frac{}{z}\left(\frac{\mathrm{\Psi }_{(\pm 1,1)}(z,\overline{z})}{wz}\right)+\mathrm{reg}$$
However, repeating the above arguments in the right sector one can easily see that $`𝒢_{\varphi _i\sigma _j}`$ is also zero. In fact, OPE is singular in both left and right sectors only for $`\mathrm{\Phi }_{(1,1)}^{\varphi _i}\overline{\mathrm{\Phi }}_{(1,1)}^{\overline{\varphi }_{\overline{j}}}`$ and $`\mathrm{\Phi }_{(1,1)}^{\sigma _i}\overline{\mathrm{\Phi }}_{(1,1)}^{\overline{\sigma }_j}`$ which correspond to the metric $`𝒢_{\varphi _i\overline{\varphi }_{\overline{j}}}`$ for chiral multiplets and the metric $`𝒢_{\sigma _i\overline{\sigma }_j}`$ for the twisted chiral multiplets, respectively. Therefore, we conclude that the target space is locally a product of the manifold $`_c(X)`$ parametrized by the chiral fields $`\varphi _i`$ and the manifold $`_𝒦(X)`$ spanned by the twisted chiral fields $`\sigma _j`$.
Appendix B. Extra derivative constraints arising from de-gauging $`𝒩=(2,2)`$ non-minimal supergravity
In this appendix we collect some more technical formulas that arise in the construction of $`𝒩=(2,2)`$ dilaton supergravity via de-gauging $`U(1)_AU(1)_V`$ non-minimal supergravity. Consistency of the de-gauging procedure requires that all terms in (4.1) with gauge symmetry generators $`𝒴`$ and $`𝒴^{}`$ vanish, so that the commutator algebra of the new covariant derivative $`\widehat{}_\alpha `$ has the form (4.1), e.g.:
$$\{\widehat{}_+,\widehat{}_+\}=\{_+,_+\}2\{_+,\lambda _+𝒴+\stackrel{~}{\lambda }_+𝒴^{}\}=$$
$$=2(_+\lambda _+)𝒴2(_+\stackrel{~}{\lambda }_+)𝒴^{}+i\lambda _+_++i\stackrel{~}{\lambda }_+_+=$$
$$=i(\lambda _++\stackrel{~}{\lambda }_+)(\widehat{}_++\lambda _+𝒴+\stackrel{~}{\lambda }_+𝒴^{})2(_+\lambda _+)𝒴2(_+\stackrel{~}{\lambda }_+)𝒴^{}=$$
$$=i(\lambda _++\stackrel{~}{\lambda }_+)\widehat{}_++[2(_+\lambda _+)+i\stackrel{~}{\lambda }_+\lambda _+]𝒴+[2(_+\stackrel{~}{\lambda }_+)+i\lambda _+\stackrel{~}{\lambda }_+]𝒴^{}=$$
$$=i(\lambda _++\stackrel{~}{\lambda }_+)\widehat{}_++[2(\widehat{}_+\lambda _+)+i\stackrel{~}{\lambda }_+\lambda _++i\stackrel{~}{\lambda }_+\lambda _+]𝒴+[2(\widehat{}_+\stackrel{~}{\lambda }_+)+i\lambda _+\stackrel{~}{\lambda }_++i\lambda _+\stackrel{~}{\lambda }_+]𝒴^{}=$$
$$=i(\lambda _++\stackrel{~}{\lambda }_+)\widehat{}_++[2(\widehat{}_+\lambda _+)+2i\stackrel{~}{\lambda }_+\lambda _+]𝒴+[2(\widehat{}_+\stackrel{~}{\lambda }_+)+2i\lambda _+\stackrel{~}{\lambda }_+]𝒴^{}$$
Therefore, we obtain the following commutation relation:
$$\{\widehat{}_+,\widehat{}_+\}=i(\lambda _++\stackrel{~}{\lambda }_+)\widehat{}_+$$
plus two constraints (Jacobi identities):
$$(\widehat{}_+\lambda _+)i\stackrel{~}{\lambda }_+\lambda _+=0,(\widehat{}_+\stackrel{~}{\lambda }_+)i\lambda _+\stackrel{~}{\lambda }_+=0.$$
Similar calculations lead to the following set of constraints:
$$\widehat{}_+\lambda _+i\stackrel{~}{\lambda }_+\lambda _+=0,,\widehat{}_+\stackrel{~}{\lambda }_+i\lambda _+\stackrel{~}{\lambda }_+=0,$$
$$\widehat{}_{}\lambda _{}+i\stackrel{~}{\lambda }_{}\lambda _{}=0,,\widehat{}_{}\stackrel{~}{\lambda }_{}i\lambda _{}\stackrel{~}{\lambda }_{}=0,$$
$$\widehat{}_+\lambda _{}+\widehat{}_{}\lambda _++i(\stackrel{~}{\lambda }_+\lambda _{}\stackrel{~}{\lambda }_{}\lambda _+)=0,$$
$$\widehat{}_+\stackrel{~}{\lambda }_\dot{}+\widehat{}_\dot{}\stackrel{~}{\lambda }_++i(\lambda _+\stackrel{~}{\lambda }_\dot{}\lambda _\dot{}\stackrel{~}{\lambda }_+)=0,$$
$$\widehat{}_+\stackrel{~}{\lambda }_{}+\widehat{}_{}\stackrel{~}{\lambda }_++2i\stackrel{~}{\lambda }_+\stackrel{~}{\lambda }_{}i(\lambda _+\stackrel{~}{\lambda }_{}+\lambda _{}\stackrel{~}{\lambda }_+)\frac{i}{2}\overline{R}=0,$$
$$\widehat{}_+\lambda _\dot{}+\widehat{}_\dot{}\lambda _++2i\lambda _+\lambda _\dot{}i(\stackrel{~}{\lambda }_+\lambda _\dot{}+\stackrel{~}{\lambda }_\dot{}\lambda _+)\frac{i}{2}\overline{F}=0,$$
$$\stackrel{~}{\lambda }_{\text{ }\text{ }\text{ }\text{ }}=i(\widehat{}_+\stackrel{~}{\lambda }_{\dot{+}}+\widehat{}_{\dot{+}}\stackrel{~}{\lambda }_+)+2\stackrel{~}{\lambda }_+\stackrel{~}{\lambda }_{\dot{+}}+(\lambda _+\stackrel{~}{\lambda }_{\dot{+}}+\stackrel{~}{\lambda }_+\lambda _{\dot{+}}),$$
$$\lambda _{\text{ }\text{ }\text{ }\text{ }}=i(\widehat{}_+\lambda _{\dot{+}}+\widehat{}_{\dot{+}}\lambda _+)+2\lambda _+\lambda _{\dot{+}}+(\lambda _+\stackrel{~}{\lambda }_{\dot{+}}+\stackrel{~}{\lambda }_+\lambda _{\dot{+}}),$$
$$\stackrel{~}{\lambda }_{\text{ }}=i(\widehat{}_{}\stackrel{~}{\lambda }_\dot{}+\widehat{}_\dot{}\stackrel{~}{\lambda }_{})2\stackrel{~}{\lambda }_{}\stackrel{~}{\lambda }_\dot{}+(\lambda _{}\stackrel{~}{\lambda }_\dot{}+\stackrel{~}{\lambda }_{}\lambda _\dot{}),$$
$$\lambda _{\text{ }}=i(\widehat{}_{}\lambda _\dot{}+\widehat{}_\dot{}\lambda _{})+2\lambda _{}\lambda _\dot{}(\lambda _{}\stackrel{~}{\lambda }_\dot{}+\stackrel{~}{\lambda }_{}\lambda _\dot{})$$
By virtue of the above relations, the Jacobi identities (4.1) are automatically satisfied. Furthermore, only half of the spinor fields are independent. In section 4 we disscuss two natural solutions to these constraints: when either $`\eta _\alpha `$ or $`\stackrel{~}{\eta }_\alpha `$ are put to zero. In both cases the remaining spinor superfields can be expressed in terms of an unconstrained real superfield.
Appendix C. Components of covariant derivatives in $`𝒩=(2,2)`$ dilaton supergravity
The expressions for the covariant derivatives in $`𝒩=(2,2)`$ dilaton supergravity evaluated at $`\theta =0`$ are, cf. :
$$\widehat{}_\alpha |=_\alpha $$
$$\widehat{}_\mu |=𝐃_\mu +\psi _\mu ^\alpha \widehat{}_\alpha |+\psi _\mu ^{\dot{\alpha }}\widehat{}_{\dot{\alpha }}|=$$
$$=𝐃_\mu +\psi _\mu ^\alpha _\alpha +\psi _\mu ^{\dot{\alpha }}_{\dot{\alpha }}$$
where $`𝐃_\mu `$ is the fully covariant gravitational derivative with the Lorentz connection $`\omega _\mu =\mathrm{\Lambda }_\mu |`$ that includes, in addition to the ordinary connection, extra terms that are bilinear in the gravitini $`\psi _\mu ^\alpha `$, $`\psi _\mu ^{\dot{\alpha }}`$. Specifically, $`𝐃_\mu `$ is defined to be <sup>11</sup> In the notations of this would correspond to $`\omega _\mu =\gamma _\mu =\phi _\mu `$.:
$$𝐃_\mu =e_\mu +\omega _\mu 𝒳$$
We also need expressions for the higher $`\theta `$ components of the covariant derivatives. The $`\theta ^\alpha `$ component of $`\widehat{}_\beta `$ is defined by :
$$\widehat{}_\alpha \widehat{}_\beta |=\frac{1}{2}\{\widehat{}_\alpha ,\widehat{}_\beta \}|$$
while the $`\theta ^\alpha `$ component of $`\widehat{}_\mu `$ is:
$$\widehat{}_\alpha \widehat{}_\mu |=[\widehat{}_\alpha ,\widehat{}_\mu ]|+\widehat{}_\mu \widehat{}_\alpha |=$$
$$=[\widehat{}_\alpha ,\widehat{}_\mu ]|+𝐃_\mu \widehat{}_\alpha |+\psi _\mu ^\beta \widehat{}_\beta \widehat{}_\alpha |+\psi _\mu ^{\dot{\beta }}\widehat{}_{\dot{\beta }}\widehat{}_\alpha |$$
From (C.1) and (4.1) we obtain the following results:
$$\widehat{}_+\widehat{}_+|=i\lambda _+_+,\widehat{}_{}\widehat{}_{}|=i\lambda _{}_{}$$
$$\widehat{}_{\dot{+}}\widehat{}_{\dot{+}}|=i\lambda _{\dot{+}}_{\dot{+}},\widehat{}_\dot{}\widehat{}_\dot{}|=i\lambda _\dot{}_\dot{}$$
$$\widehat{}_+\widehat{}_{}|=\frac{1}{4}\overline{H}𝒳,\widehat{}_+\widehat{}_\dot{}|=\frac{1}{4}\overline{G}𝒳$$
$$\widehat{}_{\dot{+}}\widehat{}_\dot{}|=\frac{1}{4}H𝒳,\widehat{}_{\dot{+}}\widehat{}_{}|=\frac{1}{4}G𝒳$$
$$\widehat{}_+\widehat{}_{\dot{+}}|=\frac{i}{2}𝐃_{\text{ }\text{ }\text{ }\text{ }}+\frac{i}{2}(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})_++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}_{}+\frac{i}{2}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)_{\dot{+}}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}_\dot{}$$
$$\widehat{}_{}\widehat{}_\dot{}|=\frac{i}{2}𝐃_{\text{ }}+\frac{i}{2}\psi _{\text{ }}^+_++\frac{i}{2}(\psi _{\text{ }}^{}+\lambda _\dot{})_{}+\frac{i}{2}\psi _{\text{ }}^{\dot{+}}_{\dot{+}}+\frac{i}{2}(\psi _{\text{ }}^\dot{}\lambda _{})_\dot{}$$
and from (C.1) we derive the series of identities that appears below:
$$\widehat{}_+\widehat{}_{\text{ }\text{ }\text{ }\text{ }}|=𝐃_{\text{ }\text{ }\text{ }\text{ }}\widehat{}_+|+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}𝐃_{\text{ }\text{ }\text{ }\text{ }}|\frac{1}{4}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\overline{H}+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\overline{G})𝒳+[i\psi _{\text{ }\text{ }\text{ }\text{ }}^+\lambda _++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})$$
$$2i\lambda _+\lambda _{\dot{+}}(\widehat{}_{\dot{+}}\lambda _+)(\widehat{}_+\lambda _{\dot{+}})]_++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}_{}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)_{\dot{+}}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}_\dot{}$$
$$\widehat{}_{}\widehat{}_{\text{ }\text{ }\text{ }\text{ }}|=𝐃_{\text{ }\text{ }\text{ }\text{ }}\widehat{}_{}|+\frac{i}{2}\psi _{\text{ }}^\dot{}𝐃_{\text{ }}|[\frac{i}{2}(\widehat{}_{\dot{+}}\overline{R})+\frac{i}{2}(\widehat{}_+F)+\frac{1}{2}(\lambda _{\dot{+}}+\frac{1}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+)H\frac{1}{2}(\lambda _+\frac{1}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}})G]𝒳+$$
$$+\frac{i}{2}(G+\psi _{\text{ }}^\dot{}\psi _{\text{ }}^+)_++\frac{i}{2}(2\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\lambda _{}+\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^{}+\lambda _\dot{}))_{}+\frac{i}{2}(\overline{H}+\psi _{\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}})_{\dot{+}}+\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^\dot{}\lambda _{})_\dot{}$$
$$\widehat{}_{\dot{+}}\widehat{}_{\text{ }\text{ }\text{ }\text{ }}|=𝐃_{\text{ }\text{ }\text{ }\text{ }}\widehat{}_{\dot{+}}|+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}𝐃_{\text{ }\text{ }\text{ }\text{ }}|\frac{1}{4}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}G+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}H)𝒳+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})_++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{}_{}+$$
$$+\left[2i\lambda _+\lambda _{\dot{+}}+(\widehat{}_{\dot{+}}\lambda _+)+(\widehat{}_+\lambda _{\dot{+}})+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _{\dot{+}}\right]_{\dot{+}}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}_\dot{}$$
$$\widehat{}_\dot{}\widehat{}_{\text{ }\text{ }\text{ }\text{ }}|=𝐃_{\text{ }\text{ }\text{ }\text{ }}\widehat{}_\dot{}|+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}𝐃_{\text{ }}|+[\frac{i}{2}(\widehat{}_+R)+\frac{i}{2}(\widehat{}_{\dot{+}}\overline{F})\frac{1}{2}\lambda _+H+\frac{1}{2}\lambda _{\dot{+}}\overline{G}\frac{1}{4}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\overline{G}\frac{1}{4}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\overline{H}]𝒳+$$
$$+\frac{i}{2}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^+H)_++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}(\psi _{\text{ }}^{}+\lambda _\dot{})_{}+\frac{i}{2}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^{\dot{+}}\overline{G})_{\dot{+}}+\frac{i}{2}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\lambda _{}2\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\lambda _\dot{})_\dot{}$$
$$\widehat{}_+\widehat{}_{\text{ }}|=𝐃_{\text{ }}\widehat{}_+|+\frac{i}{2}\psi _{\text{ }}^\dot{}𝐃_{\text{ }}|[\frac{i}{2}(\widehat{}_\dot{}\overline{R})+\frac{i}{2}(\widehat{}_{}\overline{F})+\frac{1}{2}\lambda _\dot{}\overline{H}\frac{1}{2}\lambda _{}\overline{G}+\frac{1}{4}\psi _{\text{ }}^+\overline{H}+\frac{1}{4}\psi _{\text{ }}^{\dot{+}}G]𝒳+$$
$$+\frac{i}{2}\psi _{\text{ }}^\dot{}\psi _{\text{ }}^+_++\frac{i}{2}(2\psi _{\text{ }}^{}\lambda _{}\overline{G}+\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^{}+\lambda _\dot{}))_{}+\frac{i}{2}\psi _{\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}_{\dot{+}}+\frac{i}{2}(\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^\dot{}\lambda _{})\overline{H})_\dot{}$$
$$\widehat{}_{}\widehat{}_{\text{ }}|=𝐃_{\text{ }}\widehat{}_{}|+\frac{i}{2}\psi _{\text{ }}^\dot{}𝐃_{\text{ }}|\frac{1}{4}(\psi _{\text{ }}^+\overline{H}+\psi _{\text{ }}^{\dot{+}}G)𝒳+\frac{i}{2}\psi _{\text{ }}^\dot{}\psi _{\text{ }}^+_++[i\psi _{\text{ }}^{}\lambda _{}+$$
$$+\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^{}+\lambda _\dot{})2i\lambda _{}\lambda _\dot{}(\widehat{}_\dot{}\lambda _{})(\widehat{}_{}\lambda _\dot{})]_{}+\frac{i}{2}\psi _{\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}_{\dot{+}}+\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^\dot{}\lambda _{})_\dot{}$$
$$\widehat{}_{\dot{+}}\widehat{}_{\text{ }}|=𝐃_{\text{ }}\widehat{}_{\dot{+}}|+\frac{i}{2}\psi _{\text{ }}^+𝐃_{\text{ }\text{ }\text{ }\text{ }}|+[\frac{i}{2}(\widehat{}_{}R)+\frac{i}{2}(\widehat{}_\dot{}F)\frac{1}{2}\lambda _{}H+\frac{1}{2}\lambda _\dot{}G\frac{1}{4}\psi _{\text{ }}^{}G\frac{1}{4}\psi _{\text{ }}^\dot{}H]𝒳+$$
$$+\frac{i}{2}\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})_++\frac{i}{2}(H+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{})_{}+\frac{i}{2}(\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)2\psi _{\text{ }}^{\dot{+}}\lambda _{\dot{+}})_{\dot{+}}+\frac{i}{2}(G+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{})_\dot{}$$
$$\widehat{}_\dot{}\widehat{}_{\text{ }}|=𝐃_{\text{ }}\widehat{}_\dot{}|+\frac{i}{2}\psi _{\text{ }}^{}𝐃_{\text{ }}|\frac{1}{4}(\psi _{\text{ }}^+\overline{G}+\psi _{\text{ }}^{\dot{+}}H)𝒳+\frac{i}{2}\psi _{\text{ }}^{}\psi _{\text{ }}^+_++\frac{i}{2}\psi _{\text{ }}^{}(\psi _{\text{ }}^{}+\lambda _\dot{})_{}+$$
$$+\frac{i}{2}\psi _{\text{ }}^{}\psi _{\text{ }}^{\dot{+}}_{\dot{+}}++[2i\lambda _{}\lambda _\dot{}+(\widehat{}_\dot{}\lambda _{})+(\widehat{}_{}\lambda _\dot{})+\frac{i}{2}\psi _{\text{ }}^{}(\psi _{\text{ }}^\dot{}\lambda _{})i\psi _{\text{ }}^\dot{}\lambda _\dot{}]_\dot{}$$
where we also used the commutation relations (4.1).
Furthermore, the $`\widehat{}_{\text{ }\text{ }\text{ }\text{ }}`$ component of $`\widehat{}_{\text{ }}`$ is given by:
$$\widehat{}_{\text{ }\text{ }\text{ }\text{ }}\widehat{}_{\text{ }}|=\widehat{}_{\text{ }\text{ }\text{ }\text{ }}|\widehat{}_{\text{ }}|+\psi _{\text{ }\text{ }\text{ }\text{ }}^\alpha \widehat{}_\alpha \widehat{}_{\text{ }}|+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{\alpha }}\widehat{}_{\dot{\alpha }}\widehat{}_{\text{ }}|$$
so that:
$$[\widehat{}_{\text{ }\text{ }\text{ }\text{ }},\widehat{}_{\text{ }}]|=[\widehat{}_{\text{ }\text{ }\text{ }\text{ }}|,\widehat{}_{\text{ }}|]+\psi _{\text{ }\text{ }\text{ }\text{ }}^\alpha \widehat{}_\alpha \widehat{}_{\text{ }}|+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{\alpha }}\widehat{}_{\dot{\alpha }}\widehat{}_{\text{ }}|\psi _{\text{ }}^\alpha \widehat{}_\alpha \widehat{}_{\text{ }\text{ }\text{ }\text{ }}|\psi _{\text{ }}^{\dot{\alpha }}\widehat{}_{\dot{\alpha }}\widehat{}_{\text{ }\text{ }\text{ }\text{ }}|=$$
$$=[𝐃_{\text{ }\text{ }\text{ }\text{ }},𝐃_{\text{ }}]+𝐃_{[\text{ }\text{ }\text{ }\text{ }}(\psi _{\text{ }]}^\alpha _\alpha )+𝐃_{[\text{ }\text{ }\text{ }\text{ }}(\psi _{\text{ }]}^{\dot{\alpha }}_{\dot{\alpha }})+\psi _{\text{ }\text{ }\text{ }\text{ }}^\alpha \widehat{}_\alpha \widehat{}_{\text{ }}|+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{\alpha }}\widehat{}_{\dot{\alpha }}\widehat{}_{\text{ }}|\psi _{\text{ }}^\alpha \widehat{}_\alpha \widehat{}_{\text{ }\text{ }\text{ }\text{ }}|\psi _{\text{ }}^{\dot{\alpha }}\widehat{}_{\dot{\alpha }}\widehat{}_{\text{ }\text{ }\text{ }\text{ }}|$$
Substituting the above expressions for the components of $`\widehat{}_\alpha \widehat{}_\mu |`$ into $`[\widehat{}_{\text{ }\text{ }\text{ }\text{ }},\widehat{}_{\text{ }}]|`$, we get:
$$[\widehat{}_{\text{ }\text{ }\text{ }\text{ }},\widehat{}_{\text{ }}]|=[𝐃_{\text{ }\text{ }\text{ }\text{ }},𝐃_{\text{ }}]+𝐃_{[\text{ }\text{ }\text{ }\text{ }}(\psi _{\text{ }]}^\alpha _\alpha )+𝐃_{[\text{ }\text{ }\text{ }\text{ }}(\psi _{\text{ }]}^{\dot{\alpha }}_{\dot{\alpha }})+$$
$$+\psi _{\text{ }\text{ }\text{ }\text{ }}^+(𝐃_{\text{ }}\widehat{}_{\dot{+}}|+\frac{i}{2}\psi _{\text{ }}^+𝐃_{\text{ }\text{ }\text{ }\text{ }}|+[\frac{i}{2}(\widehat{}_{}R)+\frac{i}{2}(\widehat{}_\dot{}F)\frac{1}{2}\lambda _{}H+\frac{1}{2}\lambda _\dot{}G\frac{1}{4}\psi _{\text{ }}^{}G\frac{1}{4}\psi _{\text{ }}^\dot{}H]𝒳+$$
$$+\frac{i}{2}\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})_++\frac{i}{2}(H+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{})_{}+\frac{i}{2}(\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)2\psi _{\text{ }}^{\dot{+}}\lambda _{\dot{+}})_{\dot{+}}+\frac{i}{2}(G+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{})_\dot{})+$$
$$+\psi _{\text{ }\text{ }\text{ }\text{ }}^{}(𝐃_{\text{ }}\widehat{}_{}|+\frac{i}{2}\psi _{\text{ }}^\dot{}𝐃_{\text{ }}|\frac{1}{4}(\psi _{\text{ }}^+\overline{H}+\psi _{\text{ }}^{\dot{+}}G)𝒳+\frac{i}{2}\psi _{\text{ }}^\dot{}\psi _{\text{ }}^+_++[i\psi _{\text{ }}^{}\lambda _{}+$$
$$+\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^{}+\lambda _\dot{})2i\lambda _{}\lambda _\dot{}(\widehat{}_\dot{}\lambda _{})(\widehat{}_{}\lambda _\dot{})]_{}+\frac{i}{2}\psi _{\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}_{\dot{+}}+\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^\dot{}\lambda _{})_\dot{})+$$
$$+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(𝐃_{\text{ }}\widehat{}_{\dot{+}}|+\frac{i}{2}\psi _{\text{ }}^+𝐃_{\text{ }\text{ }\text{ }\text{ }}|+[\frac{i}{2}(\widehat{}_{}R)+\frac{i}{2}(\widehat{}_\dot{}F)\frac{1}{2}\lambda _{}H+\frac{1}{2}\lambda _\dot{}G\frac{1}{4}\psi _{\text{ }}^{}G\frac{1}{4}\psi _{\text{ }}^\dot{}H]𝒳+$$
$$+\frac{i}{2}\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})_++\frac{i}{2}(H+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{})_{}+\frac{i}{2}(\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)2\psi _{\text{ }}^{\dot{+}}\lambda _{\dot{+}})_{\dot{+}}+\frac{i}{2}(G+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{})_\dot{})+$$
$$+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}(𝐃_{\text{ }}\widehat{}_\dot{}|+\frac{i}{2}\psi _{\text{ }}^{}𝐃_{\text{ }}|\frac{1}{4}(\psi _{\text{ }}^+\overline{G}+\psi _{\text{ }}^{\dot{+}}H)𝒳+\frac{i}{2}\psi _{\text{ }}^{}\psi _{\text{ }}^+_++\frac{i}{2}\psi _{\text{ }}^{}(\psi _{\text{ }}^{}+\lambda _\dot{})_{}+$$
$$+\frac{i}{2}\psi _{\text{ }}^{}\psi _{\text{ }}^{\dot{+}}_{\dot{+}}+[2i\lambda _{}\lambda _\dot{}+(\widehat{}_\dot{}\lambda _{})+(\widehat{}_{}\lambda _\dot{})+\frac{i}{2}\psi _{\text{ }}^{}(\psi _{\text{ }}^\dot{}\lambda _{})i\psi _{\text{ }}^\dot{}\lambda _\dot{}]_\dot{})$$
$$\psi _{\text{ }}^+(𝐃_{\text{ }\text{ }\text{ }\text{ }}\widehat{}_+|+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}𝐃_{\text{ }\text{ }\text{ }\text{ }}|\frac{1}{4}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\overline{H}+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\overline{G})𝒳+[i\psi _{\text{ }\text{ }\text{ }\text{ }}^+\lambda _++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})$$
$$2i\lambda _+\lambda _{\dot{+}}(\widehat{}_{\dot{+}}\lambda _+)(\widehat{}_+\lambda _{\dot{+}})]_++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}_{}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)_{\dot{+}}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}_\dot{})$$
$$\psi _{\text{ }}^{}(𝐃_{\text{ }\text{ }\text{ }\text{ }}\widehat{}_{}|+\frac{i}{2}\psi _{\text{ }}^\dot{}𝐃_{\text{ }}|[\frac{i}{2}(\widehat{}_{\dot{+}}\overline{R})+\frac{i}{2}(\widehat{}_+F)+\frac{1}{2}(\lambda _{\dot{+}}+\frac{1}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+)H\frac{1}{2}(\lambda _+\frac{1}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}})G]𝒳+$$
$$+\frac{i}{2}(G+\psi _{\text{ }}^\dot{}\psi _{\text{ }}^+)_++\frac{i}{2}(2\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\lambda _{}+\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^{}+\lambda _\dot{}))_{}+\frac{i}{2}(\overline{H}+\psi _{\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}})_{\dot{+}}+\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^\dot{}\lambda _{})_\dot{})$$
$$\psi _{\text{ }}^{\dot{+}}(𝐃_{\text{ }\text{ }\text{ }\text{ }}\widehat{}_{\dot{+}}|+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}𝐃_{\text{ }\text{ }\text{ }\text{ }}|\frac{1}{4}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}G+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}H)𝒳+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})_++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{}_{}+$$
$$+[2i\lambda _+\lambda _{\dot{+}}+(\widehat{}_{\dot{+}}\lambda _+)+(\widehat{}_+\lambda _{\dot{+}})+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _{\dot{+}}]_{\dot{+}}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}_\dot{})$$
$$\psi _{\text{ }}^\dot{}(𝐃_{\text{ }\text{ }\text{ }\text{ }}\widehat{}_\dot{}|+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}𝐃_{\text{ }}|+[\frac{i}{2}(\widehat{}_+R)+\frac{i}{2}(\widehat{}_{\dot{+}}\overline{F})\frac{1}{2}\lambda _+H+\frac{1}{2}\lambda _{\dot{+}}\overline{G}\frac{1}{4}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\overline{G}\frac{1}{4}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\overline{H}]𝒳+$$
$$+\frac{i}{2}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^+H)_++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}(\psi _{\text{ }}^{}+\lambda _\dot{})_{}+\frac{i}{2}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^{\dot{+}}\overline{G})_{\dot{+}}+\frac{i}{2}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\lambda _{}2\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\lambda _\dot{})_\dot{})$$
Now one can compare this huge formula with the leading component of the commutator $`[\widehat{}_{\text{ }\text{ }\text{ }\text{ }},\widehat{}_{\text{ }}]`$ computed directly from the Bianchi identities (4.1):
$$[\widehat{}_{\text{ }\text{ }\text{ }\text{ }},\widehat{}_{\text{ }}]=\left(\lambda _{\dot{+}}S\lambda _+\overline{S}+i(\widehat{}_{\dot{+}}S)+i(\widehat{}_+\overline{S})\right)𝒳+$$
$$+\left(\frac{i}{2}\lambda _\dot{}\overline{R}\frac{i}{2}\lambda _{}\overline{F}\frac{1}{2}(\widehat{}_\dot{}\overline{R})\frac{1}{2}(\widehat{}_{}\overline{F})\right)\widehat{}_{\dot{+}}+\left(\frac{i}{2}\lambda _{}R\frac{i}{2}\lambda _\dot{}F+\frac{1}{2}(\widehat{}_{}\overline{R})+\frac{1}{2}(\widehat{}_\dot{}F)\right)\widehat{}_++$$
$$+\left(\frac{i}{2}\lambda _{\dot{+}}\overline{R}+\frac{i}{2}\lambda _+F\frac{1}{2}(\widehat{}_{\dot{+}}\overline{R})+\frac{1}{2}(\widehat{}_+F)\right)\widehat{}_\dot{}+\left(\frac{i}{2}\lambda _{\dot{+}}\overline{F}+\frac{i}{2}\lambda _+R\frac{1}{2}(\widehat{}_{\dot{+}}\overline{F})+\frac{1}{2}(\widehat{}_+R)\right)\widehat{}_{}$$
where we denoted:
$$S=\frac{i}{2}(\widehat{}_\dot{}\overline{R})+\frac{i}{2}(\widehat{}_{}\overline{F})+\frac{1}{2}\lambda _\dot{}\overline{R}\frac{1}{2}\lambda _{}\overline{F}$$
By comparing the coefficients of $`\widehat{}_+|=_+`$ in (C.1) and (C.1), for example, we find:
$$\frac{i}{2}\lambda _{}R\frac{i}{2}\lambda _\dot{}F+\frac{1}{2}(\widehat{}_{}\overline{R})+\frac{1}{2}(\widehat{}_\dot{}F)=𝐃_{[\text{ }\text{ }\text{ }\text{ }}\psi _{\text{ }]}^++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})+$$
$$+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^\dot{}\psi _{\text{ }}^++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{}\psi _{\text{ }}^+\psi _{\text{ }}^+[i\psi _{\text{ }\text{ }\text{ }\text{ }}^+\lambda _++\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})2i\lambda _+\lambda _{\dot{+}}$$
$$(\widehat{}_{\dot{+}}\lambda _+)(\widehat{}_+\lambda _{\dot{+}})]\frac{i}{2}\psi _{\text{ }}^{}(G+\psi _{\text{ }}^\dot{}\psi _{\text{ }}^+)\frac{i}{2}\psi _{\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^++\lambda _{\dot{+}})\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^+H)$$
In the same fashion we obtain the relations defining the other components of the superfields $`R`$ and $`F`$:
$$\frac{i}{2}\lambda _{\dot{+}}\overline{F}+\frac{i}{2}\lambda _+R\frac{1}{2}(\widehat{}_{\dot{+}}\overline{F})+\frac{1}{2}(\widehat{}_+R)=𝐃_{[\text{ }\text{ }\text{ }\text{ }}\psi _{\text{ }]}^{}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+(H+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{})+\psi _{\text{ }\text{ }\text{ }\text{ }}^{}[i\psi _{\text{ }}^{}\lambda _{}+$$
$$+\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^{}+\lambda _\dot{})2i\lambda _{}\lambda _\dot{}(\widehat{}_\dot{}\lambda _{})(\widehat{}_{}\lambda _\dot{})]+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(H+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{})+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{}(\psi _{\text{ }}^{}+\lambda _\dot{})$$
$$\frac{i}{2}\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\frac{i}{2}\psi _{\text{ }}^{}(2\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\lambda _{}+\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^{}+\lambda _\dot{}))\frac{i}{2}\psi _{\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\frac{i}{2}\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}(\psi _{\text{ }}^{}+\lambda _\dot{})$$
$$\frac{i}{2}\lambda _\dot{}\overline{R}\frac{i}{2}\lambda _{}\overline{F}\frac{1}{2}(\widehat{}_\dot{}\overline{R})\frac{1}{2}(\widehat{}_{}\overline{F})=$$
$$=𝐃_{[\text{ }\text{ }\text{ }\text{ }}\psi _{\text{ }]}^{\dot{+}}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+(\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)2\psi _{\text{ }}^{\dot{+}}\lambda _{\dot{+}})+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)2\psi _{\text{ }}^{\dot{+}}\lambda _{\dot{+}})+$$
$$+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{}\psi _{\text{ }}^{\dot{+}}\frac{i}{2}\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)\frac{i}{2}\psi _{\text{ }}^{}(\overline{H}+\psi _{\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}})\psi _{\text{ }}^{\dot{+}}[2i\lambda _+\lambda _{\dot{+}}+(\widehat{}_{\dot{+}}\lambda _+)+$$
$$+(\widehat{}_+\lambda _{\dot{+}})+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)i\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _{\dot{+}}]\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^{\dot{+}}\overline{G})$$
$$\frac{i}{2}\lambda _{\dot{+}}\overline{R}+\frac{i}{2}\lambda _+F\frac{1}{2}(\widehat{}_{\dot{+}}\overline{R})+\frac{1}{2}(\widehat{}_+F)=$$
$$=𝐃_{[\text{ }\text{ }\text{ }\text{ }}\psi _{\text{ }]}^\dot{}+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+(G+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{})+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^\dot{}\lambda _{})+\frac{i}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}(G+\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{})+$$
$$+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\left[2i\lambda _{}\lambda _\dot{}+(\widehat{}_\dot{}\lambda _{})+(\widehat{}_{}\lambda _\dot{})+\frac{i}{2}\psi _{\text{ }}^{}(\psi _{\text{ }}^\dot{}\lambda _{})i\psi _{\text{ }}^\dot{}\lambda _\dot{}\right]\frac{i}{2}\psi _{\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}$$
$$\frac{i}{2}\psi _{\text{ }}^{}\psi _{\text{ }}^\dot{}(\psi _{\text{ }}^\dot{}\lambda _{})\frac{i}{2}\psi _{\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\frac{i}{2}\psi _{\text{ }}^\dot{}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\psi _{\text{ }}^\dot{}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\lambda _{}2\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\lambda _\dot{})$$
$$\lambda _{\dot{+}}S\lambda _+\overline{S}+i(\widehat{}_{\dot{+}}S)+i(\widehat{}_+\overline{S})=$$
$$=[𝐃_{\text{ }\text{ }\text{ }\text{ }},𝐃_{\text{ }}]_𝒳+\psi _{\text{ }\text{ }\text{ }\text{ }}^+[\frac{i}{2}(\widehat{}_{}R)+\frac{i}{2}(\widehat{}_\dot{}F)\frac{1}{2}\lambda _{}H+\frac{1}{2}\lambda _\dot{}G\frac{1}{4}\psi _{\text{ }}^{}G\frac{1}{4}\psi _{\text{ }}^\dot{}H]\frac{1}{4}\psi _{\text{ }\text{ }\text{ }\text{ }}^{}(\psi _{\text{ }}^+\overline{H}+$$
$$+\psi _{\text{ }}^{\dot{+}}G)+\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}[\frac{i}{2}(\widehat{}_{}R)+\frac{i}{2}(\widehat{}_\dot{}F)\frac{1}{2}\lambda _{}H+\frac{1}{2}\lambda _\dot{}G\frac{1}{4}\psi _{\text{ }}^{}G\frac{1}{4}\psi _{\text{ }}^\dot{}H]\frac{1}{4}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}(\psi _{\text{ }}^+\overline{G}+\psi _{\text{ }}^{\dot{+}}H)+$$
$$+\frac{1}{4}\psi _{\text{ }}^+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}\overline{H}+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\overline{G})+\psi _{\text{ }}^{}\left[\frac{i}{2}(\widehat{}_{\dot{+}}\overline{R})+\frac{i}{2}(\widehat{}_+F)+\frac{1}{2}(\lambda _{\dot{+}}+\frac{1}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^+)H\frac{1}{2}(\lambda _+\frac{1}{2}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}})G\right]+$$
$$+\frac{1}{4}\psi _{\text{ }}^{\dot{+}}(\psi _{\text{ }\text{ }\text{ }\text{ }}^{}G+\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}H)\psi _{\text{ }}^\dot{}\left[\frac{i}{2}(\widehat{}_+R)+\frac{i}{2}(\widehat{}_{\dot{+}}\overline{F})\frac{1}{2}\lambda _+H+\frac{1}{2}\lambda _{\dot{+}}\overline{G}\frac{1}{4}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\overline{G}\frac{1}{4}\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\overline{H}\right]$$
Appendix D. Derivation of the projection formula in $`𝒩=(2,2)`$ dilaton supergravity
Here we repeat the derivation of the local density projection formula in $`𝒩=(2,2)`$ dilaton supergravity. Namely, we start by writing the most general expression for the chiral projector with the right dimension and index structure:
$$d^2xd^4\theta E^1=d^2xe^1[\widehat{}_+\widehat{}_{}+X^+\widehat{}_++X^{}\widehat{}_{}+Y]\widehat{\overline{}}^2|$$
where the coefficients $`X^\alpha `$ and $`Y`$ are to be determined. Following , we evaluate (D.1) for the kinetic action $`=\overline{\mathrm{\Phi }}\mathrm{\Phi }`$ of a free chiral multiplet:
$$d^2xd^4\theta E^1\overline{\mathrm{\Phi }}\mathrm{\Phi }=d^2xe^1[(\widehat{}^2\widehat{\overline{}}^2\overline{\mathrm{\Phi }})\mathrm{\Phi }|+(\widehat{}_+\widehat{\overline{}}^2\overline{\mathrm{\Phi }})(\widehat{}_{}\mathrm{\Phi })|(\widehat{}_{}\widehat{\overline{}}^2\overline{\mathrm{\Phi }})(\widehat{}_+\mathrm{\Phi })|+$$
$$+(\widehat{\overline{}}^2\overline{\mathrm{\Phi }})(\widehat{}^2\mathrm{\Phi })|+X^+(\widehat{}_+\widehat{\overline{}}^2\overline{\mathrm{\Phi }})\mathrm{\Phi }|+X^+(\widehat{\overline{}}^2\overline{\mathrm{\Phi }})(\widehat{}_+\mathrm{\Phi })|$$
$$+X^{}(\widehat{}_{}\widehat{\overline{}}^2\overline{\mathrm{\Phi }})\mathrm{\Phi }|+X^{}(\widehat{\overline{}}^2\overline{\mathrm{\Phi }})(\widehat{}_{}\mathrm{\Phi })|+Y(\widehat{\overline{}}^2\overline{\mathrm{\Phi }})\mathrm{\Phi }|]$$
Clearly, the left-hand side of this formula is invariant under complex conjugation. So, the right-hand side must be invariant as well. As we will see in a moment, this condition completely determines the unknown coefficients $`X^\alpha `$ and $`Y`$. It suffices to consider only bosonic terms. Using the formulas in appendix C along with the definition (4.1), one can easily compute the relevant components:
$$\widehat{}_+\widehat{}_{\dot{+}}\widehat{}_\dot{}\overline{\mathrm{\Phi }}|\psi _{\text{ }\text{ }\text{ }\text{ }}^{}𝐃_{\text{ }}\overline{\varphi }+(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)\overline{A}$$
$$\widehat{}_{}\widehat{}_{\dot{+}}\widehat{}_\dot{}\overline{\mathrm{\Phi }}|\psi _{\text{ }}^+𝐃_{\text{ }\text{ }\text{ }\text{ }}\overline{\varphi }+(\psi _{\text{ }}^\dot{}\lambda _{})\overline{A}$$
$$\widehat{}_+\widehat{}_{}\widehat{}_{\dot{+}}\widehat{}_\dot{}\overline{\mathrm{\Phi }}|𝐃_{\text{ }}𝐃_{\text{ }\text{ }\text{ }\text{ }}\overline{\varphi }+i(\psi _{\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^+\lambda _+\psi _{\text{ }}^+)𝐃_{\text{ }\text{ }\text{ }\text{ }}\overline{\varphi }+$$
$$+i(\psi _{\text{ }}^\dot{}\lambda _{})\psi _{\text{ }\text{ }\text{ }\text{ }}^{}𝐃_{\text{ }}\overline{\varphi }\frac{i}{4}\overline{H}\overline{A}+i\psi _{\text{ }}^{\dot{+}}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\overline{A}i(\psi _{\text{ }}^\dot{}\lambda _{})(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)\overline{A}$$
Substituting these into (D.1), we find that the resulting component action is invariant under complex conjugation only if we put:
$$X^+=i(\psi _{\text{ }}^\dot{}\lambda _{})$$
$$X^{}=i(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)$$
$$Y=\frac{1}{4}\overline{H}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+(\psi _{\text{ }}^\dot{}\lambda _{})(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)$$
With these expressions for $`X^\alpha `$ and $`Y`$ the result does not depend on whether we use the chiral projector (D.1) or its complex conjugate, of course, as it should be. Specifically, the chiral projection formula in $`𝒩=(2,2)`$ dilaton supergravity takes the following form:
$$d^2xd^4\theta E^1=d^2xe^1[\widehat{}_+\widehat{}_{}+i(\psi _{\text{ }}^\dot{}\lambda _{})\widehat{}_+i(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+)\widehat{}_{}+$$
$$+(\frac{1}{4}\overline{H}\psi _{\text{ }\text{ }\text{ }\text{ }}^\dot{}\psi _{\text{ }}^{\dot{+}}+(\psi _{\text{ }}^\dot{}\lambda _{})(\psi _{\text{ }\text{ }\text{ }\text{ }}^{\dot{+}}\lambda _+))]\widehat{\overline{}}^2|$$
One can easily check that this expression is equivalent to the chiral density projector (5.1) obtained by de-gauging the corresponding projector in the non-minimal $`𝒩=(2,2)`$ supergravity.
References
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# Untitled Document
QUIVER VARIETIES AND YANGIANS
MICHELA VARAGNOLO <sup>1</sup><sup>1</sup>1 The author is partially supported by EEC grant no. ERB FMRX-CT97-0100.
Département de Mathématique, Université de Cergy-Pontoise, 2 av. A. Chauvin, 95302 Cergy-Pontoise cedex. e-mail: michela.varagnolomath.u-cergy.fr
Mathematics Subject Classification (1991) : Primary 17B37, 14L30, 19D55.
Key words : yangian, equivariant homology, convolution product.
Abstract. We prove a conjecture of Nakajima (, for type $`A`$ it was announced in ) giving a geometric realization, via quiver varieties, of the Yangian of type $`ADE`$ (and more in general of the Yangian associated to every symmetric Kac-Moody Lie algebra). As a corollary we get that the finite dimensional representation theory of the quantized affine algebra and that of the Yangian coincide.
1. The algebra $`𝕐_{\mathrm{}}(\mathrm{L}𝔤)`$.
Let $`𝔤`$ be a simple, simply laced, complex Lie algebra over $``$ with Cartan matrix $`A=(a_{kl})_{k,lI}.`$ Denote by $`L𝔤=𝔤[t,t^1]`$ the loop Lie algebra of $`𝔤`$. The Yangian $`𝕐_{\mathrm{}}(\mathrm{L}𝔤)`$ is the associative algebra, free over $`[\mathrm{}],`$ generated by $`𝕩_{k,r}^\pm ,𝕙_{k,r}`$ $`(kI,r)`$ with the following defining relations
$$[𝕙_{k,r},𝕙_{l,s}]=0,[𝕙_{k,0},𝕩_{l,s}^\pm ]=\pm a_{kl}𝕩_{l,s}^\pm ,$$
$`(1.1)`$
$$2[𝕙_{k,r+1},𝕩_{l,s}^\pm ]2[𝕙_{k,r},𝕩_{l,s+1}^\pm ]=\pm \mathrm{}a_{kl}(𝕙_{k,r}𝕩_{l,s}^\pm +𝕩_{l,s}^\pm 𝕙_{k,r}),$$
$`(1.2)`$
$$[𝕩_{k,r}^+,𝕩_{l,s}^{}]=\delta _{kl}𝕙_{k,r+s},$$
$`(1.3)`$
$$2[𝕩_{k,r+1}^\pm ,𝕩_{l,s}^\pm ]2[𝕩_{k,r}^\pm ,𝕩_{l,s+1}^\pm ]=\pm \mathrm{}a_{kl}(𝕩_{k,r}^\pm 𝕩_{l,s}^\pm +𝕩_{l,s}^\pm 𝕩_{k,r}^\pm ),$$
$`(1.4)`$
$$\underset{wS_m}{}[𝕩_{k,r_{w(1)}}^\pm ,[𝕩_{k,r_{w(2)}}^\pm ,\mathrm{},[𝕩_{k,r_{w(m)}}^\pm ,𝕩_{l,s}^\pm ]\mathrm{}]]=0,kl$$
$`(1.5)`$
for all sequences of non-negative integers $`r_1,\mathrm{},r_m`$, where $`m=1a_{kl}.`$
Set
$$[n]=\frac{q^nq^n}{qq^1}n.$$
2. Quiver varieties.
Let $`I`$ (resp. $`E`$) be the set of vertices (resp. edges) of a finite graph $`(I,E)`$ with no edge loops. For $`k,lI`$ let $`n_{kl}`$ be the number of edges joining $`k`$ and $`l`$. Put $`a_{kl}=2\delta _{kl}n_{kl}`$. The map $`(I,E)A=(a_{kl})_{k,lI}`$ is a bijection from the set of finite graphs with no loops onto the set of symmetric generalized Cartan matrices. Let $`\alpha _k`$ and $`\omega _k`$, $`kI`$, be the simple roots and fundamental weights of the symmetric Kac-Moody algebra corresponding to $`A.`$ Let $`H`$ be the set of edges of $`(I,E)`$ together with an orientation. For $`hH`$ let $`h^{}I`$ (resp. $`h^{\prime \prime }I`$) the incoming (resp. the outcoming) vertex of $`h`$. If $`hH`$ we denote by $`\overline{h}H`$ the same edge with opposite orientation. Take two collection of finite dimensional complex vector spaces $`V=(V_k)_{kI},W=(W_k)_{kI}.`$ Let us fix once for all the following convention : the dimension of the graded vector space $`V`$ is identified with the element $`𝕧=_{kI}v_k\alpha _k`$ in the root lattice (where $`v_k`$ is the dimension of $`V_k`$). Similarly the dimension of $`W`$ is identified with the weight $`𝕨=_kw_k\omega _k`$ (where $`w_k`$ is the dimension of $`W_k`$). Set
$$M(𝕧,𝕨)=\underset{hH}{}\text{Hom}(V_{h^{\prime \prime }},V_h^{})\underset{kI}{}\text{Hom}(W_k,V_k)\underset{kI}{}\text{Hom}(V_k,W_k).$$
The group $`G_𝕧=_k\text{GL}(V_k)`$ acts on $`M(𝕧,𝕨)`$ by $`g(B,i,j)=(gBg^1,gi,jg^1).`$ We denote by $`B_h`$ the component of the element $`B`$ in $`\text{Hom}(V_{h^{\prime \prime }},V_h^{})`$. Let us consider the map
$$\mu _{𝕧,𝕨}:M(𝕧,𝕨)\underset{kI}{}\text{Hom}(V_k,V_k),(B,i,j)\underset{h}{}\epsilon (h)B_hB_{\overline{h}}+ij,$$
where $`\epsilon `$ is any function $`\epsilon :H^\times `$ such that $`\epsilon (h)+\epsilon (\overline{h})=0.`$ We say that a triple $`(B,i,j)\mu _{𝕧,𝕨}^1(0)`$ is stable if there is no nontrivial $`B`$-invariant subspace of $`\text{Ker}j`$. Let $`\mu _{𝕧,𝕨}^1(0)^s`$ be the set of stable triples. The group $`G_𝕧`$ acts freely on $`\mu _{𝕧,𝕨}^1(0)^s`$. Put
$$T(𝕧,𝕨)=\mu _{𝕧,𝕨}^1(0)^s/G_𝕧,N(𝕧,𝕨)=\mu _{𝕧,𝕨}^1(0)//G_𝕧$$
and let $`\pi :T(𝕧,𝕨)N(𝕧,𝕨)`$ be the affinization map (it sends $`G_𝕧(B,i,j)`$ to the only closed $`G_𝕧`$-orbit contained in $`\overline{G_𝕧(B,i,j)}`$). It is proved in \[4, 3.10(2)\] that $`T(𝕧,𝕨)`$ is a smooth quasi-projective variety. Given $`𝕧^1,𝕧^2[I]`$ consider the fiber product $`Z(𝕧^1,𝕧^2;𝕨)=T(𝕧^1,𝕨)\times _\pi T(𝕧^2,𝕨).`$ Take $`𝕧^2=𝕧^1+\alpha _k`$ where $`\alpha _k`$ is a simple root and assume that $`V^1V^2`$ have dimension $`𝕧^1`$, $`𝕧^2`$, respectively. Consider the closed subvariety $`C_k^+(𝕧^2,𝕨)`$ of $`Z(𝕧^1,𝕧^2;𝕨)`$ consisting of the pairs of triples $`(B^1,i^1,j^1),`$ $`(B^2,i^2,j^2)`$ such that $`B_{|V^1}^2=B^1,i^2=i^1,j_{|V^1}^2=j^1.`$ Put $`C_k^{}(𝕧^2,𝕨)=\phi \left(C_k^+(𝕧^2,𝕨)\right)Z(𝕧^2,𝕧^1;𝕨)`$ where $`\phi :T(𝕧^1,𝕨)\times T(𝕧^2,𝕨)T(𝕧^2,𝕨)\times T(𝕧^1,𝕨)`$ permutes the components. The varieties $`C_k^\pm (𝕧^2,𝕨)`$ are nonsingular \[4, 5.7\]. The group $`\stackrel{~}{G}_𝕨=G_𝕨\times ^\times `$ acts on $`T(𝕧,𝕨)`$ by
$$(g,t)(B,i,j)=(tB,t^2ig^1,gj),gG_𝕨,t^\times .$$
Let $`𝒱_k=\mu _{𝕧,𝕨}^1(0)^s\times _{G_𝕧}V_k`$ and $`𝒲_k`$ be respectively the $`k`$-th tautological bundle and the trivial $`W_k`$-bundle on $`T(𝕧,𝕨).`$ The bundles $`𝒱_k,𝒲_k,`$ are $`\stackrel{~}{G}_𝕨`$-equivariant. The group $`\stackrel{~}{G}_𝕨`$ acts also on $`N(𝕧,𝕨)`$, $`C_k^\pm (𝕧^2,𝕨)`$, and $`Z(𝕧^1,𝕧^2;𝕨)`$. Let $`q`$ be the trivial line bundle on $`T(𝕧,𝕨)`$ with the degree one action of $`^\times `$. For any complex $`G`$-variety X let $`K^G(X)`$ be the Grothendieck ring of $`G`$-equivariant coherent sheaves on X. Put
$$_k(𝕧,𝕨)=q^2𝒲_k(1+q^2)𝒱_k+q^1\underset{h^{}=k}{}𝒱_{h^{\prime \prime }}K^{\stackrel{~}{G}_𝕨}(T(𝕧,𝕨)).$$
The rank of $`_k(𝕧,𝕨)`$ is $`(𝕨𝕧|\alpha _k),`$ where $`(|)`$ is the standard metric on the weight lattice of $`𝔤.`$ We fix a pair of linear maps $`𝕨𝕨_\pm `$ on the weight lattice which are adjoint with respect to $`(|)`$, and such that $`𝕨_++𝕨_{}=𝕨`$ for all $`𝕨`$.
3. Equivariant homology and convolution product.
Let $`G`$ be a complex, connected, linear algebraic group. For any complex $`G`$-variety $`X`$, let $`H_i^G(X)`$ (resp. $`H_G^i(X)`$) be the $`i`$-th space of $`G`$-equivariant complex Borel-Moore homology (resp. of $`G`$-equivariant complex cohomology). Put
$$H^G(X)=\underset{i}{}H_i^G(X),H_G(X)=\underset{i}{}H_G^i(X).$$
See for details on equivariant Borel-Moore homology. Let us only recall the following well known facts.
\- If $`Y`$ is a closed $`G`$-subvariety of $`X`$ and $`X`$ is smooth, then $`H^G(Y)=H_G(X,XY).`$ Moreover there is a natural map $`H_G(X)H^G(X).`$ Call $`\alpha ^oH^G(X)`$ the image of $`\alpha H_G(X).`$ The $``$-product in equivariant cohomology induce, via the Poincaré duality, a product, noted $``$, in equivariant homology. We will denote also by a dot the product $`H_G(X)H^G(X)H^G(X).`$ \- Any $`G`$-equivariant vector bundle $`E`$ on $`X`$ admits an equivariant Chern polynomial $`\lambda _z(E)H_G(X)[z].`$ The coefficient of $`z`$ in $`\lambda _z(E)`$ is the equivariant first Chern class $`c_1(E)H_G(X).`$ The coefficient of $`z^{\text{rk }(E)}`$ in $`\lambda _z(E)`$ is the equivariant Euler class $`\lambda (E)H_G(X).`$ If $`E`$ is invertible, then $`\lambda _z(E)=1+c_1(E)z`$. Moreover, for any $`E`$ and $`F`$ we have $`\lambda _z(EF)=\lambda _z(E)\lambda _z(F)`$. The class $`\lambda _z(E)`$ depends only on the class of $`E`$ in $`K^G(X).`$ \- If $`TG`$ is a maximal torus, put $`𝔱=\text{Lie}(T)`$. Then $`H_G^{2i}(pt)=S^{2i}(𝔱^{})^W`$, where $`S^i`$ is the $`i`$-symmetric product and $`W`$ is the Weyl group. We will use the following (see \[1, Proposition 2.6.47\]) : Lemma. Let $`X`$ be a smooth $`G`$-variety and let $`C_i`$ ($`i=1,2`$) be two smooth closed $`G`$-subvarieties. Set $`C_3=C_1C_2`$ and let $`\gamma _i:C_iX`$ ($`i=1,2,3`$) be the natural embedding. Suppose that $`C_1`$ and $`C_2`$ are transversal. Then, for all $`\alpha H_G(C_1)`$ and $`\beta H_G(C_2)`$,
$$\gamma _1(\alpha ^o)\gamma _2(\beta ^o)=\gamma _3\left((\alpha _{|C_3}\beta _{|C_3})^o\right),$$
where $`\alpha _{|C_3}`$ (resp. $`\beta _{|C_3}`$) is the restriction of $`\alpha `$ (resp. $`\beta `$) to $`C_3`$. $``$Let us recall the definition of the convolution product. Given quasi-projective $`G`$-varieties $`X_1,X_2,X_3,`$ consider the projection $`p_{ij}:X_1\times X_2\times X_3X_i\times X_j`$ for all $`1i<j3`$. Consider subvarieties $`Z_{ij}X_i\times X_j`$ such that the restriction of $`p_{13}`$ to $`p_{12}^1Z_{12}p_{23}^1Z_{23}`$ is proper and maps to $`Z_{13}`$. The convolution product is the map
$$:H^G(Z_{12})H^G(Z_{23})H^G(Z_{13}),\alpha \beta p_{13}((p_{12}^{}\alpha )(p_{23}^{}(\beta )).$$
See \[1, 2.7 and the remark (iii), page 113\] for more details on convolution product. We will essentially consider the case $`X_i=T(𝕧^i,𝕨)`$ and $`Z_{ij}=Z(𝕧^i,𝕧^j;𝕨)`$, where $`𝕧^1,𝕧^2,𝕧^3,𝕨[I]`$ and $`1i<j3`$.
4. Statement of the Result.
Let $`(I,E)`$ be a graph of type $`ADE`$. Fix $`𝕨,𝕧^1,𝕧^2[I]`$, with $`𝕧^2=𝕧^1+\alpha _k`$. For any $`k`$, denote by $`𝒱_k^1`$ (resp. $`𝒱_k^2`$) the vector bundle $`𝒱_k𝒪_{T(𝕧^2,𝕨)}`$ (resp. $`𝒪_{T(𝕧^1,𝕨)}𝒱_k`$) over $`T(𝕧^1,𝕨)\times T(𝕧^2,𝕨).`$ The restriction to $`C_k^+(𝕧^2,𝕨)`$ of the sheaf $`𝒱_k^1`$ is a subsheaf of $`𝒱_k^2.`$ The quotient sheaf $`_k^+=𝒱_k^2/𝒱_k^1`$ is $`\stackrel{~}{G}_𝕨`$-invariant and invertible. Put $`_k^{}=\phi ^{}_k^+.`$ Consider the following varieties
$$N(𝕨)=\underset{𝕧}{}N(𝕧,𝕨),T(𝕨)=\underset{𝕧}{}T(𝕧,𝕨),Z(𝕨)=\underset{𝕧^{},𝕧^{\prime \prime }}{}Z(𝕧^{},𝕧^{\prime \prime };𝕨),$$
where $`𝕧,𝕧^{},𝕧^{\prime \prime }`$ take all the possible values in $`[I]`$. Let $`\mathrm{\Delta }^\pm `$ be the two natural embeddings
$$\mathrm{\Delta }^+:C_k^+(𝕧^2,𝕨)Z(𝕧^1,𝕧^2;𝕨)\text{and}\mathrm{\Delta }^{}:C_k^{}(𝕧^2,𝕨)Z(𝕧^2,𝕧^1;𝕨).$$
If $`r0`$, put
$$x_{k,r}^\pm =\underset{𝕧^2}{}(1)^{(\alpha _k|𝕧_\pm ^2)}\mathrm{\Delta }_{}^\pm (c_1(_k^\pm )^o)^rH^{\stackrel{~}{G}_𝕨}(Z(𝕨)).$$
$`(4.1)`$
Let $`\mathrm{\Delta }:T(𝕧,𝕨)T(𝕧,𝕨)\times T(𝕧,𝕨)`$ be the diagonal embedding and set $`\mathrm{}=c_1(q^2)^o`$. Define $`h_{k,r}`$ as the coefficient of $`\mathrm{}z^{r1}`$ in
$$\left(1+\underset{𝕧}{}\mathrm{\Delta }_{}\frac{\lambda _{1/z}(_k(𝕧,𝕨))}{\lambda _{1/z}(q^2_k(𝕧,𝕨))}\right)^{},$$
$`(4.2)`$
where $``$ stands for the expansion at $`z=\mathrm{}.`$ The following result was conjectured by Nakajima (\[5, Introduction\], in the result was announced for type $`A`$). Theorem. For all $`𝕨[I]`$, the map $`𝕩_{k,r}^\pm x_{k,r}^\pm ,`$ $`𝕙_{k,r}h_{k,r}`$ extends uniquely to an algebra homomorphism $`\mathrm{\Phi }_𝕨:𝕐_{\mathrm{}}(\mathrm{L}𝔤)H^{\stackrel{~}{G}_𝕨}(Z(𝕨)).`$ $``$Remark. We can prove a similar result for any symmetric Kac-Moody algebra. Let $`A=(a_{kl})_{k,lI}`$ be a symmetric generalized Cartan matrix. In the definition of the Yangian, the relation (1.4) becomes
$$\{\begin{array}{c}[𝕩_{k,r+1}^\pm ,𝕩_{k,s}^\pm ][𝕩_{k,r}^\pm ,𝕩_{k,s+1}^\pm ]=\pm \mathrm{}(𝕩_{k,r}^\pm 𝕩_{k,s}^\pm +𝕩_{k,s}^\pm 𝕩_{k,r}^\pm )\\ \\ \eta _{a_{kl}}(z\frac{\mathrm{}}{2},w)𝕩_k^\pm (z)𝕩_l^\pm (w)=\eta _{a_{kl}}(z,w\frac{\mathrm{}}{2})𝕩_l^\pm (w)𝕩_k^\pm (z)(\text{if}kl)\end{array}$$
where
$$𝕩_k^\pm (z)=\underset{r0}{}𝕩_{k,r}^\pm z^r,\text{and}\eta _a(z,w)=\underset{j=1}{\overset{n}{}}\left(zw+(1+a2j)\mathrm{}/2\right).$$
In this case the action of $`^\times `$ on $`T(𝕧,𝕨)`$ and the complex $`_k(𝕧,𝕨)`$ has to be changed as in . In the proof of the theorem there are only minor and evident changes to do.
5. Proof of the Result.
The proof is as in \[5, sections 10 and 11\] : we check relations (1.1), (1.2), (1.5) and relations (1.3) and (1.4) in the case $`kl`$ by direct computation. Relations (1.3) and (1.4) in the case $`k=l`$ are proved by reduction to the $`𝔰𝔩_2`$-case. We insist here only on the parts which need different calculations.
Relation (1.1). It is an immediate consequence of the definition, since for all $`xH^{\stackrel{~}{G}_𝕨}(Z(𝕧,𝕧^{};𝕨))`$ we have
$$h_{k,0}x=\text{rk }_k(𝕧,𝕨)x=(𝕨𝕧|\alpha _k)x,$$
$$xh_{k,0}=\text{rk }_k(𝕧^{},𝕨)x=(𝕨𝕧^{}|\alpha _k)x.$$
Relation (1.2). We prove only the plus case, the minus being similar. Fix $`𝕧^2=𝕧^1+\alpha _l`$. We identify $`_k(𝕧^1,𝕨)`$ and $`_k(𝕧^2,𝕨)`$ with their pull-back to $`C_l^+(𝕧^2,𝕨)`$ via the 1-st and the 2-nd projection. Then, in $`\mathrm{K}^{\stackrel{~}{G}_𝕨}(C_l^+(𝕧^2,𝕨))`$, we have
$$_k(𝕧^1,𝕨)q^2_k(𝕧^1,𝕨)=_k(𝕧^2,𝕨)q^2(𝕧^2,𝕨)+[a_{kl}](q^1q)_l^+.$$
It follows that $`[h_{k,r},x_{l,s}^+]H^{\stackrel{~}{G}_𝕨}(C_l^+(𝕧^2,𝕨))`$ is the coefficient of $`\mathrm{}z^{r1}`$ in
$$\left(\lambda _{1/z}(_k(𝕧^1,𝕨)q^2_k(𝕧^1,𝕨))\lambda _{1/z}(_k(𝕧^2,𝕨)q^2_k(𝕧^2,𝕨))\right)^{}x_{l,s}^+=$$
$$=\left(\left(\lambda _{1/z}([a_{kl}](q^1q)_l^+)1\right)\lambda _{1/z}(_k(𝕧^2,𝕨)q^2_k(𝕧^2,𝕨))\right)^{}x_{l,s}^+.$$
Set
$$A_s=\lambda _{1/z}(_k(𝕧^2,𝕨)q^2_k(𝕧^2,𝕨))x_{l,s}^+,$$
$$X=\lambda _{1/z}([a_{kl}](q^1q)_l^+)=\frac{1(c_l^+a_{kl}\mathrm{}/2)z^1}{1(c_l^++a_{kl}\mathrm{}/2)z^1}.$$
Then the LHS and the RHS of the relation (1.2) are respectively equal to the coefficient of $`\mathrm{}z^{r1}`$ in
$$\left(2z(X1)A_s2(X1)A_{s+1}\right)^{}=\left(2(X1)(zc_l^+)A_s\right)^{}\text{and}\left(\mathrm{}a_{kl}(X+1)A_s\right)^{}.$$
We are then reduced to the identity, easily checked, in $`H^{\stackrel{~}{G}_𝕨}(C_l^+(𝕧^2,𝕨))`$ :
$$2(X1)(zc_l^+)=\mathrm{}a_{kl}(X+1).$$
Relation (1.3) with $`kl.`$ Fix $`𝕧^1,𝕧^2,\stackrel{~}{𝕧}^2,𝕧^3`$, such that
$$\stackrel{~}{𝕧}^2=𝕧^1\alpha _l=𝕧^3\alpha _k=𝕧^2\alpha _k\alpha _l.$$
If $`1i<j3`$, consider the projections
$$p_{ij}:T(𝕧^1,𝕨)\times T(𝕧^2,𝕨)\times T(𝕧^3,𝕨)T(𝕧^i,𝕨)\times T(𝕧^j,𝕨),$$
$$\stackrel{~}{p}_{ij}:T(𝕧^1,𝕨)\times T(\stackrel{~}{𝕧}^2,𝕨)\times T(𝕧^3,𝕨)T(\stackrel{~}{𝕧}^i,𝕨)\times T(\stackrel{~}{𝕧}^j,𝕨),$$
where we set $`\stackrel{~}{𝕧}^1=𝕧^1,\stackrel{~}{𝕧}^3=𝕧^3.`$ We have
$$x_{k,r}^+x_{l,s}^{}=(1)^{(\alpha _k|𝕧_+^2)+(\alpha _l|𝕧_{}^3)}p_{13}\left(p_{12}^{}(c_1(𝒱_k^2/𝒱_k^1)^{or})p_{23}^{}(c_1(𝒱_l^3/𝒱_l^2)^{os})\right),$$
$$x_{l,s}^{}x_{k,r}^+=(1)^{(\alpha _l|\stackrel{~}{𝕧}_{}^2)+(\alpha _k|𝕧_+^3)}\stackrel{~}{p}_{13}\left(\stackrel{~}{p}_{12}^{}(c_1(\stackrel{~}{𝒱}_l^2/𝒱_l^1)^{os})\stackrel{~}{p}_{23}^{}(c_1(𝒱_k^3/\stackrel{~}{𝒱}_k^2)^{or})\right).$$
It is proved in \[5, Lemma 10.2.1\] that the intersections
$$p_{12}^1C_k^+(𝕧^2,𝕨)p_{23}^1C_l^{}(𝕧^3,𝕨)\text{and}\stackrel{~}{p}_{12}^1C_l^{}(\stackrel{~}{𝕧}^2,𝕨)\stackrel{~}{p}_{23}^1C_k^+(𝕧^3,𝕨)$$
are transversal in $`T(𝕧^1,𝕨)\times T(𝕧^2,𝕨)\times T(𝕧^3,𝕨)`$ and $`T(𝕧^1,𝕨)\times T(\stackrel{~}{𝕧}^2,𝕨)\times T(𝕧^3,𝕨)`$ and that there exists a $`\stackrel{~}{G}_𝕨`$-equivariant isomorphisms between them which induces the isomorphisms :
$$𝒱_k^2/𝒱_k^1𝒱_k^3/\stackrel{~}{𝒱}_k^2\text{and}\stackrel{~}{𝒱}_l^2/𝒱_l^1𝒱_l^3/𝒱_l^2.$$
The result follows from the lemma in section 3.
Relation (1.4) with $`kl`$. We prove only the plus case. Fix $`𝕧^1,𝕧^2,\stackrel{~}{𝕧}^2,𝕧^3,`$ such that
$$𝕧^3=\stackrel{~}{𝕧}^2+\alpha _k=𝕧^2+\alpha _l=𝕧^1+\alpha _k+\alpha _l.$$
Consider the projections $`p_{ij}`$ and $`\stackrel{~}{p}_{ij}`$ ($`1i<j3`$) as before. The intersections
$$Z_{kl}=p_{12}^1C_k^+(𝕧^2,𝕨)p_{23}^1C_l^{}(𝕧^3,𝕨)\text{and}Z_{lk}=\stackrel{~}{p}_{12}^1C_l^{}(\stackrel{~}{𝕧}^2,𝕨)\stackrel{~}{p}_{23}^1C_k^+(𝕧^3,𝕨)$$
are transversal in $`T(𝕧^1,𝕨)\times T(𝕧^2,𝕨)\times T(𝕧^3,𝕨)`$ and $`T(𝕧^1,𝕨)\times T(\stackrel{~}{𝕧}^2,𝕨)\times T(𝕧^3,𝕨)`$ (see \[5, Lemma 10.3.1\]). Since $`kl`$, the restriction of $`p_{13}`$ and $`\stackrel{~}{p}_{13}`$ to $`Z_{kl}`$ and $`Z_{lk}`$ is an embedding into $`Z(𝕧^1,𝕧^3;𝕨).`$ Call it $`\iota _{kl}`$ and $`\iota _{lk}`$ respectively. Put $`b_k=c_1(𝒱_k^3𝒱_k^1)`$, $`b_l=c_1(𝒱_l^3𝒱_l^1).`$ We have (see the lemma in section 3)
$$x_{k,r}^+x_{l,s}^+=(1)^{(\alpha _k|𝕧_+^2)+(\alpha _l|𝕧_+^3)}\iota _{kl}\left((p_{12}^{}(b_k^r)_{|Z_{kl}}p_{23}^{}(b_l^s)_{|Z_{kl}})^o\right),$$
$$x_{l,s}^+x_{k,r}^+=(1)^{(\alpha _l|\stackrel{~}{𝕧}_+^2)+(\alpha _k|𝕧_+^3)}\iota _{lk}\left((p_{12}^{}(b_k^r)_{|Z_{lk}}p_{23}^{}(b_l^s)_{|Z_{kl}})^o\right).$$
Take $`hH`$ such that $`h^{}=l`$ and $`h^{\prime \prime }=k.`$ The map $`B_{\overline{h}}`$ may be viewed as a section of the $`\stackrel{~}{G}_𝕨`$-bundle $`_{kl}=q(𝒱_l^3/𝒱_l^1)^{}(𝒱_k^3/𝒱_k^1)`$ on $`p_{13}(Z_{kl})`$ (where we set $`_{kl}=0`$ if $`a_{kl}=0`$). Similarly $`B_h`$ is a section of the $`\stackrel{~}{G}_𝕨`$-bundle $`_{lk}=q(𝒱_k^3/𝒱_k^1)^{}(𝒱_l^3/𝒱_l^1)`$ on $`\stackrel{~}{p}_{13}(Z_{lk})`$ (where again we set $`_{lk}=0`$ if $`a_{kl}=0`$). In \[5, 10.3.9\] it is proved that $`B_{\overline{h}}`$ and $`B_h`$ are transversal to the zero section respectively. Moreover
$$p_{13}(Z_{kl})B_{\overline{h}}^1(0)=\stackrel{~}{p}_{13}(Z_{lk})B_h^1(0).$$
Then,
$$\iota _{kl}(c_1(_{kl})^o)x_{k,r}^+x_{l,s}^+=(1)^{a_{kl}}\iota _{lk}(c_1(_{lk})^o)x_{l,s}^+x_{k,r}^+,$$
i.e.
$$\iota _{kl}(b_k^ob_l^o+\mathrm{}/2)x_{k,r}^+x_{l,s}^+=\iota _{lk}(b_l^ob_k^o+\mathrm{}/2)x_{l,s}^+x_{k,r}^+.$$
The relation (1.4) follows immediately from this.
Relations (1.3) and (1.4) with $`k=l`$. Thank to the same argument than in \[5, 11.3\] we are reduce to the case of $`(I,E)`$ of type $`A_1.`$ In this case $`𝕧`$ and $`𝕨`$ are identified with natural numbers, so let us call them $`v`$ and $`w`$. Moreover we will omit everywhere the subindex 1. Let $`\text{Gr}_v(w)`$ be the variety of $`v`$-dimensional subspaces in $`W`$. It is easy to see that $`T(v,w)T^{}\text{Gr}_v(w)`$. The group $`G_w`$ acts in the obvious way on $`T(v,w)`$. The group $`^\times `$ acts by scalar multiplication on the fibers of the cotangent bundle. Fix $`T_1,\mathrm{},T_w`$ such that
$$\begin{array}{c}K^{\stackrel{~}{G}_w}(\text{Gr}_v(w))=[q^{\pm 1},T_1^{\pm 1},\mathrm{},T_w^{\pm 1}]^{S_v\times S_{wv}},\\ \\ ^i𝒱=e_i(T_1,\mathrm{},T_v),^i𝒲=e_i(T_1,\mathrm{},T_w),\end{array}$$
where $`e_i`$ is the $`i`$-th elementary symmetric function. We get
$$(v,w)=q^2𝒲(1+q^2)𝒱=q^2(T_{v+1}+\mathrm{}+T_w)(T_1+\mathrm{}+T_v).$$
Put $`t_k=c_1(T_k)^o`$. Then $`H^{\stackrel{~}{G}_w}(T(w))=_{v=0}^w[\mathrm{},t_1,\mathrm{}t_w]^{S_v\times S_{wv}}.`$ The following lemma is proved as in \[1, Claim 7.6.7\]. Lemma.The space $`H^{\stackrel{~}{G}_w}(T(w))`$ is a faithful module over $`H^{\stackrel{~}{G}_w}(Z(w))`$. $``$The operators $`x_r^\pm `$ on $`H^{\stackrel{~}{G}_w}(T(w))`$ can be written down explicitely. Put
$$O(v,w)=\{(V^1,V^2)\text{Gr}_{v1}(w)\times \text{Gr}_v(w)|V^1V^2\}.$$
The Hecke correspondence $`C^+(v,w)`$ is the conormal bundle to $`O(v,w)`$. Consider the projections $`p_1,p_2`$ from $`O(v,w)`$ to the first and the second component and let $`\pi :T(v,w)Gr_v(w)`$ be the projection. We can prove as in \[6, Lemme 5\] that if $`\alpha H_{\stackrel{~}{G}_w}(O(v,w))`$ and $`\beta H^{\stackrel{~}{G}_w}(\text{Gr}_v(w))`$, then
$$\pi ^{}(\alpha )\pi ^{}(\beta )=p_1(\lambda (q^2T^{}p_1)\alpha p_{2}^{}{}_{}{}^{}\beta ),$$
where $`T^{}p_1`$ is the relative cotangent bundle to $`p_1`$. The map $`p_1`$ is a $`^{wv}`$-fibration, then we have
$$\lambda (q^2T^{}p_1)=\underset{m=v+1}{\overset{w}{}}(t_mt_v+\mathrm{})H^{\stackrel{~}{G}_w}(O(v,w)).$$
Let us introduce the following notation. Fix $`z[1,w]=\{1,2,\mathrm{},w\}`$ and let $`I=(I_1,I_2)`$ be a partition of $`[1,w]`$ into two subset of cardinality $`z`$ and $`wz`$ respectively, say $`I_1=\{a_1,a_2,\mathrm{},a_z\},`$ $`I_2=\{b_1,b_2,\mathrm{},b_{wz}\}.`$ Then put
$$f(t_{I_1};t_{I_2})=f(t_{a_1},t_{a_2},\mathrm{},t_{a_z},t_{b_1},t_{b_2},\mathrm{},t_{b_{wz}}).$$
Thus (see \[6, Lemme 1\]), for any $`f[\mathrm{},t_1,\mathrm{},t_w]^{S_v\times S_{wv}},`$
$$x_r^+(f)(t_{[1,v1]};t_{[v,w]})=\underset{k=v}{\overset{w}{}}f(t_{[1,v1]\{k\}};t_{[v,w]\{k\}})t_k^r\underset{m[v,w]\{k\}}{}\left(1+\frac{\mathrm{}}{t_kt_m}\right),$$
$`(5.1)`$
$$x_r^{}(f)(t_{[1,v+1]};t_{[v+2,w]})=\underset{k=1}{\overset{v+1}{}}f(t_{[1,v+1]\{k\}};t_{[v+2,w]\{k\}})t_k^r\underset{m[1,v+1]\{k\}}{}\left(1+\frac{\mathrm{}}{t_mt_k}\right).$$
$`(5.2)`$
We have
$$\lambda _z(v,w)=\frac{\underset{m=v+1}{\overset{w}{}}\left(1z(t_m\mathrm{})\right)}{_{m=1}^v\left(1z(t_m+1)\right)}.$$
Thus $`h_r(f)`$ is the coefficient of $`\mathrm{}z^{r1}`$ in
$$f\left(\underset{m=1}{\overset{v}{}}\frac{zt_m\mathrm{}}{zt_m}\underset{m=v+1}{\overset{w}{}}\frac{zt_m+\mathrm{}}{zt_m}\right)^{}.$$
$`(5.3)`$
Proposition. Relations $`(1.3)`$ and $`(1.4)`$ hold in the $`𝔰𝔩_2`$-case.
Proof. Let us prove relation (1.3). Fix $`f[\mathrm{},t_1,\mathrm{},t_w]^{S_v\times S_{wv}}.`$ Using formulas (5.1) and (5.2), we have
$$\left(x_s^{}x_r^+(f)\right)(t_{[1,v]};t_{[v+1,w]})=\underset{l=1}{\overset{v}{}}\underset{k[v+1,w]\{l\}}{}f(t_{([1,v]\{l\})\{k\}};t_{([v+1,w]\{l\})\{k\}})t_l^st_k^rX_{kl},$$
$$\left(x_r^+x_s^{}(f)\right)(t_{[1,v]};t_{[v+1,w]})=\underset{l[1,v]\{k\}}{}\underset{k=v+1}{\overset{w}{}}f(t_{([1,v]\{k\})\{l\}};t_{([v+1,w]\{k\})\{l\}})t_l^st_k^rY_{kl},$$
where
$$\begin{array}{c}X_{kl}=\underset{m[1,v]\{l\}}{}\left(1+\frac{\mathrm{}}{t_mt_l}\right)\underset{n[v+1,w]\{l\}\{k\}}{}\left(1+\frac{\mathrm{}}{t_kt_n}\right),\\ \\ Y_{kl}=\underset{m[1,v]\{k\}\{l\}}{}\left(1+\frac{\mathrm{}}{t_mt_l}\right)\underset{n[v+1,w]\{k\}}{}\left(1+\frac{\mathrm{}}{t_kt_n}\right).\end{array}$$
The terms with $`kl`$ cancel out in the bracket. We get
$$[x_r^+,x_s^{}](f)=f\underset{k=v+1}{\overset{w}{}}t_k^st_k^r\underset{m[1,v]}{}\left(1+\frac{\mathrm{}}{t_mt_k}\right)\underset{n[v+1,w]\{k\}}{}\left(1+\frac{\mathrm{}}{t_kt_n}\right)$$
$$f\underset{l=1}{\overset{v}{}}t_l^st_l^r\underset{m[1,v]\{l\}}{}\left(1+\frac{\mathrm{}}{t_mt_l}\right)\underset{n[v+1,w]}{}\left(1+\frac{\mathrm{}}{t_lt_n}\right).$$
Put
$$A(z)=\underset{m=1}{\overset{w}{}}(zt_m),B(z)=\underset{m=1}{\overset{v}{}}(zt_m\mathrm{})\underset{m=v+1}{\overset{w}{}}(zt_m+\mathrm{}).$$
Then it is easy to check that
$$\mathrm{}[x_r^+,x_s^{}](f)=f\underset{k=1}{\overset{w}{}}t_k^{r+s}\frac{B(t_k)}{A^{}(t_k)}=f\text{res}_{\mathrm{}}z^{r+s}\frac{B(z)}{A(z)}.$$
This is the definition of $`h_{r+s}(f)`$ given in (5.3). As for the relation (1.4), note that, using (5.1), we get $`\left(x_s^+x_r^+(f)\right)(t_{[1,v2]};t_{[v1,w]})=`$
$$=\underset{l=v1}{\overset{w}{}}\underset{k[v1,w]\{l\}}{}f(t_{[1,v2]\{k,l\}};t_{[v1,w]\{k,l\}})t_l^st_k^rZ_{kl},$$
where
$$Z_{kl}=\underset{n[v1,w]\{l\}}{}\left(1+\frac{\mathrm{}}{t_lt_n}\right)\underset{m[v1,w]\{k,l\}}{}\left(1+\frac{\mathrm{}}{t_kt_m}\right).$$
The relation in the plus case follows now by a direct computation. $``$
Relation (1.5). The proof is exactly as in \[5, 10.4\], so we omit it .
Remark. Nakajima \[5, Theorem 9.4.1\] has proved that there exists an algebra morphism
$$\mathrm{\Psi }_𝕨:𝕌_q(\mathrm{L}𝔤)K^{\stackrel{~}{G}_𝕨}(Z(𝕨))_{[q,q^1]}(q),$$
where the algebra to the left is the quantized enveloping algebra of $`\mathrm{L}𝔤`$ and the algebra to the right is equipped with the convolution product. Using $`\mathrm{\Phi }_𝕨`$ and $`\mathrm{\Psi }_𝕨`$ we can construct the finite dimensional simple modules of $`𝕐_{\mathrm{}}(\mathrm{L}𝔤)`$ and $`𝕌_q(\mathrm{L}𝔤)`$ respectively (see \[5, section 14\]). In particular $`𝕐_{\mathrm{}}(\mathrm{L}𝔤)`$ and $`𝕌_q(\mathrm{L}𝔤)`$ have the same finite dimensional representation theory. More precisely let $``$ (resp. $`𝔇`$) be the abelian category of finite dimensional $`𝕌_q(\mathrm{L}𝔤)`$-(resp. $`𝕐_{\mathrm{}}(\mathrm{L}𝔤)`$-)modules such that the Drinfeld polynomials of the simple factors have roots in $`q^{}`$ (resp. $``$). Proposition. The characters (as $`𝕌_q(𝔤)`$-modules and $`𝕌(𝔤)`$-modules resp.) of the simple finite dimensional modules in $``$ and in $`𝔇`$ are the same. $``$
Aknowledgments. The author thanks E. Vasserot for useful discussions.
References.
1. Chriss, N., Ginzburg, V.: Representation theory and complex geometry, Birkhäuser, Boston-Basel-Berlin, 1997.
2. Ginzburg, V., Vasserot, E.: Langlands reciprocity for affine quantum groups of type $`A_n`$, Internat. Math. Res. Notices, 3, 1993, 67-85.
3. Lusztig, G.: Cuspidal local systems and graded Hecke algebras I, Inst. Hautes Études Sci. Publ. Math., 67, 1988, 145-202.
4. Nakajima, H.: Quiver varieties and Kac-Moody algebras, Duke. Math. J., 91, 1998, 515-560.
5. Nakajima, H.: Quiver varieties and finite dimensional representations of quantum affine algebras, Preprint QA/9912158.
6. Vasserot, E.: Représentations de groupes quantiques et permutations, Ann. Scient. Éc. Norm. Sup., 26, 1993, 747-773.
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# Reduction of covers and Hurwitz spaces
## Introduction
Let $`R`$ be a complete discrete valuation ring whose residue field $`k`$ is algebraically closed of characteristic $`p`$ and whose quotient field $`K`$ is of characteristic zero. Let $`f_K:Y_K_K^1`$ be a Galois cover defined over $`K`$. We ask ourselves whether $`f_K`$ has good reduction. In case $`p`$ divides the order of the Galois group, this is a hard question. We cannot expect good reduction, in general.
A recent paper of Raynaud gives a criterion for good reduction in a first case. Let $`f_K:Y_K_K^1`$ be a Galois cover branched at $`0`$, $`1`$ and $`\mathrm{}`$. Suppose that $`p`$ strictly divides the order of the Galois group $`G`$, but not the ramification indices of $`f_K`$ and that the center of $`G`$ is trivial. Let $`P/p`$ be a $`p`$-Sylow of $`G`$, and denote by $`n:=[N_G(P):C_G(P)]`$ the index of the centralizer of $`P`$ in the normalizer of $`P`$. Let $`e`$ be the absolute ramification index of $`p`$ in $`K`$. It is shown that the cover has good reduction, provided that $`en<p1`$. The idea of the proof is to suppose that $`f_K`$ has bad reduction and to study the semistable reduction of $`f_K`$. Raynaud proves general structure results on the semistable reduction. Using these techniques, he proves that bad reduction implies $`enp1`$.
In this paper, we follow the approach of Raynaud. We consider the reduction of $`G`$-covers $`f_K:Y_K_K^1`$ branched at four points. We suppose that $`p|||G|`$, but that $`p`$ does not divide the ramification indices of $`f_K`$. This is the next case to study after the result of Raynaud. However, we put a much stronger condition on the group $`G`$. In particular, we assume that $`n=2`$. The criterion for good reduction given by Raynaud extends to our situation. The stronger condition on the Galois group allows us to get a stronger result. For instance, we are able to describe the semistable model of $`f_K`$.
When passing from $`3`$ to $`4`$ branch points, a new aspect arises: the reduction of $`f_K`$ might depend on the position of the branch points. It is therefore natural to study the corresponding Hurwitz space, i.e. the moduli space of $`G`$-covers of a certain type, and its reduction to positive characteristic. Hurwitz spaces were first introduced in a purely geometric context, but have since then proved to be useful for studying arithmetic aspects of covers as well, see e.g. , . There are many variants of Hurwitz spaces. To fix ideas, let $`G`$ be a group, $`r3`$ and denote by $`H:=H_r^{\mathrm{in}}(G)`$ the Hurwitz space parameterizing $`G`$-Galois covers of $`^1`$ with $`r`$ branch points. Then $`H`$ is a smooth variety defined over $``$. It is known that $`H`$ has good reduction to characteristic $`p`$ provided $`p|G|`$.
Recently, Abramovich and Oort suggested a definition of an arithmetic compactification of Hurwitz spaces. We follow, and somewhat simplify, this approach. The idea is to take the closure of $`H`$ inside a bigger moduli space which parameterizes maps between stably marked curves. We obtain an algebraic space $`\overline{H}`$ which is proper over $``$ and contains $`H`$ as a dense open subspace. The complement $`\overline{H}^{\mathrm{bad}}:=\overline{H}H`$ is a closed subspace supported in positive characteristic and corresponds to $`G`$-covers with bad reduction. Let us denote by $`\overline{H}^{\mathrm{good}}`$ the closure of $`H𝔽_p`$ inside $`\overline{H}𝔽_p`$. In the case of $`r=4`$ branch points, there is a finite map $`\overline{H}^{\mathrm{good}}_\lambda ^1`$. Let $`d^{\mathrm{good}}`$ be its degree. For a “generic” choice of $`\lambda K\{0,1\}`$, there will be exactly $`d^{\mathrm{good}}`$ nonisomorphic $`G`$-covers $`f_K:Y_K_K^1`$ branched in $`0,1,\mathrm{}`$ and $`\lambda `$ which have good reduction. More precisely, there is a finite set $`\overline{\lambda }_1,\mathrm{},\overline{\lambda }_m\overline{𝔽}_p`$ of exceptional values such that $`d^{\mathrm{good}}`$ is the number of covers as above with good reduction provided $`\lambda \overline{\lambda }_i(modp)`$ for all $`i`$. The $`\lambda _i`$’s are the images of the points where $`\overline{H}^{\mathrm{good}}`$ intersects $`\overline{H}^{\mathrm{bad}}`$. One way to actually compute the number $`d^{\mathrm{good}}`$ and the exceptional values $`\overline{\lambda }_i`$ would be to determine the precise structure of $`\overline{H}^{\mathrm{bad}}`$. This seems a difficult problem, in general. In this paper, we are able to solve it in a special case, thanks to a surprising connection with modular curves.
Even though our definition of $`\overline{H}`$ works without any extra assumption, it seems to be a hard problem to describe the structure of $`\overline{H}`$ and $`\overline{H}^{\mathrm{bad}}`$ in general. To our knowledge, the only case that has been studied is the case of modular curves, which admit an interpretation as Hurwitz spaces. Let us sketch this correspondence in the case of $`X_1(p)`$. Let $`f_K:Y_K_K^1`$ be a $`G`$-cover branched at $`4`$ points of order $`2`$, where $`G`$ is a dihedral group of order $`2p`$. We may identify $`Y_K`$ with an elliptic curve $`E_K`$ over $`K`$; the cover $`f_K`$ factors through a $`p`$-cyclic isogeny $`\pi _K:E_KE_k^{}`$. Moreover, the choice of an element $`\sigma `$ of $`G`$ of order $`p`$ defines a $`p`$-torsion point $`P:=\sigma (0)E_K[p]`$ generating the kernel of $`\pi _K`$. The pair $`(E_K,P)`$ corresponds to a point on $`X_1(p)`$. This gives an identification $`HX_1(p)`$, for a suitable Hurwitz space $`H`$. By the results of Katz and Mazur , the reduction of $`X_1(p)`$ to characteristic $`p`$ is well understood. One can check that the subspace $`\overline{H}^{\mathrm{bad}}`$ of the arithmetic compactification of $`H`$ corresponds to the component of $`X_1(p)𝔽_p`$ parameterizing pairs $`(E,P)`$ such that $`P=0`$ and hence the isogeny $`\pi :EE^{}=E/<P>`$ is inseparable. Even without using the very precise results of , the theory of elliptic curves gives the following result on good reduction of Galois covers. Suppose the elliptic curve $`E_K^{}`$ given by the equation $`y^2=x(x1)(x\lambda )`$ has good ordinary reduction. Then there are precisely $`p1`$ nonisomorphic $`G`$-covers branched in $`0`$, $`1`$, $`\mathrm{}`$ and $`\lambda `$, with ramification index $`2`$, which have good reduction. On the other hand, there is no such cover with good reduction if $`E_K^{}`$ has supersingular reduction.
### Results
In this paper, we look at the following situation. Let $`G`$ be a finite group and $`p`$ an odd prime which strictly divides $`|G|`$. We assume that the normalizer of a $`p`$-Sylow of $`G`$ is a dihedral group. Let $`K`$ be as in the beginning, and let $`f_K:Y_K_K^1`$ be a $`G`$-Galois cover branched at $`4`$ points, of order prime-to-$`p`$. We prove that the cover $`f_K`$ has either good reduction or a very specific type of bad reduction, which we call modular reduction.
We will briefly explain what this means. Assume that $`f_K`$ has bad reduction. Following Raynaud , §3.2, we associate to the $`G`$-cover $`f_K`$ a $`\mathrm{\Delta }`$-cover $`g_K:Z_K_K^1`$, called the auxiliary cover. Here $`\mathrm{\Delta }`$ is a subgroup of $`G`$, and $`g_K`$ is branched in the same points as $`f_K`$ and has bad reduction. The statement that $`f_K`$ has modular reduction of level $`N`$ means essentially that $`\mathrm{\Delta }`$ is a dihedral group of order $`2N`$ (where $`p|N`$) and that $`g_K`$ has ramification of order $`2`$. In particular, $`g_K`$ gives rise to a $`K`$-point on $`X_1(N)`$. This is the link between our results and modular curves.
To explain the construction of $`g_K`$, we assume for simplicity that the normalizer of a $`p`$-Sylow of $`G`$ is of order $`2p`$ and that the branch points of $`f_K`$ do not coalesce modulo $`p`$. Let $`f:YX`$ be the special fiber of the semistable model of $`f_K`$. The curve $`X`$ consists of 5 components: the strict transform of the original component $`X_0`$, and 4 tails $`X_1,\mathrm{},X_4`$ containing the specializations of the branch points $`x_i`$. The cover $`f`$ is inseparable over $`X_0`$ and separable over the tails. Let $`E`$ be a component of $`Y`$ above $`X_0`$. The decomposition group $`\mathrm{\Delta }:=D(E)G`$ is dihedral of order $`2p`$, the inertia group $`I(E)`$ cyclic of order $`p`$ (see Fig. 2 in Section 2). We obtain $`g:ZX`$ by replacing, for $`i=1,\mathrm{},4`$, the (disconnected) $`G`$-cover $`f^1(X_i)X_i`$ by a $`\mathrm{\Delta }`$-cover $`Z_iX_i`$ which is locally, i.e. in an étale neighborhood of $`E`$, isomorphic to $`f^1(X_i)X_i`$ and tamely ramified above $`x_iX_i`$. This is possible in a unique way, by the Katz–Gabber Lemma, . Using formal patching one can show that $`g:ZX`$ is the reduction of a $`\mathrm{\Delta }`$-cover $`g_K:Z_K_K^1`$ which, in some sense, contains all the information about the bad reduction of $`f_K`$.
Let $`HH_4^{\mathrm{in}}(G)`$ be the subset of the Hurwitz space corresponding to $`G`$-covers with $`4`$ branch points and prime-to-$`p`$ ramification. Denote by $`\overline{H}`$ its arithmetic compactification and by $`\overline{H}^{\mathrm{bad}}\overline{H}𝔽_p`$ the subspace corresponding to bad reduction in characteristic $`p`$. The map $`f:YX`$ discussed above corresponds to a $`k`$-point on $`\overline{H}^{\mathrm{bad}}`$; the associated map $`g:ZX`$ corresponds essentially to a $`k`$-point on $`X_1(N)𝔽_p`$. The formal patching argument mentioned above can be used to show that the deformation theory of $`f`$ and $`g`$ are equivalent. This gives a strong connection between the subspace $`\overline{H}^{\mathrm{bad}}\overline{H}`$ and the component of $`X_1(N)𝔽_p`$ corresponding to inseparable isogenies. Using the results of on the reduction of $`X_1(N)`$, we prove our Reduction Theorem, which describes $`\overline{H}𝔽_p`$. It can be roughly stated as follows (see Fig. 1). The subspaces $`\overline{H}^{\mathrm{good}}`$ and $`\overline{H}^{\mathrm{bad}}`$ are smooth curves over $`𝔽_p`$; they intersect transversally in the supersingular points, i.e. the points of $`\overline{H}^{\mathrm{bad}}`$ with supersingular $`\lambda `$-value. The scheme $`\overline{H}𝔽_p`$ is not reduced, in general; the irreducible components of $`\overline{H}^{\mathrm{bad}}`$ have multiplicity $`p1`$ or $`(p1)/2`$. Moreover, each irreducible component of $`\overline{H}^{\mathrm{bad}}`$ is essentially the reduction of a modular curve.
We call the $`\overline{}`$-rational points of $`\overline{H}`$ lying above $`0`$, $`1`$ and $`\mathrm{}`$ the cusps of $`H`$. A cusp corresponds to a degenerate cover $`f_\overline{}:Y_\overline{}X_\overline{}`$, obtained from a smooth $`G`$-cover by coalescing of branch points. The curve $`X_\overline{}`$ is the union of two projective lines intersecting in one point. The cover $`f_\overline{}`$ is ramified above the singular point of $`X_\overline{}`$, say of order $`n`$. Our next result, which we call the Cusp Principle, states that this cusp has bad reduction (i.e. its closure in $`\overline{H}`$ meets $`\overline{H}^{\mathrm{bad}}`$) if and only if $`p|n`$. Essentially, this is an application of Raynaud’s result, since the degenerate cover $`f_\overline{}`$ is build up from two $`3`$ point covers. The Cusp Principle is the key result in our calculation of the number of $`G`$-covers with good reduction. The point here is that, via the Hurwitz classification and the braid action, one can explicitly compute the set of cusps of $`H`$ and decide for each of them whether they have good or bad reduction. In particular, given a finite group $`G`$ and an odd prime $`p`$ verifying all the assumptions made above, we can compute two numbers, $`d`$ and $`d^{\mathrm{bad}}`$, such that
$$|Cov(G,\lambda )^{\mathrm{good}}|=\{\begin{array}{cc}dd^{\mathrm{bad}},\hfill & \text{if }\lambda \text{ is ordinary,}\hfill \\ d\frac{p+1}{p}d^{\mathrm{bad}},\hfill & \text{if }\lambda \text{ is supersingular.}\hfill \end{array}$$
Here $`Cov(G,\lambda )^{\mathrm{good}}`$ is the set of isomorphism classes of $`G`$-covers of $`^1`$ with good reduction, branched in $`0`$, $`1`$, $`\mathrm{}`$ and $`\lambda `$.
Let $`G=PSL_2(\mathrm{})`$, where $`\mathrm{}`$ is an odd prime different from $`p`$ such that $`p`$ exactly divides the order of $`G`$, and consider $`G`$-covers branched in $`4`$ points of order $`\mathrm{}`$. We have computed the cusps of the corresponding Hurwitz spaces, using the computer program ho , for $`\mathrm{}31`$. From this information, we can deduce the complete structure of $`\overline{H}^{\mathrm{bad}}𝔽_p`$ and the number of covers with good reduction.
This paper owes a lot to Raynaud. It started as an attempt to understand a talk he gave in Oberwolfach, June 1997. In this talk Raynaud presented Example 4.3.2. In the problem session of the same conference, he gave a similar problem as an exercise. In a way, this paper is our solution of this exercise. We would also like to thank Bas Edixhoven for a helpful conversation and for sending his manuscript , and Andrew Kresch for comments on an earlier version of Section 1. The second author gratefully acknowledges financial support from the Deutsche Forschungsgemeinschaft.
### Notation
In this paper we will understand by a semistable curve a flat projective morphism $`XS`$ of schemes whose geometric fibers are reduced connected curves having at most ordinary double points as singularities. We write $`X^{\mathrm{sm}}`$ for the subset of smooth points of the morphism $`XS`$. A mark on $`X/S`$ is a closed subscheme $`DX^{\mathrm{sm}}`$ which is finite and étale over $`S`$. The pair $`(X/S,D)`$ is called a marked semistable curve. A stably marked curve is either a pointed stable curve $`(X/S;x_i)`$ in the sense of or a marked semistable curve $`(X/S,D)`$ which becomes a pointed stable curve after an étale base change $`S^{}S`$. By an algebraic stack we mean an algebraic stack in the sense of Deligne–Mumford .
## 1 Complete Hurwitz spaces
The goal of this section is to define arithmetic compactifications of Hurwitz spaces for $`G`$-covers. For a given Hurwitz space $`H`$, such a compactification $`\overline{H}`$ should be a proper model of $`H`$ over $``$ whose points in positive characteristic correspond to the reductions of the covers which are parameterized by $`H`$. To make this precise, one first has to give the definition of the reduction of a $`G`$-cover. This is done in Section 1.1.
Our definition of the complete Hurwitz space for $`G`$-covers essentially follows the approach of . We let $`\overline{}`$ be the closure of the moduli stack of $`G`$-covers inside a bigger moduli stack parameterizing certain maps between stably marked curves. Then we define $`\overline{H}`$ as the coarse moduli space associated to $`\overline{}`$. This is the content of Section 1.2. Section 1.3 discusses a technical problem that arises from our definition. The reader who is not interested in this abstract approach may wish to skip these two sections.
### 1.1 Reduction of $`G`$-covers
In this section we give some terminology and recall some general facts concerning the reduction of $`G`$-covers to positive characteristic. We closely follow , §2. However, our definition of the model of a $`G`$-cover is not exactly the same as Raynaud’s. Also, since we allow the bottom curve to degenerate, we have to consider a slightly more general situation than in , §2.
###### Definition 1.1.1
Let $`K`$ be a field and $`G`$ a finite group. A $`G`$-cover defined over $`K`$ is a finite separable morphism $`f:YX`$ of smooth projective and geometrically irreducible $`K`$-curves together with an isomorphism $`GAut(Y/X)`$ such that $`|G|=\mathrm{deg}f`$. We say that a $`G`$-cover $`f`$ is tame if it is tamely ramified.
Throughout this section, we assume the following situation. Let $`R`$ be a complete discrete valuation ring with quotient field $`K`$ of characteristic zero, and residue field $`k=\overline{k}`$ of characteristic $`p>0`$. Let $`f_K:Y_KX_K`$ be a $`G`$-cover defined over $`K`$ ($`f_K`$ is automatically tame). Write $`x_{1,K},\mathrm{},x_{r,K}X_K(\overline{K})`$ for the branch points and $`y_{1,K},\mathrm{},y_{n,K}Y_K(\overline{K})`$ for the ramification points of $`f_K`$. We assume that $`2g+r3`$, where $`g`$ is the genus of $`X_K`$.
We would like to define a model of the $`G`$-cover $`f_K`$ over the ring $`R`$. After replacing $`K`$ by a finite extension $`K^{}/K`$ and $`R`$ by its integral closure in $`K^{}`$, we may assume that the ramification points $`y_{i,K}`$ of $`f_K`$ are $`K`$-rational. After a further extension of $`K`$, we may assume that the smooth stably marked curve $`(Y_K;y_{i,K})`$ extends to a stably marked curve $`(Y_R;y_{i,R})`$ over $`R`$, . In particular, $`Y_R`$ is semistable over $`R`$ and the points $`y_{i,K}`$ specialize to pairwise distinct, smooth points $`y_i`$ on the special fiber $`Y`$ of $`Y_R`$. Since the stably marked model is unique, the action of $`G`$ on $`Y_K`$ extends to $`Y_R`$. Let $`X_R:=Y_R/G`$ be the quotient scheme and $`X`$ the special fiber of $`X_R`$. By , Appendice, $`X_R`$ is again a semistable curve over $`R`$. Since the ramification points $`y_{i,K}`$ specialize to pairwise distinct smooth points on $`Y`$, the branch points $`x_{j,K}`$ specialize to pairwise distinct smooth points $`x_jX`$. According to , the stably marked curve $`(X_K;x_{j,K})`$ over $`K`$ extends to a stably marked curve $`(X_{0,R};x_{j,R}^{})`$ over $`R`$, and we have a well defined contraction morphism $`X_RX_{0,R}`$ sending $`x_{j,R}`$ to $`x_{j,R}^{}`$. Let $`f_R:Y_RX_R`$ and $`f_{0,R}:Y_RX_{0,R}`$ be the natural maps and $`f:YX`$, $`f_0:YX_0`$ the induced maps on the special fibers.
###### Definition 1.1.2
Let $`f_R:Y_RX_R`$ and $`f_{0,R}:Y_RX_{0,R}`$ be as above. We call $`f_R:Y_RX_R`$ the quotient model and $`f_{0,R}:Y_RX_{0,R}`$ the stable model over $`R`$ of the $`G`$-cover $`f_K`$.
The quotient model and the stable model of $`f_K`$ exist after a finite extension of $`K`$. It is clear that these models are stable with respect to any further extension of $`K`$. In many places it is better to work with the quotient model, because it is a finite map. However, the stable model is easier to study the moduli of. Therefore, our definition of a complete Hurwitz space will be based on the stable model.
###### Definition 1.1.3
1. The $`G`$-cover $`f_K`$ has (potentially) good reduction if (after a finite extension of $`K`$) the special fiber $`f:YX`$ of the quotient model of $`f_K`$ is a tame $`G`$-cover.
2. The $`G`$-cover $`f_K`$ has (potentially) admissible reduction if (after a finite extension of $`K`$) the special fiber $`f:YX`$ of the quotient model of $`f_K`$ is a tame admissible cover. By this we mean that $`f`$ is finite, separable, tamely ramified over the smooth locus $`X^{\mathrm{sm}}`$ and has tame admissible ramification over the ordinary double points of $`X`$, see , §4, or , where such a cover is called kummérien.
3. The $`G`$-cover $`f_K`$ has bad reduction if it does not have potentially admissible reduction.
For the rest of this subsection we will omit the word “potentially” and assume that $`K`$ is chosen such that the quotient and stable model exist over $`R`$. Note that $`f_K`$ has good reduction if and only if it has admissible reduction and $`X_R`$ is smooth over $`R`$. If $`f_K`$ has admissible reduction then $`X_R=X_{0,R}`$, i.e. quotient and stable model are the same.
Let $`W`$ be a component of $`X`$. We call $`W`$ an original component if it is the strict transform of a component of $`X_0`$ (otherwise, the map $`XX_0`$ contracts $`W`$). We will say that $`f`$ is separable over $`W`$ if for one (and therefore for all) components $`Z`$ of $`Y`$ above $`W`$ the restriction $`f|_Z:ZW`$ is a separable morphism. Equivalently, the inertia group of $`Z`$ (the group of elements of $`G`$ acting trivially on $`Z`$) is trivial. Proposition 1.1.4 below extends , Corollaire 2.4.9, to our situation, which includes admissible reduction. The proof is essentially the same, with , Théorème 3.2, as additional ingredient.
###### Proposition 1.1.4
The $`G`$-cover $`f_K`$ has admissible reduction if and only if $`f`$ is separable over the original components of $`X`$.
From Proposition 1.1.4 we can deduce the following well known fact.
###### Corollary 1.1.5
Assume that the order of $`G`$ is prime to the characteristic of $`k`$. Then $`f_K`$ has potentially admissible reduction. If in addition the branch points $`x_{j,K}`$ of $`f_K`$ specialize to pairwise distinct points on the special fiber of a smooth model of $`X_K`$, then $`f_K`$ has potentially good reduction.
The following example plays a central role in this paper.
###### Example 1.1.6
Let $`R`$ and $`K`$ be as before. Choose four $`R`$-rational points $`x_1,\mathrm{},x_4`$ on $`^1`$ such that $`x_ix_j(𝔪)`$, where $`𝔪R`$ is the maximal ideal. Then $`(_R^1;x_i)`$ is a smooth, stably marked curve over $`R`$. Let $`G`$ be the dihedral group of order $`2p`$, where $`p`$ is an odd prime, equal to the residue characteristic of $`R`$. Let $`f_K:E_K_K^1`$ be a $`G`$-cover branched only in the $`4`$ points $`x_i`$, with ramification of order $`2`$. We have $`4p`$ ramification points $`y_{i,j}`$ on $`E_K`$, where $`1i4`$, $`1jp`$ and $`f_K(y_{i,j})=x_i`$. The curve $`E_K`$ has genus $`1`$. After extending $`K`$ we may assume that the $`y_{i,j}`$ are $`K`$-rational. Choosing e.g. $`y_{1,1}`$ as the origin gives $`E_K`$ the structure of an elliptic curve. Moreover, $`f_K`$ can be written as the composition
$$f_K:E_K\stackrel{p}{}E_K^{}\stackrel{2}{}_K^1$$
of a $`p`$-cyclic isogeny of elliptic curves and a cyclic cover of degree $`2`$. The cover $`f_K`$ extends to a finite flat morphism $`f_R:E_R_R^1`$. Moreover, $`E_R`$ is an elliptic curve and $`f_R`$ factors through an isogeny $`\pi _R:E_RE_R^{}`$.
There are two cases to consider. First, $`\pi _R`$ might be étale. In this case, the ramification points $`y_{i,j}`$ extend to disjoint sections $`y_{i,j}:SpecRE_R`$ and $`f:E_R_R^1`$ is tamely ramified along the sections $`x_i:SpecR_R^1`$. In other words, $`f_K`$ has good reduction. Now assume that $`\pi _R`$ is not étale. Then its restriction $`\pi :EE^{}`$ to the special fiber is purely inseparable, and for fixed $`i`$ the $`p`$ points $`y_{i,j}`$, $`j=1,\mathrm{},p`$ specialize to the same point of $`E`$. We see that $`(E_R;y_{i,j})`$ is not a stably marked curve. Let $`(Z_R;y_{i,j})`$ be the extension of $`(E_K;y_{i,j})`$ to a stably marked curve over $`R`$ and $`q_R:Z_RE_R`$ the contraction morphism. We can identify $`E`$ with its strict transform in $`Z_R`$. The special fiber $`Z`$ of $`Z_R`$ has exactly $`5`$ components $`E,Z_1,\mathrm{},Z_4`$. For $`i=1,\mathrm{},4`$, the curve $`Z_i`$ is smooth and of genus $`0`$, connected to $`E`$ in one point and contains the specialization of the points $`y_{i,j}`$ for $`j=1,\mathrm{},p`$. Let $`X_R:=Z_R/G`$ be the quotient. The special fiber $`X:=X_R_Rk`$ has $`5`$ components $`X_0,\mathrm{},X_4`$. In fact, $`X_0`$ is the original component, and for $`i=1,\mathrm{},4`$, the component $`X_i`$ is the image of $`Z_i`$ and contains the specialization of $`x_i`$. The restriction of $`f:ZX`$ to $`X_i`$ is a $`G`$-cover $`Z_iX_i`$ ramified in two points, with ramification of order $`2`$ and $`2p`$ (so it is not tame).
### 1.2 Complete Hurwitz stacks
In this section we define the concept of a complete Hurwitz stack, following the idea of . Since there are many different versions of Hurwitz stacks, we do this first in detail for one specific kind, namely for $`_{[r]}^{\mathrm{in}}(G)`$, the inner Hurwitz stack for $`G`$-covers of genus $`0`$ curves with unordered branch points. Then we define several variants of the above. In Section 1.2.4 we look at Hurwitz spaces as coarse moduli spaces.
In this paper we only consider Hurwitz spaces for Galois covers. Moreover, the target curve of the cover will always be of genus $`0`$ and is considered “up to isomorphism”. Thus, we only consider “reduced” Hurwitz spaces in the terminology used by Fried . The genus zero assumption is made only to simplify the notation. It is easy to extend all our definitions to nonreduced Hurwitz spaces. It seems much less trivial to do the same for Hurwitz spaces parameterizing non-Galois covers. For instance, in a completion of the classical Hurwitz stack for simple covers is constructed, using the moduli space of stable maps as an ambient space. This construction is more involved than the one we give here. Another problem, discussed in , Section 4.1, is that “taking quotients by finite groups does not commute with base change”. Therefore, the method proposed in of studying complete moduli of non-Galois covers by going to the Galois closure probably does not work very well with our definition of complete Hurwitz stacks. A related problem is discussed in Section 1.3.
#### 1.2.1 The ambient stack
Let $`(X/S,C)`$ and $`(Y/S,D)`$ be stably marked curves, defined over the same scheme $`S`$. A morphism of stably marked curves is an $`S`$-morphism $`f:YX`$ such that $`f(D)C`$. To ease notation, we will usually write $`X`$ and $`Y`$ instead of $`(X/S,C)`$ and $`(Y/S,D)`$. We fix an integer $`r>0`$ and let $`𝒮_{\left[r\right]}`$ be the following category. Objects of $`𝒮_{\left[r\right]}`$ are morphisms $`f:YX`$ between stably marked curves such that $`X`$ has genus $`0`$ and is stably $`r`$-marked. A morphism between an $`S`$-object $`f:(Y,D)(X,C)`$ and an $`S^{}`$-object $`f^{}:(Y^{},D^{})(X^{},C^{})`$ of $`𝒮_{\left[r\right]}`$ consists of a Cartesian diagram
(1)
$$\begin{array}{ccc}Y^{}& & Y\\ f^{}& & f& & \\ X^{}& & X\\ & & & & \\ S^{}& & S\end{array}$$
such that $`D^{}=D\times _SS^{}`$ and $`C^{}=C\times _SS^{}`$. It is clear that $`𝒮_{\left[r\right]}`$ is a stack.
Let $`G`$ be a finite group and $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$ the following category. Objects of $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$ (over a scheme $`S`$) are pairs $`(f,\sigma )`$, where $`f:YX`$ is an object of $`𝒮_{\left[r\right]}`$ defined over $`S`$ and $`\sigma :GAut(Y/X)`$ is an action of $`G`$ on $`Y`$ commuting with $`f`$ such that the induced action on every geometric fiber of $`f`$ is faithful (equivalently, $`\sigma `$ induces a closed immersion $`\stackrel{~}{\sigma }:G_S\underset{¯}{\mathrm{Aut}}(Y/X)`$ of group schemes). Mostly we will omit the map $`\sigma `$ and simply write $`f:YX`$ for an object of $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$. A morphism between two objects $`f^{}:Y^{}X^{}`$ and $`f:YX`$ of $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$ is an $`𝒮_{\left[r\right]}`$-morphism (1) such that the top arrow $`Y^{}Y`$ is $`G`$-equivariant. Again it is clear that $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$ is a stack.
###### Proposition 1.2.1
The stacks $`𝒮_{\left[r\right]}`$ and $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$ are algebraic, separated and locally of finite type over $``$.
###### Proof.
The stack $`\overline{}_{g,[n]}`$ classifying stably $`n`$-marked curves of fixed genus $`g`$ is algebraic, separated and of finite type over $``$. A standard Hilbert scheme argument (see e.g. , Chap. 0.5) shows that $`𝒮_{\left[r\right]}`$ is algebraic, separated and locally of finite type over $``$. Let $`f:YX`$ be an object of $`𝒮_{\left[r\right]}`$ defined over a scheme $`S`$. Since $`\underset{¯}{\mathrm{Aut}}(Y/X)`$ is a finite $`S`$-group scheme, the functor $`Hom_S(G_S,\underset{¯}{\mathrm{Aut}}(Y/X))`$ is represented by a finite $`S`$-scheme. It is clear that the natural morphism $`S\times _{𝒮_{[r]}}𝒮_{\left[r\right]}^{\mathrm{in}}(G)Hom_S(G_S,Aut(Y/X))`$ is a locally closed immersion. Therefore, the forgetful morphism $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)𝒮_{\left[r\right]}`$ is relatively representable, separated and of finite type. This completes the proof of the proposition. $`\mathrm{}`$
#### 1.2.2 The complete inner Hurwitz stack
Let $`r`$ and $`G`$ be as before. We define the Hurwitz stack $`_{\left[r\right]}^{\mathrm{in}}(G)`$ as follows. Objects of $`_{\left[r\right]}^{\mathrm{in}}(G)`$ over a scheme $`S`$ are morphisms $`f:YX`$ between smooth $`S`$-curves, together with an action of $`G`$ on $`Y`$, commuting with $`f`$, such that the following holds. The curve $`X/S`$ has genus $`0`$ and the geometric fibers of $`f`$ are tame $`G`$-covers (see Definition 1.1.1) with exactly $`r`$ branch points. Morphisms in $`_{\left[r\right]}^{\mathrm{in}}(G)`$ are Cartesian diagrams of the form (1) such that the top horizontal arrow is $`G`$-equivariant. We call an object $`f:YX`$ of $`_{\left[r\right]}^{\mathrm{in}}(G)`$ a tame $`G`$-cover, defined over $`S`$. It is proved e.g. in that $`_{\left[r\right]}^{\mathrm{in}}(G)`$ is an algebraic stack, smooth and of finite type over $``$.
Let $`f:YX`$ be an $`S`$-object of $`_{\left[r\right]}^{\mathrm{in}}(G)`$. Then $`f`$ is finite and tamely ramified along a divisor $`CX`$ which is finite étale of degree $`r`$ over $`S`$. Hence $`(X/S,C)`$ is a (smooth) stably $`r`$-marked curve. Moreover, $`(Y/S,D)`$ is a (smooth) stably marked curve, where $`D:=f^1(C)Y`$ is the (reduced) inverse image of $`C`$. Therefore, we obtain a natural monomorphism
(2)
$$_{\left[r\right]}^{\mathrm{in}}(G)𝒮_{\left[r\right]}^{\mathrm{in}}(G),$$
identifying $`_{\left[r\right]}^{\mathrm{in}}(G)`$ with a full subcategory of $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$. We will show in Proposition 1.2.4 below that (2) is a locally closed immersion. Note however that we do not need this fact to make the following definition.
###### Definition 1.2.2
The complete Hurwitz stack $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ is the closure of $`_{\left[r\right]}^{\mathrm{in}}(G)`$ inside $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$ (i.e. the smallest closed substack of $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$ containing $`_{\left[r\right]}^{\mathrm{in}}(G)`$ as a full subcategory).
Proposition 1.2.1 shows that $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ is an algebraic stack, separated and locally of finite type over $``$. Let $`k`$ be an algebraically closed field and $`f:YX`$ an object of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ defined over $`k`$. Choose an étale neighborhood $`U\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ of the point $`s:Speck\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ corresponding to $`f`$ and let $`\eta :SpecKU`$ be a generic point of some irreducible component of $`U`$. By , Exercise II.4.11, $`\eta `$ extends to a morphism $`\eta _R:SpecRU`$, where $`R`$ is a discrete valuation ring of $`K`$ with residue field $`k`$, and the restriction of $`\eta _R`$ to the special point is equal to $`s`$. The morphism $`SpecR\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ corresponds to an $`R`$-object $`f_R:Y_RX_R`$ of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ with special fiber $`f:YX`$. Since $`_{\left[r\right]}^{\mathrm{in}}(G)`$ is dense in $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$, the generic fiber $`f_K:Y_KX_K`$ of $`f_R`$ is actually an object of $`_{\left[r\right]}^{\mathrm{in}}(G)`$, i.e. a tame $`G`$-cover. Moreover, $`K`$ has characteristic $`0`$. We are essentially (modulo taking the completion of $`R`$) in the situation of Section 1.1. It is clear that $`f_R:Y_RX_R`$ is the stable model of the $`G`$-cover $`f_K`$. In particular, $`f_K`$ has good reduction if and only if $`f`$ is an object of $`_{\left[r\right]}^{\mathrm{in}}(G)`$. We say that $`f`$ is a bad cover if $`f_K`$ has bad reduction (Definition 1.1.3). By Proposition 1.1.4, $`f`$ is a bad cover if and only if some irreducible component of $`Y`$ has a nontrivial inertia group (with respect to the action of $`G`$).
###### Lemma 1.2.3
There is a unique closed reduced substack $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{bad}}\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ characterized by the following property. For an algebraically closed field $`k`$, a $`k`$-object $`f:YX`$ of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ is a bad cover if and only if it is an object of the substack $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{bad}}`$.
###### Proof.
By , Lemma II.4.5, it suffices to show that the subset of bad covers is stable under specialization. More precisely, let $`R`$ be a discrete valuation ring and $`f_R:Y_RX_R`$ an object of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ defined over $`R`$. Assume that the generic fiber $`f_K:Y_KX_K`$ of $`f_R`$ is a bad cover. We have to show that the special fiber $`f:YX`$ is a bad cover, too. As remarked above, $`f_K`$ (resp. $`f`$) is a bad cover if and only if some irreducible component of $`Y_K`$ (resp. $`Y`$) has a nontrivial inertia group. Clearly, this property is stable under specialization. $`\mathrm{}`$
By the lemma, $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{adm}}:=\overline{}_{\left[r\right]}^{\mathrm{in}}(G)\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{bad}}`$ is a dense open substack of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$. The discussion before Lemma 1.2.3 shows that the geometric points of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{adm}}`$ correspond to tame admissible covers $`f:YX`$ over $`k`$ which arise as the reduction of tame $`G`$-covers. It follows that an object $`f:YX`$ of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{adm}}`$ (defined over an arbitrary scheme $`S`$) lies in the full subcategory $`_{\left[r\right]}^{\mathrm{in}}(G)`$ if and only if $`X/S`$ is smooth. Since smoothness of $`X/S`$ is an open condition on $`S`$, $`_{\left[r\right]}^{\mathrm{in}}(G)`$ is an open substack of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{adm}}`$, and hence an open substack of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$.
###### Proposition 1.2.4
The algebraic stack $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ is reduced, proper and of finite type over $``$ and contains $`_{\left[r\right]}^{\mathrm{in}}(G)`$ and $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{adm}}`$ as dense open substacks.
###### Proof.
It only remains to show that $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ is reduced and proper. We know that the dense open substack $`_{\left[r\right]}^{\mathrm{in}}(G)`$ is reduced, therefore $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ is reduced as well. To prove properness, we apply , Corollaire 7.3.10. We are immediately reduced to the following situation. Let $`R`$ be a complete discrete valuation ring with residue field $`k=\overline{k}`$ of characteristic $`p`$ and quotient field $`K`$ of characteristic $`0`$. Let $`f_K:Y_KX_K`$ be an object of $`_{\left[r\right]}^{\mathrm{in}}(G)`$, i.e. a $`G`$-cover. This is the situation of Section 1.1. After a finite extension of $`K`$, $`f_K`$ has a stable model $`f_R:Y_RX_R`$ over $`R`$ (Definition 1.1.2). Obviously, $`f_R`$ is an object of $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$. Therefore, $`f_R`$ is an object of $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$, by Definition 1.2.2. This proves that $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ is proper. $`\mathrm{}`$
The stack $`_{\left[r\right]}^{\mathrm{in}}(G)`$ is called the inner Hurwitz stack for $`G`$-covers with $`r`$ (unordered) branch points. We will say that the stack $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$ is the ambient stack of $`_{\left[r\right]}^{\mathrm{in}}(G)`$. We will call the stack $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)`$ the completion of $`_{\left[r\right]}^{\mathrm{in}}(G)`$.
###### Remark 1.2.5
It is easy to check that the stack $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{adm}}`$ can be identified with the compactification of $`_{\left[r\right]}^{\mathrm{in}}(G)`$ constructed in or . It follows from loc.cit. that $`\overline{}_{\left[r\right]}^{\mathrm{in}}(G)^{\mathrm{adm}}`$ is smooth over $``$ and proper over $`[1/|G|]`$.
###### Variant 1.2.6
Let $`_{\left[r\right]}^{\mathrm{ab}}(G)`$ be the stack whose objects are $`S`$-morphisms $`f:YX`$ which are locally on $`S`$ tame $`G`$-covers. More precisely, after an étale localization $`S^{}S`$, there exists an action of $`G`$ on $`Y`$ such that $`f:YX`$ becomes an object of $`_{\left[r\right]}^{\mathrm{in}}(G)`$. We call $`_{\left[r\right]}^{\mathrm{ab}}(G)`$ the absolute Hurwitz stack for $`G`$-covers with $`r`$ (unordered) branch points, see also . We embed $`_{\left[r\right]}^{\mathrm{ab}}(G)`$ into an ambient stack $`𝒮_{\left[r\right]}^{\mathrm{ab}}(G)`$. Objects of $`𝒮_{\left[r\right]}^{\mathrm{ab}}(G)`$ are pairs $`(f,𝒢)`$, where $`f:YX`$ is an object of $`𝒮_{\left[r\right]}`$ defined over a scheme $`S`$ and $`𝒢\underset{¯}{\mathrm{Aut}}(Y/X)`$ is an étale subgroup scheme which becomes isomorphic to the constant $`S`$-group scheme $`G_S`$ after an étale localization of $`S`$. We define the completion $`\overline{}_{\left[r\right]}^{\mathrm{ab}}(G)`$ as the closure of $`_{\left[r\right]}^{\mathrm{ab}}(G)`$ inside $`𝒮_{\left[r\right]}^{\mathrm{ab}}(G)`$.
###### Variant 1.2.7
Let $`_r^{\mathrm{in}}(G)`$ and $`_r^{\mathrm{ab}}(G)`$ be the inner resp. absolute Hurwitz stack for $`G`$-covers with $`r`$ ordered branch points. We can embed $`_r^{\mathrm{in}}(G)`$ into an ambient stack $`𝒮_r^{\mathrm{in}}(G)`$. Objects of $`𝒮_r^{\mathrm{in}}(G)`$ are objects $`f:YX`$ of $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$ together with $`r`$ sections $`x_1,\mathrm{},x_r:SX`$ such that $`C=_ix_i(S)`$ is the mark of the stably marked curve $`X`$. We define the completion $`\overline{}_r^{\mathrm{in}}(G)`$ as the closure of $`_r^{\mathrm{in}}(G)`$ inside $`𝒮_r^{\mathrm{in}}(G)`$. Similar for $`\overline{}_r^{\mathrm{ab}}(G)`$.
We obtain a diagram
(3)
$$\begin{array}{ccc}\overline{}_r^{\mathrm{in}}(G)& & \overline{}_{\left[r\right]}^{\mathrm{in}}(G)\\ & & & & \\ \overline{}_r^{\mathrm{ab}}(G)& & \overline{}_{\left[r\right]}^{\mathrm{ab}}(G).\end{array}$$
The vertical arrows in (3) are principal $`Aut(G)`$-bundles, the horizontal arrows are principal $`𝔖_r`$-bundles. All stacks in (3) are algebraic, reduced and proper and of finite type over $``$.
#### 1.2.3 Complete Hurwitz stacks for a given type
With $`G`$ and $`r`$ as before, let $`\underset{¯}{C}=(C_1,\mathrm{},C_r)`$ be an $`r`$-tuple of conjugacy classes of $`G`$. We denote by $`(\underset{¯}{C})(\zeta _n)`$ the field generated by the values $`\chi (\tau )`$, where $`\chi `$ runs over all irreducible characters of $`G`$ and $`\tau C_i`$, for $`i=1,\mathrm{},r`$. In other words, $`(\underset{¯}{C})`$ is the minimal number field over which every class $`C_i`$ becomes rational. Let us make a “choice of $`n`$th root of unity over $`(\underset{¯}{C})`$”, i.e. we choose an orbit under $`Gal(\overline{}/(\underset{¯}{C}))`$ of primitive $`n`$th roots of unity, see , Section 8.2.1. Let $`\mathrm{\Lambda }(\underset{¯}{C})`$ be a Dedekind domain with fraction field $`(\underset{¯}{C})`$. We define an open substack
(4)
$$_r^{\mathrm{in}}(\underset{¯}{C})_\mathrm{\Lambda }_r^{\mathrm{in}}(G)_{}\mathrm{\Lambda },$$
corresponding to $`G`$-covers with inertia type $`\underset{¯}{C}`$. To be more precise, let $`f:Y(X;x_i)`$ be an object of $`_r^{\mathrm{in}}(G)\mathrm{\Lambda }`$, defined over an algebraically closed field $`k`$. We say that $`f`$ has inertia type $`\underset{¯}{C}`$ if $`C_i`$ is the conjugacy class associated to the branch point $`x_i`$ (with respect to a canonical choice of $`m_i`$th root of unity, induced by the natural map $`\mathrm{\Lambda }k`$), compare , Section 2.2.1. An object of $`_r^{\mathrm{in}}(G)\mathrm{\Lambda }`$, defined over an arbitrary scheme $`S`$, is said to have inertia type $`\underset{¯}{C}`$ if all its geometric fibers have inertia type $`\underset{¯}{C}`$. We define the completion $`\overline{}_r^{\mathrm{in}}(\underset{¯}{C})_\mathrm{\Lambda }`$ as the closure of $`_r^{\mathrm{in}}(\underset{¯}{C})_\mathrm{\Lambda }`$ inside $`𝒮_r^{\mathrm{in}}(G)_{}\mathrm{\Lambda }`$. Clearly, $`\overline{}_r^{\mathrm{in}}(𝒞)_\mathrm{\Lambda }`$ is an algebraic stack, proper and of finite type over $`\mathrm{\Lambda }`$. If the choice of $`\mathrm{\Lambda }`$ is understood, we will omit it from the notation, and we say that $`\overline{}_r^{\mathrm{in}}(\underset{¯}{C})`$ is defined over $`\mathrm{\Lambda }`$, or that $`\mathrm{\Lambda }`$ is the domain of definition of $`\overline{}_r^{\mathrm{in}}(\underset{¯}{C})`$.
In a similar way, we can define further variants of complete Hurwitz stacks: $`\overline{}_r^{\mathrm{ab}}(\underset{¯}{C})`$, $`\overline{}_{[r]}^{\mathrm{in}}(\underset{¯}{C})`$ and $`\overline{}_{[r]}^{\mathrm{ab}}(\underset{¯}{C})`$. The domains of definition of these stacks are Dedekind domains whose fraction fields are suitable subfields of $`(\underset{¯}{C})`$. For instance, $`\overline{}_{[r]}^{\mathrm{in}}(\underset{¯}{C})`$ can be defined over the ring of integers of the smallest field over which $`\underset{¯}{C}`$ as a tuple is rational, compare , Definition 3.15.
#### 1.2.4 Complete Hurwitz spaces as coarse moduli spaces
Let $`\overline{}:=\overline{}_r^{\mathrm{in}}(\underset{¯}{C})`$ be the complete Hurwitz stack over $`\mathrm{\Lambda }`$, as defined in Section 1.2.3. We denote by $`\overline{H}:=\overline{H}_r^{\mathrm{in}}(\underset{¯}{C})`$ the associated coarse moduli space, and call it the complete Hurwitz space over $`\mathrm{\Lambda }`$ for $`G`$-covers with inertia type $`\underset{¯}{C}`$. By construction , $`\overline{H}`$ is an algebraic space, proper and of finite type over $`\mathrm{\Lambda }`$, and contains the usual Hurwitz space $`H=H_r^{\mathrm{in}}(\underset{¯}{C})`$ as a dense open subspace. Actually, $`H`$ is a scheme and is smooth over $`\mathrm{\Lambda }`$, see . We define a closed subspace $`\overline{H}^{\mathrm{bad}}`$ as the image of the natural morphism $`\overline{}^{\mathrm{bad}}\overline{H}`$. Thus, $`\overline{H}^{\mathrm{adm}}:=\overline{H}\overline{H}^{\mathrm{bad}}`$ is the coarse moduli space associated to $`\overline{}^{\mathrm{adm}}=\overline{}\overline{}^{\mathrm{bad}}`$ and contains $`H`$ as a dense open subset. According to , $`\overline{H}^{\mathrm{adm}}`$ is a normal scheme.
Let $`\overline{}\overline{}_{0,r}`$ be the natural forgetful morphism. It is known that $`\overline{}_{0,r}`$ is represented by a smooth projective scheme over $``$, see . By the universal property of the coarse moduli space, we obtain a morphism $`\overline{H}\overline{}_{0,r}`$.
###### Proposition 1.2.8
Assume that the complete Hurwitz stack $`\overline{}`$ is normal and that the forgetful morphism $`\overline{}\overline{}_{0,r}\mathrm{\Lambda }`$ is relatively representable and finite. Then the complete Hurwitz space $`\overline{H}`$ is a normal scheme, finite over $`\overline{}_{0,r}\mathrm{\Lambda }`$.
###### Proof.
If $`\overline{}`$ is normal and $`\overline{}\overline{}_{0,r}\mathrm{\Lambda }`$ finite, then the coarse moduli space $`\overline{H}`$ is the normalization of $`\overline{}_{0,r}\mathrm{\Lambda }`$ in $`H`$, see , Proposition IV.3.10. This proves the proposition. $`\mathrm{}`$
###### Variant 1.2.9
In the same manner, we define complete Hurwitz spaces $`\overline{H}_{\left[r\right]}^{\mathrm{in}}(G)`$, $`\overline{H}_r^{\mathrm{ab}}(G)`$, etc. Proposition 1.2.8 applies as well. There are natural maps between all these variants. For instance, diagram (3) induces an analogous diagram of finite morphisms between algebraic spaces. But unlike the maps in (3), these maps are in general not étale.
### 1.3 Quotient model versus stable model
Let us fix a finite group $`G`$ and an integer $`r0`$. Let $`\overline{}:=\overline{}_r^{\mathrm{in}}(G)`$ be the complete inner Hurwitz stack defined in Section 1.2. In Section 1.1 we have defined two different models of a $`G`$-cover $`f_K:Y_KX_K`$, where $`K`$ is a complete discrete valued field: the quotient model $`f_R:Y_RX_R`$ and the stable model $`f_{0,R}:Y_RX_{0,R}`$. The definition of $`\overline{}`$ was made such that the stable model $`f_{0,R}:Y_RX_{0,R}`$ of the $`G`$-cover $`f_K`$ is an object of $`\overline{}`$. One drawback of the stable model is that it is in general not a finite map. Over the discrete valuation ring $`R`$, one can recover the quotient model from the stable model by taking the quotient scheme $`X_R:=Y_R/G`$. The problem is that taking quotients does not commute with arbitrary base change. So it is not clear how to define a quotient model of an object $`f_0:YX_0`$ of $`\overline{}`$ over an arbitrary scheme $`S`$. In this section we propose a definition that works well when $`S`$ is either the spectrum of an algebraically closed field or a normal scheme $`S`$ with function field of characteristic $`0`$. This will be enough for our purposes.
###### Definition 1.3.1
Let $`S`$ be a scheme and $`f_0:YX_0`$ an object of $`\overline{}`$ defined over $`S`$. Let $`f:YX`$ be a finite morphism between marked semistable curves commuting with the action of $`G`$ on $`Y`$. We say that $`f`$ is a quotient model of $`f_0`$ if $`f_0`$ is the composition of $`f`$ with an $`S`$-morphism $`XX_0`$ and if for every geometric fiber $`f_s:Y_sX_s`$ of $`f`$ the following holds.
1. The natural morphism $`Y_s/GX_s`$ induces a bijection on geometric points.
2. Let $`Y_i`$ be a component of $`Y_s`$ and $`X_i`$ the component of $`X_s`$ under $`Y_i`$. Then the degree of $`Y_i`$ over $`X_i`$ is equal to the order of the stabilizer $`D(Y_i)G`$ of $`Y_i`$.
###### Proposition 1.3.2
Let $`S`$ be either the spectrum of an algebraically closed field $`k`$ or a normal scheme, generically of characteristic $`0`$. Let $`f_0:YX_0`$ be an object of $`\overline{}`$ over $`S`$. Then there exists a unique quotient model $`f:YX`$ of $`f_0`$.
###### Proof.
Let us first assume that $`S`$ is a normal scheme, generically of characteristic $`0`$. We may assume that $`S=SpecR`$ is local, with algebraically closed residue field $`k`$. The generic fiber $`f_K:Y_KX_K`$ is an admissible cover. In particular, $`f_K`$ is a quotient model of itself. Let $`X:=Y/G`$. It follows from , Proposition 4.2, and , Corollary A.7.2.2, that $`f:YX`$ is a quotient model of $`f_0`$. To show that it is unique, suppose we have another quotient model $`f^{}:YX^{}`$ of $`f_0`$. Since $`f^{}`$ commutes with the action of $`G`$ on $`Y`$, there exists a morphism $`\kappa :XX^{}`$ of $`S`$-schemes such that $`\kappa f=f^{}`$. We claim that $`\kappa `$ is an isomorphism of $`X_0`$-schemes. Since $`X`$ and $`X^{}`$ are flat over $`S`$ and $`\kappa `$ is the identity on the generic fiber, it suffices to prove that the restriction of $`\kappa `$ to the special fiber is an isomorphism of $`X_0`$-schemes. Hence we have reduced the proposition to the case $`S=Speck`$. Let us now prove this case. We have seen before that every $`k`$-object $`f_0:YX_0`$ of $`\overline{}`$ is the reduction of a $`G`$-cover $`f_K:Y_KX_K`$, where $`K`$ is the quotient field of a discrete valuation ring $`R`$ with residue field $`k`$. Therefore, the special fiber $`f:YX`$ of the quotient model of $`f_K`$ is a quotient model of $`f_0`$. It remains to prove its uniqueness. The quotient $`X^{}:=Y/G`$ is a semistable curve over $`k`$ and the natural map $`X^{}X`$ is a bijection on geometric points, by assumption. Let $`Y_i`$ be a component of $`Y`$ and $`X_i`$ resp. $`X_i^{}`$ the component of $`X`$ resp. of $`X^{}`$ under $`Y_i`$. It follows from , Proposition IV.2.5, that $`X_i`$ is the $`n`$th Frobenius twist $`(X_i^{})^{F^n}`$ of $`X_i^{}`$, where $`p^n`$ is the order of the inertia group $`I(Y_i)G`$ of $`Y_i`$. Clearly, this characterizes the map $`X^{}X`$ and hence the quotient model $`f:YX`$ uniquely. $`\mathrm{}`$
## 2 Semistable reduction
In this section the notations and conventions are as in Section 1.1. Let us recall them briefly. Let $`R`$ be a complete discrete valuation ring with quotient field $`K`$ of characteristic zero and residue field $`k=\overline{k}`$ of characteristic $`p`$. Let $`f_K:Y_KX_K`$ be a $`G`$-cover of smooth curves defined over $`K`$. Let $`g=g(X_K)`$ and $`r`$ the number of branch points of $`f_K`$. Let $`f_R:Y_RX_R`$ be the quotient model of $`f_K`$ defined in Definition 1.1.2. Write $`f:YX`$ for its special fiber. The branch points of $`f_K`$ specialize to distinct points $`x_1,\mathrm{},x_r`$ of the smooth locus of $`X`$. Let $`m_i`$ be the ramification index of $`x_i`$ in $`f_K`$. Recall that we also defined a different model for $`f_K`$, called the stable model $`f_{0,R}:Y_RX_{0,R}`$. It is obtained from $`f_R`$ by composing with a contraction map $`X_RX_{0,R}`$.
We will determine the structure of the reduction $`f:YX`$ in case the cover $`f_K`$ has bad reduction, under suitable assumptions (Condition 2.2.2 below). All results of this section rely on the results of and . Results from these papers are recalled very briefly and the reader is referred to these papers for more details.
Condition 2.2.2 plays an essential role in the rest of the paper. It will enable us to compute the structure of $`f:YX`$. Most importantly, we will assume that the normalizer of a $`p`$-Sylow group $`P`$ of $`G`$ is a dihedral group. We will define an auxiliary cover $`g:ZX`$. The condition on the normalizer $`N_G(P)`$ will imply that this auxiliary cover is equivariant under a (dihedral) subgroup of $`N_G(P)`$. In Section 3 we will relate the deformation theory of $`f`$ to the deformation theory of $`g`$.
In Section 2.1 we will recall some results on inertia and decomposition groups of points and components of $`Y`$. In Section 2.2 a technical lemma is proved. This enables us to extend some constructions and notations from to the case that $`X_0`$ is not smooth in Section 2.3. In Section 2.4 we show that $`f:YX`$ has a fairly simple structure, which we call modular reduction. In Section 2.5 we study the reduction of the degenerate covers in characteristic zero, i.e. the covers corresponding to cusps of the Hurwitz space.
### 2.1 Inertia and decomposition groups
###### Lemma 2.1.1
A singular point of $`Y`$ maps to a singular point of $`X`$, i.e. $`G`$ acts on $`Y`$ without inversions.
###### Proof.
The ramification points $`y_{i,K}`$ specialize to distinct smooth points of $`Y`$, by construction. It follows from , Proposition 2.3.2.b, that there are no inversions. $`\mathrm{}`$
Note in particular that the above lemma implies that the irreducible components of $`Y`$ do not intersect themselves. For an irreducible component $`Z`$ of $`Y`$, we denote by $`I(Z)`$ its inertia group, and $`D(Z)`$ its decomposition group. Analogously, for a closed point $`y`$ of $`Y`$, we will write $`I(y)`$ for its inertia group $`I(y)`$. It is equal to the decomposition group of $`y`$, since $`k`$ is algebraically closed.
###### Lemma 2.1.2
* Let $`Z`$ be an irreducible component of $`Y`$. Then $`I(Z)`$ is a $`p`$-group and a normal subgroup of the decomposition group $`D(Z)`$.
* Let $`y_R:Spec(R)Y_R`$ be a section whose image is contained in the smooth locus of $`Y_R`$. Let $`Z`$ be the irreducible component of $`Y`$ on which $`y:=y_Rk`$ lies. Write $`m=p^\alpha n`$, with $`\mathrm{gcd}(n,p)=1`$, for the ramification index of $`y_K`$ in $`f`$. Then $`I(Z)`$ is normal in $`I(y)`$ and $`I(y)/I(Z)`$ is cyclic of order $`n`$.
* Any branch point $`x_{i,K}`$ of $`f_K`$ whose ramification index $`m_i`$ is prime-to-$`p`$, specializes to a component of $`X`$ over which $`f`$ is separable.
###### Proof.
Part (a) and (b) follow from , Lemme 6.3.3. Let $`x_{i,R}`$ be a branch point such that the ramification index $`m_i`$ of $`x_{i,K}`$ is prime-to-$`p`$. Suppose it specializes to a component $`W`$ of $`X`$ over which $`f`$ is inseparable. Then the inertia group of any component $`Z`$ of $`Y`$ mapping to $`W`$ is nontrivial. Hence $`x_i:=x_{i,R}k`$ will be branched of order $`p^am_i`$ with $`a>0`$, by Part (b). This is in contradiction with the assumption that the ramification points specialize to distinct points on $`Y`$. This proves (c). $`\mathrm{}`$
###### Lemma 2.1.3
Let $`y`$ be a singular point of $`Y`$. Let $`Z_1,Z_2`$ be the two irreducible components of $`Y`$ passing through $`y`$.
* The inertia group $`I(y)`$ is an extension of a cyclic group of order prime-to-$`p`$ by a $`p`$-group.
* The groups $`I(Z_1)`$ and $`I(Z_2)`$ are normal subgroups of $`I(y)`$ and $`<I(Z_1),I(Z_2)>`$ is the $`p`$-Sylow group of $`I(y)`$.
###### Proof.
Part (a) follows from , Proposition 2.3.2.a, since $`G`$ acts without inversions. Part (b) is proved analogous to , Lemme 6.3.6.iii. The assumptions in that lemma differ from the assumptions in the present case, but the proof goes through. $`\mathrm{}`$
### 2.2 First properties
We will suppose now that $`f_K:Y_KX_K`$ is a $`G`$-cover branched at four points $`x_1,\mathrm{},x_4`$, where $`X_K`$ has genus zero. Let $`m_i`$ be the ramification index of a point above $`x_i`$. We suppose that $`f_K`$ has bad reduction, and denote by $`f:YX`$ the special fiber of the quotient model of $`f_K`$.
###### Notation 2.2.1
Let $`U`$ be the union of the components $`W`$ of $`X`$ such that $`f`$ is inseparable over $`W`$. Let $`P`$ be a $`p`$-Sylow group of $`G`$. We denote by $`N_G(P)`$ the normalizer of $`P`$ in $`G`$ and by $`C_G(P)`$ the centralizer of $`P`$ in $`G`$.
###### Condition 2.2.2
In the rest of the paper we will assume the following conditions to hold:
| (a) $`p2`$, | (c) $`p|||G|`$, |
| --- | --- |
| (b) $`m_1,\mathrm{},m_4`$ prime-to-$`p`$, | (d) $`N_G(P)`$ is a dihedral group. |
###### Example 2.2.3
Here are some examples of groups for which Condition 2.2.2.(d) is satisfied. Let $`G=PSL_2(\mathrm{}^\alpha )`$, where $`\mathrm{}>2`$ is a prime. Suppose that $`p\mathrm{}`$ is a prime exactly dividing $`|G|=(\mathrm{}^\alpha 1)\mathrm{}^\alpha (\mathrm{}^\alpha +1)/2.`$ Then $`N_G(P)`$ is a dihedral group of order $`\mathrm{}^\alpha 1`$ or order $`\mathrm{}^\alpha +1`$, , Abschnitt II.8. Special cases are $`G=PSL_2(5)=A_5`$ and $`G=PSL_2(9)=A_6`$.
Recall that there is a map $`XX_0`$ which contracts some components; the strict transforms of the components of $`X_0`$ in $`X`$ are called the original components. Since we assumed that $`r=4`$, there are two possibilities for $`X_0`$: it can be smooth or not. In case $`X_0`$ is singular, it will consist of two genus zero components which meet in a unique point. We will call these components $`W_1`$ and $`W_2`$ and will also write $`W_1,W_2`$ for their strict transform in $`X`$. Similarly, we will write $`X_0`$ for the strict transform of $`X_0`$ in $`X`$, in case $`X_0`$ is smooth.
Suppose $`X_0`$ is not smooth. Since $`X_0`$ is stably marked, there will be exactly two branch points specializing to each of the components $`W_i`$. In $`X`$ the two components $`W_1`$ and $`W_2`$ are connected by a chain of $`^1`$’s, since the dual graph of $`X`$ is a tree. Let $`\mathrm{\Lambda }`$ the union of $`W_1`$, $`W_2`$ and the components connecting the two. We will say that a component $`W`$ of $`X`$ is a tail if it is not contained in $`\mathrm{\Lambda }`$ and meets the rest of $`X`$ in a unique point.
In case $`X_0`$ is smooth we just put $`\mathrm{\Lambda }=\{X_0\}`$. The definition of tail then becomes the usual definition of tail, i.e. one views the dual graph of $`X`$ to be oriented from $`X_0`$.
###### Lemma 2.2.4
Suppose that $`f_K`$ has bad reduction. Let $`W`$ be an irreducible component of $`X`$. Then $`f|_W`$ is separable if and only if $`W`$ is a tail.
In case $`X_0`$ is smooth, the statement of this lemma is the same as , Lemme 3.1.2, except that the model we are looking at is slightly different from the one considered in that paper. The definition of tail is made in such a way that $`f`$ will be separable exactly over the tails of $`X`$. The proof of Lemma 2.2.4 relies on Condition 2.2.2. If one does not assume the condition, it will not be true in general that $`f`$ will be exactly separable over the tails. For counter-examples see for example for $`p=2`$ or for general $`p`$.
###### Proof.
We split the proof up in two parts. First we consider the statement of the lemma for components $`W`$ which are not contained in $`\mathrm{\Lambda }`$. Then we show that $`f`$ is inseparable for all components contained in $`\mathrm{\Lambda }`$.
Let $`W`$ be a component of $`X`$ such that $`f|_W`$ is separable and $`W`$ is not contained in $`\mathrm{\Lambda }`$. The proof of , Proposition 2.4.8, carries over to this situation and shows that $`W`$ is connected to the rest of $`X`$ in a single point, i.e. $`W`$ is a tail.
Now suppose that $`W`$ is a tail of $`X`$. Let $`Z`$ be a component of $`Y`$ which maps to $`W`$. In case there is a branch point specializing to $`W`$, the cover is separable over $`W`$ by Lemma 2.1.2.(c) and Condition 2.2.2.(b). Suppose there are no branch points specializing to $`W`$ and $`ZW`$ is inseparable. Then $`I(Z)`$ is a nontrivial $`p`$-group which is a normal subgroup of $`D(Z)`$. Let $`Z^{}=Z/I(Z)`$. Since $`p`$ exactly divides the order of $`G`$, it follows that $`D(Z)/I(Z)`$ is of order prime-to-$`p`$. Then $`Z^{}W`$ is Galois of prime-to-$`p`$ order and branched at at most one point, hence trivial. Since $`ZZ^{}=W`$ is purely inseparable, $`Z`$ is a component of $`Y`$ of genus zero which meets the rest of $`Y`$ in a single point. By assumption, there is no ramification point specializing to $`Z`$. This contradicts the minimal character of $`Y`$. Hence $`f|_W`$ is separable.
We are now going to prove the lemma for components contained in $`\mathrm{\Lambda }`$, i.e. we have to show that $`f`$ is inseparable over the components contained in $`\mathrm{\Lambda }`$.
Suppose that $`X_0`$ is smooth. The cover $`f`$ is inseparable over $`X_0`$, by Proposition 1.1.4. This finishes the proof of the lemma for $`X_0`$ smooth.
Suppose that $`X_0`$ is singular and suppose that there exists a component of $`\mathrm{\Lambda }`$ over which $`f`$ is separable. Let $`U`$ be the union of the components of $`X`$ over which the cover is inseparable. If $`f`$ is separable over both $`W_1`$ and $`W_2`$, then the reduction is admissible by Proposition 1.1.4. It is no restriction to suppose that $`f`$ is inseparable over $`W_1`$. Let $`U^{}`$ be the connected component of $`U`$ containing $`W_1`$. The assumption on $`\mathrm{\Lambda }`$ implies that $`U^{}`$ does not contain $`W_2`$. The number of branch points specializing to $`U^{}`$ is at most two. Let $`\overline{U}^{}`$ be the union of $`U^{}`$ and the components of $`X`$ which are adjacent to $`U^{}`$.
We orient the dual graph of $`\overline{U}^{}`$ starting from $`W_1`$. Let $`𝔹^{}`$ be the set of tails of $`\overline{U^{}}`$. Let $`b_0`$ be the unique component of $`\overline{U}^{}`$ on the geodesic connecting $`W_1`$ and $`W_2`$ over which $`f`$ is separable. By assumption, such a component exist. For $`b𝔹^{}\{b_0\}`$, we say that the corresponding tail $`X_b`$ is primitive if one of the branch points specializes to this component. Otherwise, we will call the tail new. Denote the set of primitive (resp. new) tails by $`𝔹_{\mathrm{prim}}^{}`$ (resp. $`𝔹_{\mathrm{new}}^{}`$). The tail $`X_b`$ of $`\overline{U}^{}`$ meets the rest of $`\overline{u}^{}`$ in a unique point $`x_b`$. Let $`y_b`$ be a point of $`Y`$ mapping to $`x_b`$. By Lemma 2.1.2, $`I(y_b)`$ is an extension of a cyclic group of order $`n_b`$ prime-to-$`p`$, by a cyclic group of order $`p`$. Denote the conductor by $`h_b`$, and put $`\sigma _b=h_b/n_b`$. Recall that the ramification at $`y_b`$ is wild. This means that $`I(y_b)`$ is a subgroup of $`N_G(P)`$. We have assumed (Condition 2.2.2) that $`N_G(P)`$ is a dihedral group. It follows from , Lemme 1.1.2, that $`\sigma _b1/2(mod)`$. In particular, for $`b𝔹^{}`$, we have $`\sigma _b1/2`$.
Analogous to , Section 3.4, one proves the following vanishing cycle formula:
$$\underset{b𝔹_{\mathrm{new}}^{}}{}(\sigma _b1)=2+\underset{b𝔹_{\mathrm{prim}}^{}}{}(1\sigma _b)+(1\sigma _{b_0}).$$
(One proves this formula by constructing an “auxiliary cover” $`Z^{}\overline{U}^{}`$ and showing that it can be lifted to characteristic zero. The vanishing cycle formula follows then from the Riemann–Hurwitz formula applied to the generic fiber of the lift.) Furthermore, for $`b𝔹_{\mathrm{new}}^{}`$ we have that $`\sigma _b11/2`$. This follows from a genus consideration, , Proposition 3.3.5. The inequalities for $`\sigma _b`$ together with the fact that $`|𝔹_{\mathrm{prim}}^{}|2`$, implies that
$$\frac{|𝔹_{\mathrm{new}}^{}|}{2}2+\frac{|𝔹_{\mathrm{prim}}^{}|+1}{2}\frac{1}{2}.$$
Which is impossible. This concludes the proof. $`\mathrm{}`$
###### Remark 2.2.5
The notation and results used in the above proof will be introduced and explained in more detail in the next section. The reason for introducing the notation twice is that now that we proved Lemma 2.2.4, the notation can be simplified considerably.
### 2.3 The auxiliary cover
In this subsection we will suppose that Condition 2.2.2 is satisfied. Furthermore, we suppose that the cover $`f_K`$ has bad reduction. We start by introducing some more notation. Similar notation is used in the proof of Lemma 2.2.4.
Let $`𝔹`$ be the set of tails of $`X`$. Every tail $`X_b`$ of $`X`$ contains a unique singular point $`x_b`$ of $`X`$. Choose a singular point $`y_b`$ of $`Y`$ mapping to $`x_b`$ and let $`Y_b`$ be the component of $`Y`$ through $`y_b`$ which is mapping to $`X_b`$. Denote by $`h_b`$ the conductor of $`Y_bX_b`$ at $`y_b`$, and by $`n_b`$ the order of the prime-to-$`p`$ ramification. Put $`\sigma _b=h_b/n_b`$; this is the jump in the higher ramification groups of $`D(y_b)`$, in the upper numbering.
Let $`P`$ be a $`p`$-Sylow group of $`G`$. Let $`U`$ be the union of all components of $`X`$ over which $`f`$ is inseparable. Let $`V`$ be any connected component of $`f^1(U)`$. Let $`y`$ be a singular point of $`V`$ and $`Z_1`$ and $`Z_2`$ be the components passing through $`y`$. By Lemma 2.2.4, the inertia groups $`I(Z_1),I(Z_2)`$ are cyclic of order $`p`$, since $`Z_1`$ and $`Z_2`$ do not map to tails of $`X`$. Therefore, by Lemma 2.1.3, we have that $`I(Z_1)=I(Z_2)`$ and both are equal to the $`p`$-part of $`I(y)`$. Here we use that $`p`$ exactly divides the order of $`G`$. It follows that all irreducible components of $`V`$ have the same $`p`$-Sylow subgroup of $`G`$ as inertia group. Therefore, there is a connected component $`V`$ of $`f^1(U)`$ such that $`I(V)=P`$. We will always assume $`V`$ to be chosen like this.
As a consequence of Lemma 2.2.4, the construction of auxiliary covers as in , Section 3.2, goes through in our slightly different context. Since our notation differs from the notation of , we will recall the result.
###### Proposition 2.3.1
There exists a cover $`g_K:Z_KX_K`$ which is Galois with group $`D(V)`$ and has a quotient model $`g_R:Z_RX_R`$. It is uniquely characterized by the following properties:
* There exists a suitable open $`\mathrm{\Omega }`$ of $`X_R`$, which contains $`X_b\{x_b\}`$ for $`b𝔹`$, such that $`Z_RX_R`$ is tamely ramified over $`\mathrm{\Omega }`$ and unramified outside the sections $`x_i`$ ($`i=1\mathrm{}4`$).
* There exists an étale neighborhood $`X_R^{}X_R`$ of $`UXX_R`$ such that $`Y_R^{}X_R^{}Ind_{D(V)}^G(Z_R^{}X_R^{})`$, where $`Z_R^{}=Z_R\times _{X_R}X_R^{}`$ and $`Y_R^{}=Y_R\times _{X_R}X_R^{}`$.
We call $`g_K`$ the auxiliary cover associated to $`f_K`$.
###### Proof.
, Proposition 3.2.6. $`\mathrm{}`$
The (special fiber of the) auxiliary cover looks as follows. As above, we let $`V`$ be a connected component of $`f^1(U)`$ with inertia group $`P`$. Restricted to $`U`$, the auxiliary cover is just $`VU`$. Let $`X_b`$ be a tail of $`X`$ and $`x_b`$ the unique point of $`X_b`$ which is singular in $`X`$; we may suppose $`x_b=\mathrm{}`$. Let $`\mathrm{\Delta }_b`$ be the inertia group of a point $`y_b`$ above $`x_b`$ which lies on $`V`$. Then, by the Katz–Gabber Lemma , there exist a cover $`Z_bX_b`$ unbranched outside $`0,\mathrm{}`$ and at most tamely ramified at 0 which locally around $`\mathrm{}`$, if we induce it up to a $`G`$-cover, agrees with $`YX`$. Now $`ZX|_{X_b}=Ind_{\mathrm{\Delta }_b}^{D(V)}(Z_bX_b)`$. Note that by construction, the $`\sigma _b`$ for $`b𝔹`$ for the cover $`g:ZX`$, are the same as for the original cover.
###### Lemma 2.3.2
For $`b𝔹_{\mathrm{new}}`$, we have $`\sigma _b11/2.`$
###### Proof.
This follows from , Proposition 3.3.5 and Lemme 1.1.2, and the assumption that $`N_G(P)`$ is a dihedral group. $`\mathrm{}`$
The following formula reflects the condition that the genus of $`Y`$ has to be equal to the genus of the generic fiber $`Y_K`$. It is proved in , Section 3, using the auxiliary cover $`g_R:Z_RX_R`$.
###### Proposition 2.3.3 (Vanishing cycle formula)
$$\underset{b𝔹_{\mathrm{new}}}{}(\sigma _b1)=2+\underset{b𝔹_{\mathrm{prim}}}{}(1\sigma _b).$$
### 2.4 Modular reduction
Let $`f_K:Y_KX_K`$ be a $`G`$-cover for which Condition 2.2.2 holds. In this subsection we will show that $`f_K`$ has either good reduction or modular reduction. Essentially, the property of having modular reduction means that the auxiliary cover introduced in the previous section is a cover $`Z_K_K^1`$ with Galois group a dihedral group and ramification of order 2 as discussed in Example 1.1.6.
###### Definition 2.4.1
Suppose $`f_K`$ has bad reduction. We will say that $`f_K`$ has modular reduction if the following conditions are satisfied.
* The curve $`X`$ has four primitive tails and no new tails. Every irreducible component of $`X`$ is either an original component or a tail.
* Let $`E`$ be a connected component of $`f^1(X_0)`$. Then $`D(E)`$ is a dihedral group of order $`2N`$ for some $`N`$ divisible by $`p`$.
* Let $`X_b`$ be a tail of $`X`$. Let $`x_b`$ be the unique singular point of $`X_b`$ in $`X`$ and let $`y_b`$ be a point of $`Y`$ mapping to $`x_b`$. Then the inertia group $`I(y_b)`$ is a dihedral group of order $`2p`$ and $`\sigma _b=1/2`$.
If $`f_K:Y_KX_K`$ has modular reduction, then we will say that the special fiber $`f:YX`$ of the quotient model of $`f_K`$ is of modular type.
The integer $`N`$ defined above will be called the level of $`f`$. This terminology reflects the relation between $`f`$ of modular type and modular curves.
###### Remark 2.4.2
Suppose that $`f_K:Y_KX_K`$ has modular reduction of level $`N`$. The auxiliary cover $`g_K:Z_KX_K`$ corresponding to $`f_K`$ is a Galois cover with Galois group the dihedral group $`\mathrm{\Delta }`$ of order $`2N`$, branched at $`x_1,\mathrm{},x_4`$ of order 2. The reduction $`g:ZX`$ is inseparable over $`X_0`$; it is separable over the tails. Note that $`g^1(X_0)`$ may be identified with $`E`$ in the definition; it is, after choice of a base point, a generalized elliptic curve.
###### Proposition 2.4.3
Let $`f_K:Y_KX_K`$ be as before, in particular we suppose that Condition 2.2.2 is satisfied and that $`f_K`$ has bad reduction. Then $`f_K`$ has modular reduction.
###### Proof.
Suppose that $`f_K`$ has bad reduction. Proposition 2.3.3 implies that
$$\underset{b𝔹_{\mathrm{new}}}{}\sigma _b1=2+\underset{b𝔹_{\mathrm{prim}}}{}1\sigma _b.$$
Recall that $`\sigma _b3/2`$ for $`b𝔹_{\mathrm{new}}`$ and $`\sigma 1/2`$ for $`b𝔹_{\mathrm{prim}}`$, Lemma 2.3.2. Hence $`|𝔹_{\mathrm{new}}|/20`$. This implies that $`|𝔹_{\mathrm{new}}|=0`$ and $`\sigma _b=1/2`$ for $`b𝔹_{\mathrm{prim}}`$. Moreover, all $`x_i`$ specialize to components over which $`f`$ is separable, so $`|𝔹_{\mathrm{prim}}|=4`$.
The decomposition group of a singular point $`y_b`$ of $`Y`$ contains a dihedral group of order $`2p`$, since $`\sigma _b=1/2`$. (It cannot be Abelian, since then $`\sigma _b`$ would be an integer by the Hasse–Arf Theorem.) Since $`D(y_b)=I(y_b)`$ is a cyclic-by-$`p`$ group, $`D(y_b)`$ is isomorphic to a dihedral group of order $`2p`$. This proves Part (c) of Definition 2.4.1.
The decomposition group $`D(V)`$ is a subgroup of $`N_G(P)`$ which is a dihedral group, by assumption. Note that $`D(V)`$ is not cyclic, since it contains (a conjugate of) the inertia group of some point $`y_b`$, which is dihedral. This proves Part (b).
The only thing left to show is the second part of Part (a). Let $`W`$ be a component of $`X`$ which is neither a tail nor an original component and let $`Z`$ be a component of $`Y`$ which maps to $`W`$. Above we have shown that there are exactly four tails. This implies that $`W`$ meets the rest of $`X`$ in two points. There is no branch point specializing to $`W`$, so the maximal separable subcover $`Z^{}W`$ of $`ZW`$ is branched at at most two points. Moreover, $`I_Z`$ is a cyclic group of order $`p`$, since $`W`$ is not a tail. This implies that the degree of $`Z^{}W`$ is prime-to-$`p`$, so it is a cyclic cover of $`^1`$ branched at two points. But then $`Z`$ has genus zero and meets the rest of $`Y`$ in exactly two points. This contradicts the minimality of $`Y`$.
This shows that the cover has modular reduction. $`\mathrm{}`$
### 2.5 Reduction of degenerate covers
In this section we will study the reduction behavior of the degenerate covers corresponding to the cusps of the Hurwitz space. Note that the theory of semistable reduction we developed so far does not apply here, since we always assumed that the generic fibers of our curves were smooth.
The situation in this section is somewhat different from that in the previous sections. Let $`R`$ be a complete discrete valuation ring whose quotient field $`K`$ is of characteristic zero and whose residue field $`k=\overline{k}`$ has characteristic $`p`$. Let $`(X_K;x_1,\mathrm{},x_4)`$ be a stably marked $`K`$-curve of genus zero, which we suppose to be singular. Let $`X_{1,K}`$ and $`X_{2,K}`$ be the two irreducible components of $`X_K`$. Let $`\tau `$ be the singular point of $`X_K`$. Let $`f_K:Y_KX_K`$ be an admissible $`G`$-Galois cover ramified at $`x_1,\mathrm{},x_4`$. Let $`\rho `$ be a point of $`Y_K`$ mapping to $`\tau `$ and let $`Y_{1,K}`$ and $`Y_{2,K}`$ be the components of $`Y_K`$ passing through $`\rho `$, where we suppose that $`Y_{i,K}`$ maps to $`X_{i,K}`$. Let $`G_i`$ be the decomposition group $`Y_{i,K}`$. Let $`X_{i,R}`$ and $`Y_{i,R}`$ be the closure of $`X_{i,K}`$ and $`Y_{i,K}`$ in $`X_R`$. Since $`\overline{}_4^{\mathrm{in}}(G)`$ is proper, $`f_K`$ extends uniquely to a map $`f_{0,R}:Y_RX_{0,R}`$ between stably marked curves over $`SpecR`$. Let $`f_i:Y_{i,R}X_{i,R}`$ be the corresponding morphism. We will denote the special fibers of $`Y_R,X_R,Y_{i,R},X_{i,R}`$ by $`Y,X,Y_i,X_i`$, respectively. The mark $`C`$ on $`X_R`$ can be written as a union $`C=C_1^{}C_2^{}`$, where $`C_i^{}`$ is a mark on $`X_{i,R}`$. Let $`\tau `$ be the unique singular point of $`X_K`$, we denote its (unique) extension to $`X_R`$ also by $`\tau `$. Let $`C_i=C_i^{}\{\tau \}`$. Let $`D_i^{}`$ be the restriction of the mark of $`Y_R`$ to $`Y_{i,R}`$ and write $`D_i`$ for the union of $`D_i^{}`$ with the points of $`Y_{i,R}`$ which are singular in $`Y_R`$. The next lemma follows immediately.
###### Lemma 2.5.1
The morphism $`f_{i,R}:(Y_{i,R},D_i)(X_{i,R},C_i)`$ constructed above is a morphism of stably marked curves.
The covers $`f_{i,K}`$ constructed above are covers of a projective line branched at three points. So for these covers we can apply the criterion for good reduction proved by Raynaud . Actually, since here we put a stronger condition on $`G`$ than in , we can show that the covers $`f_{i,K}`$ have good reduction iff $`p`$ does not divide the ramification indices, Lemma 2.5.2. In other words, the condition on the field of definition in the Theorem of Raynaud is not needed in this case.
Proposition 2.5.3 is the key result in this section. It describes the reduction of the degenerate covers to characteristic $`p`$. The idea is that we can understand the reduction of such a degenerate cover, because we understand the reduction of the two three point covers of which it is made. Proposition 2.5.3 will be used in Section 5 to explicitly describe the bad part of the Hurwitz space.
###### Lemma 2.5.2
Let $`f_K:Y_KX_K_K^1`$ be a $`G`$-cover branched at three points $`0,1,\mathrm{}`$. Suppose that $`p`$ exactly divides the order of $`G`$ and the normalizer of a $`p`$-Sylow group of $`G`$ is dihedral. Denote the ramification indices of $`f_K`$ by $`m_1,m_2,m_3`$. Let $`f_R:Y_RX_R`$ be the quotient model corresponding to $`X_{0,R}=_R^1`$. Then $`f_R`$ has good reduction if and only if $`pm_i`$ for $`i=1,2,3`$.
###### Proof.
This follows immediately from the vanishing cycle formula , Section 3.4, combined with the estimates for $`\sigma _b`$ from Lemma 2.3.2. $`\mathrm{}`$
###### Proposition 2.5.3
Let $`f_K:Y_KX_K`$ be as above and let $`f:YX`$ be its reduction. Let $`n`$ be the order of the ramification of $`f_K`$ above the singular point $`\tau `$. We denote by $`\tau _k`$ the image of $`\tau `$ on the special fiber. Let $`\rho _k`$ be a point on $`Y`$ above $`\tau _k`$.
* The cover $`f_K`$ has admissible reduction iff $`pn`$.
* The inertia group $`I(\rho _k)`$ has order $`n`$
* If $`f_K`$ has bad reduction, then $`f`$ is of modular type.
* Suppose $`f_K`$ has bad reduction, then $`n`$ divides the level $`N`$ of $`f`$. Let $`Z_1`$ and $`Z_2`$ be the irreducible components of $`Y`$ passing through $`\rho _k`$. Then $`D(Z_1)`$ and $`D(Z_2)`$ are dihedral groups of order $`2n`$.
###### Proof.
The cover $`f:YX`$ is the reduction of a degenerate cover, by assumption. However, there will be covers representing points in the interior of the Hurwitz stack in characteristic zero which specialize to $`f`$. Therefore, in case $`f_K`$ has bad reduction, $`f`$ will be of modular type by Proposition 2.4.3. This proves Part (c).
This implies that in case $`f_K`$ has bad reduction, $`f_0:YX_0`$ does not contract any components to $`\tau `$. This is clearly also the case if $`f_K`$ does not have bad reduction. Since $`f_i:Y_{i,R}X_{i,R}`$ for $`i=1,2`$ are flat, we have $`g(Y_i)=g(Y_{i,K})`$. Moreover, $`g(Y)=g(Y_K)`$. Now
$$g(Y)=(|G|/|G_1|)g(Y_1)+(|G|/|G_2|)g(Y_2)+1|G|/|G_1||G|/|G_2|+|G|/|I(\rho _k)|$$
and
$$g(Y_K)=(|G|/|G_1|)g(Y_{1,K})+(|G|/|G_2|)g(Y_{2,K})+1|G|/|G_1||G|/|G_2|+|G|/n.$$
This shows that $`|I(\rho _k)|=n`$.
Suppose $`p|n`$. Then Part (b) of Lemma 2.1.2 implies that the covers $`f_{i,K}`$ have bad reduction for $`i=1,2`$. It follows that $`f_K`$ has bad reduction. Conversely, suppose that $`f_K`$ has bad reduction. Part (c) implies that $`f_K`$ has modular reduction. In particular it follows that $`I(Z_i)`$ is $`p`$-cyclic. Part (b) of Lemma 2.1.2 implies that $`p|n`$.
Suppose that $`f_K`$ has bad reduction. The decomposition groups $`D(Z_i)`$ are subgroups of $`N_G(P)`$, since the corresponding inertia groups are nontrivial. Moreover, they contain a dihedral group of order $`2p`$, by the definition of modular reduction (Definition 2.4.1). It follows that the $`D(Z_i)`$ are dihedral groups. The maximal separable subcover of $`Z_iX_i`$ is $`f_i^{}:Z_i^{}:=Z_i/I(Z_i)X_i`$. Note that $`f_i^{}`$ is branched at three points of order $`2,2,n/p`$ hence $`g(Z_i^{})=0`$. It follows that the degree of $`f_i^{}`$ is $`2n/p`$, hence $`D(Z_i)`$ is a dihedral group of order $`2n`$. By definition of the level $`N`$ of $`f`$, the subgroup of $`N_G(P)`$ generated by $`D(Z_1)`$ and $`D(Z_2)`$ is a dihedral group of order $`2N`$; this implies that $`n|N`$. $`\mathrm{}`$
## 3 A Reduction Theorem
In this section we prove a theorem about the structure of $`\overline{H}𝔽_p`$, where $`\overline{H}`$ is a complete Hurwitz space for $`G`$-covers satisfying Condition 2.2.2. As explained in the introduction, this theorem relies on, and in some sense extends, the results of Katz and Mazur on the reduction of the modular curve $`X_1(p)`$. This is somewhat surprising. Like modular curves, the Hurwitz spaces we look at are curves, defined over small number fields, and arise as quotients of the upper half plane by discrete subgroups of $`\mathrm{GL}_2()`$, see . However, these groups are non-congruence subgroups, in general.
### 3.1 Statement of the main results
#### 3.1.1
Let $`G`$ be a finite group, $`\underset{¯}{C}=(C_1,C_2,C_3,C_4)`$ a class vector in $`G`$ of length $`4`$ and $`p`$ an odd prime. We denote by $`m_i`$ the order of the elements of $`C_i`$. We assume that Condition 2.2.2 holds, with respect to $`G`$, $`p`$ and $`m_i`$.
Let $`K:=K(\underset{¯}{C})(\zeta _m)`$ be the minimal field over which the classes $`C_i`$ are rational, see Section 1.2.3. We choose a prime ideal $`𝔭`$ of $`K`$ dividing $`p`$, and denote by $`\mathrm{\Lambda }:=𝒪_{K,𝔭}`$ its local ring. We let $`\overline{H}:=\overline{H}_4^{\mathrm{in}}(\underset{¯}{C})`$ be the complete Hurwitz space over $`\mathrm{\Lambda }`$, as defined in Section 1.2.4. By construction, $`\overline{H}`$ is an algebraic space, proper and of finite type over $`\mathrm{\Lambda }`$. It contains the Hurwitz space $`H:=H_4^{\mathrm{in}}(\underset{¯}{C})`$ as a dense open subscheme. The scheme $`H`$ is smooth over $`\mathrm{\Lambda }`$; its generic fiber $`HK`$ is the reduced Hurwitz curve studied e.g. in .
Let $`\overline{H}^{\mathrm{bad}}\overline{H}`$ be the closed subspace corresponding to bad covers. Since bad covers occur only in positive characteristic, $`\overline{H}^{\mathrm{bad}}`$ is a closed subspace of $`\overline{H}𝔽_q`$, where $`𝔽_q`$ is the residue field of $`\mathrm{\Lambda }`$. The complement $`\overline{H}^{\mathrm{adm}}:=\overline{H}\overline{H}^{\mathrm{bad}}`$ is a dense open subscheme and corresponds to admissible covers. Let $`\overline{H}^{\mathrm{good}}`$ be the closure of $`\overline{H}^{\mathrm{adm}}𝔽_q`$ in $`\overline{H}𝔽_q`$. Note that this is an abuse of notation, since $`\overline{H}^{\mathrm{good}}`$ has nontrivial intersection with $`\overline{H}^{\mathrm{bad}}`$, in general. There is a natural map $`\overline{H}_\lambda ^1\mathrm{\Lambda }`$. We say that an $`\overline{𝔽}_p`$-rational point $`s`$ of $`\overline{H}^{\mathrm{bad}}`$ is supersingular if the corresponding value $`\lambda (s)\overline{𝔽}_p`$ is supersingular.
###### Theorem 3.1.1
1. The complete Hurwitz space $`\overline{H}`$ is a normal scheme of dimension $`2`$. The natural map $`\overline{H}_\lambda ^1\mathrm{\Lambda }`$ is finite and flat.
2. The subspaces $`\overline{H}^{\mathrm{bad}}`$ and $`\overline{H}^{\mathrm{good}}`$ are smooth projective curves over $`𝔽_q`$. They intersect transversally in the supersingular points.
In the rest of Section 3.1 we state a number of results on the local structure of $`\overline{H}`$ and explain how Theorem 3.1.1 can be deduced from them. Let us give a brief outline. On the open subset $`\overline{H}^{\mathrm{adm}}\overline{H}`$, Theorem 3.1.1 is known to hold, see . Therefore, it suffices to look at $`\overline{H}`$ in a neighborhood of a point corresponding to a bad cover. Theorem 3.1.2 below describes the universal deformation ring of such a bad cover. Essentially, this theorem implies that the complete Hurwitz stack $`\overline{}`$ associated to $`\overline{H}`$ is regular and that a version of Theorem 3.1.1 holds for $`\overline{}`$. To finish the proof of Theorem 3.1.1, we have to study the monodromy action, i.e. the action of the group of automorphisms of a bad cover on its universal deformation ring. Propositions 3.1.4 describes this action, and Theorem 3.1.1 follows. Under some extra hypotheses (Condition 3.1.6), we can improve our results on the monodromy action, and we can actually show that $`\overline{H}`$ is regular. The relevant statement is made in Proposition 3.1.7.
The proofs of Theorem 3.1.2, Proposition 3.1.4 and Proposition 3.1.7 are postponed to Section 3.3.
#### 3.1.2 The universal deformation ring
Let $`\overline{}:=\overline{}_4^{\mathrm{in}}(\underset{¯}{C})`$ be the complete Hurwitz stack over $`\mathrm{\Lambda }`$, associated to $`\overline{H}`$, and let $`\overline{}^{\mathrm{bad}},\overline{}^{\mathrm{good}}\overline{}`$ be the closed substacks corresponding to the closed subspaces $`\overline{H}^{\mathrm{bad}},\overline{H}^{\mathrm{good}}\overline{H}`$. We look at the following situation. Let $`k`$ be an algebraically closed field of characteristic $`p`$ and
$$f_0:Y(X_0;x_i)$$
an object of $`\overline{}^{\mathrm{bad}}`$, defined over $`k`$. In the rest of this section we will mostly write “$`Y`$” instead of “$`f_0:Y(X_0;x_i)`$” for this object; we understand that the curve $`Y`$ is equipped with an action of the group $`G`$, a mark $`DY`$ and a map $`f_0`$ to the stably marked curve $`(X_0;x_i)`$.
Let $`R_Y`$ be the strict complete local ring of $`\overline{}`$ at the $`k`$-point corresponding to $`Y`$. Let $`Y_{\mathrm{univ}}`$ be the object of $`\overline{}_4^{\mathrm{in}}(\underset{¯}{C})`$ corresponding to the tautological morphism $`SpecR_Y\overline{}`$. By a general property of algebraic stacks, $`Y_{\mathrm{univ}}`$ is the universal deformation of $`Y`$ as object of $`\overline{}`$. This means the following. Let $`W(k)`$ denote the ring of Witt vectors over $`k`$ and $`\widehat{𝒞}_k`$ the category of complete local Noetherian $`W(k)`$-algebras with residue field $`k`$. We let $`Def(Y)`$ be the functor which assigns to $`R\widehat{𝒞}_k`$ the set of isomorphism classes of deformations of $`Y`$ over $`R`$. Then $`Y_{\mathrm{univ}}`$ defines an equivalence
$$Hom_{\widehat{𝒞}_k}(R_Y,)\stackrel{}{}Def(Y).$$
The closed substack $`\overline{}^{\mathrm{bad}}\overline{}`$ (resp. $`\overline{}^{\mathrm{good}}\overline{}`$) corresponds to a subfunctor $`Def(Y)^{\mathrm{bad}}Def(Y)`$ (resp. $`Def(Y)^{\mathrm{good}}Def(Y)`$), which is represented by a quotient ring $`R_Y^{\mathrm{bad}}`$ (resp. $`R_Y^{\mathrm{good}}`$) of $`R_Y`$.
We denote by $`Def(X_0;x_i)`$ the deformation functor for $`(X_0;x_i)`$ as object of $`\overline{}_{0,4}`$ and by $`R_{X_0}`$ the universal deformation ring. The morphism $`\overline{}\overline{}_{0,4}`$ induces a transformation $`Def(Y)Def(X;x_i)`$, hence a morphism $`R_{X_0}R_Y`$. The ring $`R_{X_0}`$ is of the form $`R_{X_0}=W(k)[[w]]`$. For instance, if $`X_0`$ is smooth we may assume that $`X_0=^1`$, $`x_1=0`$, $`x_2=1`$, $`x_3=\mathrm{}`$, $`x_4=\lambda _0k\{0,1\}`$ and $`w=\lambda \lambda _0`$. We say that $`Y`$ is supersingular (resp. ordinary) if $`\lambda _0`$ is supersingular (resp. ordinary). If $`X_0`$ is singular, we will also say that $`Y`$ is ordinary.
###### Theorem 3.1.2
The ring $`R_Y`$ is regular of dimension $`2`$ and a finite flat extension of $`R_{X_0}`$. Moreover, there exists a regular sequence $`(t,\pi )`$ for $`R_Y`$ such that
1. $`R_Y^{\mathrm{bad}}=R_Y/(\pi )k[[t]]`$,
2. The induced (finite) morphism $`R_{X_0}_{W(k)}kR_Y^{\mathrm{bad}}`$ is inseparable of degree $`p`$. Its separable part $`R_{X_0}_{W(k)}k(R_Y^{\mathrm{bad}})^pk[[t^p]]`$ is tamely ramified (of degree $`d`$); if $`X_0`$ is smooth then $`d=1`$.
3. If $`Y`$ is ordinary, then $`p=u\pi ^{p1}`$ for a unit $`uR_Y^\times `$ and $`R_Y^{\mathrm{good}}=0`$.
4. If $`Y`$ is supersingular, then $`p=u\pi ^{p1}`$, where $`uR_Y`$ is a local equation for $`\overline{}_4^{\mathrm{in}}(\underset{¯}{C})^{\mathrm{good}}`$. More precisely, $`R_Y^{\mathrm{good}}=R_Y/(u)k[[\overline{\pi }]]`$, and $`(u,\pi )`$ is a regular sequence for $`R_Y`$.
In Section 3.2 we will prove this theorem in the case that the group $`G`$ is dihedral and the conjugacy classes $`C_i`$ represent reflections. In this case, the $`G`$-cover $`f:YX_0`$ corresponds essentially to a point on the modular curve $`X_1(N)`$, and Theorem 3.1.2 follows from the results of . We will prove the general case in Section 3.3.2 by reduction to the dihedral case.
###### Corollary 3.1.3
The complete Hurwitz stack $`\overline{}`$ is regular of dimension $`2`$. The natural morphism $`\overline{}\overline{}_{0,4}`$ is finite and flat. The closed substacks $`\overline{}^{\mathrm{bad}}`$ and $`\overline{}^{\mathrm{good}}`$ are smooth over $`𝔽_q`$, and intersect transversally in the supersingular point.
#### 3.1.3 The monodromy action
Corollary 3.1.3 together with Proposition 1.2.8 implies Part (i) of Theorem 3.1.1. To prove Part (ii) of Theorem 3.1.1, we use the general fact that the coarse moduli scheme is locally the quotient of (an étale cover of) the corresponding algebraic stack by a finite group action. Actually, it suffices to look at the strict complete local rings. Hence we are led to study the monodromy action of the automorphisms of $`Y`$ on the universal deformation ring $`R_Y`$.
Let $`Y`$ be as in Section 3.1.2. We write $`Aut_k(Y)`$ for the group of $`k`$-linear automorphisms of $`Y`$, considered as object of $`\overline{}`$. Since $`(X_0;x_i)`$ has no nontrivial automorphism, an element $`\sigma Aut_k(Y)`$ is a $`k`$-automorphism of $`Y`$ such that $`f_0\sigma =f_0`$ and $`g\sigma =\sigma g`$ for all $`gG`$. By the universal property of $`Y_{\mathrm{univ}}`$, for each $`\sigma `$ there exists a unique automorphism $`\gamma :R_Y\stackrel{}{}R_Y`$ such that $`\sigma `$ lifts to a unique $`\gamma `$-semilinear automorphism $`\sigma _{\mathrm{univ}}:Y_{\mathrm{univ}}\stackrel{}{}Y_{\mathrm{univ}}`$. We call $`\gamma Aut_{\widehat{𝒞}_k}(R_Y)`$ the monodromy action of $`\sigma `$. We define the monodromy group $`\mathrm{\Gamma }`$ of $`Y`$ as the image of the homomorphism $`Aut_k(Y)Aut_{\widehat{𝒞}_k}(R_Y)`$. Let $`Y_{\overline{\eta }}`$ be the geometric generic fiber of $`Y_{\mathrm{univ}}`$ (note that $`R_Y`$ is a domain, by Theorem 3.1.2). Identifying the group of $`\overline{\eta }`$-automorphisms of $`Y_{\overline{\eta }}`$ with $`C_G`$, the center of $`G`$, we obtain a natural exact sequence
(5)
$$1C_GAut_k(Y)\mathrm{\Gamma }\mathrm{\hspace{0.33em}1}.$$
Let $`s:Speck\overline{H}`$ be the geometric point of the Hurwitz space corresponding to $`Y`$. The strict complete local ring of $`\overline{H}`$ is the ring of $`\mathrm{\Gamma }`$-invariants of $`R_Y`$:
(6)
$$\widehat{𝒪}_{\overline{H},s}=R_Y^\mathrm{\Gamma }.$$
#### 3.1.4 The absolute monodromy action
To study the action of $`\mathrm{\Gamma }`$ on $`R_Y`$, it will turn out to be useful to enlarge $`\mathrm{\Gamma }`$ and look at the absolute monodromy group $`\mathrm{\Gamma }^{\mathrm{ab}}`$.
Let $`\overline{}_4^{\mathrm{ab}}(\underset{¯}{C})`$ be the absolute version of the complete Hurwitz stack $`\overline{}`$. Let us for the moment consider $`Y`$ and $`Y_{\mathrm{univ}}`$ as objects of $`\overline{}_4^{\mathrm{ab}}(\underset{¯}{C})`$. In other words, we forget the embedding $`GAut(Y/X)`$, and only retain its image, see Variant 1.2.6. It is still true that $`Y_{\mathrm{univ}}`$ is the universal deformation of $`Y`$. But we obtain a bigger automorphism group, in general. We denote by $`Aut_k^{\mathrm{ab}}(Y)`$ the group of $`k`$-automorphisms of $`Y`$ as object of $`\overline{}_4^{\mathrm{ab}}(\underset{¯}{C})`$. An element $`\sigma Aut_k^{\mathrm{ab}}(Y)`$ is a $`k`$-automorphism of $`Y`$ such that $`f_0\sigma =f_0`$ and $`\sigma g\sigma ^1G`$ for all $`gG`$. We find that $`G`$ is a normal subgroup of $`Aut_k^{\mathrm{ab}}(Y)`$, and $`Aut_k(Y)`$ is the centralizer of $`G`$ in $`Aut_k^{\mathrm{ab}}(Y)`$. The group $`Aut_k^{\mathrm{ab}}(Y)`$ acts on $`R_Y`$; this action extends the action of $`Aut_k(Y)`$. We write $`\mathrm{\Gamma }^{\mathrm{ab}}`$ for the image of $`Aut_k^{\mathrm{ab}}(Y)`$ in $`Aut_{\widehat{𝒞}_k}(R_Y)`$. Clearly, $`\mathrm{\Gamma }\mathrm{\Gamma }^{\mathrm{ab}}`$, and the exact sequence (5) becomes
(7)
$$1GAut_k^{\mathrm{ab}}(Y)\mathrm{\Gamma }^{\mathrm{ab}}\mathrm{\hspace{0.33em}1}.$$
###### Proposition 3.1.4
There exist two characters $`\chi ^{\mathrm{bad}}:\mathrm{\Gamma }^{\mathrm{ab}}𝔽_p^\times `$ and $`\chi ^{\mathrm{adm}}:\mathrm{\Gamma }^{\mathrm{ab}}\mu _d`$ (here $`\mu _dR_Y`$ denotes the set of $`d`$th roots of unity, for $`d`$ as in Theorem 3.1.2 (ii)) with the following properties. The parameters $`t,\pi R_Y`$ in Theorem 3.1.2 can be chosen such that for all $`\gamma \mathrm{\Gamma }^{\mathrm{ab}}`$
$$\gamma (t)=\chi ^{\mathrm{adm}}(\gamma )t,\gamma (\pi )\chi ^{\mathrm{bad}}(\gamma )\pi (mod\pi ^2).$$
Moreover, the homomorphism $`(\chi ^{\mathrm{bad}},\chi ^{\mathrm{adm}}):\mathrm{\Gamma }^{\mathrm{ab}}𝔽_p^\times \times \mu _d`$ is injective.
We will prove this proposition in Section 3.3.3. In the rest of this subsection we will show that Theorem 3.1.2 and Proposition 3.1.4 together imply Theorem 3.1.1. Since $`\mathrm{\Gamma }\mathrm{\Gamma }^{\mathrm{ab}}`$, Proposition 3.1.4 holds also for $`\mathrm{\Gamma }`$.
The closed subscheme $`\overline{H}^{\mathrm{bad}}\overline{H}`$ is the image of the natural morphism $`\overline{}^{\mathrm{bad}}\overline{H}`$, by definition. It follows that the strict complete local ring $`\widehat{𝒪}_{\overline{H}^{\mathrm{bad}},s}`$ is the image of the natural morphism $`R_Y^\mathrm{\Gamma }R_Y^{\mathrm{bad}}`$. By Proposition 3.1.4, the order of $`\mathrm{\Gamma }`$ is prime-to-$`p`$. Therefore, the natural map $`R_Y^\mathrm{\Gamma }(R_Y^{\mathrm{bad}})^\mathrm{\Gamma }`$ is surjective, hence $`\widehat{𝒪}_{\overline{H}^{\mathrm{bad}},s}=(R_Y^{\mathrm{bad}})^\mathrm{\Gamma }`$, see , A 7. Now Theorem 3.1.2 (i) and Proposition 3.1.4 imply
(8)
$$\widehat{𝒪}_{\overline{H}^{\mathrm{bad}},s}=(R_Y^{\mathrm{bad}})^\mathrm{\Gamma }k[[t^d^{}]],$$
where $`d^{}|d`$ is the order of $`\chi ^{\mathrm{adm}}(\mathrm{\Gamma })`$. We have shown that $`\overline{H}^{\mathrm{bad}}`$ is a smooth curve. The same argument shows that
(9)
$$\widehat{𝒪}_{\overline{H}^{\mathrm{good}},s}=(R_Y^{\mathrm{good}})^\mathrm{\Gamma }\{\begin{array}{cc}0\hfill & \text{if }Y\text{ is ordinary},\hfill \\ k[[\overline{\pi }^\mu ]]\hfill & \text{if }Y\text{ is supersingular},\hfill \end{array}$$
where $`\mu |(p1)`$ is the order of $`\chi ^{\mathrm{bad}}(\mathrm{\Gamma })`$. We conclude that $`\overline{H}^{\mathrm{good}}`$ is a smooth curve and that $`\overline{H}^{\mathrm{good}}_\lambda ^1`$ is ramified of order $`(p1)/\mu `$ in the supersingular points. Let us assume that $`Y`$ is supersingular and denote by $`\overline{H}_p^{\mathrm{red}}`$ (resp. by $`\overline{}_p^{\mathrm{red}}`$) the closed subscheme $`(\overline{H}𝔽_p)^{\mathrm{red}}\overline{H}`$ (resp. the closed substack $`(\overline{}𝔽_p)^{\mathrm{red}}\overline{}`$). By Theorem 3.1.2 (iv), $`\overline{}_p^{\mathrm{red}}\times _\overline{}SpecR_Y=R_Y/(\pi u)k[[\overline{\pi },\overline{u}\overline{\pi }\overline{u}=0]]`$. Taking invariants and arguing as before, we get
(10)
$$\widehat{𝒪}_{\overline{H}_p^{\mathrm{red}},s}=(R_Y/(\pi u))^\mathrm{\Gamma }k[[\overline{\pi }^\mu ,\overline{u}\overline{\pi }^\mu \overline{u}=0]]$$
(note that $`\chi ^{\mathrm{adm}}=1`$, since $`X_0`$ is smooth). This completes the proof of Theorem 3.1.1, modulo the proofs of Theorem 3.1.2 and Proposition 3.1.4.
Translating the statements of Theorem 3.1.2 into geometric properties of the map $`\overline{H}_\lambda ^1`$, we obtain the following corollary.
###### Corollary 3.1.5
Let $`s`$ be the point on $`\overline{H}^{\mathrm{bad}}`$ corresponding to $`Y`$.
1. The natural map $`\overline{H}^{\mathrm{bad}}_\lambda ^1𝔽_q`$ is finite, with inseparability degree $`p`$. Its separable part $`(\overline{H}^{\mathrm{bad}})^{(p)}_\lambda ^1𝔽_q`$ is tamely ramified at $`\lambda =0,1,\mathrm{}`$ and étale everywhere else. More precisely, its ramification index in $`s^{(p)}`$ is equal to $`[Im(\chi ^{\mathrm{adm}}):\mu _d]`$.
2. The natural map $`\overline{H}^{\mathrm{good}}_\lambda ^1𝔽_q`$ is tamely ramified at $`\lambda =0,1,\mathrm{}`$ and the supersingular values of $`\lambda `$, and is étale everywhere else. Its ramification index in $`s\overline{H}^{\mathrm{good}}\overline{H}^{\mathrm{bad}}`$ is $`[\chi ^{\mathrm{bad}}(\mathrm{\Gamma }):𝔽_p^\times ]`$ (this happens if and only if $`Y`$ is supersingular).
3. Let $`m_s`$ be the multiplicity of $`\overline{H}^{\mathrm{bad}}`$ in a neighborhood of $`s`$ in $`\overline{H}𝔽_q`$, i.e. the length of the Artinian local ring $`\widehat{𝒪}_{\overline{H},\eta }/p\widehat{𝒪}_{\overline{H},\eta }`$, where $`\eta :Speck((t^d^{}))\overline{H}`$ comes from equation (8). Then $`m_s=[\chi ^{\mathrm{bad}}(\mathrm{\Gamma }):𝔽_p^\times ]`$.
#### 3.1.5 Regularity
We have a better control of the monodromy group $`\mathrm{\Gamma }`$, if we assume that — in addition to Condition 2.2.2 — the following holds.
###### Condition 3.1.6
* The center $`C_G`$ of $`G`$ is trivial.
* Let $`G^{}`$ be a subgroup of $`G`$ which contains an element of order $`p`$ and an element of one of the classes $`C_i`$, $`i=1,\mathrm{},4`$. Then $`G=G^{}`$.
###### Proposition 3.1.7
Assume that Condition 3.1.6 holds. Let $`N`$ be the level of $`Y`$. Then $`\chi ^{\mathrm{adm}}|_\mathrm{\Gamma }=1`$ and
$$\mathrm{\Gamma }\{\begin{array}{ccc}/2& \text{if}\hfill & N=p,\hfill \\ 1& \text{if}\hfill & N>p.\hfill \end{array}$$
We will prove Proposition 3.1.7 in Section 3.3.4.
###### Corollary 3.1.8
If Condition 3.1.6 holds, then $`\overline{H}`$ is regular.
###### Proof.
It follows from Theorem 3.1.2 and Propositions 3.1.4 and 3.1.7 that $`R_Y^\mathrm{\Gamma }`$ has a regular sequence $`(t,\pi ^{})`$, where $`\pi ^{}\pi ^\mu `$ and $`\mu =1,2`$. $`\mathrm{}`$
###### Remark 3.1.9
We do not know of any example in which either $`\overline{H}`$ or $`\overline{H}_4^{\mathrm{ab}}(\underset{¯}{C})`$ is not regular. On the other hand, the Hurwitz spaces $`\overline{H}_{[4]}^{\mathrm{in}}(\underset{¯}{C})`$ and $`\overline{H}_{[4]}^{\mathrm{ab}}(\underset{¯}{C})`$ behave like modular curves. We have to replace the $`\lambda `$-line with the $`j`$-line. For the special values $`j=0`$ and $`j=1728`$, the monodromy action may become more complicated then in Proposition 3.1.4. Therefore, $`\overline{H}_{[4]}^{\mathrm{in}}(\underset{¯}{C})`$ and $`\overline{H}_{[4]}^{\mathrm{ab}}(\underset{¯}{C})`$ are very often not regular. However, Theorem 3.1.1 remains true.
### 3.2 Dihedral covers and generalized elliptic curves
In this section we show that Theorem 3.1.2 is true in the case the group $`G`$ is a dihedral group of order $`2N`$, where $`p||N`$, and $`\underset{¯}{C}`$ consists of $`4`$ times the conjugacy class of reflections. In order to prove this, we relate the deformation theory of a bad cover, which arises as the reduction of a $`G`$-cover of type $`\underset{¯}{C}`$, to the deformation theory of a generalized elliptic curve endowed with a certain level structure. Once this is achieved, Theorem 3.1.2 follows from the results of on the reduction of the modular curve $`X_1(p)`$.
#### 3.2.1 Generalized elliptic curves
We start by recalling some definitions from , and . A generalized elliptic curve over a scheme $`S`$ is a semistable curve $`E/S`$ of genus $`1`$ together with a section $`0:SE^{\mathrm{sm}}`$ and an $`S`$-morphism $`+:E^{\mathrm{sm}}\times _SEE`$, verifying the following properties. The geometric fibers of $`E/S`$ are either smooth or “$`n`$-gons”. The restriction of $`+`$ to $`E^{\mathrm{sm}}`$ gives $`E^{\mathrm{sm}}`$ the structure of a commutative group scheme with identity $`0`$. Moreover, $`E^{\mathrm{sm}}`$ acts on $`E`$ by “rotation”, see , Definition II.1.12. If $`E/S`$ is smooth of genus $`1`$ and $`0:SE`$ a section, then there exists one and only one such group law $`+:E\times EE`$, and we call $`(E,0)`$ an elliptic curve over $`S`$.
Let $`A`$ be a finite Abelian group and $`E/S`$ a generalized elliptic curve. A weak $`A`$-structure on $`E`$ is a group homomorphism $`\varphi :AE^{\mathrm{sm}}(S)`$ such that the Cartier divisor $`\varphi (A):=_{aA}\varphi (a)`$ is a subgroup-scheme of $`E^{\mathrm{sm}}`$. A weak $`A`$-structure $`\varphi `$ is called an $`A`$-structure if $`\varphi (A)`$ meets every irreducible component of every geometric fiber of $`E/S`$. The following two examples are classical. For $`A=/n`$, an $`A`$-structure is called a $`\mathrm{\Gamma }_1(n)`$-structure. For $`A=/n\times /n`$, an $`A`$-structure is called a $`\mathrm{\Gamma }(n)`$-structure.
#### 3.2.2 $`\mathrm{\Gamma }_2(N)`$-structures
We fix an integer $`N>0`$ and an odd prime number $`p`$ such that $`p||N`$. We define the Abelian group
$$A:=/2N\times /2.$$
Similarly, we let $`A^{}:=/2N^{}\times /2`$, where $`N=pN^{}`$, and let $`\tau :AA^{}`$ be the natural projection.
###### Definition 3.2.1
A $`\mathrm{\Gamma }_2(N)`$-structure on a generalized elliptic curve $`E/S`$ is an $`A`$-structure $`\varphi :AE(S)`$, with $`A`$ as above. We say that $`\varphi `$ is étale if the subscheme $`\varphi (A)E`$ is étale over $`S`$ (equivalently, the induced map $`\varphi _s:AE_s(K)`$ is injective for all geometric points $`s:SpeckS`$). We say that $`\varphi `$ is $`p`$-local if $`Ker\varphi A`$ has order $`p`$.
The following is obvious from :
###### Remark 3.2.2
1. If $`2N`$ is invertible on $`S`$ then $`\varphi `$ is étale.
2. If $`\varphi `$ is $`p`$-local then $`S`$ is an $`𝔽_p`$-scheme. Moreover, $`\varphi =\varphi ^{}\tau `$, where $`\varphi ^{}:A^{}E(S)`$ is an étale $`\mathrm{\Gamma }_2(N^{})`$-structure.
###### Proposition 3.2.3
Let $`G`$ be a dihedral group of order $`2N`$ and $`f:E_K^1`$ a $`G`$-cover, branched in $`4`$ points with ramification index $`2`$. After a finite extension of $`K`$, there exists a map $`\varphi :AE(K)`$, whose image $`\varphi (A)Y`$ is the set of ramification points of $`f`$, such that $`(E,\varphi )`$ is an elliptic curve with an étale $`\mathrm{\Gamma }_2(N)`$-structure.
###### Proof.
The cover $`f`$ factors through an étale $`N`$-cyclic cover $`\pi :EE^{}`$, and $`E`$ and $`E^{}`$ are smooth projective curves of genus $`1`$. We may assume that all ramification points of $`f`$ are $`K`$-rational. Let us choose one ramification point $`0E(K)`$, and set $`0^{}:=\pi (0)`$. Now we can regard $`\pi `$ is an $`N`$-cyclic isogeny between elliptic curves. Moreover, the branch points of $`f`$ are precisely the points of $`E`$ lying above the $`2`$-torsion points of $`E^{}`$. Therefore, the set of branch points is a subgroup of $`E[2N](K)`$, of order $`4N`$, and contains a point of order $`2N`$. It is clear that this subgroup is isomorphic to $`A`$. $`\mathrm{}`$
Consider the exact sequence
(11)
$$0/NA/2\times /2\mathrm{\hspace{0.33em}0}$$
of Abelian groups, where $`a/N`$ is send to $`(2a,0)A`$ and $`(a,b)A`$ is send to $`(\overline{a},\overline{b})`$. Let $`(E,\varphi )`$ be a generalized elliptic curve with $`\mathrm{\Gamma }_2(N)`$-structure. The image $`\varphi (/N)E`$ is a subgroup scheme of $`E^{\mathrm{sm}}`$, finite and flat over $`S`$. Since $`E^{\mathrm{sm}}`$ acts on $`E`$, we can form the quotient scheme $`E^{}:=E/\varphi (/N)`$. One checks fiber by fiber that $`E^{}/S`$ is again a generalized elliptic curve. Via the exact sequence (11), $`E^{}`$ is endowed with a $`\mathrm{\Gamma }(2)`$-structure $`\overline{\varphi }:/2\times /2E^{}(S)`$. Let $`[1]:E^{}\stackrel{}{}E^{}`$ be the canonical involution (see , Chapitre II) and $`X_0:=E^{}/<[1]>`$ the quotient. We write $`f_\varphi :EX_0`$ for the natural map. We choose a bijection $`\alpha :\{1,2,3,4\}/2\times /2`$ and let $`x_jX_0(S)`$ be the image of $`\overline{\varphi }(\alpha (j))`$, for $`j=1,\mathrm{},4`$.
###### Proposition 3.2.4
Suppose $`2`$ is invertible on $`S`$.
1. The curve $`(X_0;x_j)`$ is stably marked, of genus $`0`$.
2. If $`\varphi `$ is étale, then $`f_\varphi :EX_0`$ is an admissible cover, ramified of order $`2`$ along the sections $`x_j`$. There is a natural $`G`$-action on $`E`$, where $`G`$ is dihedral of order $`2N`$, such that $`X_0=E/G`$. If in addition $`E/S`$ is smooth then $`f_\varphi `$ is a tame $`G`$-cover.
###### Proof.
It suffices to prove the proposition in the case $`S=Speck`$, where $`k`$ is an algebraically closed field. Carrying a $`\mathrm{\Gamma }(2)`$-structure, $`E^{}`$ is either smooth or a $`2`$-gon. Since $`[1]`$ is the identity on $`E^{}[2]/2\times /2`$, the points $`x_1,\mathrm{},x_4X_0`$ are distinct and smooth. If $`E`$ is smooth then $`X_0_k^1`$. Otherwise, $`[1]`$ restricts to an involution on each component of $`E^{}`$ and interchanges the two singular point. In this case, $`X_0`$ is the union of two projective lines meeting transversally in one point, and each component of $`X_0`$ contains two of the points $`x_1,\mathrm{},x_4`$. This proves (i).
We have seen that $`E^{}X_0`$ is ramified at $`x_1,\mathrm{},x_4`$ of order $`2`$ and étale everywhere else. Assume that $`E/S`$ is smooth and $`\varphi (A)E`$ is étale. Then $`\pi :EE^{}`$ is an étale $`N`$-cyclic isogeny. It follows that $`f_\varphi `$ is a dihedral Galois cover, ramified at $`x_1,\mathrm{},x_4`$ of order $`2`$. In case $`E`$ is singular, the map $`f_\varphi `$ may be ramified in the singular points. However, using the description of the group law on a Néron polygon given in , Chapitre II, one checks that $`f_\varphi `$ is admissible. $`\mathrm{}`$
#### 3.2.3 Deformation
Let $`k`$ be an algebraically closed field of characteristic $`p`$. We fix a generalized elliptic curve $`E`$ over $`k`$ and a $`\mathrm{\Gamma }_2(N)`$-structure $`\varphi :AE(k)`$. We assume that $`\varphi `$ is $`p`$-local. Let $`Def(E,\varphi )`$ denote the deformation functor classifying isomorphism classes of deformations $`(E_R,\varphi _R)`$ of $`(E,\varphi )`$ over complete local $`W(k)`$-algebras $`R`$ with residue field $`k`$. Let $`Def(E,\varphi )^{\mathrm{loc}}`$ be the subfunctor corresponding to deformations $`(E_R,\varphi _R)`$ where $`\varphi _R`$ is $`p`$-local.
###### Proposition 3.2.5
The functor $`Def(E,\varphi )`$ has a universal deformation ring $`R_\varphi `$. The ring $`R_\varphi `$ is regular of dimension $`2`$; there exists a regular sequence $`(t,\pi )`$ with the following properties.
1. The ring $`R_\varphi ^{\mathrm{loc}}:=R_\varphi /(\pi )k[[t]]`$ is the universal deformation ring for $`Def(E,\varphi )^{\mathrm{loc}}`$.
2. If $`E`$ is ordinary, then $`p=u\pi ^{p1}`$, where $`uR_\varphi ^\times `$.
3. If $`E`$ is supersingular, then $`p=u\pi ^{p1}`$, and $`(u,\pi )`$ is another regular sequence for $`R_\varphi `$.
###### Proof.
Let us first assume that $`E`$ is smooth. We write $`Def(E,\varphi |_{/p})`$ for the functor classifying deformations of $`E`$ together with the $`\mathrm{\Gamma }_1(p)`$-structure $`\varphi |_{/p}`$. Since $`AA^{}\times /p`$ and the order of $`A^{}`$ is prime-to-$`p`$, the morphism $`Def(E,\varphi )Def(E,\varphi |_{/p})`$ that sends $`(E_R,\varphi _R)`$ to $`(E_R,\varphi _R|_{/p})`$ is an equivalence. Therefore, if $`E`$ is smooth, the proposition is a direct consequence of , Section 13.5. Namely, the universal deformation ring $`R_\varphi `$ can be identified with the strict complete local ring of the moduli stack $`(\mathrm{\Gamma }_1(p))`$ at the point corresponding to $`(E,\varphi |_{/p})`$.
It is clear from how to extend this to the general case. Actually, since $`p||N`$, the situation here is somewhat easier than in . Since $`\varphi `$ is $`p`$-local, the number of components of $`E`$ is prime-to-$`p`$ (more precisely, if $`E`$ is singular, the number of components is $`2n`$, where $`n|N^{}`$). By , Théorème III.1.2, the generalized elliptic curve $`E`$ admits a universal deformation $`E_{R_0}`$ over $`R_0=W(k)[[t]]`$. As in the smooth case, the morphism $`Def(E,\varphi )Def(E,\varphi |_{/p})`$ is an equivalence, if we regard $`\varphi |__p`$ as a weak $`/p`$-structure. It follows from that $`Def(E,\varphi )`$ admits a universal deformation $`(E_{\mathrm{univ}},\varphi _{\mathrm{univ}})`$ over a ring $`R_\varphi `$. In fact, $`SpecR_\varphi `$ is a closed subscheme of $`E_{R_0}[p]^\times `$ and $`E_{\mathrm{univ}}=E_{R_0}_{R_0}R_\varphi `$. We regard $`t`$ as an element of $`R_\varphi `$ via the natural morphism $`R_0=W(k)[[t]]R_\varphi `$. Let us choose a formal parameter $`T`$ of $`E_{\mathrm{univ}}`$ along the $`0`$-section (see , Section 2.2.3). The point $`P_{\mathrm{univ}}:=\varphi _{\mathrm{univ}}(2N^{},0)`$ is a point of exact order $`p`$ on $`E_{\mathrm{univ}}`$. Since $`\varphi `$ is $`p`$-local, $`P_{\mathrm{univ}}|_E=\varphi (2N^{},0)=0`$. Hence we may regard $`\pi :=T(P_{\mathrm{univ}})`$ as an element of $`R_\varphi `$. We claim that $`(t,\pi )`$ is a regular sequence for $`R_\varphi `$ such that (i), (ii) and (iii) hold. We have already mentioned that, if $`E`$ is smooth, this is proved in . If $`E`$ is singular, the situation is essentially the same as for $`E`$ smooth and ordinary, see . Namely, we may choose $`T`$ such that $`E_{R_0}[p]^\times =SpecR_0[T\mathrm{\Phi }_p(1+T)=0]`$, where $`\mathrm{\Phi }_p(X)=(X^p1)/(X1)`$. Therefore, $`R_\varphi =W(k)[\zeta _p][[t]]`$ and $`\pi =\zeta _p1`$. This completes the proof of the proposition. $`\mathrm{}`$
Let $`(E_R,\varphi _R)`$ be a deformation of $`(E,\varphi )`$. Following Section 3.2.2 we associate to every deformation $`(E_R,\varphi _R)`$ of $`(E,\varphi )`$ a finite map $`f_{\varphi _R}:E_RX_{0,R}`$ and sections $`x_{1,R},\mathrm{},x_{4,R}:SpecRX_{0,R}`$ such that $`(X_{0,R};x_{j,R})`$ is a stably marked curve of genus $`0`$. This gives rise to a morphism
(12)
$$Def(E,\varphi )Def(X_0;x_j)$$
of deformation functors and hence to a $`W(k)`$-algebra morphism $`R_{X_0}R_\varphi `$.
###### Proposition 3.2.6
The ring $`R_\varphi `$ is a finite and flat extension of $`R_{X_0}W(k)[[w]]`$. Modulo $`\pi `$, we obtain a finite extension
$$R_{X_0}_{W(k)}kk[[w]]R_\varphi ^{\mathrm{loc}}k[[t]]$$
with inseparability degree $`p`$. Its separable part $`k[[w]](R_\varphi ^{\mathrm{loc}})^p=k[[t^p]]`$ is tamely ramified of degree $`d`$. Here
$$d=\{\begin{array}{cc}1& \text{if }E\text{ is smooth,}\hfill \\ N^{}/n& \text{if }E\text{ is a }2n\text{-gon.}\hfill \end{array}$$
###### Proof.
Let $`E^{\prime \prime }:=E/\varphi (/p)`$ be the quotient by the subgroup scheme $`\varphi (/p)E^{\mathrm{sm}}`$ and $`\varphi ^{\prime \prime }:A^{}E^{\prime \prime }(k)`$ the induced $`\mathrm{\Gamma }_2(N^{})`$-structure. Clearly, the map $`f_\varphi :EX_0`$ induced by $`\varphi `$ factors through the projection $`EE^{\prime \prime }`$. The resulting map $`f_{\varphi ^{\prime \prime }}:E^{\prime \prime }X_0`$ is the map induced by $`\varphi ^{\prime \prime }`$. We see that the morphism (12) can be written as the composition of two morphisms, as follows:
(13)
$$Def(E,\varphi )Def(E^{\prime \prime },\varphi ^{\prime \prime })Def(X_0;x_j).$$
Let $`R_{\varphi ^{\prime \prime }}`$ be the universal deformation ring of $`Def(E^{\prime \prime },\varphi ^{\prime \prime })`$. By Proposition 3.2.4 (ii), the map $`f_{\varphi ^{\prime \prime }}:E^{\prime \prime }X_0`$ is admissible, and is “Galois” with dihedral Galois group of order $`2N^{}`$. If $`E`$ is singular, then $`f_{\varphi ^{\prime \prime }}`$ is ramified of order $`d`$ over the unique singular point of $`X_0`$, where $`d`$ is as in the statement of the proposition. It is not hard to see that any deformation of the admissible cover $`f_{\varphi ^{\prime \prime }}`$ together with the group action corresponds to a unique deformation of $`(E^{\prime \prime },\varphi ^{\prime \prime })`$. Therefore, it follows from that $`R_{\varphi ^{\prime \prime }}=R_{X_0}[zz^d=w]W(k)[[z]]`$. Here we identify $`R_{X_0}`$ with $`W(k)[[w]]`$, such that, if $`X_0`$ is singular, $`w`$ is the deformation parameter of the singular point. We are reduced to showing that $`R_{\varphi ^{\prime \prime }}R_\varphi `$ is finite and flat, and purely inseparable of degree $`p`$ modulo $`\pi `$.
Let $`B:=R_{\varphi ^{\prime \prime }}_{W(k)}kk[[z]]`$ and $`\stackrel{~}{B}:=B^{1/p}k[[z^{1/p}]]`$. Let $`(E_{\stackrel{~}{B}}^{\prime \prime },\varphi _{\stackrel{~}{B}}^{\prime \prime })`$ be the deformation of $`(E^{\prime \prime },\varphi ^{\prime \prime })`$ corresponding to the natural morphism $`R_{\varphi ^{\prime \prime }}\stackrel{~}{B}`$. We can define a generalized elliptic curve $`E_{\stackrel{~}{B}}`$ over $`\stackrel{~}{B}`$ such that $`E_{\stackrel{~}{B}}^{\prime \prime }=E_{\stackrel{~}{B}}^{(p)}`$ is the $`p`$th power Frobenius twist of $`E_{\stackrel{~}{B}}`$. The relative Frobenius $`F:E_{\stackrel{~}{B}}E_{\stackrel{~}{B}}^{\prime \prime }`$ is a $`p`$-cyclic, purely inseparable “isogeny” whose kernel is generated by the $`0`$-section of $`E_{\stackrel{~}{B}}`$ (which is a point of exact order $`p`$). Moreover, there exists a unique $`p`$-local $`\mathrm{\Gamma }_2(N)`$-structure $`\varphi _{\stackrel{~}{B}}:AE_{\stackrel{~}{B}}(\stackrel{~}{B})`$ such that $`\varphi _{\stackrel{~}{B}}^{\prime \prime }=F\varphi _{\stackrel{~}{B}}`$. Conversely, let $`(E_R,\varphi _R)`$ be any $`p`$-local deformation of $`(E,\varphi )`$. By , we can canonically identify $`E_R^{\prime \prime }:=E_R/\varphi _R(/p)`$ with $`E_R^{(p)}`$ and the quotient map $`E_RE_R^{\prime \prime }`$ with the Frobenius $`F:E_RE_R^{(p)}`$. It follows that $`\stackrel{~}{B}=R_\varphi ^{\mathrm{loc}}`$. Hence $`R_{\varphi ^{\prime \prime }}R_\varphi `$ is purely inseparable of degree $`p`$ modulo $`\pi `$. Nakayama’s Lemma shows that $`R_\varphi `$ is finite over $`R_{\varphi ^{\prime \prime }}`$. Since a finite morphism between two regular local rings of the same dimension is automatically flat (see , V.3.8), the proposition is proved. $`\mathrm{}`$
#### 3.2.4 Stabilization
Let $`(E_R,\varphi _R)`$ be a deformation of $`(E,\varphi )`$. A stabilization of $`(E_R,\varphi _R)`$ is a morphism $`q_R:Z_RE_R`$ between semistable $`R`$-curves together with a map $`\psi _R:AZ_R(R)`$ such that (i) $`(Z_R,\psi _R(A))`$ is a stably marked curve, (ii) $`\varphi _R=q_R\psi _R`$, and (iii) $`q_R`$ is the contraction of the marked semistable curve $`(Z_R,\psi _R(A^{}))`$ (see ).
###### Lemma 3.2.7
For every deformation $`(E_R,\varphi _R)`$, there exists a stabilization $`(Z_R,\psi _R)`$. Assume that there exists a dense open subset $`USpecR`$ such that $`\varphi _U`$ is étale. Then $`(Z_R,\psi _R)`$ is unique up to unique isomorphism.
###### Proof.
If $`\varphi _U`$ is étale, then $`(E_U,\varphi _U(A))`$ is stably marked, so necessarily $`Z_U=E_U`$. Therefore, the uniqueness follows from the fact that the moduli stack of stably marked curves is separated, see .
The notion of stabilization is certainly compatible with base change $`RR^{}`$. Hence it suffices to prove the existence of stabilization for the universal deformation $`(E_{\mathrm{univ}},\varphi _{\mathrm{univ}})`$. Recall that $`\pi =T(\varphi _{\mathrm{univ}}(2N^{},0))R_\varphi `$, where $`T`$ is a formal parameter of $`E_{\mathrm{univ}}`$ along the $`0`$-section. Relative to $`T`$, the formal group of $`E_{\mathrm{univ}}^{\mathrm{sm}}`$ is given by a power series $`\mathrm{\Phi }(T_1,T_2)=T_1+T_2+\mathrm{}R[[T_1,T_2]]`$. Since $`\varphi _{\mathrm{univ}}`$ is a group homomorphism, we have
(14)
$$T(\varphi _{\mathrm{univ}}(2N^{}m,0))m\pi (mod\pi ^2),m/p.$$
Let $`q_{\mathrm{univ}}:Z_{\mathrm{univ}}E_{\mathrm{univ}}`$ be the blowup of $`E_{\mathrm{univ}}`$ along the closed subscheme $`\varphi _{\mathrm{univ}}(A^{})(\pi )E_{\mathrm{univ}}`$. Since $`\varphi _{\mathrm{univ}}(A^{})`$ consists of $`|A^{}|=4N^{}`$ pairwise disjoint sections $`SpecR_\varphi E_{\mathrm{univ}}^{\mathrm{sm}}`$, the blowup $`Z_{\mathrm{univ}}`$ is a semistable curve over $`R_\varphi `$. Denote its special fiber by $`Z`$. Then
$$Z=E\underset{bA^{}}{}Z_b,$$
where $`Z_b_k^1`$ is connected to $`E`$ in the point $`z_b:=\varphi (b)`$. For $`bA^{}`$, the translate $`T_b:=T\varphi _{\mathrm{univ}}(b)`$ is a formal parameter of $`E`$ along the section $`\varphi _{\mathrm{univ}}(b)`$. By the definition of $`Z_{\mathrm{univ}}`$, $`\stackrel{~}{T}_b:=T_b/\pi `$ is a regular function in a neighborhood of $`Z_b\{z_b\}Z_{\mathrm{univ}}`$ and defines an isomorphism $`Z_b_k^1`$ mapping $`z_b`$ to $`\mathrm{}`$. For $`aA`$, let $`\psi _{\mathrm{univ}}(a)`$ be the closure of $`\varphi _{\mathrm{univ}}(a)|_KE_{\mathrm{univ}}(K)=Z_{\mathrm{univ}}(K)`$ in $`Z_{\mathrm{univ}}`$. This defines a map $`\psi _{\mathrm{univ}}:AZ_{\mathrm{univ}}(R_\varphi )`$ such that $`\varphi _{\mathrm{univ}}=q_{\mathrm{univ}}\psi _{\mathrm{univ}}`$. Write $`a=b+c`$, with $`bA^{}`$ and $`c/p`$. Using (14) one finds that $`\psi _{\mathrm{univ}}(a)`$ meets the special fiber in the point $`cZ_b(𝔽_p)𝔽_p\{\mathrm{}\}`$. By construction, $`(Z_{\mathrm{univ}},\psi _{\mathrm{univ}})`$ is a stabilization of $`(E_{\mathrm{univ}},\varphi _{\mathrm{univ}})`$. This proves the lemma. $`\mathrm{}`$
As in the proof of the lemma, let $`Z_{\mathrm{univ}}`$ be the stabilization of the universal deformation $`(E_{\mathrm{univ}},\varphi _{\mathrm{univ}})`$. Let $`f_{0,_{\mathrm{univ}}}:Z_{\mathrm{univ}}X_{0,\mathrm{univ}}`$ be the composition of $`q_{\mathrm{univ}}`$ with the map $`f_{\varphi _{\mathrm{univ}}}:E_{\mathrm{univ}}X_{0,\mathrm{univ}}`$ induced by $`\varphi _{\mathrm{univ}}`$. Let $`K`$ be the fraction field of $`R_\varphi `$. Since $`\varphi _K`$ is étale, $`Z_K=E_K`$. By Proposition 3.2.3, $`f_{0,K}:Z_KX_{0,K}`$ is a $`G`$-cover, branched at $`4`$ points of order $`2`$, where $`G`$ is dihedral of order $`2N`$. In other words, $`Z_K`$ is a $`K`$-object of the Hurwitz stack $`_4^{\mathrm{in}}(G)`$. By the uniqueness of stabilization, the action of $`G`$ extends to $`Z_{\mathrm{univ}}`$. Therefore, $`Z_{\mathrm{univ}}`$, together with the mark $`\psi _{\mathrm{univ}}(A)`$, the action of $`G`$ and the map $`f_{0,\mathrm{univ}}`$, is an $`R_\varphi `$-object of the complete Hurwitz stack $`\overline{}_4^{\mathrm{in}}(G)`$. In particular, the special fiber $`Z`$ of $`Z_{\mathrm{univ}}`$ is a $`k`$-object of $`\overline{}_4^{\mathrm{in}}(G)^{\mathrm{bad}}`$. Via pullback, we obtain a morphism
(15)
$$Def(E,\varphi )Def(Z)$$
of deformation functors, compatible with the morphisms $`Def(E,\varphi )Def(X_0;x_j)`$ and $`Def(Z)Def(X_0;x_j)`$.
###### Proposition 3.2.8
The morphism (15) is an isomorphism; it induces an isomorphism between $`Def(E,\varphi )^{\mathrm{loc}}`$ and $`Def(Z)^{\mathrm{bad}}`$.
###### Proof.
The first statement is equivalent to the assertion that $`Z_{\mathrm{univ}}`$ is the universal deformation of $`Z`$. By construction, the maximal subset $`USpecR_\varphi `$ such that $`\varphi _U`$ is étale is precisely the maximal subset such that $`Z_UX_{0,U}`$ is an admissible cover. Therefore, the first statement of the proposition implies the second.
Let $`Z_R`$ be a deformation of $`Z`$. We denote by $`\psi :AZ(k)`$ the restriction of $`\psi _{\mathrm{univ}}`$ to $`Z`$. Clearly, $`\psi `$ lifts uniquely to a map $`\psi _R:AZ_R(R)`$ such that $`\psi _R(A)`$ is the mark of the stably marked curve $`Z_R`$. Let $`q_R:Z_RE_R`$ be the contraction of the marked semistable curve $`(Z_R,\psi _R(A^{}))`$, and let $`\varphi _R:=q_R\psi _R`$. We claim that the assignment $`Z_R(E_R,\varphi _R)`$ defines a morphism $`Def(Z)Def(E,\varphi )`$, which is the inverse of (15). This is clear from the construction and the following lemma.
###### Lemma 3.2.9
There exists a unique morphism $`+_R:E_R^{\mathrm{sm}}\times E_RE_R`$ such that $`(E_R,+_R,\varphi _R(0))`$ is a generalized elliptic curve and $`\varphi _R`$ is a $`\mathrm{\Gamma }_2(N)`$-structure.
###### Proof.
By construction, $`E_R`$ is a semistable $`R`$-curve of genus $`1`$. We claim that every geometric fiber $`E_s`$ of $`E_R`$ is either smooth or a Néron polygon, and that $`\varphi _R(A)`$ meets every irreducible component of $`E_s`$. In fact, this is an open condition on $`SpecR`$, and it is true for the special fiber $`E`$.
Let us first prove the uniqueness of $`+_R`$. Suppose a morphism $`+_R`$ satisfying the conditions of the lemma exists. Since $`\varphi _R:AE_R^{\mathrm{sm}}(R)`$ is a group homomorphism (with respect to the group law induced by $`+_R`$), it induces an action of $`A`$ on $`E_R`$ via translation. Since $`\varphi _R`$ is a $`\mathrm{\Gamma }_2(N)`$-structure, $`A`$ acts transitively on the set of irreducible components of every geometric fiber. By definition of this action, $`\varphi _R`$ is equivariant (here $`A`$ acts on itself by translation). But $`(E_R,\psi _R(A^{}))`$ is stably marked, so there can exists at most one $`A`$-action on $`E_R`$ with this property. Therefore, we can apply , Théorème II.3.2, to conclude that there exists at most one morphism $`+_R`$ with the claimed properties.
The assignment $`Z_R(E_R,\varphi _R)`$ is clearly compatible with base change $`RR^{}`$. Hence it suffices the prove the existence of $`+_R`$ in the case $`R=R_Z`$, i.e. when $`Z_R`$ is the universal deformation of $`Z`$. Let $`USpecR`$ be the maximal subset such that $`Z_U`$ is a $`G`$-cover. By the construction of the complete Hurwitz stack $`\overline{}_4^{\mathrm{in}}(G)`$, $`U`$ is open and dense. Moreover, $`Z_U=E_U`$ is a smooth curve of genus $`1`$. By , Proposition II.2.7, there exists a unique structure of elliptic curve on $`(E_U,\varphi _U(0))`$. As in the proof of Proposition 3.2.3 one shows that $`\varphi _U`$ is a $`\mathrm{\Gamma }_2(N)`$-structure. In particular, $`A`$ acts on $`E_U`$ such that $`\varphi _U`$ is equivariant. Since $`(E_R,\varphi _R(A^{}))`$ is stably marked and $`USpecR`$ is dense, this action extends uniquely to $`E_R`$, and $`\varphi _R`$ is equivariant. By , Proposition II.2.7, the induced action of $`A`$ on $`\mathrm{Pic}^0E_R/R`$ is trivial. Therefore, we can apply , Théorème II.3.2, to show that there exists a structure of generalized elliptic curve on $`E_R`$ such that the action of $`aA`$ on $`E_R`$ is given by translation with $`\varphi _R(a)`$. It remains to show that $`\varphi _R`$ is a weak $`A`$-structure. But this is a closed condition on $`SpecR`$ (see ) and it is true on $`USpecR`$. Hence it is true on $`SpecR`$. This concludes the proof of Lemma 3.2.9 and Proposition 3.2.8. $`\mathrm{}`$
By Proposition 3.2.8, we can identify the universal deformation rings $`R_\varphi `$ and $`R_Z`$, and regard $`Z_{\mathrm{univ}}`$ as the universal deformation of $`Z`$. In view of Proposition 3.2.5 and Proposition 3.2.6, we obtain:
###### Corollary 3.2.10
Theorem 3.1.2 is true for $`Y=Z`$.
###### Remark 3.2.11
Let $`X_0(N)_{_{(p)}}`$ and $`X_1(N)_{_{(p)}}`$ be the arithmetic models over $`_{(p)}`$ of the modular curves $`X_0(N)`$ and $`X_1(N)`$, as defined in and . The results of this section imply that $`X_0(N)_{_{(p)}}\overline{H}_{[4]}^{\mathrm{ab}}(\underset{¯}{C})`$ and $`X_1(N)_{_{(p)}}\overline{H}_{[4]}^{\mathrm{in}}(\underset{¯}{C})`$, where
(16)
$$\underset{¯}{C}:=\{\begin{array}{cc}(2A,2A,2A,2A)& \text{if }N\text{ is odd,}\hfill \\ (2A,2A,2B,2B)& \text{if }N\text{ is even,}\hfill \end{array}$$
and $`2A`$, $`2B`$ denote the conjugacy classes of the “reflections” in a dihedral group of order $`2N`$. Let $`X_2(N)`$ be the coarse moduli space for generalized elliptic curves with $`\mathrm{\Gamma }_2(N)`$-structure. One can show that $`X_2(N)_{_{(p)}}\overline{H}_4^{\mathrm{in}}(\underset{¯}{\overset{~}{C}})`$, where $`\underset{¯}{\overset{~}{C}}=(2\stackrel{~}{A},2\stackrel{~}{A},2\stackrel{~}{B},2\stackrel{~}{B})`$ is the tuple of conjugacy classes in a dihedral group of order $`4N`$, as in (16).
### 3.3 Proof of the Reduction Theorem
We are now ready to complete the proof of the Reduction Theorem. The main argument is given in Section 3.3.2, where we compare the deformation theory of $`Y`$ to the deformation theory of the special fiber $`Z`$ of the associated auxiliary cover. Since we have modular reduction, Theorem 3.1.2 follows from the results of Section 3.2. In Section 3.3.3 and Section 3.3.4 we prove Proposition 3.1.4 and Proposition 3.1.7, using the same method: we first reduce the statements to the dihedral case and then use the results of Section 3.2.
#### 3.3.1 The auxiliary cover
Let $`f_0:YX_0`$ be as in Section 3.1.2. As we have seen in the paragraph following Definition 1.2.2, there exists a complete discrete valuation ring $`R`$ with residue field $`k`$ and quotient field $`K`$ of characteristic $`0`$ such that $`f_0:YX_0`$ is the reduction of a $`G`$-cover $`f_K:Y_K_K^1`$. More precisely, the $`G`$-cover $`f_K`$ has a stable model $`f_{0,R}:Y_RX_{0,R}`$ over $`R`$, with special fiber $`f_0`$. We denote by $`f_R:Y_RX_R`$ the quotient model of $`f_K`$ and by $`f:YX`$ its special fiber (see Section 1.1). Note that $`f_0`$ factors through $`f`$ and that $`f`$ does not depend on the choice of the lift $`f_K`$, by Proposition 1.3.2.
We are assuming that Condition 2.2.2 holds. Therefore, it follows from Proposition 2.4.3 that the $`G`$-cover $`f_K`$ has modular reduction of level $`N`$, where $`N`$ is some integer diving $`|G|`$ and divisible by $`p`$ (see Definition 2.4.1). Let $`g_K:Z_K_K^1`$ be the auxiliary cover associated to $`f_K`$. The cover $`g_K:Z_K_K^1`$ is a $`\mathrm{\Delta }`$-cover, where $`\mathrm{\Delta }G`$ is a dihedral group of order $`2N`$; it has the same branch locus as $`f_K`$, but with ramification of order $`2`$. Let $`\underset{¯}{C}_{\mathrm{aux}}`$ be the inertia type of $`g_K`$, and let $`\overline{}_{\mathrm{aux}}:=\overline{}_4^{\mathrm{in}}(\underset{¯}{C}_{\mathrm{aux}})_{_{(p)}}`$ be the complete Hurwitz stack over $`_{(p)}`$ for $`\mathrm{\Delta }`$-covers with inertia type $`\underset{¯}{C}_{\mathrm{aux}}`$. By definition, $`g_K`$ is a $`K`$-object of $`\overline{}_{\mathrm{aux}}`$. Let $`g_R:Z_RX_R`$ and $`g_{0,R}:Z_RX_{0,R}`$ be the quotient and the stable model of $`g_K`$. By definition, $`g_{0,R}`$ is an $`R`$-object of $`\overline{}_{\mathrm{aux}}`$ extending $`g_K`$; its special fiber $`g_0:ZX_0`$ is a $`k`$-object of $`\overline{}_{\mathrm{aux}}^{\mathrm{bad}}`$. In the sequel, we will denote this object simply by $`Z`$.
According to Section 3.2.2, there exists a map $`\varphi _K:AZ_K(K)`$ such that $`(Z_K,\varphi _K)`$ is an elliptic curve with $`\mathrm{\Gamma }_2(N)`$-structure, and the cover $`g_K:Z_KX_K`$ is induced by $`\varphi _K`$. We extend $`\varphi _K`$ to a map $`\psi _R:AZ_R(R)`$; let $`q_R:Z_RE_R`$ be the contraction of $`(Z_R,\psi _R(A^{}))`$ and let $`\varphi _R:=q_R\psi _R`$. By Lemma 3.2.9, $`(E_R,\varphi _R)`$ is a generalized elliptic curve with $`\mathrm{\Gamma }_2(N)`$-structure. Let $`(E,\varphi )`$ be the special fiber of $`(E_R,\varphi _R)`$. Clearly, we are in the situation of Section 3.2.4. In particular, Theorem 3.1.2 is true for $`Y=Z`$, by Corollary 3.2.10.
#### 3.3.2 Proof of Theorem 3.1.2
By the construction of the auxiliary cover, there exists an étale map $`U^{(1)}X`$, covering $`X_0`$, and a $`G`$-equivariant isomorphism
(17)
$$Y\times _XU^{(1)}Ind_\mathrm{\Delta }^G(Z\times _XU^{(1)})$$
of $`X`$-schemes. Over the open subset $`U^{(2)}:=XX_0`$, the map $`g`$ is tamely ramified along the mark $`CX`$. Let $`X_{\mathrm{univ}}:=Z_{\mathrm{univ}}/\mathrm{\Delta }`$. Since $`R_Z`$ is regular, the natural morphism $`g_{\mathrm{univ}}:Z_{\mathrm{univ}}X_{\mathrm{univ}}`$ is a quotient model of $`Z_{\mathrm{univ}}`$, by Proposition 1.3.2. In particular, $`X_{\mathrm{univ}}`$ is a semistable curve and carries a natural mark $`C_{\mathrm{univ}}X_{\mathrm{univ}}`$.
###### Construction 3.3.1
Let $`R𝒞_k`$ be Artinian, and let $`Z_R`$ be a deformation of $`Z`$ over $`R`$. There exist a unique morphism $`R_ZR`$ such that $`Z_R=Z_{\mathrm{univ}}_{R_Z}R`$. Let $`X_R:=X_{\mathrm{univ}}_{R_Z}R`$. Then $`Z_RX_R`$ is a quotient model of $`Z_R`$. Note that the pair $`(U^{(1)},U^{(2)})`$, where $`U^{(1)}X`$ and $`U^{(2)}X`$ are as above, is an étale covering of $`X`$. For $`i=1,2`$, let $`Y^{(i)}:=Y\times _XU^{(i)}`$ and $`Z^{(i)}:=Z\times _XU^{(i)}`$. For $`i,j=1,2`$, let $`U^{(i,j)}:=U^{(i)}\times _XU^{(j)}`$, $`Y^{(i,j)}:=Y\times _XU^{(i,j)}`$ and $`Z^{(i,j)}:=Z\times _XU^{(i,j)}`$. Since $`R`$ is Artinian, the étale covering $`(U^{(1)},U^{(2)})`$ of $`X`$ extends uniquely to an étale covering $`(U_R^{(1)},U_R^{(2)})`$ of $`X_R`$. Actually, $`U_R^{(2)}=X_RX_0`$. Over $`U^{(2)}`$, the finite map $`f:YX`$ is tamely ramified along $`C`$. By a theorem of Grothendieck and Murre, there exists a unique extension of $`Y^{(2)}U^{(2)}`$ to a finite morphism $`Y_R^{(2)}U_R^{(2)}`$ which is tamely ramified along $`C_RX_R`$. Define
(18)
$$Y_R^{(1)}:=Ind_\mathrm{\Delta }^G(Z_R\times _{X_R}U_R^{(1)}).$$
By (17), $`Y_R^{(1)}_Rk=Y^{(1)}`$. We claim that there exist $`G`$-equivariant isomorphisms
(19)
$$\alpha _R^{(i,j)}:Y_R^{(i)}\times _{U_R^{(i)}}U_R^{(i,j)}\stackrel{}{}Y_R^{(j)}\times _{U_R^{(j)}}U_R^{(i,j)},i,j=1,2$$
of $`U_R^{(i,j)}`$-schemes which extend the identity on $`Y^{(i,j)}`$. For $`(i,j)(1,1)`$, this follows again from Grothendieck–Murre. For $`(i,j)=(1,1)`$, we obtain (19) via the canonical identification
(20)
$$Y_R^{(1)}\times _{U_R^{(1)}}U_R^{(1,1)}Ind_H^G(Z_R\times _{X_R}U_R^{(1,1)})Y_R^{(1)}\times _{U_R^{(1)}}U_R^{(1,1)}$$
(on the left hand side we use the first projection $`U_R^{(1,1)}U_R^{(1)}`$, on the right hand side we use the second projection). It is clear that the isomorphisms (19) verify the obvious cocycle condition. Therefore, there exist a finite morphism $`f_R:Y_RX_R`$ such that $`Y_R^{(i)}=Y_R\times _{X_R}U_R^{(i)}`$, for $`i=1,2`$.
Construction 3.3.1 associates to any deformation $`Z_RDef(Z)(R)`$ over an Artinian ring $`R𝒞_k`$ a finite map $`f_R:Y_RX_R`$ which extends $`f:YX`$. Moreover, the $`G`$-action on $`Y`$ and the mark $`DY`$ extend to $`Y_R`$. Since $`X_R`$ is projective over $`R`$, we can apply Grothendieck’s Existence Theorem and extend this construction to the case of an arbitrary ring $`R\widehat{𝒞}_k`$. We claim that the constructed curve $`Y_R`$, together with the $`G`$-action, the natural morphism $`f_{0,R}:Y_RX_{0,R}`$ and the mark $`D_RY_R`$, is an object of $`\overline{}`$. It suffices to prove this in the case $`R=R_Z`$, $`Z_R=Z_{\mathrm{univ}}`$. So let $`Y_{R_Z}X_{\mathrm{univ}}`$ be the map associated to $`Z_{\mathrm{univ}}`$ by Construction 3.3.1. Since the generic fiber of $`g_{\mathrm{univ}}:Z_{\mathrm{univ}}X_{\mathrm{univ}}`$ is a $`\mathrm{\Delta }`$-cover, the generic fiber of the map $`Y_{R_Z}X_{\mathrm{univ}}`$ is a $`G`$-cover, i.e. an object of $``$. By construction, $`Y_{R_Z}X_{\mathrm{univ}}`$ is an object of $`𝒮_{\left[r\right]}^{\mathrm{in}}(G)`$. Therefore, $`Y_{R_Z}X_{\mathrm{univ}}`$ is an object of $`\overline{}`$. We have shown that Construction 3.3.1 defines a morphism of functors
(21)
$$Def(Z)Def(Y).$$
To complete the proof of Theorem 3.1.2, we have to show that (21) is an isomorphism. We need two lemmas.
###### Lemma 3.3.2
Let $`R𝒞_k`$ be an Artinian $`k`$-algebra. Let $`Z_R`$ be a deformation of $`Z`$ over $`R`$ and $`Y_R`$ its image under (21). If $`Y_RY_kR`$ is a trivial deformation, then $`Z_RZ_kR`$ is trivial, too.
###### Proof.
It suffices to show that the deformation $`(E_R,\varphi _R)`$ of $`(E,\varphi )`$ corresponding to $`Z_R`$ is trivial. Since $`Y_RY_kR`$, the closed embedding $`EY`$ extends to a closed embedding $`E_kRY_R`$. By Construction 3.3.1 and descent, we obtain a closed embedding $`E_kRZ_R`$ extending the closed embedding $`EZ`$. Composition with the contraction morphism $`q_R:Z_RE_R`$ yields a morphism $`E_kRE_R`$ which restricts to the identity on the special fiber. Since both $`E_kR`$ and $`E_R`$ are flat and of finite type over $`R`$, it is an isomorphism. By construction, this isomorphism identifies $`\varphi _R`$ with $`\varphi _kR`$. This proves the lemma. $`\mathrm{}`$
###### Lemma 3.3.3
Let $`R\widehat{𝒞}_k`$ be a normal domain, with fraction field of characteristic $`0`$. Let $`Y_R`$ be a deformation of $`Y`$ over $`R`$. Then there exists a deformation $`Z_R`$ of $`Z`$ over $`R`$ whose image under (21) is isomorphic to $`Y_R`$.
###### Proof.
Let $`X_R^{}:=Y_R/G`$. By Proposition 1.3.2, the natural map $`f_R:Y_RX_R^{}`$ is a quotient model of $`Y_R`$. In particular, $`X_R^{}`$ is a semistable curve with special fiber $`X`$. The image of $`D_RY_R`$ is a mark $`C_R^{}X_R^{}`$. There exists a maximal open subset $`U_R^{}X_R^{}`$, containing $`C_R^{}`$, over which $`f_R`$ is tamely ramified along $`C_R^{}`$. Clearly, $`U_R^{}`$ contains the generic fiber and $`U_R^{}_Rk=XX_0`$.
We claim that there exists a finite morphism $`g_R:Z_RX_R^{}`$ extending $`g:ZX`$, with $`\mathrm{\Delta }`$ acting on $`Z_R`$, characterized by the following two properties: (i) over $`U_R^{}`$, the map $`g_R`$ is tamely ramified along $`C_R^{}`$, (ii) over an étale neighborhood of $`X_R^{}U_R^{}`$, $`Y_R`$ is isomorphic to $`Ind_\mathrm{\Delta }^G(Z_R)`$. In fact, one can construct $`g_R`$ using the same method as in Construction 3.3.1. It follows that the generic fiber of $`g_R`$ is a $`\mathrm{\Delta }`$-cover, hence an object of $`_{\mathrm{aux}}`$. Therefore, $`Z_RDef(Z)(R)`$. By construction, $`g_R:Z_RX_R^{}`$ is a quotient model of $`Z_R`$. The uniqueness of the quotient model implies $`X_R^{}=X_{\mathrm{univ}}_{R_Z}R`$. A formal verification shows that $`Y_R`$ is the image of $`Z_R`$ under (21). $`\mathrm{}`$
We are now going to complete the proof of Theorem 3.1.2. The morphism (21) induces a homomorphism $`R_YR_Z`$ of local rings. We have to show that it is an isomorphism. Let $`𝔪_YR_Y`$, $`𝔪_ZR_Z`$ be the maximal ideals. Lemma 3.3.2 implies that for every $`N>0`$ the $`k`$-module $`R_Z/𝔪_YR_Z`$ is generated by the images of $`1`$ and $`(𝔪_Z)^N`$. It follows that $`R_Z/𝔪_YR_Z=k`$. So $`R_YR_Z`$ is surjective, by Nakayama’s Lemma.
Let $`𝔭SpecR_Y`$ be a generic point. The quotient $`A:=R_Y/𝔭`$ is a complete local domain with residue field $`k`$ and fraction field $`K`$ of characteristic $`0`$. The integral closure $`\stackrel{~}{A}`$ of $`A`$ in $`K`$ is again a complete local domain with residue field $`k`$. So by Lemma 3.3.3, the morphism $`R_Y\stackrel{~}{A}`$ factors via $`R_YR_Z`$. Therefore, $`I:=Ker(R_YR_Z)𝔭`$. It follows that $`I`$ is contained in the nilradical of $`R_Y`$. But $`R_Y`$ is reduced, so $`I=0`$ and (21) is an isomorphism. This completes the proof of Theorem 3.1.2. $`\mathrm{}`$
#### 3.3.3 Proof of Proposition 3.1.4
We fix an element $`\gamma \mathrm{\Gamma }^{\mathrm{ab}}`$ and choose an automorphism $`\sigma :Y\stackrel{}{}YAut_k^{\mathrm{ab}}(Y)`$ which induces $`\gamma `$. Such an automorphism $`\sigma `$ is unique up to composition with an element of $`G`$ (later in the proof we will give a “canonical” choice). Recall that the generalized elliptic curve $`E`$ is a closed subscheme of $`Y`$. Let $`E_0`$ be the identity component of $`E`$. Since $`E_0`$ is an irreducible component of $`Y`$ above $`X_0`$, we may assume $`\sigma |_{E_0}=Id_{E_0}`$, after composing $`\sigma `$ with an appropriate element of $`G`$. Since $`E`$ is either irreducible or a Néron polygon, it follows that $`\sigma `$ induces an automorphism $`\sigma |_E:E\stackrel{}{}E`$ of the $`k`$-curve $`E`$ which normalizes the action of $`\mathrm{\Delta }=D(E)`$ and commutes with $`f_0|_E:EX_0`$. By the construction of $`Z`$, $`\sigma |_E`$ extends uniquely to an automorphism $`\sigma _Z:Z\stackrel{}{}Z`$ such that $`\sigma _Z`$ and $`\sigma `$ agree in an étale neighborhood of $`E`$. Note that $`\sigma _Z`$ normalizes the action of $`\mathrm{\Delta }`$ and commutes with $`g_0:ZX_0`$. In other words, $`\sigma _ZAut_k^{\mathrm{ab}}(Z)`$. Therefore, $`\sigma _Z`$ lifts to a $`\gamma ^{}`$-semilinear automorphism $`\sigma _{Z,\mathrm{univ}}:Z_{\mathrm{univ}}\stackrel{}{}Z_{\mathrm{univ}}`$, for some $`W(k)`$-algebra automorphism $`\gamma ^{}`$ of $`R_Y=R_Z`$. Let $`\widehat{E}`$ be the formal completion of $`Z_{\mathrm{univ}}`$ along $`EZ_{\mathrm{univ}}`$. It follows from Construction 3.3.1 that we can identify $`\widehat{E}`$ with the formal completion of $`Y_{\mathrm{univ}}`$ along $`EY_{\mathrm{univ}}`$, and that $`\sigma _{Z,\mathrm{univ}}|_{\widehat{E}}=\sigma _{\mathrm{univ}}|_{\widehat{E}}`$. In particular, $`\gamma ^{}=\gamma `$. Therefore, to prove Proposition 3.1.4, we may assume that $`Y=Z`$ and $`G=\mathrm{\Delta }`$.
Recall that there exists a map $`\psi _{\mathrm{univ}}:AZ_{\mathrm{univ}}(R_Z)`$ such that the following holds: (i) $`\psi _{\mathrm{univ}}(A)`$ is the mark of the stably marked curve $`Z_{\mathrm{univ}}`$, (ii) $`(Z_{\mathrm{univ}},\psi _{\mathrm{univ}})`$ is the stabilization of $`(E_{\mathrm{univ}},\varphi _{\mathrm{univ}})`$ and (iii) $`(E_{\mathrm{univ}},\varphi _{\mathrm{univ}})`$ is the universal deformation of its special fiber $`(E,\varphi )`$. Let $`\psi :AZ(k)`$ be the restriction of $`\psi _{\mathrm{univ}}`$ to $`Z`$ and let $`Z_0`$ be the component of $`Z`$ which meets $`E`$ in $`0E`$. By the proof of Lemma 3.2.7, we may identify $`Z_0`$ with $`_k^1`$ such that $`\mathrm{}Z_0`$ is the point where $`Z_0`$ meets $`E`$ and such that $`\psi (2N^{}a,0)=a𝔽_pZ_0(k)`$. It is clear that $`\sigma :Z\stackrel{}{}Z`$ restricts to an automorphism of $`Z_0`$ which fixes $`\mathrm{}`$ and permutes the points $`\psi (2N^{}a,0)`$. Therefore, $`\sigma `$ acts on $`Z_0`$ as $`zcz+b`$, with $`b,c𝔽_p`$, $`c0`$. Composing $`\sigma `$ by an element of the decomposition group $`D(Z_0)G`$, we may assume that $`b=0`$. This determines $`\sigma `$ uniquely. Moreover,
$$\chi ^{\mathrm{bad}}(\gamma ):=c^1$$
defines a homomorphism $`\chi ^{\mathrm{bad}}:\mathrm{\Gamma }^{\mathrm{ab}}𝔽_p^\times `$. Since $`\sigma _{\mathrm{univ}}`$ is an automorphism of $`Z_{\mathrm{univ}}`$ as stably marked curve, it descents to an automorphism $`\stackrel{~}{\sigma }_{\mathrm{univ}}:E_{\mathrm{univ}}\stackrel{}{}E_{\mathrm{univ}}`$ such that $`q_{\mathrm{univ}}\sigma _{\mathrm{univ}}=\stackrel{~}{\sigma }_{\mathrm{univ}}q_{\mathrm{univ}}`$, where $`q_{\mathrm{univ}}:Z_{\mathrm{univ}}E_{\mathrm{univ}}`$ is the projection morphism. Let $`P_{\mathrm{univ}}:=\varphi _{\mathrm{univ}}(2N^{},0)E_{\mathrm{univ}}`$. Since $`\varphi _{\mathrm{univ}}`$ is a group homomorphism, we have $`aP_{\mathrm{univ}}=\varphi _{\mathrm{univ}}(2N^{}a,0)`$, for $`a/p`$. Using the definition of $`\chi ^{\mathrm{bad}}`$, we obtain
(22)
$$\stackrel{~}{\sigma }_{\mathrm{univ}}(aP_{\mathrm{univ}})=\chi ^{\mathrm{bad}}(\gamma )^1aP_{\mathrm{univ}},a/p.$$
In particular, $`\stackrel{~}{\sigma }_{\mathrm{univ}}`$ fixes the $`0`$-section of $`E_{\mathrm{univ}}`$. Similar to the proof of Lemma 3.2.9, one shows that $`\stackrel{~}{\sigma }`$ is a $`\gamma `$-semilinear automorphism of the generalized elliptic curve $`E_{\mathrm{univ}}`$, i.e. is compatible with the “group law” on $`E_{\mathrm{univ}}`$. In particular, $`\sigma |_E:E\stackrel{}{}E`$ is a $`k`$-linear automorphism of the generalized elliptic curve $`E`$.
Recall that $`E_{\mathrm{univ}}=E_{R_0}_{R_0}R_Z`$, where $`E_{R_0}`$ is the universal deformation of $`E`$, defined over $`R_0=W(k)[[t]]R_Z`$. It follows that $`\stackrel{~}{\sigma }_{\mathrm{univ}}`$ is induced by a $`\gamma _0`$-linear automorphism $`\sigma _{R_0}:E_{R_0}\stackrel{}{}E_{R_0}`$, where $`\gamma _0:=\gamma |_{R_0}`$ is the monodromy action of $`\sigma |_E`$ on $`R_0`$. If $`E`$ is smooth then $`\sigma |_E=Id_E`$. So in this case we have $`\gamma (t)=t`$, as claimed in Proposition 3.1.4. If $`E`$ is not smooth, it is isomorphic to the standard $`2n`$-gon, for some $`n|N^{}`$, see , Section II.1.1. So we identify the group of components $`E^{\mathrm{sm}}/E_0^{\mathrm{sm}}`$ with $`/2n`$ and the individual components $`E_i`$, $`i/2n`$, with $`_k^1`$. According to , Proposition II.1.10, the automorphism $`\sigma |_E`$ is given by the formula
$$\sigma |_E(x,i)=(\zeta ^ix,i),x^1,i/2n,$$
for some $`2n`$th root of unity $`\zeta `$. But $`\sigma |_{E_i}:E_i\stackrel{}{}E_i`$ lifts to an element of the decomposition group $`D(E_i)G`$, which is dihedral of order $`2N/n`$, by Proposition 2.5.3 (d). We conclude that $`\zeta `$ is a $`d`$th root of unity, where $`d`$ is the greatest common divisor of $`2n`$ and $`N^{}/n`$. We set $`\chi ^{\mathrm{adm}}(\gamma ):=\zeta `$. It is obvious that this defines a group homomorphism $`\chi ^{\mathrm{adm}}:\mathrm{\Gamma }^{\mathrm{ab}}\mu _d`$. If we identify $`E_{\mathrm{univ}}`$ with the Tate elliptic curve $`𝔾_m^t/q^{}`$, where $`q=t^{2n}`$, then we find $`\gamma (t)=\zeta t`$, see , Chapitre VII.
It remains to prove the formula $`\gamma (\pi )\chi ^{\mathrm{bad}}(\gamma )\pi (mod\pi ^2)`$. Recall that $`\pi =T(P_{\mathrm{univ}})`$, where $`T`$ is a formal parameter of $`E_{\mathrm{univ}}`$ along the $`0`$-section. Actually, $`T`$ was chosen as a formal parameter of $`E_{R_0}`$, see the proof of Proposition 3.2.5. It follows from the preceding discussion that we may assume $`\stackrel{~}{\sigma }_{\mathrm{univ}}^{}(T)=T`$. Using (22) and the formal group law on $`E_{\mathrm{univ}}`$, we get
(23)
$$\pi ^{}:=T(\stackrel{~}{\sigma }_{\mathrm{univ}}(P_{\mathrm{univ}}))\chi ^{\mathrm{bad}}(\gamma )^1\pi (mod\pi ^2).$$
By definition, we have
(24)
$$T\gamma (\pi ^{})=\stackrel{~}{\sigma }_{\mathrm{univ}}^{}(T\pi ^{})=u(T\pi ),uR_Z[[T]]^\times .$$
Comparing coefficients, we find $`u1(mod\pi )`$, and therefore $`\gamma (\pi ^{})\pi (mod\pi ^2)`$. With (23), we conclude that $`\gamma (\pi )\chi ^{\mathrm{bad}}(\gamma )\pi (mod\pi ^2)`$. $`\mathrm{}`$
#### 3.3.4 Proof of Proposition 3.1.7
We assume that Condition 3.1.6 holds. Recall that the curve $`X`$ consists of five components $`X_0,\mathrm{},X_4`$, where $`X_i`$ contains the specialization of the branch point $`x_i`$, for $`i=1,\mathrm{},4`$. Fix $`i\{1,\mathrm{},4\}`$ and let $`W_i`$ be a component of $`f^1(X_i)`$. The restriction of $`f:YX`$ to $`W_i`$ is a $`D(W_i)`$-Galois cover $`W_iX_i`$, branched at $`2`$ points. Over $`x_i`$, the cover $`W_iX_i`$ is tamely ramified, with inertia type $`C_i`$. Over the point where $`X_i`$ meets $`X_0`$, we have wild ramification, of order $`2p`$. It follows from Condition 3.1.6 (f) that $`D(W_i)=G`$, i.e. $`W_i=f^1(X_i)X_i`$ is a $`G`$-cover. Let $`wW_i`$ be a point where $`W_i`$ meets $`EY`$. The inertia group $`I(w)`$ is dihedral of order $`2p`$.
Since $`C_G=1`$, by Condition 3.1.6 (e), every element $`\gamma \mathrm{\Gamma }`$ is induced by a unique element $`\sigma Aut_k(Y)`$. It is clear that $`\sigma `$ fixes the component $`W_i`$ and that $`\sigma (w)`$ is again a singular point of $`Y`$. Moreover, there exists an element $`\tau G`$ such that $`\sigma ^{}:=\tau ^1\sigma Aut_k^{\mathrm{ab}}(Y)`$ fixes $`w`$. The element $`\tau `$ is unique up to composition with an element of $`I(w)`$ and normalizes $`I(w)`$. For $`hI(w)`$, we have $`\sigma ^{}h(\sigma ^{})^1=\tau ^1h\tau `$. By Condition 2.2.2 (d), we may assume that $`\sigma ^{}`$ centralizes $`I(w)`$. According to the Katz–Gabber Lemma, there exists a unique $`I(w)`$-cover $`Z_iX_i`$ which is isomorphic to $`W_iX_i`$ in an étale neighborhood of $`w`$ and tamely ramified above $`x_iX_i`$. In fact, $`Z_i`$ is a component of $`g^1(X_i)`$, where $`g:ZX`$ is the auxiliary cover associated to $`f`$. The automorphism $`\sigma ^{}|_{W_i}`$ induces an isomorphism $`\sigma _{Z_i}^{}`$ of $`Z_i`$ centralizing $`I(w)`$. Recall that there exists an isomorphism $`Z_i_k^1`$ such that the action of $`I(w)`$ on $`Z_i`$ is generated by the translation $`zz+1`$ and the reflection $`zz`$. Using this identification it is easy to see that $`\sigma _{Z_i}^{}=Id_{Z_i}`$, hence $`\sigma |_{W_i}=\tau |_{W_i}`$. It follows that $`\tau G`$ centralizes the action of $`D(W)=G`$ on $`W`$. But $`G`$ has trivial center by Condition 3.1.6 (e), so $`\sigma |_{W_i}=Id_{W_i}`$.
We have shown that $`\sigma `$ is the identity on all components of $`Y`$ except possibly on those that lie above $`X_0`$. Therefore, $`\sigma `$ restricts to an automorphism $`\sigma |_E:E\stackrel{}{}E`$ of $`E`$ which fixes all points where $`E`$ meets another component of $`Y`$. Recall that the set of points where $`E`$ meets another component is the image of a $`\mathrm{\Gamma }_2(N^{})`$-structure $`\varphi ^{}:A^{}=/2N^{}\times /2E(k)`$, with $`N^{}=N/p`$. In particular, $`\sigma `$ fixes $`0E`$ and is thus an isomorphism of $`E`$ as generalized elliptic curve (see the proof of Proposition 3.1.4 above).
We can now finish the proof of Proposition 3.1.7 by applying the classification of automorphism groups of generalized elliptic curves, see . Assume that $`N^{}>1`$. Using that $`\varphi ^{}(A^{})`$ meets every component of $`E`$ and contains points of order $`>2`$ we conclude that $`\sigma |_E=Id_E`$. Therefore, $`\chi ^{\mathrm{bad}}(\gamma )=\chi ^{\mathrm{adm}}(\gamma )=1`$, so $`\gamma =1`$.
Suppose $`N^{}=1`$. Then $`\sigma |_E<[1]>`$, where $`[1]:E\stackrel{}{}E`$ is the canonical involution of $`E`$. Therefore, $`\chi ^{\mathrm{bad}}(\gamma )=\pm 1`$ and $`\chi ^{\mathrm{adm}}(\gamma )=1`$. It remains to be shown that $`[1]:E\stackrel{}{}E`$ lifts to an element $`\sigma Aut_k(Y)`$ inducing $`\gamma \mathrm{\Gamma }`$ with $`\chi ^{\mathrm{bad}}(\gamma )=1`$. We can set $`\sigma |_{W_i}:=Id_{W_i}`$ for $`W_i:=f^1(X_i)`$, $`i=1,\mathrm{},4`$, and define $`\sigma |_{\tau (E)}`$ as the canonical involution on the generalized elliptic curve $`\tau (E)`$, for all $`\tau G`$. This completes the proof. $`\mathrm{}`$
## 4 Applications to good reduction
In this section we apply the results obtained in the previous sections to questions of good reduction of Galois covers. We extend Raynaud’s criterion for good reduction to our situation (Theorem 4.1.2). Since we are in a very special situation, we get a somewhat sharper bound. For covers with $`4`$ branch points, this result is not useful to produce covers with good reduction, in practice. The rigidity method, which can be used to construct covers over fields with low ramification, hardly ever works for $`4`$ branch points.
However, our results on the reduction of the Hurwitz space allow us to compute the number of covers with good reduction, for given type and position of the branch points (Theorem 4.2.5). The rough idea is this. In characteristic $`0`$, the structure of the Hurwitz space $`H`$ as a cover of the $`\lambda `$-line is known, and has a nice description in terms of the braid action. Our Reduction Theorem describes the structure of $`\overline{H}𝔽_p`$. The Cusp Principle links these two results. It states that a cusp of $`\overline{H}`$ has good reduction if and only if it corresponds to an admissible cover with prime-to-$`p`$ ramification over the singular point.
### 4.1 Good and bad reduction
#### 4.1.1
Let $`G`$ be a finite group and $`\underset{¯}{C}=(C_1,C_2,C_3,C_4)`$ be a $`4`$-tuple of conjugacy classes of $`G`$. Let $`K_0`$ be a field of characteristic $`0`$ such that the individual conjugacy classes $`C_i`$ are rational over $`K_0`$, i.e. $`(\underset{¯}{C})K_0`$. We fix an algebraic closure $`\overline{K}_0`$ of $`K_0`$. We choose an element $`\lambda K_0\{0,1\}`$ and define $`Cov(\underset{¯}{C},\lambda )`$ as the set of isomorphism classes of $`G`$-covers $`f:Y^1`$, defined over $`\overline{K}_0`$, of type $`(\underset{¯}{C};0,1,\mathrm{},\lambda )`$. In other words, $`Cov(\underset{¯}{C},\lambda )=\pi ^1(\lambda )`$, where
(25)
$$\pi :H_4^{\mathrm{in}}(\underset{¯}{C})_\lambda ^1\{0,1,\mathrm{}\}$$
is the natural map from the inner Hurwitz space of $`G`$-covers of type $`\underset{¯}{C}`$ to the $`\lambda `$-line, and we see $`\lambda `$ as a $`\overline{K}_0`$-rational point on $`_\lambda ^1`$.
Since the domain of definition of $`H_4^{\mathrm{in}}(\underset{¯}{C})`$ is contained in $`K_0`$ (by assumption), we obtain a natural action of $`Gal(\overline{K}_0/K_0)`$ on $`Cov(\underset{¯}{C},\lambda )`$. In more concrete terms, this action is given as follows. For $`\sigma Gal(\overline{K}_0/K_0)`$ and $`fCov(\underset{¯}{C},\lambda )`$, we can form the twisted cover $`{}_{}{}^{\sigma }𝑓:{}_{}{}^{\sigma }𝑌^1`$ by applying $`\sigma `$ to the coefficients of the equations defining $`f`$ and the action of $`G`$ on $`Y`$. To each $`fCov(\underset{¯}{C},\lambda )`$ we can associate the field of moduli of $`f`$ (relative to $`K_0`$), i.e. the fixed field of all $`\sigma Gal(\overline{K}_0/K_0)`$ such that $`{}_{}{}^{\sigma }𝑓f`$. Equivalently, $`K`$ is the field of rationality of the point on $`H_4^{\mathrm{in}}(\underset{¯}{C})`$ corresponding to $`f`$. If the center of $`G`$ is trivial then $`f`$ has a unique model $`f_K:Y_K_K^1`$ over $`K`$. See e.g. .
#### 4.1.2
In the situation of 4.1.1, we will now make the following assumptions. The field $`K_0`$ is complete with respect to a discrete valuation $`v`$. We denote by $`R_0`$ the corresponding valuation ring. The residue field $`k_0`$ of $`v`$ is assumed to be algebraically closed of characteristic $`p`$, where $`p`$ is an odd prime number. We assume that Conditions 2.2.2 and 3.1.6 hold, with respect to the class vector $`\underset{¯}{C}`$ and the prime $`p`$. Finally, we assume that
(26)
$$\lambda 0,1,\mathrm{}(modv).$$
For a $`G`$-cover $`f:Y^1`$ in $`Cov(\underset{¯}{C},\lambda )`$ we can ask whether it has good or bad reduction at $`v`$, in the sense of Section 1.1 (under Condition (26), good reduction is equivalent to admissible reduction). If a cover $`fCov(\underset{¯}{C},\lambda )`$ has good reduction, its reduction $`f_k:Y_k_k^1`$ is a $`G`$-cover over the field $`k`$ of positive characteristic, of type $`(\underset{¯}{C};0,1,\mathrm{},\overline{\lambda })`$ ($`\overline{\lambda }k`$ denotes the residue of $`\lambda K_0`$). By a theorem of Grothendieck, all $`G`$-covers over $`k`$ of type $`(\underset{¯}{C};0,1,\mathrm{},\overline{\lambda })`$ arise as the reduction of a unique $`G`$-cover $`fCov(\underset{¯}{C},\lambda )`$ with good reduction. This motivates the following question.
###### Question 4.1.1
How many covers $`fCov(\underset{¯}{C},\lambda )`$ have good reduction at $`v`$?
Theorem 4.2.5 below answers this question in an explicit way. Its proof relies on the Reduction Theorem 3.1.1 and the Cusp Principle, Proposition 4.2.1.
#### 4.1.3
Let $`\overline{H}=\overline{H}_4^{\mathrm{in}}(\underset{¯}{C})`$ be the completion of the Hurwitz space $`H=H_4^{\mathrm{in}}(\underset{¯}{C})`$, see Section 1. We understand that $`\overline{H}`$ is defined over $`\mathrm{\Lambda }:=𝒪_{(\underset{¯}{C}),𝔭}`$. Here $`(\underset{¯}{C})K_0`$ is as in Section 1.2.3 and $`𝔭`$ is the prime ideal corresponding to the restriction of $`v`$ to $`(\underset{¯}{C})`$. By Theorem 3.1.1 (i), $`\overline{H}`$ is a normal scheme of dimension $`2`$, proper and flat over $`\mathrm{\Lambda }`$. Let $`fCov(\underset{¯}{C},\lambda )`$ and be $`K`$ its field of moduli (relative to $`K_0`$). Since $`K/K_0`$ is finite, $`v`$ extends uniquely to $`K`$. We denote by $`R`$ the corresponding valuation ring. Since $`\overline{H}`$ is proper over $`\mathrm{\Lambda }`$, the morphism $`SpecKH`$ corresponding to $`f`$ extends uniquely to a morphism $`\varphi :SpecR\overline{H}`$, giving rise to a $`k`$-rational point $`s:Speck\overline{H}`$. By the definition of $`\overline{H}^{\mathrm{bad}}`$, $`f`$ has good (resp. bad) reduction if and only if $`s\overline{H}^{\mathrm{bad}}`$ (resp. $`s\overline{H}^{\mathrm{bad}}`$).
We can be more precise. Since we assume the center of $`G`$ to be trivial, $`f`$ descents to a unique $`G`$-cover $`f_K:Y_K_K^1`$ defined over $`K`$. Let $`K^{}/K`$ be the minimal extension of $`K`$ over which $`f_KK^{}`$ has a stable model $`f_{0,R^{}}:Y_R^{}_R^{}^1`$ over the valuation ring $`R^{}K^{}`$ (see Section 1.1). We denote by $`f_{0,k}:Y_k_k^1`$ the special fiber of $`f_{0,R^{}}`$. The morphism $`f_{0,k}`$ (together with the induced marks on $`Y_k`$ and $`_k^1`$ and the $`G`$-action on $`Y_k`$) is an object of the Hurwitz stack $`\overline{}`$ associated to $`\overline{H}`$ and corresponds to the $`k`$-point $`s:Speck\overline{H}`$. The stable model $`f_{0,R^{}}`$ is a deformation of $`f_{0,k}`$. Hence it corresponds to a unique morphism
(27)
$$R_{Y_k}R^{}$$
of $`W(k)`$-algebras, where $`R_{Y_k}`$ is the universal deformation ring of $`f_{0,k}`$, see Section 3.1.2. The extension $`K^{}/K`$ is Galois, and its Galois group is a subgroup of $`\mathrm{\Gamma }`$, the universal monodromy group of $`f_{0,k}`$, see Section 3.1.3. The morphism (27) is equivariant with respect to the injection $`Gal(K^{}/K)\mathrm{\Gamma }`$. The morphism
(28)
$$\widehat{𝒪}_{\overline{H},s}=R_{Y_k}^\mathrm{\Gamma }R$$
obtained from (27) by taking invariants corresponds to the morphism $`\varphi :SpecR\overline{H}`$. If $`f`$ has good reduction, then $`\mathrm{\Gamma }=1`$ and $`K^{}=K`$. On the other hand, if $`f`$ has bad reduction, it has modular reduction of level $`N`$, where $`N`$ is an integer strictly divisible by $`p`$, see Proposition 2.4.3.
We say that $`\lambda K_0`$ is ordinary (resp. supersingular) if the elliptic curve $`y^2=x(x1)(x\lambda )`$ has ordinary (resp. supersingular) reduction modulo $`v`$. By , Corollary IV.4.22, $`\lambda `$ is supersingular if and only if $`h_p(\lambda )0(modv)`$, where $`h_p(X)[X]`$ is an explicit polynomial of degree $`(p1)/2`$. By a theorem of Igusa, $`h_p(X)`$ is separable modulo $`p`$, i.e. there are exactly $`(p1)/2`$ supersingular values $`\overline{\lambda }`$.
###### Theorem 4.1.2
Assume that $`fCov(\underset{¯}{C},\lambda )`$ has bad reduction of level $`N`$. Let $`K`$ be the field of moduli of $`f`$, relative to $`K_0`$, and denote by $`e`$ the ramification index of $`p`$ in $`K`$. Then
(29)
$$e\{\begin{array}{ccc}(p1)/2& \text{if}\hfill & N=p\hfill \\ p1& \text{if}\hfill & N>p.\hfill \end{array}$$
Moreover, if $`\lambda `$ is supersingular then the inequality (29) is strict.
###### Proof.
As explained above, the cover $`fCov(\underset{¯}{C},\lambda )`$ gives rise to a local ring homomorphism $`R_{Y_k}R^{}`$, where $`R_{Y_k}`$ is the universal deformation ring of the reduction of $`f`$ and $`R^{}/R`$ is the minimal extension over which $`f`$ has a stable model. We denote by $`\pi R^{}`$ the image of the element of $`R_{Y_k}`$ with the same name, given by Theorem 3.1.2. Clearly, $`v(\pi )>0`$. Since $`\pi ^{p1}|p`$ (Theorem 3.1.2 (iii) and (iv)), the ramification index of $`p`$ in $`K^{}`$ is at least $`p1`$. But $`Gal(K^{}/K)\mathrm{\Gamma }`$, so the inequality (29) follows from Proposition 3.1.7. If $`\lambda `$ is supersingular then $`p=\pi ^{p1}u`$, where $`v(u)>0`$, showing that the inequality (29) is strict. $`\mathrm{}`$
###### Remark 4.1.3
The inequality $`e(p1)/2`$ can also be deduced from , Corollaire 4.2.5 and Théorème 5.1.1.
### 4.2 The Hurwitz space as cover of the $`\lambda `$-line
#### 4.2.1 The Hurwitz classification and braid action
Let us for the moment consider the morphism $`\pi `$ in (25) as an unramified cover of Riemann surfaces. We choose a point $`\lambda _0(\mathrm{},0)`$ on the negative real line and a point $`x_0\{xImx>0\}`$ on the upper half plane. Let $`\gamma _i`$ be the unique element of $`\pi _1(_{}^1\{0,1,\mathrm{},\lambda _0\},x_0)`$ represented by a simple loop which crosses the real line exactly twice, turning around the $`i`$th point in the list $`(0,1,\mathrm{},\lambda _0)`$, in counterclockwise orientation. The group $`\pi _1(_{}^1\{0,1,\mathrm{},\lambda _0\},x_0)`$ is generated by the $`\gamma _i`$, $`i=1,\mathrm{},4`$, with the only relation $`_i\gamma _i=1`$. With these choices made, there is a canonical bijection
(30)
$$\pi ^1(\lambda _0)Ni_4^{\mathrm{in}}(\underset{¯}{C}):=\{𝐠=(g_1,\mathrm{},g_4)G=<g_i>,g_iC_i,\underset{i}{}g_i=1\}/G,$$
(here $`G`$ acts on the set of tuples $`𝐠`$ by diagonal conjugation).
The cover $`\pi `$ induces an action of $`\mathrm{\Pi }:=\pi _1(^1\{0,1,\mathrm{}\},\lambda _0)`$ on $`\pi ^1(\lambda _0)`$, hence on $`Ni_4^{\mathrm{in}}(\underset{¯}{C})`$ via (30). To describe this action explicitly, we denote by $`_4`$ the Hurwitz braid group on $`4`$ strings, with generators $`Q_1,Q_2,Q_3`$ and relations , (3.1.a-c). We identify $`_4`$ with the fundamental group of $`𝒰_4:=\{𝐱^1()|𝐱|=4\}`$, with base point $`𝐱_0=\{0,1,\mathrm{},\lambda _0\}`$. Thus we obtain an embedding
(31)
$$\mathrm{\Pi }:=\pi _1(^1\{0,1,\mathrm{}\},\lambda _0)_4.$$
The elements
(32)
$$a_0:=Q_3Q_2Q_1^2Q_2^1Q_3^1,a_1:=Q_3Q_2^2Q_3^1,a_{\mathrm{}}:=Q_3^2,$$
lie in the image of (31) and define standard generators of $`\mathrm{\Pi }`$. In particular, $`a_0a_1a_{\mathrm{}}=1`$, and $`a_w`$ is represented by a simple loop around $`w`$, for $`w\{0,1,\mathrm{}\}`$. One can check using the formula , (3.1.d), that the induced action of $`\mathrm{\Pi }`$ on $`Ni_4^{\mathrm{in}}(\underset{¯}{C})`$ is given by
(33)
$$[g_1,g_2,g_3,g_4]a_w=\{\begin{array}{ccc}[g_1^\gamma ,g_2,g_3,g_4^\gamma ],\hfill & \gamma =g_4g_1,\hfill & \text{if}w=0,\hfill \\ [g_1,g_2^\gamma ,g_3^{[g_2^1,g_4^1]},g_4^\gamma ],\hfill & \gamma =g_2g_4,\hfill & \text{if}w=1,\hfill \\ [g_1,g_2,g_3^\gamma ,g_4^\gamma ],\hfill & \gamma =g_3g_4,\hfill & \text{if}w=\mathrm{}\hfill \end{array}$$
(we use the notation $`g^\gamma =\gamma ^1g\gamma `$).
#### 4.2.2 The cusps
Let $`\overline{\pi }:\overline{H}_\lambda ^1`$ denote the canonical map, which extends $`\pi `$. We say that a $`\overline{}`$-point $`c`$ on $`\overline{H}`$ is a cusp if $`\overline{\pi }(c)\{0,1,\mathrm{}\}`$. We write $`Cusps(\underset{¯}{C},w)`$ for the set of cusps above $`w\{0,1,\mathrm{}\}`$. The cusps are the ramification points of the finite, tamely ramified morphism $`\overline{\pi }_{}`$. By Section 4.2.1 we can identify cusps with certain braid orbits:
(34)
$$Cusps(\underset{¯}{C},w)Ni_4^{\mathrm{in}}(\underset{¯}{C})/<a_w>,w\{0,1,\mathrm{}\}.$$
Let us fix a cusp $`c`$ above $`w\{0,1,\mathrm{}\}`$, represented by a class $`[𝐠]Ni_4^{\mathrm{in}}(\underset{¯}{C})`$. The conjugacy class of the element $`\gamma G`$ associated to $`[𝐠]`$ in (33) does only depend on the orbit of $`[𝐠]`$ under the action of $`a_w`$, and is thus canonically associated to the cusp $`c`$. We call $`n:=ord(\gamma )`$ the order of the cusp $`c`$. Note that the length of the $`a_w`$-orbit of $`[𝐠]`$ divides $`n`$.
As a point on $`\overline{H}`$, the cusp $`c`$ corresponds to an admissible cover
$$f_K:Y_KX_K$$
between stably marked curves over a number field $`K`$, together with an action of $`G`$ on $`Y_K`$. The bottom curve $`X_K`$ is singular, consisting of two components $`X_{1,K},X_{2,K}_K^1`$, intersecting in one point. It is shown e.g. in , Section 4.3.3, that the order $`n`$ is the ramification index of $`f`$ above the singular point of $`X_K`$.
We choose once and for all a valuation $`\overline{v}`$ of $`\overline{}`$ extending the valuation of $`(\underset{¯}{C})`$ corresponding to the prime ideal $`𝔭`$. Since $`\overline{H}`$ is proper over $`\mathrm{\Lambda }`$, the $`\overline{}`$-valued point $`c`$ reduces (with respect to the valuation $`\overline{v}`$) to an $`\overline{𝔽}_p`$-valued point $`\overline{c}`$ on $`\overline{H}`$. We say that the cusp $`c`$ has good (resp. bad) reduction if $`\overline{c}\overline{H}^{\mathrm{bad}}`$ (resp. $`\overline{c}\overline{H}^{\mathrm{bad}}`$). If $`c`$ has bad reduction then the point $`\overline{c}\overline{H}^{\mathrm{bad}}`$ corresponds to a bad cover $`f:YX`$ of modular type of level $`N`$, for some integer $`N`$ strictly divisible by $`p`$. We say that $`c`$ has bad reduction of level $`N`$.
###### Proposition 4.2.1 (The Cusp Principle)
A cusp $`c`$ of order $`n`$ has bad reduction if and only if $`p|n`$. In this case, $`n`$ divides the level $`N`$.
###### Proof.
After a finite extension of $`K`$, the admissible cover $`f_K`$ corresponding to $`c`$ extends to an $`R`$\- object $`f_R:Y_RX_R`$ of the Hurwitz stack $`\overline{}`$, where $`R`$ is the valuation ring of $`K`$ corresponding to $`\overline{v}|_K`$. The special fiber $`f:YX`$ of $`f_R`$ is the object of $`\overline{}`$ which induces the $`\overline{𝔽}_p`$-point $`\overline{c}`$ on $`\overline{H}`$. We are exactly in the situation of Section 2.5. The Cusp Principle follows from Proposition 2.5.3. $`\mathrm{}`$
#### 4.2.3 The bad components
According to Theorem 3.1.1, $`\overline{H}^{\mathrm{bad}}\overline{𝔽}_p`$ is a smooth projective curve over $`\overline{𝔽}_p`$. We pick a connected component $`W\overline{H}^{\mathrm{bad}}\overline{𝔽}_p`$. We call $`W`$ a bad component. The geometric points of $`W`$ correspond to bad covers of modular type. It is clear that the level $`N`$ of these covers is constant on $`W`$. Therefore, we may call $`N`$ the level of $`W`$. We let $`\eta `$ be the generic point of $`W`$, and we denote by $`m`$ the multiplicity of $`W`$ inside $`\overline{H}^{\mathrm{bad}}`$. By definition, $`m`$ is the length of the local ring of $`\eta `$ on $`\overline{H}\overline{𝔽}_p`$.
###### Proposition 4.2.2
If $`N=p`$ then $`m=(p1)/2`$. Otherwise, $`m=p1`$.
###### Proof.
This follows directly from Corollary 3.1.5 (iii) and Proposition 3.1.7. $`\mathrm{}`$
By Corallary 3.1.5 (i), the natural map $`W_\lambda ^1\overline{𝔽}_p`$ is inseparable, with inseparability degree $`p`$. Moreover, the induced map $`W^{(p)}_\lambda ^1\overline{𝔽}_p`$ is tamely ramified in $`0`$, $`1`$, $`\mathrm{}`$, and étale everywhere else. We are interested in describing it in more detail. We write $`N=pN^{}`$, and let $`X_2(N^{})`$ be the coarse moduli space for generalized elliptic curves with $`\mathrm{\Gamma }_2(N^{})`$-structure, see Section 3.2. We fix a bijection $`\alpha :\{1,2,3,4\}/2\times /2`$; this determines a finite map $`X_2(N^{})_\lambda ^1`$, by Proposition 3.2.4 (i). Since $`N^{}`$ is prime-to-$`p`$, $`X_2(N^{})\overline{𝔽}_p`$ is a smooth projective curve over $`\overline{𝔽}_p`$ and the map $`X_2(N^{})\overline{𝔽}_p_\lambda ^1\overline{𝔽}_p`$ is finite, tamely ramified in $`0`$, $`1`$, $`\mathrm{}`$ and étale everywhere else.
###### Proposition 4.2.3
There exists a finite map $`h:WX_2(N^{})\overline{𝔽}_p`$, compatible with the maps to $`_\lambda ^1\overline{𝔽}_p`$. This map is the composition of a purely inseparable map of degree $`p`$ and an étale map.
###### Proof.
Let $`s:SpeckW`$ be a geometric point, corresponding to a cover $`f_0:YX_0`$ over $`k`$. Following Section 3, we associate to $`f_0`$ a generalized elliptic curve $`E`$ over $`k`$, together with a $`\mathrm{\Gamma }_2(N)`$-structure $`\varphi `$. We may assume that the ordering of the branch points $`x_1,\mathrm{},x_4X_0`$ is compatible with the bijection $`\alpha :\{1,2,3,4\}/2\times /2`$ we have chosen above. As in the proof of Proposition 3.2.6, we let $`E^{\prime \prime }:=E/\varphi (/p)`$ and $`\varphi ^{\prime \prime }`$ be the induced $`\mathrm{\Gamma }_2(N^{})`$-structure. The pair $`(E^{\prime \prime },\varphi ^{\prime \prime })`$ gives rise to a geometric point $`s^{}:SpeckX_2(N^{})\overline{𝔽}_p`$. One easily checks that $`(E^{\prime \prime },\varphi ^{\prime \prime })`$ is unique up to isomorphism. In other words, $`sh(s):=s^{}`$ is a well defined map on geometric points. We claim that this map is induced by a finite morphism $`h:WX_2(N^{})\overline{𝔽}_p`$, as in the statement of the proposition. In order to prove this claim, it suffices to show the following. Let $`R_s`$ (resp. $`R_s^{}`$) be the complete local ring of $`s`$ on $`W`$ (resp. of $`s^{}`$ on $`X_2(N^{})\overline{𝔽}_p`$). Then there exists a finite morphism $`h_s^{}:R_s^{}R_s`$ of local $`\overline{𝔽}_p`$-algebras, purely inseparable of degree $`p`$, with the following property. Let $`\overline{K}`$ be an algebraic closure of the fraction field of $`R_s`$, $`\stackrel{~}{s}:Spec\overline{K}W`$ the tautological point and $`\stackrel{~}{s}^{}:Spec\overline{K}X_2(N^{})\overline{𝔽}_p`$ the point induced by $`h_s^{}`$. Then $`h(\stackrel{~}{s})=\stackrel{~}{s}^{}`$.
Following Section 3.1.4, we identify $`R_s=\widehat{𝒪}_{W,s}`$ with $`(R_Y^{\mathrm{bad}})^\mathrm{\Gamma }=k[[t]]^\mathrm{\Gamma }`$. By Proposition 3.1.7, $`\mathrm{\Gamma }`$ acts trivially on $`k[[t]]`$, so $`R_s=k[[t]]`$. In particular, the point $`\stackrel{~}{s}`$ corresponds to the generic fiber of $`Y_{\mathrm{univ}}_{R_Y}k[[t]]`$. In the same way, we can identify $`R_s^{}`$ with $`R_{\varphi ^{\prime \prime }}^{\mathrm{\Gamma }^{\prime \prime }}`$, where $`R_{\varphi ^{\prime \prime }}`$ is the universal deformation ring of $`(E^{\prime \prime },\varphi ^{\prime \prime })`$ and $`\mathrm{\Gamma }^{\prime \prime }`$ the corresponding monodromy group. By the proof of Proposition 3.2.6, we have $`R_{\varphi ^{\prime \prime }}=k[[t^p]]`$. Moreover, $`\mathrm{\Gamma }^{\prime \prime }`$ is trivial, because
$$Aut(E^{\prime \prime },\varphi ^{\prime \prime })=\{\begin{array}{cc}/2& \text{if }N^{}=1,\hfill \\ 1& \text{otherwise.}\hfill \end{array}$$
If we define $`h_s^{}`$ as the natural injection $`k[[t^p]]k[[t]]`$, the proposition follows. $`\mathrm{}`$
###### Remark 4.2.4
One can show that the étale part $`h^{(p)}:W^{(p)}X_2(N^{})\overline{𝔽}_p`$ of $`h`$ is in fact an isomorphism. In this sense, bad components are “modular”.
#### 4.2.4 The number of covers with good reduction
Let us denote by $`Cov(\underset{¯}{C},\lambda )^{\mathrm{good}}`$ the subset of $`Cov(\underset{¯}{C},\lambda )`$ containing the covers with good reduction. Define
$$d:=|Ni_4^{\mathrm{in}}(\underset{¯}{C})|,d^{\mathrm{bad}}:=|\{[𝐠]Ni_4^{\mathrm{in}}(\underset{¯}{C})p|ord(g_3g_4)\}|.$$
We know from (30) that $`d=\mathrm{deg}\overline{\pi }=|Cov(\underset{¯}{C},\lambda )|`$. Using the Cusp Principle and the Reduction Theorem, we can show:
###### Theorem 4.2.5
We have
$$|Cov(\underset{¯}{C},\lambda )^{\mathrm{good}}|=\{\begin{array}{cc}dd^{\mathrm{bad}},\hfill & \text{if }\lambda \text{ is ordinary,}\hfill \\ d\frac{p+1}{p}d^{\mathrm{bad}},\hfill & \text{if }\lambda \text{ is supersingular.}\hfill \end{array}$$
In particular, if $`\lambda `$ is ordinary and $`Ni_4^{\mathrm{in}}(\underset{¯}{C})\mathrm{}`$ then $`Cov(\underset{¯}{C},\lambda )^{\mathrm{good}}\mathrm{}`$.
###### Proof.
According to Theorem 3.1.1, $`\overline{H}^{\mathrm{good}}`$ is a smooth curve over $`𝔽_q`$ and the natural map $`\overline{\pi }^{\mathrm{good}}:\overline{H}^{\mathrm{good}}_\lambda ^1𝔽_q`$ is finite. Let $`S:=\overline{H}^{\mathrm{good}}\overline{H}^{\mathrm{bad}}`$. By Theorem 3.1.1 (iii), $`S`$ contains exactly the points on $`\overline{H}^{\mathrm{bad}}`$ with a supersingular $`\overline{\lambda }`$-value. By definition, we have $`\overline{H}^{\mathrm{good}}S=\overline{H}^{\mathrm{adm}}𝔽_q`$. It follows that $`\overline{H}^{\mathrm{good}}S_\lambda ^1𝔽_q`$ is tamely ramified in $`0`$, $`1`$ and $`\mathrm{}`$ and étale everywhere else. Comparing the degrees of $`\overline{\pi }^{\mathrm{good}}`$ and of $`\overline{H}^{\mathrm{adm}}_\lambda ^1`$ above $`\mathrm{}`$, we get
(35)
$$\mathrm{deg}\overline{\pi }^{\mathrm{good}}=\underset{c}{}e_c=dd^{\mathrm{bad}}.$$
Here $`c`$ runs over the set of cusps above $`\mathrm{}`$ with good reduction, and $`e_c`$ denotes the ramification index of $`\overline{\pi }`$ in $`c`$ (which is equal to the length of the $`a_{\mathrm{}}`$-orbit of $`Ni_4^{\mathrm{in}}(\underset{¯}{C})`$). The second equality in (35) is a consequence of the Cusp Principle, Proposition 4.2.1.
In case $`\lambda `$ is ordinary, the statement of the theorem follows directly from (35). Assume that $`\lambda `$ is supersingular, and let $`S_\lambda `$ be the set of points in $`S`$ above $`\overline{\lambda }`$. For $`sS_\lambda `$, we denote by $`m_s`$ the ramification index of $`\pi ^{\mathrm{good}}`$ in $`s`$. We obtain
(36)
$$|Cov(\underset{¯}{C},\lambda )^{\mathrm{good}}|=\mathrm{deg}\overline{\pi }^{\mathrm{good}}\underset{sS_\lambda }{}m_s.$$
According to Corollary 3.1.5 (ii) and (iii), $`m_s`$ equals the multiplicity of the bad component $`W_s`$ meeting $`\overline{H}^{\mathrm{good}}`$ in $`s`$. By Corollary 3.1.5 (i), the natural map $`W_s_\lambda ^1\overline{𝔽}_p`$ is the composition of a purely inseparable map of degree $`p`$ and a map which is étale away from $`0`$, $`1`$ and $`\mathrm{}`$. Therefore,
(37)
$$d=\mathrm{deg}\overline{\pi }^{\mathrm{good}}+p\underset{sS_\lambda }{}m_s.$$
The equations (35), (36) and (37) together imply Theorem 4.2.5 in the supersingular case. $`\mathrm{}`$
### 4.3 Examples
In this section, we explain how the results we obtained can be used to compute the reduction of the Hurwitz space, in an explicit example. We take $`G=PSL_2(\mathrm{})`$, where $`\mathrm{}p`$ is a prime such that $`p`$ exactly divides $`|G|=(\mathrm{}^21)\mathrm{}/2`$. Recall that $`G`$ has two conjugacy classes of order $`\mathrm{}`$, which we denote by $`\mathrm{}A`$ and $`\mathrm{}B`$. We take $`\underset{¯}{C}=(\mathrm{}A,\mathrm{}A,\mathrm{}A,\mathrm{}A)`$ or $`\underset{¯}{C}=(\mathrm{}A,\mathrm{}A,\mathrm{}B,\mathrm{}B)`$. The normalizer of a $`p`$-Sylow group of $`G`$ for $`p\mathrm{}`$ is a dihedral group of order $`\mathrm{}+1`$ or $`\mathrm{}1`$, depending on whether $`p|\mathrm{}+1`$ or $`p|\mathrm{}1`$, . Note that Condition 3.1.6 is also satisfied. We are interested in computing the reduction to characteristic $`p`$ of $`\overline{H}:=\overline{H}_4^{\mathrm{in}}(\underset{¯}{C})`$.
We will explain the algorithm for computing the reduction of $`\overline{H}`$ in the special case $`\mathrm{}=11`$ and $`\underset{¯}{C}=(\mathrm{}A,\mathrm{}A,\mathrm{}B,\mathrm{}B)`$. After that, we give a table with the reduction of $`\overline{H}`$ for $`\mathrm{}31`$, to all odd primes $`p\mathrm{}`$ exactly dividing the order of $`G`$.
Take $`\mathrm{}=11`$ and $`\underset{¯}{C}=(\mathrm{}A,\mathrm{}A,\mathrm{}B,\mathrm{}B)`$. Since $`|G|=660`$, we know that $`\overline{H}`$ has good reduction to characteristic $`p2,3,5,11`$. The reduction to characteristic 2 and 11 we cannot compute using our methods. Let us first take $`p=3`$. The normalizer of a 3-Sylow group is a dihedral group of order 12, so the possible levels associated to the bad components are 3 and 6. To decide which ones occur, we want to apply the Cusp Principle 4.2.1. For this we need to know the order of the cusps in characteristic zero.
Using the program ho, , we compute the irreducible components $`Z`$ of $`\overline{H}\overline{}`$, and the ramification indices of $`Z_\lambda ^1`$, for each irreducible component. Let us relate these ramification indices to the order of the cusps. Let $`s`$ be a cusp of $`\overline{H}\overline{}`$ defined over $`K`$, and let $`f_K:Y_KX_K`$ be the corresponding admissible cover. Let $`\tau `$ be the unique singular point of $`X_K`$ and $`\rho `$ a singular point of $`Y_K`$. Let $`n`$ be the ramification index of $`\rho `$ and $`G_1`$ and $`G_2`$ the decomposition groups of the two components of $`Y_K`$ passing through $`\rho `$. Using the description of the normalizers of elements in $`G=PSL_2(\mathrm{})`$ given in , Abschnitt II.8, one easily checks that the ramification index $`e`$ of $`s`$ in $`\overline{H}_\lambda ^1`$ is equal to $`n`$, unless $`n=\mathrm{}`$ and $`G_i/\mathrm{}`$ for some $`i`$, see. Section 4.2.1. In particular, $`p|n`$ iff $`p|e`$.
$$\begin{array}{cccc}& & & \\ \multicolumn{4}{c}{\mathrm{}=11,\underset{¯}{C}=(\mathrm{}A,\mathrm{}A,\mathrm{}B,\mathrm{}B)}\\ & & & \\ \text{ramification}\hfill & \hfill \text{deg}& \hfill g& \hfill \text{num}\\ & & & \\ 2^1;;1^2\hfill & \hfill 2& \hfill 0& \hfill 1\\ & & & \\ 2^26^2;;1^43^4\hfill & \hfill 16& \hfill 1& \hfill 1\\ & & & \\ 2^31^56^25^2;;5^211^13^4\hfill & \hfill 33& \hfill 2& \hfill 1\end{array}$$
The notation is as follows. Each row corresponds to an irreducible component; the last entry of each row gives the number of isomorphic components. The first entry gives the ramification of $`\overline{H}_\lambda ^1`$ over $`0,1,\mathrm{}`$. Here $`2^1`$ means one ramification point of order 2, and $`1^2`$ means two ramification points of order 1. A “$``$” indicated that over this point, the ramification indices are the same as over the previous point. The next entries give the genus of the components and its degree over the $`\lambda `$-line.
The Cusp Principle implies that the first component $`W_1`$ has good reduction to characteristic 3, since 3 does not divide the order of any of the cusps. The second and third component, which we will denote by $`W_2`$ and $`W_3`$, have bad reduction to characteristic 3. Since both these components have a cusp of order $`n=6`$ and the order of the cusp divides the level of the bad component it reduces to, we see that in both cases there will be a bad component of level $`N=6`$. A bad component of level 6 is a cover of $`X_2(N^{})\overline{𝔽}_3`$, purely inseparable of degree $`3`$, with $`N^{}=N/p=2`$, see Remark 4.2.4. The curve $`X_2(2)\overline{𝔽}_3`$ is a cover of degree 2 of the $`\lambda `$-line, branched at 0 and 1 and unbranched at $`\mathrm{}`$. In Proposition 4.2.2 we computed the multiplicity of the bad components. A bad component of level 6 has multiplicity $`p1=2`$. We conclude that the reduction of the irreducible components $`W_2`$ and $`W_3`$ each have one bad component, and it is of level 6.
Now let us have a look at the good components. The good and the bad components intersect over the supersingular $`\lambda `$’s. In characteristic 3, there is only one supersingular $`\lambda `$, namely $`\lambda =1(mod3)`$. This means that the good and bad components meet in two points. Since the multiplicity of the bad component is two, the good components will be ramified of order two in these intersection points. From this we can compute the number of covers with good reduction, for each value of $`\lambda `$.
$$|\mathrm{Cov}(W_2,\lambda )^{\mathrm{good}}|=\{\begin{array}{cc}4\hfill & \text{ if }\lambda 1(mod3),\hfill \\ 0\hfill & \text{ if }\lambda 1(mod3).\hfill \end{array}$$
$$|\mathrm{Cov}(W_3,\lambda )^{\mathrm{good}}|=\{\begin{array}{cc}21\hfill & \text{ if }\lambda 1(mod3),\hfill \\ 17\hfill & \text{ if }\lambda 1(mod3).\hfill \end{array}$$
In general we are not able to calculate the number of good components. However, since the degree of $`W_2`$ is sufficiently small, we can describe what its good part looks like. As remarked before, the degree of $`W_2^{\mathrm{good}}`$ over $`_\lambda ^1`$ is 4 and it is ramified over $`1`$ at two points of order two. Outside the supersingular $`\lambda `$’s, the ramification is as in characteristic zero. So over $`0`$ and 1, there are two ramification points of order two, and over $`\mathrm{}`$ it is unramified. From this it follows that $`W_2^{\mathrm{good}}`$ is connected.
Now let us have a look at $`p=5`$. In this case, the components $`W_1`$ and $`W_2`$ both have good reduction, since 5 does not divide the order of any of the cusps. The component $`W_3`$ has bad reduction. Note that the normalizer of a 5-Sylow group of $`G`$ is a dihedral group of order 10, so the only possibility for the level is 5. A bad component of level 5 is a cover of $`X_2(1)\overline{𝔽}_5`$, purely inseparable of degree $`5`$. The curve $`X_2(1)\overline{𝔽}_5`$ is isomorphic to the $`\lambda `$-line. It has multiplicity $`(p1)/2=2`$. In characteristic 5, there are two supersingular $`\lambda `$-values: the primitive sixth roots of unity. In these points the good part will have an “extra” ramification of order two. So as before we are able to compute the number of covers with good reduction.
$$|\mathrm{Cov}(W_3,\lambda )^{\mathrm{good}}|=\{\begin{array}{cc}23\hfill & \text{ if }\lambda \zeta _6,\zeta _6^5(mod5),\hfill \\ 21\hfill & \text{ if }\lambda \zeta _6,\zeta _6^5(mod5).\hfill \end{array}$$
The results are summarized in the following lemma.
###### Lemma 4.3.1
Let $`W_1,W_2,W_3`$ be the three irreducible components of $`H_4^{\mathrm{in}}(PSL_2(11))\overline{}`$, as described above.
* Then $`W_1`$ has good reduction to characteristic $`p2,11`$.
* The component $`W_2`$ has good reduction to characteristic $`p2,3,11`$. In characteristic 3, it as two irreducible components: a bad component of level 6 and a good component. The degree of the good component over the $`\lambda `$-line is 4.
* The component $`W_3`$ has good reduction to characteristic $`p2,3,5,11.`$ In characteristic 3, there is one bad component, of level 6, and the degree of the good part over the $`\lambda `$-line is 21. In characteristic 5, there is one bad component, of level 5, and the degree of the good part over the $`\lambda `$-line is 23.
The Hurwitz space might have irreducible components of large degree having good reduction at many primes. For example, take $`\mathrm{}=31`$ and $`\underset{¯}{C}=(\mathrm{}A,\mathrm{}A,\mathrm{}B,\mathrm{}B)`$. The Hurwitz space in characteristic zero has an irreducible component of genus 37, whose degree over the $`\lambda `$-line is 128. The cusps of this component all have ramification index a power of 2. We conclude from the Cups Principle that this component has good reduction to all primes $`p2,31`$. For details, see the table below.
###### Example 4.3.2 (Raynaud)
Let $`H`$ be the inner Hurwitz space parameterizing Galois covers of $`^1`$ with Galois group $`A_5`$ which are branched at four points of order 3. Let $`\overline{H}`$ be its completion, over $`_{(5)}`$. (Raynaud considered the absolute Hurwitz space; it is easy to make the adaption to that case.) In characteristic zero, $`\overline{H}\overline{}`$ is connected and has degree 18 over the $`\lambda `$-line. Over $`0,1,\mathrm{}`$ this cover has ramification $`3^25^21^2`$. We conclude as above that the reduction of the Hurwitz space to characteristic 5 has one bad component of level 5. So the good degree is 8. If $`\lambda `$ reduces to a supersingular value, there are 6 covers with good reduction.
The above example was presented by Raynaud in his talk in Oberwolfach, June 1997. In the problem session of the same conference he proposed the exercise of computing the number of covers with good reduction to characteristic 3 for $`G=A_5`$ and ramification of order 5. The answer to this exercise appears in the first rows of the table below.
The following table describes the reduction of the Hurwitz space $`\overline{H}:=\overline{H}_4^{\mathrm{in}}(\underset{¯}{C})`$, where $`G=PSL_2(\mathrm{})`$ and $`\underset{¯}{C}`$ is either $`(\mathrm{}A,\mathrm{}A,\mathrm{}B,\mathrm{}B)`$ or $`(\mathrm{}A,\mathrm{}A,\mathrm{}A,\mathrm{}A)`$. Every row corresponds to an isomorphism class of irreducible components in $`\overline{H}`$; the entry “num” gives the number of isomorphic components. The first column gives the $`\mathrm{}`$, the second column gives the class vector. The third column gives the ramification of $`\overline{H}_\lambda ^1`$ in characteristic zero over $`0,1,\mathrm{}`$. Here $`a^b`$ means: $`b`$ ramification points of order $`a`$ and $``$ means: the same as the previous point. The entries “deg” and $`g`$ give the degree over the $`\lambda `$-line and the genus of the component (in characteristic zero). The last three entries describe the reduction for odd primes different from $`\mathrm{}`$ which exactly divide the order of $`G`$. No statement is made for other primes. A dash means: the component has good reduction to all such primes. Each prime $`p`$ such that the component has bad reduction to characteristic $`p`$, is listed on a separate row. Under “bad components”, for each prime, all the bad components are listed. The last entry gives the degree of the good part over the $`\lambda `$-line. This is the number of covers with good reduction, for $`\lambda `$ ordinary. The number of covers with good reduction for supersingular $`\lambda `$ can be computed by Theorem 4.2.5. A component has good reduction to characteristic $`p`$ for odd primes $`p\mathrm{}`$ which are not listed and which strictly divide the order of $`G`$. The prime $`\mathrm{}=13`$ is missing from the table because our computer refused to run ho for $`PSL_2(13)`$. The prime $`\mathrm{}=17`$ is missing because there are no primes $`p17`$ which exactly divide the order of $`PSL_2(17)`$.
$$\begin{array}{ccccccccc}& & & & & & & & \\ \hfill \mathrm{}& \text{ Ni}\hfill & \text{ramification}\hfill & \hfill \text{deg}& \hfill g& \hfill \text{num}& p\hfill & \text{ bad comp}\hfill & \hfill \text{gdeg}\\ & & & & & & & & \\ & & & & & & & & \\ \hfill 5& AABB\hfill & 2^1;;1^2\hfill & \hfill 2& \hfill 0& \hfill 1& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 5& AABB\hfill & 3^11^2;3^12^1;\hfill & \hfill 5& \hfill 0& \hfill 1& 3\hfill & 1\times N=3\hfill & \hfill 2\\ & & & & & & & & \\ \hfill 5& AAAA\hfill & 5^13^11^2;;\hfill & \hfill 10& \hfill 0& \hfill 1& 3\hfill & 1\times N=3\hfill & \hfill 7\\ & & & & & & & & \\ \hfill 7& AABB\hfill & 4^2;;1^42^2\hfill & \hfill 8& \hfill 0& \hfill 1& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 7& AABB\hfill & 1^34^23^1;;7^13^12^2\hfill & \hfill 14& \hfill 0& \hfill 1& 3\hfill & 1\times N=3\hfill & \hfill 11\\ & & & & & & & & \\ \hfill 7& AAAA\hfill & 1^2;2^2;\hfill & \hfill 2& \hfill 0& \hfill 3& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 7& AAAA\hfill & 2^23^1;;\hfill & \hfill 7& \hfill 0& \hfill 1& 3\hfill & 1\times N=3\hfill & \hfill 4\\ & & & & & & & & \\ \hfill 11& AABB\hfill & 2^1;;1^2\hfill & \hfill 2& \hfill 0& \hfill 1& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 11& AABB\hfill & 2^26^2;;1^43^4\hfill & \hfill 16& \hfill 1& \hfill 1& 3\hfill & 1\times N=6\hfill & \hfill 4\\ & & & & & & & & \\ \hfill 11& AABB\hfill & 2^31^56^25^2;;5^211^13^4\hfill & \hfill 33& \hfill 2& \hfill 1& 3\hfill & 1\times N=6\hfill & \hfill 21\\ & & & & & & 5\hfill & 1\times N=5\hfill & \hfill 23\\ & & & & & & & & \\ \hfill 11& AAAA\hfill & 1^13^1;;\hfill & \hfill 4& \hfill 0& \hfill 4& 3\hfill & 1\times N=3\hfill & \hfill 1\\ & & & & & & & & \\ \hfill 11& AAAA\hfill & 3^45^2;;\hfill & \hfill 22& \hfill 3& \hfill 1& 3\hfill & 4\times N=3\hfill & \hfill 10\\ & & & & & & 5\hfill & 1\times N=5\hfill & \hfill 12\\ & & & & & & & & \\ \hfill 19& AABB\hfill & 2^1;;1^2\hfill & \hfill 2& \hfill 0& \hfill 1& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 19& AABB\hfill & 2^410^4;;1^85^8\hfill & \hfill 48& \hfill 9& \hfill 1& 5\hfill & 1\times N=10\hfill & \hfill 16\\ & & & & & & & & \\ \hfill 19& AABB\hfill & 2^51^93^39^310^4;;5^89^33^319^1\hfill & \hfill 95& \hfill 17& \hfill 1& 5\hfill & 1\times N=10\hfill & \hfill 55\\ & & & & & & & & \\ \hfill 19& AAAA\hfill & 1^25^2;;\hfill & \hfill 12& \hfill 1& \hfill 4& 5\hfill & 1\times N=5\hfill & \hfill 2\\ & & & & & & & & \\ \hfill 19& AAAA\hfill & 3^35^89^3;;\hfill & \hfill 76& \hfill 18& \hfill 1& 5\hfill & 4\times N=5\hfill & \hfill 36\\ & & & & & & & & \\ \hfill 23& AABB\hfill & 4^2;;1^42^2\hfill & \hfill 8& \hfill 0& \hfill 1& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 23& AABB\hfill & 4^412^4;;1^86^43^82^4\hfill & \hfill 64& \hfill 13& \hfill 1& 3\hfill & 1\times N=12\hfill & \hfill 16\\ & & & & & & & & \\ \hfill 23& AABB\hfill & 1^{11}4^611^512^4;;\hfill & \hfill 138& \hfill 32& \hfill 1& 3\hfill & 1\times N=12\hfill & \hfill 90\\ & & 6^411^53^82^623^1\hfill & & & & 11\hfill & 1\times N=11\hfill & \hfill 83\\ & & & & & & & & \\ \hfill 23& AAAA\hfill & 1^2;2^1;\hfill & \hfill 2& \hfill 0& \hfill 3& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 23& AAAA\hfill & 3^11^1;;\hfill & \hfill 4& \hfill 0& \hfill 4& 3\hfill & 1\times N=3\hfill & \hfill 1\\ & & & & & & & & \\ \hfill 23& AAAA\hfill & 2^26^2;;3^41^4\hfill & \hfill 16& \hfill 1& \hfill 2& 3\hfill & 1\times N=6\hfill & \hfill 4\\ & & & & & & & & \\ \hfill 23& AAAA\hfill & 2^63^811^56^4;;\hfill & \hfill 115& \hfill 24& \hfill 1& 3\hfill & 3\times N=6,2\times N=3\hfill & \hfill 67\\ & & & & & & 11\hfill & 1\times N=11\hfill & \hfill 60\\ & & & & & & & & \\ \hfill 29& AABB\hfill & 2^1;;1^2\hfill & \hfill 2& \hfill 0& \hfill 1& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 29& AABB\hfill & 2^614^6;;1^{12}7^{12}\hfill & \hfill 96& \hfill 25& \hfill 1& 7\hfill & 1\times N=14\hfill & \hfill 12\\ & & & & & & & & \\ \hfill 29& AABB\hfill & 5^63^52^714^615^4;;\hfill & \hfill 203& \hfill 54& \hfill 1& 3\hfill & 1\times N=15,1\times N=3\hfill & \hfill 128\\ & & 1^415^47^{12}5^63^5\hfill & & & & 5\hfill & 1\times N=15,1\times N=5\hfill & \hfill 113\\ & & & & & & 7\hfill & 1\times N=14\hfill & \hfill 119\\ & & & & & & & & \\ \hfill 29& AAAA\hfill & 1^37^3;;\hfill & \hfill 24& \hfill 4& \hfill 4& 7\hfill & 1\times N=7\hfill & \hfill 3\\ & & & & & & & & \\ \hfill 29& AAAA\hfill & 1^{14}3^57^{12}5^615^429^1;;\hfill & \hfill 232& \hfill 54& \hfill 1& 3\hfill & 1\times N=15,1\times N=3\hfill & \hfill 157\\ & & & & & & 5\hfill & 1\times N=15,1\times N=5\hfill & \hfill 142\\ & & & & & & 7\hfill & 2\times N=7\hfill & \hfill 148\\ & & & & & & & & \\ \hfill 31& AABB\hfill & 16^8;;1^{16}2^84^88^8\hfill & \hfill 128& \hfill 37& \hfill 1& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 31& AABB\hfill & 1^{15}3^516^85^615^4;;\hfill & \hfill 248& \hfill 67& \hfill 1& 3\hfill & 1\times N=15,1\times N=3\hfill & \hfill 173\\ & & 5^615^431^{11}2^84^83^58^8\hfill & & & & 5\hfill & 1\times N=15,1\times N=5\hfill & \hfill 158\\ & & & & & & & & \\ \hfill 31& AAAA\hfill & 2^1;;1^2\hfill & \hfill 2& \hfill 0& \hfill 3& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 31& AAAA\hfill & 1^42^2;4^2;\hfill & \hfill 8& \hfill 0& \hfill 3& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 31& AAAA\hfill & 1^82^44^4;8^4;\hfill & \hfill 32& \hfill 0& \hfill 3& \hfill & \hfill & \hfill \\ & & & & & & & & \\ \hfill 31& AAAA\hfill & 3^52^815^45^64^88^8;;\hfill & \hfill 217& \hfill 51& \hfill 1& 3\hfill & 1\times N=15,1\times N=3\hfill & \hfill 142\\ & & & & & & 5\hfill & 1\times N=15,1\times N=5\hfill & \hfill 127\end{array}$$
DRL, University of Pennsylvania
209 South 33rd Street
Philadelphia, PA 19104-6395
bouw@math.upenn.edu
wewers@math.upenn.edu
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# Abelian Extensions of Algebras in Congruence-Modular Varieties
## Introduction
The theory of abelian extensions of algebras has a long history, going from abelian extensions of groups and modules to abelian extensions of algebras in arbitrary varieties. Extensions of modules and abelian extensions of groups and Lie algebras are standard topics in treatises on Homological Algebra. The similarity of the situation for groups and Lie algebras points to a common generalization, which is our goal in this paper. Previous attempts at such a generalization, and related work, can be found in , , , , , , , , .
The most general treatment of which we are aware, in , treats abelian extensions in the generality of an algebra in an arbitrary variety of algebras, and a Beck module over that algebra. (A *Beck module* over $`A`$ in a variety $`𝐕`$ is an abelian group object in the category $`(𝐕A)`$ of algebras over $`A`$.)
Our treatment is less general because we require $`𝐕`$ to be congruence-modular. However, the apparatus of commutator theory, in particular, the existence of a difference term, allows us to prove an important lemma (lemma 3.2) which lets us work with a simpler and more conceptual definition of an extension than that in . We are also able to treat module extensions with the same theory, in a fairly natural way.
In both cases, there is an associated cohomology theory, the cohomology group in dimension one being the group of equivalence classes of extensions. We have shown that these groups are isomorphic for $`𝐕`$ congruence-modular, but will not give the proof in this paper as it is quite tedious. The theory in (and also and ) is called *comonadic cohomology*, and we call ours *clone cohomology*.
We frame our theory not in terms of algebras over $`A`$, and Beck modules over $`A`$, but of objects in equivalent categories we call the category of $`A`$-overalgebras and the category of abelian group $`A`$-overalgebras. We will briefly explain our reasons for doing so at the end of §1, where these categories are defined.
The theory of abelian extensions is interesting because it organizes, from a certain point of view, the possible structures of a class of algebras, related to $`A`$ (or an $`A`$-overalgebra $`Q`$) and an abelian group $`A`$-overalgebra $`M`$, into an abelian group, functorial in a way that we will discuss in §9. Also, that abelian group is a cohomology group for a suitable cohomology theory derived from $`A`$ and $`M`$.
After a section of preliminaries, §1 of the paper defines the categories of $`A`$-overalgebras and abelian group $`A`$-overalgebras. §2 sketches the theory of enveloping ringoids very briefly, and is included to show how that theory can be applied to constructing abelian group $`A`$-overalgebras free on an $`A`$-tuple of sets of generators. §3 defines abelian extensions and performs some preliminary analyses of them. §4 defines and explores the formalism of factor sets of extensions. §5 introduces the definition of equivalence of extensions, and defines the set $`𝐄_𝐕(A,M)`$ of equivalence classes of extensions. §6 then shows how the set of equivalence classes of extensions can be seen as a cohomology group. §7 explores composition operations between abelian extensions and homomorphisms, and §8 discusses the group law in the set of equivalence classes of extensions. §9 contains a reformulation of the definition that defines $`𝐄_𝐕(Q,M)`$ for an $`A`$-overalgebra $`Q`$ and abelian group $`A`$-overalgebra $`M`$, functorially in $`Q`$ and $`M`$. §10 shows how module extensions can be treated. §11 presents a cohomology theory we call *clone cohomology* because its definition intimately involves the clone of the variety $`𝐕`$ to which $`A`$, $`Q`$, and $`M`$ all belong. §12 explores varying the variety $`𝐕`$ used in defining clone cohomology, giving *relative clone cohomology*. Finally, we pose a number of questions that seem important to ask about this theory, and about the relationship of clone cohomology and comonadic cohomology.
## 0. Preliminaries
### Category theory
We follow in terminology and notation.
### Homological algebra
We assume a familiarity with concepts and conventions of homological algebra, such as can be found in , , and .
### Universal algebra
We assume the basic definitions of universal algebra, such as can be found in , are known to the reader. Unlike some authors, we admit the possibility that an algebra can have an empty underlying set.
The kernel congruence of a homomorphism $`f`$ will be denoted by $`\mathrm{ker}f`$. The other sort of kernel, $`f^1(0)`$, will be denoted by $`\mathrm{Ke}f`$.
We denote the greatest and least congruences of $`A`$ by $`_A`$ and $`_A`$, and the identity homomorphism of $`A`$ by $`1_A`$. Basic operations or term operations of $`A`$ will occasionally be denoted by $`\omega ^A`$ or $`t^A`$, but we almost always drop the superscript.
If $`A`$ stands for an algebra, we use $`U(A)`$ to stand for the underlying set of the algebra.
### Clones
A *clone* is an $``$-tuple of sets $`V_n`$ (the $`n`$-ary elements of the clone $`V`$) such that for each $`n`$, and each $`i`$ with $`1in`$, there is an element $`\pi _{in}^VV_n`$, called the *$`i^{\text{th}}`$ of $`n`$ projection*, and such that, for each $`n`$, each $`n^{}`$-tuple $`𝐯`$ of elements of $`V_n`$, and each $`v^{}V_n^{}`$, there is an element $`v^{}𝐯V_n`$, called the *clone composite* of $`v^{}`$ and $`𝐯`$, satisfying
1. $`\pi _{in}^V𝐯=v_i`$,
2. $`v\pi _{1n}^V,\mathrm{},\pi _{nn}^V=v`$, and
3. $`u(\mathrm{𝐯𝐰})=(u𝐯)𝐰`$,
whenever the relevant compositions are defined. ($`\mathrm{𝐯𝐰}`$ stands for $`v_1𝐰,\mathrm{},v_n𝐰`$ if $`𝐯`$ is an $`n`$-tuple.)
As an example of a clone, given a set $`S`$, we have the *clone of $`S`$*, denoted by $`\mathrm{Clo}S`$. $`\mathrm{Clo}_nS`$ is the set of $`n`$-ary functions from $`S`$ to $`S`$, $`\pi _{in}^{\mathrm{Clo}S}`$ is the $`n`$-ary function on $`S`$ choosing the $`i^{\text{th}}`$ of its $`n`$ arguments, and given $`n`$, $`n^{}`$, an $`n^{}`$-tuple of $`n`$-ary functions $`𝐟`$, and an $`n^{}`$-ary function $`f^{}`$, the clone composite $`f^{}𝐟`$ is the function defined by
$$𝐬f^{}(f_1(𝐬),\mathrm{},f_n^{}(𝐬)).$$
Another example of a clone is the clone of a variety $`𝐕`$, denoted by $`\mathrm{Clo}𝐕`$. Elements of $`\mathrm{Clo}_n𝐕`$ are equivalence classes of $`n`$-ary term operations of the algebras in $`𝐕`$, where terms $`t`$ and $`t^{}`$ are equivalent if $`t(𝐱)=t^{}(𝐱)`$ is an identity of $`𝐕`$.
If $`V`$, $`V^{}`$ are clones, a *homomorphism of clones* from $`V`$ to $`V^{}`$ is an $``$-tuple $`f`$ of functions $`f_n:V_nV_n^{}`$, such that for all $`i`$ and $`n`$, $`f_n(\pi _{in}^V)=\pi _{in}^V^{}`$, and for all $`n`$ and $`n^{}`$, $`𝐯V_n^n^{}`$, and $`v^{}V_n^{}`$, we have $`f_n^{}(v^{})f_n(𝐯)=f_n(v^{}𝐯)`$.
An algebra $`A`$ in a variety $`𝐕`$ is the same as a clone homomorphism from $`\mathrm{Clo}V`$ to $`\mathrm{Clo}S`$, where $`S`$ is the underlying set of $`A`$.
A clone $`V`$ can be viewed as a category with one object for each natural number. The arrows from $`n`$ to $`n^{}`$ are $`n^{}`$-tuples of elements of $`V_n`$, and the identity of $`n`$ is $`\pi _{1n}^V,\mathrm{},\pi _{nn}^V`$. In the resulting category, $`n`$ is the $`n`$-fold direct power of $`1`$. Thus, the category is what is often called a *theory*.
### The modular commutator
In a congruence-modular variety $`𝐕`$ (i.e., such that for all $`A𝐕`$, $`\mathrm{Con}A`$ is a modular lattice) the congruence lattices admit a binary operation, called the *commutator*, that generalizes some well-known operations such as the commutator of two normal subgroups of a group. A comprehensive treatment of the theory of this operation, and related matters, can be found in .
The commutator of two congruences $`\alpha `$, $`\beta \mathrm{Con}A`$ is denoted by $`[\alpha ,\beta ]`$. A congruence $`\alpha `$ is said to be *abelian* if $`[\alpha ,\alpha ]=_A`$.
One definition of the commutator is as follows: If $`A𝐕`$, a congruence-modular variety of algebras, and $`\theta `$, $`\psi \mathrm{Con}A`$, then $`[\theta ,\psi ]`$ is the least congruence such that for all $`n`$, for all $`(n+1)`$-ary terms $`t`$, for all $`a`$, $`a^{}A`$ such that $`a𝜃a^{}`$, and for all $`𝐛`$, $`𝐜A^n`$ such that $`b_i𝜓c_i`$ for all $`i`$, we have
$$t(a,𝐛)[\theta ,\psi ]t(a,𝐜)$$
implies
$$t(a^{},𝐛)[\theta ,\psi ]t(a^{},𝐜).$$
### Difference terms
If $`𝐕`$ is a congruence-modular variety of algebras, a ternary term $`d`$ is called a *difference term* for $`𝐕`$ if
1. $`d(x,x,y)=y`$ is an identity of $`𝐕`$, and
2. for all $`A𝐕`$, $`\theta \mathrm{Con}A`$, and $`x`$, $`yA`$ such that $`x𝜃y`$, we have $`d(x,y,y)[\theta ,\theta ]x`$.
At least one such term exists for any congruence-modular variety.
### $`𝐕`$-objects in a category
If $`𝐂`$ is a category, and $`𝐕`$ is a variety of algebras, a $`𝐕`$-object in $`𝐂`$ is a pair $`c,F`$, consisting of an object $`c𝐂`$, and a contravariant functor $`F:𝐂𝐕`$, such that $`UF=𝐂(,c)`$, where $`U:𝐕\text{Set}`$ is the forgetful functor. If $`c,F`$ and $`c^{},F^{}`$ are $`𝐕`$-objects, a *homomorphism of $`𝐕`$-objects from $`c,F`$ to $`c^{},F^{}`$* is an arrow $`f:cc^{}`$ such that for each object $`d𝐂`$, the function $`𝐂(d,f):𝐂(d,c)𝐂(d,c^{})`$ is the underlying function of a (necessarily unique, since $`U`$ is faithful) arrow $`\overline{f}:F(d)F^{}(d)`$.
$`𝐕`$-objects in $`𝐂`$, and the homomorphisms between them, form a category in an obvious manner, which we denote by $`𝐕[𝐂]`$.
## 1. The Categories $`A`$-Set, $`\text{Ov}[A,𝐕]`$, and $`\text{Ab}[A,𝐕]`$
### The category of $`A`$-sets
Let $`A`$ be an algebra. We define an *$`A`$-set* to be a $`U(A)`$-tuple of sets. If $`S`$ is an $`A`$-set, we will denote the member of the tuple corresponding to an element $`aA`$ by $`{}_{a}{}^{}S`$. If $`S`$, $`S^{}`$ are $`A`$-sets, an *$`A`$-function* from $`S`$ to $`S^{}`$ is a $`U(A)`$-tuple $`f`$ such that the element corresponding to each $`aA`$, which we denote by $`{}_{a}{}^{}f`$, is a function from $`{}_{a}{}^{}S`$ to $`{}_{a}{}^{}S_{}^{}`$. We write $`f:SS^{}`$. $`A`$-sets and $`A`$-functions form a category, $`A`$-Set, in an obvious manner.
### Overalgebras
If $`A`$ is an algebra, an *$`A`$-overalgebra* is an $`A`$-set $`Q`$, such that for each $`n`$, and each $`n`$-ary basic operation $`\omega `$ of the type of $`A`$, there is a $`U(A)^n`$-tuple of functions $`\omega _𝐚^Q:{}_{a_1}{}^{}Q\times \mathrm{}\times {}_{a_n}{}^{}Q{}_{\omega (𝐚)}{}^{}Q`$.
In what follows, we will write $`{}_{𝐚}{}^{}Q`$ for the product $`{}_{a_1}{}^{}Q\times \mathrm{}\times {}_{a_n}{}^{}Q`$. Thus, $`\omega _𝐚^Q:{}_{𝐚}{}^{}Q{}_{\omega (𝐚)}{}^{}Q`$.
If $`Q`$, $`Q^{}`$ are $`A`$-overalgebras, a *homomorphism of $`A`$-overalgebras from $`Q`$ to $`Q^{}`$* is an $`A`$-function $`f:QQ^{}`$ such that for each $`n`$, each $`n`$-ary basic operation $`\omega `$, each $`𝐚A^n`$, and each $`𝐪{}_{𝐚}{}^{}Q`$, we have
$${}_{\omega (𝐚)}{}^{}f(\omega _𝐚^Q(𝐪))=\omega _𝐚^Q^{}({}_{𝐚}{}^{}f(𝐪)),$$
where $`{}_{𝐚}{}^{}f(𝐪)`$ stands for $`{}_{a_1}{}^{}f(q_1),\mathrm{},{}_{a_n}{}^{}f(q_n)`$.
$`A`$-overalgebras and their homomorphisms form a category in an obvious manner, which we denote by $`\text{Ov}[A]`$.
For an example of an $`A`$-overalgebra, let $`B,\pi `$ be an object of the “comma category” $`(\mathrm{\Omega }`$-$`\text{Alg}A)`$ of algebras over $`A`$. That is, let $`B`$ be an algebra of the same type $`\mathrm{\Omega }`$ as $`A`$, and let $`\pi :BA`$ be a homomorphism. Then we define the $`A`$-overalgebra $`[[B,\pi ]]`$ by $`{}_{a}{}^{}[[B,\pi ]]=\pi ^1(a)`$ and $`\omega _𝐚^{[[B,\pi ]]}(𝐛)=\omega ^B(𝐛)`$.
If $`Q`$ is an $`A`$-overalgebra, we define the *total algebra* of $`Q`$, denoted by $`AQ`$, to be the set of pairs $`\{a,q:aA,q{}_{a}{}^{}Q\}`$, provided with operations defined by
$$\omega (a_1,q_1,\mathrm{},a_n,q_n)=\omega (𝐚),\omega _𝐚^Q(𝐪);$$
we say that an $`A`$-overalgebra $`Q`$ is *totally in $`𝐕`$*, where $`V`$ is a variety of algebras of the type of $`A`$, if $`AQ𝐕`$.
If $`Q`$ is an $`A`$-overalgebra, then accompanying the total algebra $`AQ`$ there is a homomorphism $`\pi _Q:AQA`$, defined by $`a,qa`$. It is clear that $`\text{Ov}[A,𝐕]`$ is equivalent as a category to the category $`(𝐕A)`$. One leg of an equivalence takes an $`A`$-overalgebra $`Q`$ to $`AQ,\pi _Q`$. The other leg takes an algebra over $`A`$, $`B,\pi `$, to $`[[B,\pi ]]`$.
There is an evident forgetful functor from $`\text{Ov}[A,𝐕]`$ to $`A`$-Set. To construct an $`A`$-overalgebra free on an $`A`$-set $`S`$ (i.e., the value of the corresponding left adjoint functor) form a free algebra $`F`$ on a disjoint union of the $`{}_{a}{}^{}S`$, (a free algebra in $`𝐕`$, that is, also known as a *relatively free* algebra), and map the generators to $`A`$ in the obvious way, giving a homomorphism $`\pi :FA`$. The $`A`$-overalgebra $`[[F,\pi ]]`$ is then free on $`S`$.
### One-one and onto homomorphisms of $`A`$-overalgebras
If $`f`$ is a homomorphism of $`A`$-overalgebras, then we say that $`f`$ is *one-one* if each $`{}_{a}{}^{}f`$ is one-one, and we say that $`f`$ is *onto* if each $`{}_{a}{}^{}f`$ is onto. If $`f`$ is both one-one and onto, then it is an isomorphism in the category of $`A`$-overalgebras.
### $`A`$-operations
Let $`A`$ be a set, and $`\omega `$ an $`n`$-ary operation on $`A`$. If $`S`$ is an $`A`$-set, then an *$`A`$-operation on $`S`$, over $`\omega `$*, is a $`U(A)^n`$-tuple $`\omega ^{}`$ of functions $`\omega _𝐚^{}:{}_{𝐚}{}^{}S{}_{\omega (𝐚)}{}^{}S`$. We can now rephrase our definition of an $`A`$-overalgebra: it is an $`A`$-set $`Q`$, together with an $`A`$-operation $`\omega ^Q`$ over $`\omega ^A`$ for every basic operation $`\omega `$ on $`A`$.
If $`Q`$ is an $`A`$-overalgebra totally in $`𝐕`$, and $`v\mathrm{Clo}_n𝐕`$, then we can define an $`A`$-operation $`v^Q`$ over $`v^A`$ by letting $`v_𝐚^Q(𝐪)`$ be the second component of the pair
$$v^{AQ}(a_1,q_1,\mathrm{},a_n,q_n);$$
we will thus use $`v^Q`$ to denote that $`A`$-operation, regardless of whether $`v`$ is an $`n`$-ary basic operation, $`n`$-ary term operation, or $`n`$-ary element of $`\mathrm{Clo}𝐕`$.
If $`A`$ is a set, and $`S`$ an $`A`$-set, then we can define $`\mathrm{Clo}^AS`$, the *clone of $`A`$-operations on $`S`$* to be, for each $`n`$, the set of pairs $`\omega ,\omega ^{}`$ where $`\omega `$ is an $`n`$-ary operation on $`A`$, and $`\omega ^{}`$ is an $`n`$-ary $`A`$-operation on $`S`$ over $`\omega `$. Projections and clone composition are easy to define, and there is a clone homomorphism $`\pi :\mathrm{Clo}^AS\mathrm{Clo}A`$ given by taking the first component of each pair.
If $`A`$ is an algebra in a variety $`𝐕`$, $`\varphi ^A`$ is the corresponding clone homomorphism from $`\mathrm{Clo}𝐕`$ to $`\mathrm{Clo}A`$, and $`Q`$ is an $`A`$-set, then an $`A`$-overalgebra structure on $`Q`$, totally in $`𝐕`$, just amounts to a clone homomorphism $`\varphi ^Q:\mathrm{Clo}𝐕\mathrm{Clo}^AQ`$, such that $`\pi \varphi ^Q=\varphi ^A`$. Accordingly, to define an $`A`$-overalgebra, it suffices to define $`v^A`$ for all $`v`$, with the condition that the given definitions provide a well-defined clone homomorphism.
### Pointed overalgebras
If $`A`$ is an algebra, a *pointed $`A`$-overalgebra* is an $`A`$-overalgebra $`P`$, such that each $`{}_{a}{}^{}P`$ has a basepoint $`{}_{a}{}^{}_{}^{P}`$, or simply $`{}_{a}{}^{}`$, such that for each basic operation $`\omega `$, $`n`$-ary, and each $`𝐚A^n`$, $`\omega _𝐚^P({}_{a_1}{}^{},\mathrm{},{}_{a_n}{}^{})={}_{\omega (𝐚)}{}^{}`$.
If $`P`$, $`P^{}`$ are pointed $`A`$-overalgebras, a *homomorphism from $`P`$ to $`P^{}`$* is a homomorphism of $`A`$-overalgebras $`f:PP^{}`$ such that for each $`a`$, $`{}_{a}{}^{}f({}_{a}{}^{}_{}^{P})={}_{a}{}^{}_{}^{P^{}}`$.
Pointed overalgebras totally in $`𝐕`$, and the homomorphisms between them, form a category $`\text{Pnt}[A,𝐕]`$ in an obvious manner. As the notation suggests, it is precisely the category $`\text{Pnt}[\text{Ov}[A,𝐕]]`$ of pointed set objects in the category $`\text{Ov}[A,𝐕]`$.
As an example of a pointed $`A`$-overalgebra, let $`\alpha \mathrm{Con}A`$. We define a pointed $`A`$-overalgebra $`\alpha ^{}`$ by $`{}_{a}{}^{}\alpha _{}^{}=\{a^{}:a𝛼a^{}\}`$, $`{}_{a}{}^{}=a`$, and $`\omega _𝐚^\alpha ^{}(𝐜)=\omega (𝐜)`$ for $`𝐜{}_{𝐚}{}^{}\alpha _{}^{}`$.
If $`P`$ is a pointed $`A`$-overalgebra, then accompanying $`AP`$ and $`\pi _P`$ there is a homomorphism $`\iota _P:AAP`$, defined by $`\iota _P:aa,{}_{a}{}^{}`$. The triple $`AP,\pi _P,\iota _P`$ can be viewed as a commutative diagram
$$\begin{array}{ccc}A& \stackrel{\iota _P}{}& AP\\ & & \pi _P;& & \\ A& =& A\end{array}$$
in slightly different terms, it is a pointed set object in the category of algebras over $`A`$.
Given such a diagram, or, a triple $`B,\pi ,\iota `$ with $`\pi :BA`$, $`\iota :AB`$, and $`\pi \iota =1_A`$, we denote by $`[[B,\pi ,\iota ]]`$ the pointed $`A`$-overalgebra with underlying $`A`$-overalgebra $`[[B,\pi ]]`$ and basepoints $`\iota (a){}_{a}{}^{}[[B,\pi ]]`$.
The constructions $`PAP,\pi _P,\iota _P`$ and $`B,\pi ,\iota [[B,\pi ,\iota ]]`$ are clearly two legs of an equivalence between the categories $`\text{Pnt}[A,𝐕]`$ and $`\text{Pnt}(𝐕A)`$.
There is an obvious forgetful functor from $`\text{Pnt}[A,𝐕]`$ to $`\text{Ov}[A,𝐕]`$. A corresponding free functor can be defined as follows: given an $`A`$-overalgebra $`Q`$, form the algebra $`B=A(AQ)`$. Define $`\pi :BA`$ by applying the universal property of the coproduct to the homomorphisms $`1_A`$ and $`\pi _Q`$. Define $`\iota :AB`$ as the insertion of $`A`$ into the coproduct. Then $`[[B,\pi ,\iota ]]`$ is a pointed $`A`$-overalgebra free on $`Q`$.
### Abelian group overalgebras
An *abelian group $`A`$-overalgebra* is an $`A`$-overalgebra $`M`$ such that on each $`{}_{a}{}^{}M`$ there is the structure of an abelian group, in such a way that the functions $`\omega _𝐚^M:{}_{𝐚}{}^{}M{}_{\omega (𝐚)}{}^{}M`$ are abelian group homomorphisms. If $`M`$, $`M^{}`$ are abelian group $`A`$-overalgebras, a *homomorphism of abelian group $`A`$-overalgebras from $`M`$ to $`M^{}`$* is an $`A`$-overalgebra homomorphism $`f:MM^{}`$ such that each $`{}_{a}{}^{}f`$ is an abelian group homomorphism.
Abelian group $`A`$-overalgebras totally in $`𝐕`$, and the homomorphisms between them, form a category in an obvious manner, which we denote by $`\text{Ab}[A,𝐕]`$. It is the category of abelian group objects in $`\text{Ov}[A,𝐕]`$.
Categorical algebraists have given the term *Beck module over $`A`$* to an abelian group object in the category $`(𝐕A)`$. Clearly, $`\text{Ab}[A,𝐕]`$ is equivalent to the category of Beck modules over $`A`$.
There is an obvious forgetful functor from $`\text{Ab}[A,𝐕]`$ to $`\text{Pnt}[A,𝐕]`$.
###### Theorem 1.1.
() Let $`𝐕`$ be a congruence-modular variety of algebras, and $`A𝐕`$. Let $`P`$ be a pointed $`A`$-overalgebra which is the underlying pointed $`A`$-overalgebra of an abelian group overalgebra. Then there is a unique assignment of abelian group structures to the pointed sets $`{}_{a}{}^{}P`$, such that $`{}_{a}{}^{}`$ is the zero element of each $`{}_{a}{}^{}P`$ and the functions $`\omega _𝐚^P`$ are abelian group homomorphisms. The group operations in $`{}_{a}{}^{}P`$ satisfy $`pp^{}+p^{\prime \prime }=d_{a,a,a}^P(p,p^{},p^{\prime \prime })`$.
### Abelian group $`A`$-overalgebras free on a pointed $`A`$-overalgebra
If $`𝐕`$ is congruence-modular, a free functor (left adjoint) for the forgetful functor from $`\text{Ab}[A,𝐕]`$ to $`\text{Pnt}[A,𝐕]`$ is as follows: Given a pointed $`A`$-overalgebra $`P`$, draw the diagram
$$\begin{array}{ccccc}A& \stackrel{\iota _P}{}& AP& \stackrel{\mathrm{nat}[\kappa ,\kappa ]}{}& (AP)/[\kappa ,\kappa ]\\ & & \pi _P& & \pi ,& & \\ A& =& A& =& A\end{array}$$
where $`\kappa =\mathrm{ker}\pi _P`$ and $`\pi `$ is the unique homomorphism making the diagram commute.
$$M=[[(AP)/[\kappa ,\kappa ],\pi ,\mathrm{nat}[\kappa ,\kappa ]\iota _P]]$$
is then a pointed overalgebra, which is an abelian group overalgebra in a unique way, by theorem 1.1. $`M`$ is an abelian group $`A`$-overalgebra free on $`P`$.
### $`𝐕^{}`$ $`A`$-overalgebras
Suppose $`𝐕^{}`$ is another variety of algebras than $`𝐕`$, perhaps of a different type. We define a $`𝐕^{}`$ $`A`$-overalgebra to be an $`A`$-overalgebra $`M`$ such that the $`{}_{a}{}^{}M`$ are algebras in $`𝐕^{}`$ and the $`\omega _𝐚^M`$ are homomorphisms of $`𝐕^{}`$, and a homomorphism of $`𝐕^{}`$ $`A`$-overalgebras to be an $`A`$-overalgebra homomorphism $`f`$ such that the $`{}_{a}{}^{}f`$ are homomorphisms of algebras in $`𝐕^{}`$. We denote the category of $`𝐕^{}`$ $`A`$-overalgebras totally in $`𝐕`$, and homomorphisms betwen them, by $`𝐕^{}[A,𝐕]`$.
The category $`𝐕^{}[A,𝐕]`$ generalizes $`\text{Pnt}[A,𝐕]`$ and $`\text{Ab}[A,𝐕]`$ in an obvious way, and we have $`𝐕^{}[A,𝐕]=𝐕^{}[\text{Ov}[A,𝐕]]`$, the category of $`𝐕^{}`$-objects in the category $`\text{Ov}[A,𝐕]`$.
Although $`𝐕^{}[A,𝐕]`$ is the category of $`𝐕^{}`$-objects in the category $`\text{Ov}[A,𝐕]`$, it is also true that given $`M𝐕^{}[A,𝐕]`$, and an $`A`$-set $`S`$, the $`A`$-functions from $`S`$ to $`M`$ form an algebra of $`𝐕^{}`$ in a natural way.
If $`M𝐕^{}[A,𝐕]`$, and $`u`$ is an $`n`$-ary basic operation, term operation, or element of $`\mathrm{Clo}𝐕^{}`$, we will denote by $`{}_{a}{}^{}u`$, or occasionally by $`{}_{a}{}^{}u_{}^{M}`$, that operation on the algebra $`{}_{a}{}^{}M`$.
### Restriction and induction functors
For each category $`𝐕^{}[A,𝐕]`$, algebra $`X𝐕`$, and homomorphism $`f:XA`$, there is a functor $`{}_{f}{}^{}\mathrm{Res}:𝐕^{}[A,𝐕]𝐕^{}[X,𝐕]`$, defined by $`{}_{x}{}^{}({}_{f}{}^{}\mathrm{Res}M)={}_{f(x)}{}^{}M`$ and $`\omega _𝐱^{{}_{f}{}^{}\mathrm{Res}M}=\omega _{f(𝐱)}^M`$, where $`f(𝐱)`$ stands for $`f(x_1),\mathrm{},f(x_n)`$.
The restriction functors all have left adjoints, which are constructed in . We call these *induction* functors.
We will have occasion to use the functor of induction of abelian group overalgebras in §10.
### Products
Let $`A`$ be an algebra. If $`S_1`$, $`\mathrm{}`$, $`S_n`$ are $`A`$-sets, then a product $`\mathrm{\Pi }_iS_i`$ in the category of $`A`$-sets is given by $`{}_{a}{}^{}(\mathrm{\Pi }_iS_i)={}_{a}{}^{}S_{1}^{}\times \mathrm{}\times {}_{a}{}^{}S_{n}^{}`$. Similarly, if $`M_1`$, $`\mathrm{}`$, $`M_n𝐕^{}[A,𝐕]`$, then a product $`\mathrm{\Pi }_iM_i𝐕^{}[A,𝐕]`$ is given by $`{}_{a}{}^{}(\mathrm{\Pi }_iM_i)={}_{a}{}^{}M_{1}^{}\times \mathrm{}\times {}_{a}{}^{}M_{n}^{}`$, by
$$v_𝐚^{\mathrm{\Pi }_iM_i}(𝐦_1,\mathrm{},𝐦_k)=v_𝐚^{M_1}(m_{11},\mathrm{},m_{k1}),\mathrm{},v_𝐚^{M_n}(m_{1n},\mathrm{},m_{kn}),$$
for $`v\mathrm{Clo}_k𝐕`$, and by
$${}_{a}{}^{}u(𝐦_1,\mathrm{},𝐦_k)={}_{a}{}^{}u_{}^{M_1}(m_{11},\mathrm{},m_{k1}),\mathrm{},{}_{a}{}^{}u_{}^{M_n}(m_{11},\mathrm{},m_{k1}),$$
for $`u\mathrm{Clo}_k𝐕^{}`$.
Suppose, on the other hand, that we are given algebras $`A_i𝐕`$, and objects $`M_i𝐕^{}[A_i,𝐕]`$, for $`i=1`$, $`\mathrm{}`$, $`n`$. Let $`A=\mathrm{\Pi }_iA_i`$. Define the *outer product* $`_iM_i𝐕^{}[A,𝐕]`$ by $`{}_{𝐚}{}^{}(_iM_i)=\mathrm{\Pi }_i({}_{a_i}{}^{}M_{i}^{})`$, and by
$`v_{𝐚_1,\mathrm{},𝐚_k}^{_iM_i}(𝐦_1,`$ $`\mathrm{},𝐦_k)`$
$`=v_{a_{11},\mathrm{},a_{k1}}^{M_1}(m_{11},\mathrm{},m_{k1}),\mathrm{},v_{a_{1n},\mathrm{},a_{kn}}^{M_n}(m_{1n},\mathrm{},m_{kn}),`$
for $`v\mathrm{Clo}_k𝐕`$, and
$${}_{𝐚}{}^{}u(𝐦_1,\mathrm{},𝐦_i)={}_{a_1}{}^{}u_{}^{M_1}(m_{11},\mathrm{},m_{k1}),\mathrm{},{}_{a_n}{}^{}u_{}^{M_n}(m_{1n},\mathrm{},m_{kn}),$$
for $`u\mathrm{Clo}_k𝐕^{}`$. That is, $`M=\mathrm{\Pi }_i({}_{\pi _{A,i}}{}^{}\mathrm{Res}M_i)`$, where the $`\pi _{A,i}:AA_i`$ are the projections to the factors.
###### Theorem 1.2.
If all the $`A_i`$ are the same algebra $`A`$, we have $`{}_{\mathrm{\Delta }_A}{}^{}\mathrm{Res}(_iM_i)=M^n`$, where the homomorphism $`\mathrm{\Delta }_A:A\mathrm{\Pi }_iA`$ is defined by $`aa,\mathrm{},a`$.
### Advantages of the overalgebra formalism
There are two main advantages for using the formalism of $`A`$-overalgebras rather than that of algebras over $`A`$. One is that given an object $`M𝐕^{}[A,𝐕]`$, the $`{}_{a}{}^{}M`$, which are algebras of $`𝐕^{}`$, are important objects of study, and a formalism that provides for this is useful. If we use the formalism of algebras over $`A`$, then we will end up considering the $`{}_{a}{}^{}M`$ anyway, in a different form, as $`\pi ^1(a)`$ for $`aA`$.
The other main advantage is that the formalism of $`A`$-overalgebras allows the $`{}_{a}{}^{}M`$ to be nondisjoint. One way this helps is to facilitate defining and using the pointed $`A`$-overalgebra $`\alpha ^{}`$ for $`\alpha \mathrm{Con}A`$. If we don’t allow nondisjoint sets, then we must work with the algebra often denoted by $`A(\alpha )=A\alpha ^{}`$, along with the projection $`\pi _\alpha ^{}`$ to $`A`$.
The other way allowing the $`{}_{a}{}^{}M`$ to be nondisjoint is helpful, and this is probably the most important advantage, is in defining and using the restriction functors. Using the formalism of algebras over $`A`$ requires that the restrictions be defined as pullbacks, and then frequent use must be made of the universal property of the pullback. This can be done, of course, but it is tedious, and tends to obscure the situation.
## 2. Enveloping Ringoids
### Ringoids
A *ringoid* is a small additive category. If $`𝐗`$ is a ringoid, with set of objects $`A`$, we write $`{}_{a^{}}{}^{}𝐗_{a}^{}`$ for $`𝐗(a,a^{})`$. A left *$`𝐗`$-module* is an additive functor from $`𝐗`$ to Ab, the category of abelian groups. If $`M`$ is a left $`𝐗`$-module, we write $`{}_{a}{}^{}M`$ instead of $`M(a)`$, and if $`m{}_{a}{}^{}M`$ and $`r{}_{a^{}}{}^{}𝐗_{a}^{}`$, we write $`rm`$ rather than $`M(r)(m)`$.
If $`A`$ is an algebra in a variety $`𝐕`$, the *enveloping ringoid for $`A`$, with respect to $`𝐕`$*, is a certain ringoid denoted by $`[A,𝐕]`$, which has the underlying set of $`A`$ as its set of objects, and such that the category of left $`[A,𝐕]`$-modules is isomorphic to the category $`\text{Ab}[A,𝐕]`$. Enveloping ringoids are treated in detail in and .
### Construction of the enveloping ringoid
Given an algebra $`A`$ in a variety $`𝐕`$, we construct an object $`M_a\text{Ab}[A,𝐕]`$ free on an $`A`$-set consisting of a singleton at $`a`$ and the null set elsewhere, for each $`aA`$. (An abelian group $`A`$-overalgebra free on an $`A`$-set $`S`$ can be constructed by constructing an $`A`$-overalgebra free on $`S`$, a pointed $`A`$-overalgebra free on that, and an abelian group $`A`$-overalgebra free on that.)
The enveloping ringoid can then be given as $`{}_{a^{}}{}^{}[A,𝐕]_a={}_{a^{}}{}^{}(M_a)`$ for all $`a`$, $`a^{}A`$.
### Abelian group $`A`$-overalgebras free on an $`A`$-set
Once the enveloping ringoid is defined, there is another method available for defining abelian group $`A`$-overlaying algebras free on $`A`$-sets. Let $`A`$ be an algebra in the variety $`𝐕`$, and let $`S`$ be an $`A`$-set. For each $`aA`$, let $`{}_{a}{}^{}M`$ be the set of finite formal linear combinations
$$\underset{i}{}r_is_i,$$
where $`r_i{}_{a}{}^{}[A,𝐕]_{b_i}`$ and $`s_i{}_{b_i}{}^{}S`$. This defines an object $`M\text{Ab}[A,𝐕]`$ free on $`S`$.
## 3. Abelian Extensions
Let $`A`$ be an algebra of $`𝐕`$, and let $`M\text{Ab}[A,𝐕]`$. We define an extension in $`𝐕`$ of $`A`$ by $`M`$ to be a triple $`\chi ,E,\pi `$ such that $`E`$ is an algebra of $`𝐕`$, $`\pi :EA`$ is an onto homomorphism, and $`\chi :{}_{\pi }{}^{}\mathrm{Res}P\kappa ^{}`$ is an isomorphism of pointed overalgebras, where $`P`$ is the underlying pointed $`A`$-overalgebra of $`M`$ and $`\kappa =\mathrm{ker}\pi `$.
For example, given $`A𝐕`$, and $`M`$, totally in $`𝐕`$, we can form the total algebra $`AM`$. Let $`\kappa =\mathrm{ker}\pi _M`$. Let $`\chi :{}_{\pi }{}^{}\mathrm{Res}M\kappa ^{}`$ be the $`(AM)`$-function given by $`{}_{a,m}{}^{}\chi (m^{})=a,m+m^{}`$.
###### Theorem 3.1.
$`\chi ,AM,\pi _P`$ is an extension in $`𝐕`$ of $`A`$ by $`M`$.
###### Proof.
Each $`{}_{a,m}{}^{}\chi `$ is clearly one-one and onto, and $`\chi `$ is an homomorphism of pointed $`(AM)`$-overalgebras because $`{}_{a,m}{}^{}\chi (0)=a,m={}_{a,m}{}^{}`$, and for each $`v\mathrm{Clo}_n𝐕`$, each $`𝐚,𝐦=a_1,m_1,\mathrm{},a_n,m_n(AM)^n`$, and each $`𝐦^{}=m_1^{},\mathrm{},m_n^{}{}_{𝐚,𝐦}{}^{}({}_{\pi }{}^{}\mathrm{Res}M)`$, we have
$`{}_{v𝐚,𝐦}{}^{}\chi (v_{𝐚,𝐦}^{{}_{\pi }{}^{}\mathrm{Res}M}(𝐦^{}))`$ $`={}_{v(𝐚,𝐦)}{}^{}\chi (v_𝐚^M(𝐦^{}))`$
$`=v(𝐚),v_𝐚^M(𝐦)+v_𝐚^M(𝐦^{})`$
$`=v^{AM}(a_1,m_1+m_1^{},\mathrm{},a_n,m_n+m_n^{})`$
$`=v_{𝐚,𝐦}^\kappa ^{}(a_1,m_1+m_1^{},\mathrm{},a_n,m_n+m_n^{})`$
$`=v_{𝐚,𝐦}^\kappa ^{}({}_{a_1,m_1}{}^{}\chi (m_1^{}),\mathrm{},{}_{a_n,m_n}{}^{}\chi (m_n^{})),`$
In this example, the onto homomorphism $`\pi =\pi _M`$ is split by the homomorphism $`\iota _M`$. We say that an abelian extension $`\chi ,E,\pi `$ is *split* if $`\pi `$ splits.
In general, given an extension $`\chi ,E,\pi `$ in $`𝐕`$ of $`A`$ by $`M`$, the pointed $`E`$-overalgebra $`\kappa ^{}=(\mathrm{ker}\pi )^{}`$ is isomorphic to the abelian group overalgebra $`{}_{\pi }{}^{}\mathrm{Res}M`$, and so is itself an abelian group overalgebra with abelian group operations as given in theorem 1.1.
### A lemma using the properties of the difference term
The following lemma will be very useful in proving properties of abelian extensions:
###### Lemma 3.2.
If $`\chi ,E,\pi `$ is an abelian extension in $`𝐕`$ of $`A`$ by $`M`$, and $`e`$, $`e^{}`$, $`e^{\prime \prime }E`$ are such that $`\pi (e)=\pi (e^{})=\pi (e^{\prime \prime })`$, then
$${}_{e}{}^{}\chi _{}^{1}(e^{})+{}_{e^{}}{}^{}\chi _{}^{1}(e^{\prime \prime })={}_{e}{}^{}\chi _{}^{1}(e^{\prime \prime }).$$
###### Proof.
Let $`\kappa =\mathrm{ker}\pi `$, and let $`a=\pi (e)`$. By the properties of $`d`$ and theorem 1.1, we have
$`{}_{e^{}}{}^{}\chi _{}^{1}(e^{\prime \prime })`$ $`=d_{a,a,a}^M({}_{e^{}}{}^{}\chi _{}^{1}(e^{\prime \prime }),{}_{a}{}^{}0,{}_{a}{}^{}0)`$
$`=d_{e^{},e^{},e}^{{}_{\pi }{}^{}\mathrm{Res}M}({}_{e^{}}{}^{}\chi _{}^{1}(e^{\prime \prime }),{}_{e^{}}{}^{}\chi _{}^{1}({}_{e^{}}{}^{}0),{}_{e}{}^{}\chi _{}^{1}({}_{e}{}^{}0))`$
$`={}_{e}{}^{}\chi _{}^{1}(d_{e^{},e^{},e}^\kappa ^{}(e^{\prime \prime },{}_{e^{}}{}^{}0,{}_{e}{}^{}0))`$
$`={}_{e}{}^{}\chi _{}^{1}(d^E(e^{\prime \prime },e^{},e))`$
$`={}_{e}{}^{}\chi _{}^{1}(d_{e,e,e}^\kappa ^{}(e^{\prime \prime },e^{},{}_{e}{}^{}0))`$
$`={}_{e}{}^{}\chi _{}^{1}(e^{\prime \prime }){}_{e}{}^{}\chi _{}^{1}(e^{}).`$
### Sections of extensions
Let $`=\chi ,E,\pi `$ be an abelian extension in $`𝐕`$ of $`A`$ by $`M`$. Recall that we say that $`=\chi ,E,\pi `$ splits if there is a homomorphism $`\sigma :AE`$ right inverse to $`\pi `$. Not every abelian extension of $`A`$ by $`M`$ splits, as we know, because not every onto homomorphism has a right inverse. (We do know that there is always at least one split extension, given previously.) At any rate, by the axiom of choice, there is always a function $`\sigma `$ such that $`\pi \sigma =1_A`$, whether or not the extension $``$ splits.
A function $`\sigma :AE`$, not necessarily a homomorphism, such that $`\pi \sigma =1_A`$, will be called a *section* of $`\pi `$ (or, of the extension $``$).
### $`E`$ and $`M`$
For each $`eE`$, $`{}_{e}{}^{}\chi _{}^{1}:\pi ^1(\pi (e)){}_{\pi (a)}{}^{}M`$ is a one-one and onto function. A section $`\sigma `$ chooses one representative $`\sigma (a)`$ for each $`\kappa `$-equivalence class $`\pi ^1(a)`$, and allows us to select functions $`{}_{\sigma (a)}{}^{}\chi _{}^{1}`$ which provide an isomorphism of pointed $`A`$-sets between $`[[E,\pi ,\sigma ]]`$ and $`M`$. We will use this pointed $`A`$-function extensively in what follows.
### Some $`A`$-sets connected with abelian extensions
Let $`𝐕`$ be a variety of algebras of type $`\mathrm{\Omega }`$, and let $`A`$ be an algebra in $`𝐕`$. We will define some $`A`$-sets that will be useful in the study of abelian extensions of $`A`$.
Let $`X_𝐕^0(A)`$, or simply $`X^0`$, denote the $`A`$-set $`[[A,1_A]]`$, that is, the underlying set of $`A`$, viewed as a set over $`A`$. We will refer to the element $`a{}_{a}{}^{}X_{}^{0}`$, for any $`aA`$, by $`[a]`$. Let $`X_𝐕^1(A)`$, or simply $`X^1`$, denote the $`A`$-set $`[[S^1,h]]`$, where $`S^1`$ is the set of of pairs (written as follows) $`[v;𝐚]`$, where $`v\mathrm{Clo}_n𝐕`$ for some $`n`$, and $`𝐚A^n`$, and $`h[v;𝐚]=v(𝐚)`$. Let $`X_𝐕^2(A)`$, or simply $`X^2`$, denote the $`A`$-set $`[[S^2,k]]`$, where $`S^2`$ is the set of triples $`[v^{},𝐯;𝐚]`$, where $`𝐚`$ is an element of $`A^n`$ for some $`n`$, $`𝐯`$ is an $`n^{}`$-tuple of elements of $`\mathrm{Clo}_n𝐕`$, for some $`n^{}`$, and $`v^{}\mathrm{Clo}_n^{}𝐕`$, and where $`k[v^{},𝐯;𝐚]=v^{}(𝐯(𝐚))`$.
Note that if $`n=0`$, $`1`$, or $`2`$, then $`X^n`$ satisfies the property that if $`aa^{}`$, the sets $`{}_{a}{}^{}X_{}^{n}`$ and $`{}_{a^{}}{}^{}X_{}^{n}`$ are disjoint. We will often take advantage of this property in formulas which involve $`A`$-functions with domain $`X^n`$, by dropping the subscript $`a`$. Thus, if $`x{}_{a}{}^{}X_{}^{n}`$ and $`\varphi `$ is an $`A`$-function from $`X^n`$ to another $`A`$-set, we will write $`{}_{a}{}^{}\varphi (x)`$ simply as $`\varphi (x)`$. This is unambiguous because $`x`$ determines $`a`$.
### $`\delta _{\sigma ,\sigma ^{}}^{}`$
As a first example of an $`A`$-function with domain one of the $`A`$-sets just defined, we define $`\delta _{\sigma ,\sigma ^{}}^{}`$ to be the $`A`$-function defined by
$$\delta _{\sigma ,\sigma ^{}}^{}([a])={}_{\sigma (a)}{}^{}\chi _{}^{1}(\sigma ^{}(a)),$$
where $``$ is an extension and $`\sigma `$, $`\sigma ^{}`$ are two sections for $``$.
###### Lemma 3.3.
If $``$ is an abelian extension and $`\sigma `$ is a section for $``$, then $`\delta _{\sigma ,\sigma }^{}0`$.
###### Proof.
$`\sigma (a)={}_{\sigma (a)}{}^{}0_{}^{\kappa ^{}}`$. Thus, $`\delta _{\sigma ,\sigma }^{}={}_{\sigma (a)}{}^{}\chi _{}^{1}(\sigma (a))={}_{\sigma (a)}{}^{}\chi _{}^{1}(0)=0`$. ∎
## 4. Factor Sets
### The factor set of an extension relative to a section
Let $`=\chi ,E,\pi `$ be an abelian extension of $`A`$ by $`M`$, and $`\sigma `$ a section. The functions $`{}_{\sigma (a)}{}^{}\chi _{}^{1}`$, for $`aA`$, allow us to express the structure of $`E`$ in terms of $`M`$. To begin we “measure the failure of $`\sigma `$ to be a homomorphism” by defining, for each $`[v;𝐚]X_𝐕^1(A)`$, the element
$`\mathrm{f}^{,\sigma }[v;𝐚]`$ $`={}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma (𝐚))){}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(\sigma (v(𝐚)))`$
$`={}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma (𝐚))).`$
It is easy to see that $`\mathrm{f}^{,\sigma }`$ is an $`A`$-function from $`X_𝐕^1(A)`$ to $`M`$. We call $`\mathrm{f}^{,\sigma }`$ the *factor set for $``$, with respect to the section $`\sigma `$*.
###### Theorem 4.1.
If $`\sigma `$ is a splitting for $``$, then $`\mathrm{f}^{,\sigma }0`$.
###### Proof.
In that case, $`\mathrm{f}^{,\sigma }[v;𝐚]={}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma (𝐚)))={}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(\sigma (v(𝐚)))={}_{v(𝐚)}{}^{}0`$ for all $`v`$ and $`𝐚`$. ∎
###### Theorem 4.2.
Together with $`\sigma `$, the factor set $`\mathrm{f}^{,\sigma }`$ determines the algebra structure of $`E`$.
###### Proof.
Let $`v\mathrm{Clo}_n𝐕`$, $`𝐞E^n`$, and $`𝐚=\pi (𝐞)`$. We have
(4.1) $`{}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(𝐞))`$ $`={}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v_{\sigma (𝐚)}^\kappa ^{}(𝐞))`$
$`={}_{v(\sigma (𝐚))}{}^{}\chi _{}^{1}(v_{\sigma (𝐚)}^\kappa ^{}(𝐞))+{}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma (𝐚)))`$
$`=v_{\sigma (𝐚)}^{{}_{\pi }{}^{}\mathrm{Res}M}({}_{\sigma (𝐚)}{}^{}\chi _{}^{1}(𝐞))+\mathrm{f}^{,\sigma }[v;𝐚]`$
$`=v_𝐚^M({}_{\sigma (𝐚)}{}^{}\chi _{}^{1}(𝐞))+\mathrm{f}^{,\sigma }[v;𝐚].`$
Thus, if the section $`\sigma `$ is given, $`v^E`$ and $`\mathrm{f}^{,\sigma }`$ determine each other, because the mappings $`{}_{\sigma (a)}{}^{}\chi _{}^{1}`$ together make up an isomorphism of $`A`$-sets from $`[[E,\pi ]]`$ to $`M`$. ∎
### Abstract factor sets
Because $`E`$ is an algebra in $`𝐕`$, sending each $`v\mathrm{Clo}_n𝐕`$ to the operation $`v^E`$, for all $`n`$, is a clone homomorphism from $`\mathrm{Clo}𝐕`$ to the clone $`\mathrm{Clo}U(E)`$ of all finitary operations on the set $`U(E)`$. This means that for each $`𝐯(\mathrm{Clo}_n𝐕)^n^{}`$, and each $`v^{}\mathrm{Clo}_n^{}𝐕`$, we have $`v_{}^{}{}_{}{}^{E}𝐯^E=(v^{}𝐯)^E`$, where the clone composition on the left takes place in $`\mathrm{Clo}U(E)`$, and that on the right takes place in $`\mathrm{Clo}𝐕`$. Using equation $`(1)`$ above, we obtain, for all $`𝐚A^n`$,
$`\mathrm{f}^{,\sigma }[v^{}𝐯;𝐚]`$ $`={}_{\sigma (v^{}(𝐯(𝐚)))}{}^{}\chi _{}^{1}(v^{}(𝐯(\sigma (𝐚))))`$
$`=v_{}^{}{}_{𝐯(𝐚)}{}^{M}({}_{\sigma (𝐯(𝐚))}{}^{}\chi _{}^{1}(𝐯(\sigma (𝐚)))+\mathrm{f}^{,\sigma }[v^{};𝐯(𝐚)]`$
$`=v_{}^{}{}_{𝐯(𝐚)}{}^{M}𝐯_𝐚^M({}_{\sigma (𝐚)}{}^{}\chi _{}^{1}(\sigma (𝐚)))+v_{}^{}{}_{𝐯(𝐚)}{}^{M}\mathrm{f}^{,\sigma }[𝐯;𝐚]+\mathrm{f}^{,\sigma }[v^{};𝐯(𝐚)]`$
$`=v_{}^{}{}_{𝐯(𝐚)}{}^{M}\mathrm{f}^{,\sigma }[𝐯;𝐚]+\mathrm{f}^{,\sigma }[v^{};𝐯(𝐚)],`$
for all such $`𝐯`$, $`v^{}`$, and $`𝐚`$, where $`𝐯(𝐚)`$ stands for $`v_1(𝐚),\mathrm{},v_n^{}(𝐚)`$. We will call an $`A`$-function $`\mathrm{f}:X_𝐕^1(A)M`$ satisfying the family of equations
$$\mathrm{f}[v^{}𝐯;𝐚]=v_{}^{}{}_{𝐯(𝐚)}{}^{M}\mathrm{f}[𝐯;𝐚]+\mathrm{f}[v^{};𝐯(𝐚)],$$
for each $`[v^{},𝐯;𝐚]X_𝐕^2(A)`$, a *factor set for $`A`$ and $`M`$*.
Let $`A𝐕`$ and let $`M\text{Ab}[A,𝐕]`$. Given an arbitrary $`A`$-function $`\mathrm{f}:X_𝐕^1(A)M`$, we define an $`A`$-function $`\mathrm{f}:X_𝐕^2(A)M`$ by the equation
$$(\mathrm{f})[v^{},𝐯;𝐚]=v_{}^{}{}_{𝐯(𝐚)}{}^{M}\mathrm{f}[𝐯;𝐚]\mathrm{f}[v^{}𝐯;𝐚]+\mathrm{f}[v^{};𝐯(𝐚)]$$
where $`𝐯(𝐚)`$ is the $`n^{}`$-tuple $`v_1(𝐚),\mathrm{},v_n^{}(𝐚)`$ and $`[𝐯;𝐚]`$ is the $`n^{}`$-tuple $`[v_1;𝐚],\mathrm{},[v_n^{};𝐚]`$.
###### Theorem 4.3.
An $`A`$-function $`\mathrm{f}:X_𝐕^1(A)M`$ is a factor set for $`A`$ and $`M`$ iff $`\mathrm{f}`$ is identically zero, i.e., if for each $`[v^{},𝐯;𝐚]X_𝐕^2(A)`$, $`(\mathrm{f})[v^{},𝐯;𝐚]={}_{v^{}(𝐯(𝐚))}{}^{}0`$.
### Abstract factor sets and abelian extensions
###### Theorem 4.4.
Every factor set for $`A`$ and $`M`$ arises as $`\mathrm{f}^{,\sigma }`$ for some extension $``$ of $`A`$ by $`M`$ and some section $`\sigma `$ of that extension.
###### Proof.
Given a factor set $`\mathrm{f}`$ for $`A`$ and $`M`$, we define $`E=U(AM)`$, and for each $`v\mathrm{Clo}_n(𝐕)`$, the $`n`$-ary operation $`v^E:E^nE`$, by the equation
$$v(a_1,m_1,\mathrm{},a_n,m_n)=v(𝐚),v_𝐚^M(𝐦)+\mathrm{f}[v;𝐚].$$
The fact that $`\mathrm{f}`$ is a factor set makes the mapping $`vv^E`$ a clone homomorphism, i.e., makes $`E`$ an algebra of $`𝐕`$. For, let $`e_i=a_i,m_i`$ for $`i=1,\mathrm{},n`$, let $`𝐯`$ be an $`n^{}`$-tuple of elements of $`\mathrm{Clo}_n𝐕`$, and let $`v^{}\mathrm{Clo}_n^{}𝐕`$. We have
$`(v^{}𝐯)(𝐞)`$ $`=(v^{}𝐯)(𝐚),(v^{}𝐯)_𝐚^M(𝐦)+\mathrm{f}[v^{}𝐯;𝐚]`$
$`=(v^{}𝐯)(𝐚),(v^{}𝐯)_𝐚^M(𝐦)+v_{}^{}{}_{𝐯(𝐚)}{}^{M}\mathrm{f}[𝐯;𝐚]+\mathrm{f}[v^{};𝐯(𝐚)],`$
while
$`v^{}(𝐯(𝐞))`$ $`=v^{}(v_1(𝐚),(v_1)_𝐚^M(𝐦)+\mathrm{f}[v_1;𝐚],\mathrm{}v_n^{}(𝐚),(v_n^{})_𝐚^M(𝐦)+\mathrm{f}[v_n^{};𝐚])`$
$`=v^{}(𝐯(𝐚)),v_{}^{}{}_{𝐯(𝐚)}{}^{M}𝐯_𝐚^M(𝐦)+v_{}^{}{}_{𝐯(𝐚)}{}^{M}\mathrm{f}[𝐯;𝐚]+\mathrm{f}[v^{};𝐯(𝐚)],`$
and the desired equality of these two elements follows from the analogous facts for $`A`$ and $`M`$.
Note $`\pi _M:EA`$ is a homomorphism with respect to the algebra $`E`$ just defined. Let us denote $`\mathrm{ker}\pi _M`$ by $`\kappa `$, and $`{}_{\pi _M}{}^{}\mathrm{Res}M`$ by $`\overline{M}`$. We have $`a,m𝜅a^{},m^{}`$ iff $`a=a^{}`$. We define $`\chi :\overline{M}\kappa ^{}`$ by the equation
$${}_{a,m}{}^{}\chi (m^{})=a,m^{}+m{}_{a,m}{}^{}\kappa _{}^{}.$$
Each $`{}_{a,m}{}^{}\chi `$ preserves the distinguished element because
$${}_{a,m}{}^{}\chi ({}_{a,m}{}^{}0)={}_{a,m}{}^{}\chi ({}_{a}{}^{}0)=a,m={}_{a,m}{}^{}0{}_{a,m}{}^{}\kappa _{}^{}.$$
Also, $`\chi `$ preserves the $`E`$-operations because for each $`v\mathrm{Clo}_n(𝐕)`$, $`e_i=a_i,m_iE`$ for $`i=1,\mathrm{},n`$, and $`m_i{}_{e_i}{}^{}\overline{M}`$ for $`i=1,\mathrm{},n`$, we have
$`v_𝐞^\kappa ^{}({}_{𝐞}{}^{}\chi (𝐦^{}))`$ $`=v_𝐞^\kappa ^{}(a_1,m_1^{}+m_1,\mathrm{},a_n,m_n^{}+m_n)`$
$`=v^E(a_1,m_1^{}+m_1,\mathrm{},a_n,m_n^{}+m_n)`$
$`=v(𝐚),v_𝐚^M(𝐦^{})+v_𝐚^M(𝐦)+\mathrm{f}[v;𝐚],`$
while on the other hand,
$`{}_{v(𝐞)}{}^{}\chi (v_𝐞^{\overline{M}}(𝐦^{}))`$ $`={}_{v(𝐚),v_𝐚^M(𝐦)+\mathrm{f}[v;𝐚]}{}^{}\chi (v_𝐚^M(𝐦^{}))`$
$`=v(𝐚),v_𝐚^M(𝐦^{})+v_𝐚^M(𝐦)+\mathrm{f}[v;𝐚].`$
Thus, $`\chi `$ is a homomorphism of pointed $`E`$-overalgebras, and it is clear that each $`{}_{e}{}^{}\chi `$ is 1-1 and onto. It follows that $`\chi `$ is an isomorphism of pointed $`E`$-overalgebras.
Now, let $`\sigma =\iota _M`$ (so that $`\sigma (a)=a,{}_{a}{}^{}0`$) and let us show that the factor set we obtain is $`\mathrm{f}`$. We have
$${}_{\sigma (\pi _Ma,m)}{}^{}\chi _{}^{1}a,m={}_{a,{}_{a}{}^{}0}{}^{}\chi _{}^{1}a,m,$$
so that the corresponding factor set is given by
$`\mathrm{f}^{,\sigma }[v;𝐚]`$ $`={}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma (𝐚)))`$
$`={}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(a_1,{}_{a_1}{}^{}0,\mathrm{},a_n,{}_{a_n}{}^{}0))`$
$`={}_{v(𝐚),{}_{v(𝐚}{}^{}0}{}^{}\chi _{}^{1}v(𝐚),\mathrm{f}[v;𝐚]`$
$`=\mathrm{f}[v;𝐚],`$
since by definition,
$${}_{v(𝐚),{}_{v(𝐚)}{}^{}0}{}^{}\chi (\mathrm{f}[v;𝐚])=v(𝐚),\mathrm{f}[v;𝐚].$$
### Effect of choice of section on the corresponding factor set
Let us see how factor sets for the same extension $`=\chi ,𝐄,\pi `$ in $`𝐕`$ of $`A`$ by $`M`$ differ when they are derived from different sections $`\sigma `$ and $`\sigma ^{}`$. First, a definition:
Let $`\delta :X_𝐕^0(A)M`$ be an $`A`$-function. We define an $`A`$-function $`\delta :X_𝐕^1(A)M`$ by the equation
$$(\delta )[v;𝐚]=v_𝐚^M(\delta [𝐚])\delta [v(𝐚)],$$
for each $`[v;𝐚]X_𝐕^1(A)`$.
###### Theorem 4.5.
Under these assumptions, we have
$$\mathrm{f}^{,\sigma ^{}}\mathrm{f}^{,\sigma }=\delta _{\sigma ,\sigma ^{}}^{}.$$
###### Proof.
Using equation $`(1)`$ for $`𝐞=\sigma ^{}(𝐚)`$ to expand $`\mathrm{f}^{,\sigma }[v;𝐚]`$, we have
$$\mathrm{f}^{,\sigma ^{}}[v;𝐚]\mathrm{f}^{,\sigma }[v;𝐚]={}_{\sigma ^{}(v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma ^{}(𝐚)))+v_𝐚^M({}_{\sigma (𝐚)}{}^{}\chi _{}^{1}(\sigma ^{}(𝐚))){}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma ^{}(𝐚))).$$
But,
$${}_{\sigma ^{}(v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma ^{}(𝐚)))={}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma ^{}(𝐚))){}_{\sigma (v(a))}{}^{}\chi _{}^{1}(\sigma ^{}(v(𝐚)))$$
by lemma 3.1; combining these two results, we obtain
$`\mathrm{f}^{,\sigma ^{}}[v;𝐚]\mathrm{f}^{,\sigma }[v;𝐚]`$ $`=v_𝐚^M({}_{\sigma (𝐚)}{}^{}\chi _{}^{1}(\sigma ^{}(𝐚))){}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(\sigma ^{}(v(𝐚)))`$
$`=v_𝐚^M(\delta _{\sigma ,\sigma ^{}}^{}[𝐚])\delta _{\sigma ,\sigma ^{}}^{}[v(𝐚)]`$
$`=(\delta _{\sigma ,\sigma ^{}}^{})[v;𝐚],`$
as was to be proved. ∎
## 5. Equivalence of Abelian Extensions
Let $`E`$ and $`\stackrel{~}{E}`$ be algebras of $`𝐕`$, a variety of algebras of type $`\mathrm{\Omega }`$. Let $`\gamma :E\stackrel{~}{E}`$ be a homomorphism, and let $`\alpha `$ and $`\stackrel{~}{\alpha }`$ be congruences of $`E`$ and $`\stackrel{~}{E}`$, respectively, such that $`\gamma (\alpha )\stackrel{~}{\alpha }`$. Then we define an $`E`$-function, $`\gamma ^{}:\alpha ^{}{}_{\gamma }{}^{}\mathrm{Res}(\stackrel{~}{\alpha }^{})`$ by the equation
$${}_{e}{}^{}\gamma _{}^{}(e^{})=\gamma (e^{}).$$
###### Theorem 5.1.
$`\gamma ^{}`$ is a homomorphism of pointed $`E`$-overalgebras.
### Equivalent extensions
Let $`=\chi ,E,\pi `$ and $`\stackrel{~}{}=\stackrel{~}{\chi },\stackrel{~}{E},\stackrel{~}{\pi }`$ be abelian extensions in $`𝐕`$ of $`A`$ by $`M`$, where $`M\text{Ab}[A,𝐕]`$. We define an *equivalence of extensions* from $``$ to $`\stackrel{~}{}`$ to be a homomorphism $`\gamma :E\stackrel{~}{E}`$, such that
1. $`\pi =\stackrel{~}{\pi }\gamma `$, and
2. $`\gamma ^{}\chi ={}_{\gamma }{}^{}\mathrm{Res}\stackrel{~}{\chi }`$.
If $``$ and $`\widehat{}`$ are equivalent via an equivalence $`\gamma `$, we write $`\gamma :\widehat{}`$.
Bearing in mind that, because of condition (1), $`{}_{\gamma }{}^{}\mathrm{Res}{}_{\stackrel{~}{\pi }}{}^{}\mathrm{Res}={}_{\pi }{}^{}\mathrm{Res}`$, we express these conditions in the following interrelated diagrams:
$$\begin{array}{ccccccccccccc}M& & \text{ }& & {}_{\pi }{}^{}\mathrm{Res}M& \stackrel{\chi }{}& \kappa ^{}& & \text{ }& & E& \stackrel{\pi }{}& A\\ & & & & & & \gamma ^{}& & & & \gamma & & & & \\ M& & \text{ }& & {}_{\pi }{}^{}\mathrm{Res}M& \underset{{}_{\gamma }{}^{}\mathrm{Res}\stackrel{~}{\chi }}{}& {}_{\gamma }{}^{}\mathrm{Res}(\stackrel{~}{\kappa }^{})& & \text{ }& & \stackrel{~}{E}& \underset{\stackrel{~}{\pi }}{}& A\\ \\ & & \text{ }& & {}_{\stackrel{~}{\pi }}{}^{}\mathrm{Res}M& \underset{\stackrel{~}{\chi }}{}& \stackrel{~}{\kappa }^{}\end{array}$$
where $`\kappa =\mathrm{ker}\pi `$ and $`\stackrel{~}{\kappa }=\mathrm{ker}\stackrel{~}{\pi }`$.
###### Theorem 5.2.
Equivalence of extensions is an equivalence relation on extensions in $`𝐕`$ of $`A`$ by $`M`$.
###### Proof.
If $`=\chi ,E,\pi `$, $`\stackrel{~}{}=\stackrel{~}{\chi },\stackrel{~}{E},\stackrel{~}{\pi }`$, and $`\overline{}=\overline{\chi },\overline{E},\overline{\pi }`$ are extensions in $`𝐕`$ of $`A`$ by $`M`$, and $`\gamma _1:\stackrel{~}{}`$, $`\gamma _2:\stackrel{~}{}\overline{}`$ are equivalences, then we first observe that
$$(\gamma _2\gamma _1)^{}=({}_{\gamma _1}{}^{}\mathrm{Res}\gamma _2^{})\gamma _1^{}.$$
Then, we have
$`(\gamma _2\gamma _1)^{}\chi `$ $`=({}_{\gamma _1}{}^{}\mathrm{Res}\gamma _2^{})\gamma _1^{}\chi `$
$`=({}_{\gamma _1}{}^{}\mathrm{Res}\gamma _2^{}){}_{\gamma _1}{}^{}\mathrm{Res}\stackrel{~}{\chi }`$
$`={}_{\gamma _1}{}^{}\mathrm{Res}(\gamma _2^{}\stackrel{~}{\chi })`$
$`={}_{\gamma _1}{}^{}\mathrm{Res}({}_{\gamma _2}{}^{}\mathrm{Res}\overline{\chi })`$
$`={}_{\gamma _2\gamma _1}{}^{}\mathrm{Res}\overline{\chi };`$
we also have $`\overline{\pi }\gamma _2\gamma _1=\pi `$, whence $`\gamma _2\gamma _1:\overline{}`$. ∎
###### Lemma 5.3.
If $`\gamma `$ is an equivalence from $``$ to $`\stackrel{~}{}`$, then $`\gamma `$ is an isomorphism, and is the unique equivalence from $``$ to $`\stackrel{~}{}`$.
###### Proof.
$`\chi `$ and $`{}_{\gamma }{}^{}\mathrm{Res}\chi `$ are isomorphisms, whence $`\gamma ^{}`$ is also by condition (2). Thus, $`\gamma `$ maps each $`\kappa `$-class onto a $`\stackrel{~}{\kappa }`$-class, in a one-one fashion. Since each $`\stackrel{~}{\kappa }`$-class is the image of a unique $`\kappa `$-class by condition (1), these mappings paste together to give $`\gamma `$ as an isomorphism.
$`\gamma `$ is unique, because it is determined by $`\gamma ^{}`$, which is determined by condition (2). ∎
We denote by $`E_𝐕(A,M)`$ the set of equivalence classes of extensions in $`𝐕`$ of $`A`$ by $`M`$. It is easy to see that if $`\stackrel{~}{}`$, then $``$ splits iff $`\stackrel{~}{}`$ splits. Thus, the split extensions of $`A`$ by $`M`$ form one or more equivalence classes. We will see below, in corollary 5.6, that all split extensions are equivalent.
### Equivalence and factor sets
We wish to relate factor sets and equivalence.
###### Lemma 5.4.
If $`\gamma :\stackrel{~}{}`$, then given a section $`\sigma `$ of $``$, we have $`\mathrm{f}^{,\sigma }=\mathrm{f}^{\stackrel{~}{},\gamma \sigma }`$.
###### Proof.
For each $`[v;𝐚]`$, we have
$`\mathrm{f}^{\stackrel{~}{},\gamma \sigma }[v;𝐚]`$ $`={}_{\gamma (\sigma (v(𝐚)))}{}^{}\stackrel{~}{\chi }_{}^{1}(v(\gamma (\sigma (𝐚))))`$
$`={}_{\sigma (v(𝐚))}{}^{}({}_{\gamma }{}^{}\mathrm{Res}\stackrel{~}{\chi })_{}^{1}({}_{\sigma (v(𝐚))}{}^{}\gamma _{}^{}(v(\sigma (𝐚))))`$
$`={}_{\sigma (v(𝐚))}{}^{}\chi _{}^{1}(v(\sigma (𝐚)))`$
$`=\mathrm{f}^{,\sigma }[v;𝐚].`$
###### Theorem 5.5.
Let $`M`$ be an $`A`$-module. Let $`=\chi ,E,\pi `$ and $`\stackrel{~}{}=\stackrel{~}{\chi },\stackrel{~}{E},\stackrel{~}{\pi }`$ be two abelian extensions in $`𝐕`$ of $`A`$ by $`M`$, with sections $`\sigma `$ and $`\tau `$, respectively. Then $``$ and $`\stackrel{~}{}`$ are equivalent if and only if the factor sets $`\mathrm{f}^{,\sigma }`$, $`\mathrm{f}^{\stackrel{~}{},\tau }:X_𝐕^1(A)M`$ differ by a function of the form $`\delta `$ where $`\delta :X_𝐕^0(A)M`$ is an $`A`$-function.
###### Proof.
Suppose $``$ and $`\stackrel{~}{}`$ are equivalent via an equivalence $`\gamma `$. By lemma 5.4, $`\mathrm{f}^{,\sigma }=\mathrm{f}^{\stackrel{~}{},\gamma \sigma }`$. On the other hand, $`\mathrm{f}^{\stackrel{~}{},\gamma \sigma }\mathrm{f}^{\stackrel{~}{},\tau }=\delta _{\tau ,\gamma \sigma }^\stackrel{~}{}`$ by theorem 4.5. The desired result follows.
Conversely, if the factor set difference $`\mathrm{f}^{\stackrel{~}{},\tau }\mathrm{f}^{,\sigma }=\delta `$ for $`\delta :X_𝐕^0(A)M`$ some $`A`$-function, then we can produce a new section $`\widehat{\tau }`$ for $`\stackrel{~}{}`$, namely
$$\widehat{\tau }(a)={}_{\tau (a)}{}^{}\stackrel{~}{\chi }(\delta [a]),$$
and we have $`\mathrm{f}^{\stackrel{~}{},\widehat{\tau }}=\mathrm{f}^{,\sigma }`$. For,
$`\mathrm{f}^{\stackrel{~}{},\widehat{\tau }}[v;𝐚]`$ $`={}_{\widehat{\tau }(v(𝐚))}{}^{}\stackrel{~}{\chi }_{}^{1}(v(\widehat{\tau }(𝐚)))`$
$`={}_{\widehat{\tau }(v(𝐚))}{}^{}\stackrel{~}{\chi }_{}^{1}(v({}_{\tau (𝐚)}{}^{}\stackrel{~}{\chi }(\delta [𝐚])))`$
$`={}_{\tau (v(𝐚))}{}^{}\stackrel{~}{\chi }_{}^{1}(v({}_{\tau (𝐚)}{}^{}\stackrel{~}{\chi }(\delta [𝐚]))){}_{\tau (v(𝐚))}{}^{}\stackrel{~}{\chi }_{}^{1}(\widehat{\tau }(v(𝐚)))`$
$`=v_𝐚^M({}_{\tau (𝐚)}{}^{}\chi _{}^{1}({}_{\tau (𝐚)}{}^{}\chi (\delta [𝐚])))+\mathrm{f}^{\stackrel{~}{},\tau }[v;𝐚]{}_{\tau (v(𝐚))}{}^{}\stackrel{~}{\chi }_{}^{1}({}_{\tau (v(𝐚))}{}^{}\stackrel{~}{\chi }(\delta [v(𝐚)]))`$
$`=\mathrm{f}^{\stackrel{~}{},\tau }[v;𝐚]+v_𝐚^M(\delta [𝐚])\delta [v(𝐚)]`$
$`=\mathrm{f}^{\stackrel{~}{},\tau }[v;𝐚]+(\delta )[v;𝐚]`$
$`=\mathrm{f}^{,\sigma }[v;𝐚],`$
as desired.
Now we will construct an isomorphism $`\gamma :E\stackrel{~}{E}`$, which is an equivalence from $``$ to $`\stackrel{~}{}`$. We define the one-one and onto function $`\gamma :E\stackrel{~}{E}`$ by
$$\gamma (e)={}_{\widehat{\tau }(\pi (e))}{}^{}\stackrel{~}{\chi }({}_{\sigma (\pi (e))}{}^{}\chi _{}^{1}(e)).$$
We have $`\pi =\stackrel{~}{\pi }\gamma `$, because $`\gamma `$ maps each $`\pi ^1(a)`$ to $`\stackrel{~}{\pi }^1(a)`$, and we have $`\gamma \sigma =\widehat{\tau }`$, because of lemma 3.3. If we are given $`𝐞E^n`$, then letting $`𝐚=\pi (𝐞)=\stackrel{~}{\pi }(\gamma (𝐞))`$, we have for each $`v\mathrm{Clo}_n(𝐕)`$,
$`\gamma (v(𝐞))`$ $`={}_{\widehat{\tau }(v(𝐚))}{}^{}\stackrel{~}{\chi }(v_𝐚^M({}_{\sigma (𝐚)}{}^{}\chi _{}^{1}(𝐞))+\mathrm{f}^{,\sigma }[v;𝐚])`$
$`={}_{\widehat{\tau }(v(𝐚))}{}^{}\stackrel{~}{\chi }(v_𝐚^M({}_{\widehat{\tau }(𝐚)}{}^{}\stackrel{~}{\chi }_{}^{1}(\gamma (𝐞)))+\mathrm{f}^{\stackrel{~}{},\widehat{\tau }}[v;𝐚])`$
$`={}_{\widehat{\tau }(v(𝐚))}{}^{}\stackrel{~}{\chi }({}_{\widehat{\tau }(v(𝐚))}{}^{}\stackrel{~}{\chi }_{}^{1}(v(\gamma (𝐞)))`$
$`=v(\gamma (𝐞)),`$
whence the 1-1 and onto function $`\gamma `$ is an isomorphism. Finally, for each $`eE`$, with $`a=\pi (e)`$, and all $`m{}_{a}{}^{}M`$, we have
$`{}_{e}{}^{}({}_{\gamma }{}^{}\mathrm{Res}\stackrel{~}{\chi })(m)`$ $`={}_{\gamma (e)}{}^{}\stackrel{~}{\chi }(m)`$
$`={}_{\widehat{\tau }(a)}{}^{}\stackrel{~}{\chi }({}_{\widehat{\tau }(a)}{}^{}\stackrel{~}{\chi }_{}^{1}({}_{\gamma (e)}{}^{}\stackrel{~}{\chi }(m)))`$
$`={}_{\widehat{\tau }(a)}{}^{}\stackrel{~}{\chi }({}_{\widehat{\tau }(a)}{}^{}\stackrel{~}{\chi }_{}^{1}(\gamma (e))+{}_{\gamma (e)}{}^{}\stackrel{~}{\chi }_{}^{1}({}_{\gamma (e)}{}^{}\stackrel{~}{\chi }(m)))`$
$`={}_{\widehat{\tau }(a)}{}^{}\stackrel{~}{\chi }({}_{\sigma (a)}{}^{}\chi _{}^{1}(e)+m)`$
$`={}_{\widehat{\tau }(a)}{}^{}\stackrel{~}{\chi }({}_{\sigma (a)}{}^{}\chi _{}^{1}(e)+{}_{e}{}^{}\chi _{}^{1}({}_{e}{}^{}\chi (m)))`$
$`={}_{\widehat{\tau }(a)}{}^{}\stackrel{~}{\chi }({}_{\sigma (a)}{}^{}\chi _{}^{1}({}_{e}{}^{}\chi (m)))`$
$`={}_{e}{}^{}(\gamma ^{}\chi )(m);`$
thus, $`\gamma ^{}\chi ={}_{\gamma }{}^{}\mathrm{Res}\stackrel{~}{\chi }`$, so that condition (2) is satisfied, and $`\gamma :\stackrel{~}{}`$. ∎
###### Corollary 5.6.
The split extensions of $`A`$ by $`M`$ form a single equivalence class.
###### Proof.
If $``$ is an extension in $`𝐕`$ of $`A`$ by $`M`$, and $`\sigma `$ is a splitting, then $`\mathrm{f}^{,\sigma }0`$ by theorem 4.1. If $`\tau `$ is another section of $``$, then $`\mathrm{f}^{,\tau }=\delta `$ for some $`A`$-function $`\delta :X^0M`$, by theorem 4.5. Thus, every factor set of a split extension has the form $`\delta `$. It follows by the theorem that all split extensions are equivalent. ∎
## 6. $`𝐄_𝐕(A,M)`$ as a Homology Object
### A fragment of a cochain complex
We previously introduced $`A`$-sets $`X^0=X_𝐕^0(A)`$, $`X^1=X_𝐕^1(A)`$, and $`X^2=X_𝐕^2(A)`$, and, with the additional data of an object $`M\text{Ab}[A,𝐕]`$, abelian group homomorphisms $``$, so that the diagram
$$A\text{-}\text{Set}(X^0,M)\stackrel{}{}A\text{-}\text{Set}(X^1,M)\stackrel{}{}A\text{-}\text{Set}(X^2,M)$$
is a fragment of a cochain complex in the abelian category Ab.
### The cohomology group
An $`A`$-function $`\mathrm{f}:X^1M`$ is a factor set for $`A`$ and $`M`$ iff $`\mathrm{f}=0`$, by definition. We showed (theorem 4.4) that each such abstract factor set is a factor set $`\mathrm{f}^{,\sigma }`$ for some extension $``$ and section $`\sigma `$. On the other hand, two factor sets $`\mathrm{f}`$, $`\mathrm{f}^{}`$ correspond to equivalent extensions iff $`\mathrm{f}\mathrm{f}^{}=\delta `$ for some $`A`$-function $`\delta :X^0M`$. It follows that
###### Theorem 6.1.
The equivalence classes of extensions of $`A`$ by $`M`$, i.e., the elements of the set $`𝐄_𝐕(A,M)`$, correspond to the elements of the cohomology group of the above fragment of a cochain complex in the abelian category Ab.
### The equivalence class of split extensions
###### Theorem 6.2.
The class of split extensions is the zero element of $`𝐄_𝐕(A,M)`$.
###### Proof.
Choosing a split extension $``$, and a splitting $`\sigma `$ as the section, we obtain $`\mathrm{f}^{,\sigma }0`$, by theorem 4.1. ∎
### $`𝐄_𝐕(A,M)`$ as an abelian group
The correspondence of theorem 6.1 allows us to place the structure of an abelian group on the set $`𝐄_𝐕(A,M)`$: If $`_1`$, $`_2`$ are represented by factor sets $`\mathrm{f}_1`$ and $`\mathrm{f}_2`$, then $`[_1]+[_2]`$ is represented by $`\mathrm{f}_1+\mathrm{f}_2`$.
###### Theorem 6.3.
If $`=\chi ,E,\pi `$ is an extension of $`A`$ by $`M`$, then $`[]=[\chi ,E,\pi ]`$.
###### Proof.
$`\chi `$ is also an isomorphism, and the formula for factor sets shows that
$$\mathrm{f}^{\chi ,E,\pi ,\sigma }=\mathrm{f}^{\chi ,E,\pi ,\sigma },$$
for any section $`\sigma `$. ∎
## 7. Compositions of Extensions and Homomorphisms
For this section, let $`=\chi ,E,\pi `$ be an abelian extension of $`A`$ by $`M`$. We will describe ways of composing $``$ with homomorphisms in $`𝐕`$ and $`\text{Ab}[A,𝐕]`$. Conceptually, these composition operations are best seen as operations on the *extension class* $`[]`$, the equivalence class of extensions in $`𝐕`$ of $`A`$ by $`M`$ represented by $``$. However, it is not hard to define the composition of an extension and an algebra homomorphism as a specific extension, and sometimes useful, so we will give the definition of that composition in that form.
### The extension $`g`$ for a homomorphism $`g:A^{}A`$
Suppose that we have another algebra $`A^{}`$ and a homomorphism $`g:A^{}A`$. We will show how to create from $``$ an extension of $`A^{}`$ by $`{}_{g}{}^{}\mathrm{Res}M`$. The first step is to consider the fibered product $`E^{}`$ of $`A^{}`$ and $`E`$, with the associated homomorphisms which we label $`\pi ^{}`$ and $`\gamma `$:
$$\begin{array}{ccc}E^{}& \stackrel{\pi ^{}}{}& A^{}\\ \gamma & & g.& & \\ E& \underset{\pi }{}& A\end{array}$$
###### Lemma 7.1.
In this situation, if $`\kappa =\mathrm{ker}\pi `$ and $`\kappa ^{}=\mathrm{ker}\pi ^{}`$, then $`\gamma ^{}:\kappa _{}^{}{}_{}{}^{}{}_{\gamma }{}^{}\mathrm{Res}\kappa ^{}`$ is an isomorphism of pointed $`E^{}`$-overalgebras.
###### Proof.
If $`e,a^{}`$ is an element of $`E^{}`$, then the elements of $`{}_{e,a^{}}{}^{}\kappa _{}^{}{}_{}{}^{}`$ are the elements of $`E^{}`$ of the form $`e^{},a^{}`$, and such a pair will belong to $`E^{}`$ iff $`\pi (e^{})=g(a^{})`$. On the other hand, $`{}_{e,a^{}}{}^{}({}_{\gamma }{}^{}\mathrm{Res}\kappa ^{})={}_{e}{}^{}\kappa _{}^{}`$ is the set of $`e^{}`$ such that $`\pi (e)=\pi (e^{})`$. However, we have $`g(a^{})=\pi (e)`$, because $`e,a^{}E^{}`$. Thus, the mapping $`e^{},a^{}e^{}`$ is one-one and onto. Since this applies to each $`{}_{e,a^{}}{}^{}\gamma _{}^{}`$, $`\gamma ^{}`$ is an isomorphism. ∎
Now, we define $`\chi ^{}=(\gamma ^{})^1{}_{\gamma }{}^{}\mathrm{Res}\chi `$. Since $`\gamma ^{}`$ and $`\chi `$ are isomorphisms, so is $`\chi ^{}`$. Also, we have $`\pi \gamma =g\pi ^{}`$, so that $`\chi ^{}:{}_{\pi ^{}}{}^{}\mathrm{Res}({}_{g}{}^{}\mathrm{Res}M)\kappa _{}^{}{}_{}{}^{}`$. Thus, we obtain an extension $`g=\chi ^{},E^{},\pi ^{}`$ of $`A^{}`$ by $`{}_{g}{}^{}\mathrm{Res}M`$. We will denote the extension class $`[g]`$ by $`[]g`$.
If $`\sigma `$ is any section for $``$, let $`\sigma ^{}:A^{}E^{}`$ be defined by $`\sigma ^{}(a^{})=\sigma (g(a^{})),a^{}`$. Then, we have
###### Lemma 7.2.
If $`\mathrm{f}=\mathrm{f}^{,\sigma }`$, then there is an extension $`\overline{}[]g`$, and a section $`\overline{\sigma }`$, such that or each $`[v;𝐜]X_𝐕^1(A^{})`$, $`\mathrm{f}^{\overline{},\overline{\sigma }}[v;𝐜]=\mathrm{f}[v;g(𝐜)]`$.
###### Proof.
Let $`\overline{}=g`$ and $`\overline{\sigma }=\sigma ^{}`$ as defined above. We have
$`\mathrm{f}^{g,\sigma ^{}}[v;𝐜]`$ $`={}_{\sigma ^{}(v(𝐜))}{}^{}\chi _{}^{}{}_{}{}^{1}(v(\sigma ^{}(𝐜)))`$
$`={}_{\sigma ^{}(v(𝐜))}{}^{}(\gamma _{}^{}{}_{}{}^{1}{}_{\gamma }{}^{}\mathrm{Res}\chi )_{}^{1}(v(\sigma ^{}(𝐜))`$
$`={}_{\sigma ^{}(v(𝐜))}{}^{}({}_{\gamma }{}^{}\mathrm{Res}\chi ^1)(\gamma ^{}(v(\sigma ^{}(𝐜))))`$
$`={}_{\sigma (v(g(𝐜)))}{}^{}\chi _{}^{1}(v(\sigma (g(𝐜))))`$
$`=\mathrm{f}^{,\sigma }[v;g(𝐜)].`$
###### Corollary 7.3.
If in addition to $``$, $`A^{}`$, and $`g`$ we have an algebra $`A^{\prime \prime }`$ and homomorphism $`h:A^{\prime \prime }A^{}`$, then $`(gh)(g)h`$.
###### Corollary 7.4.
The mapping $`[][]g`$ is a group homomorphism from $`𝐄_𝐕(A,M)`$ to the group $`𝐄_𝐕(A^{},{}_{g}{}^{}\mathrm{Res}M)`$.
### The extension class $`\dot{g}[]`$ for a homomorphism $`\dot{g}:MM^{}`$
Suppose again that $`=\chi ,E,\pi `$ is given, and that $`M^{}`$ is another object of $`\text{Ab}[A,𝐕]`$, and $`\dot{g}:MM^{}`$ a homomorphism. $``$ determines an extension class $`[]`$. We will specify an extension class $`\dot{g}[]`$ of $`A`$ by $`M^{}`$ by giving a representative factor set.
Let $`\mathrm{f}^{,\sigma }`$ be a factor set for $``$, where $`\sigma `$ is some section. Since $`\mathrm{f}^{,\sigma }`$ is an $`A`$-function from $`X^1`$ to $`M`$, and $`\dot{g}`$ is a homomorphism from $`M`$ to $`M^{}`$, we can compose them and obtain an $`A`$-function $`\mathrm{f}^{}=\dot{g}\mathrm{f}^{,\sigma }:X^1M^{}`$.
###### Theorem 7.5.
$`\mathrm{f}^{}`$ is a factor set for $`A`$ and $`M^{}`$.
###### Proof.
We have
$`(\mathrm{f}^{})[v^{},𝐯;𝐚]`$ $`=v_{}^{}{}_{𝐯(𝐚)}{}^{M^{}}\mathrm{f}^{}[𝐯;𝐚]\mathrm{f}^{}[v^{}𝐯;𝐚]+\mathrm{f}^{}[v^{};𝐯(𝐚)]`$
$`=\dot{g}v_{}^{}{}_{𝐯(𝐚)}{}^{M}\mathrm{f}[𝐯;𝐚]\dot{g}\mathrm{f}[v^{}𝐯;𝐚]+\dot{g}\mathrm{f}[v^{};𝐯(𝐚)]`$
$`=\dot{g}(\mathrm{f})[v^{},𝐯;𝐚]`$
$`=\dot{g}({}_{v^{}𝐯(𝐚)}{}^{}0_{}^{M})`$
$`={}_{v^{}𝐯(𝐚)}{}^{}0_{}^{M^{}}.`$
Thus, by the construction of theorem 4.4, $`\mathrm{f}^{}`$ is the factor set of some extension, which we shall for the moment denote by $`^{\sigma ,\dot{g}}`$.
###### Theorem 7.6.
If $`\sigma ^{}`$ is another section of $``$, then $`^{\sigma ,\dot{g}}^{\sigma ^{},\dot{g}}`$. I.e., $`[^{\sigma ,\dot{g}}]`$ does not depend on $`\sigma `$.
###### Proof.
$`\mathrm{f}^{,\sigma }\mathrm{f}^{,\sigma ^{}}=\delta `$ for some $`A`$-function $`\delta :X^0M`$. Then, $`\dot{g}\mathrm{f}^{,\sigma }\dot{g}\mathrm{f}^{,\sigma ^{}}=\dot{g}(\delta )`$. However, $`\dot{g}(\delta )=(\dot{g}\delta )`$. For,
$`\dot{g}(\delta )[v;𝐚]`$ $`=\dot{g}(v_𝐚^M(\delta [𝐚])\delta [v(𝐚)])`$
$`=v_𝐚^M^{}(\dot{g}\delta [𝐚])\dot{g}\delta [v(𝐚)]`$
$`=(\dot{g}\delta )[v;𝐚];`$
Thus, $`^{\sigma ,\dot{g}}^{\sigma ^{},\dot{g}}`$ by theorem 5.5. ∎
We denote $`[^{\sigma ,\dot{g}}]`$ by $`\dot{g}[]`$.
###### Theorem 7.7.
If $`M^{\prime \prime }`$ is another object of $`\text{Ab}[A,𝐕]`$, and $`\dot{h}:M^{}M^{\prime \prime }`$ is another homomorphism, then $`(\dot{h}\dot{g})[]=\dot{h}(\dot{g}[])`$.
###### Theorem 7.8.
The mapping $`[]\dot{g}[]`$ is an abelian group homomorphism from $`𝐄_𝐕(A,M)`$ to $`𝐄_𝐕(A,M^{})`$.
###### Proof.
Follows from the definition of $`\dot{g}`$ being a homomorphism of abelian group objects in the category $`\text{Ov}[A,𝐕]`$. ∎
### Relationship of these two compositions
###### Theorem 7.9.
If $``$ is an extension of $`A`$ by $`M`$, and $`g`$ and $`\dot{g}`$ are given, then $`(\dot{g}[])g=({}_{g}{}^{}\mathrm{Res}\dot{g})([]g)`$.
###### Proof.
We will show that representative factor sets for these two extension classes are equal.
Let $`\mathrm{f}`$ be a factor set representing the homology class corresponding to the extension class of $``$. Then $`\dot{g}\mathrm{f}`$ represents $`\dot{g}[]`$. A factor set $`\mathrm{f}^{}`$ for $`(\dot{g}[])g`$ can then be defined by
$$\mathrm{f}^{}:[v;𝐜](\dot{g}\mathrm{f})[v;g(𝐜)].$$
On the other hand, a factor set $`\overline{\mathrm{f}}`$ for $`[]g`$ can be defined by
$$\overline{\mathrm{f}}:[v;𝐜]\mathrm{f}[v;g(𝐜)]$$
and then a factor set for $`({}_{g}{}^{}\mathrm{Res}\dot{g})([]g)`$ can be given by $`({}_{f}{}^{}\mathrm{Res}\dot{g})\overline{\mathrm{f}}`$. For each $`[v;𝐜]`$, we have
$`({}_{g}{}^{}\mathrm{Res}\dot{g})\overline{\mathrm{f}}[v;𝐜]`$ $`={}_{g(v(𝐜))}{}^{}\dot{g}(\mathrm{f}[v;g(𝐜)])`$
$`=\dot{g}\mathrm{f}[v;g(𝐜)]`$
$`=\mathrm{f}^{}[v;𝐜].`$
## 8. The Addition Operation on $`𝐄_𝐕(Q,M)`$
Now we will describe a method for constructing the result of adding a pair $`_1`$, $`_2`$ of extensions of $`A`$ by $`M`$. This method closely parallels the method of adding two module extensions called the Baer sum. However, while the Baer sum uses only universal properties to construct an extension having the desired properties, we can do this only up to a point, and must complete the construction of a representative of $`[_1]+[_2]`$ by using factor sets and theorem 4.4.
### Outer product of a finite tuple of extensions
Let $`_i=\chi _i,E_i,\pi _i`$ be an extension of $`A_i`$ by $`M_i`$ for $`i=1`$, $`\mathrm{}`$, $`n`$, where $`A_i𝐕`$ and $`M_i\text{Ab}[A_i𝐕]`$ for all $`i`$. Define $`E=\mathrm{\Pi }_iE_i`$ and $`A=\mathrm{\Pi }_iA_i`$. Define $`\pi :EA`$ by the equation $`\pi (𝐞)=\pi _1(e_1),\mathrm{},\pi _n(e_n)`$. We will construct an extension $`=\chi ,E,A`$ of $`A`$ by $`_iM_i`$ from the extensions $`_i`$.
If $`\kappa _i=\mathrm{ker}\pi _i`$ for all $`i`$, $`\kappa =\mathrm{ker}\pi `$, and $`\pi _{E,i}:EE_i`$ are the projections to the factors, then we have
###### Lemma 8.1.
$`\kappa ^{}`$ is naturally isomorphic to $`\mathrm{\Pi }_i({}_{\pi _{E,i}}{}^{}\mathrm{Res}\kappa _i^{})`$.
We define $`\chi :{}_{\pi }{}^{}\mathrm{Res}M\kappa ^{}`$ by the equation
$${}_{𝐞}{}^{}\chi (𝐦)={}_{e_1}{}^{}\chi _{1}^{}(m_1),\mathrm{},{}_{e_n}{}^{}\chi _{n}^{}(m_n),$$
and finally, we have
###### Theorem 8.2.
$`\chi ,E,\pi `$ is an extension of $`A`$ by $`M`$.
We denote this extension by $`_i_i`$.
### A factor set of $`_i_i`$
Let $`\sigma _i`$ be a section of $`_i`$ for each $`i`$, and let $`\sigma :AE`$ be the section of $`_i_i`$ defined by $`\sigma (𝐚)=\sigma _1(a_1),\mathrm{},\sigma _n(a_n)`$.
###### Theorem 8.3.
For each $`[v;𝐚_1,\mathrm{},𝐚_{\mathrm{}}]X_𝐕^1(A)`$, we have
$$\mathrm{f}^{_i_i,\sigma }[v;𝐚_1,\mathrm{},𝐚_{\mathrm{}}]=\mathrm{f}^{_1,\sigma _1}[v;a_{11},\mathrm{},a_\mathrm{}1],\mathrm{},\mathrm{f}^{_n,\sigma _n}[v;a_{1n},\mathrm{},a_\mathrm{}n].$$
###### Proof.
We have
$`\mathrm{f}^{_i_i,\sigma }`$ $`[v;𝐚_1,\mathrm{},𝐚_{\mathrm{}}]`$
$`={}_{\sigma (v(𝐚_1,\mathrm{},𝐚_{\mathrm{}}))}{}^{}\chi _{}^{1}(v(\sigma (𝐚_1),\mathrm{},\sigma (𝐚_{\mathrm{}})))`$
$`={}_{\sigma _1(v(a_{11},\mathrm{},a_\mathrm{}1))}{}^{}\chi _{1}^{1}(v(\sigma _1(a_{11},\mathrm{},a_\mathrm{}1))),`$
$`\mathrm{},{}_{\sigma _n(v(a_{1n},\mathrm{},a_\mathrm{}n))}{}^{}\chi _{1}^{1}(v(\sigma _n(a_{1n},\mathrm{},a_\mathrm{}n)))`$
$`=\mathrm{f}^{_1,\sigma _1}[v;a_{11},\mathrm{},a_\mathrm{}1],\mathrm{},\mathrm{f}^{_n,\sigma _n}[v;a_{1n},\mathrm{},a_\mathrm{}n].`$
### The addition operation on $`𝐄_𝐕(Q,M)`$
We observe that $`+^M:M^2M`$ is an arrow of $`\text{Ab}[A,𝐕]`$. Given a pair $`_1,_2`$ of extensions in $`𝐕`$ of $`A`$ by $`M`$, $`(_1_2)\mathrm{\Delta }_A`$ is an extension of $`A`$ by $`{}_{\mathrm{\Delta }_A}{}^{}\mathrm{Res}M^2=M^2`$. Thus, we can form the composite $`+^M[(_1_2)\mathrm{\Delta }_A]`$, which is a class of extensions in $`𝐕`$ of $`A`$ by $`M`$.
###### Theorem 8.4.
We have
$$+^M[(_1_2)\mathrm{\Delta }_A]=[_1]+[_2].$$
###### Proof.
Let sections $`\sigma _1`$, $`\sigma _2`$ of the extensions $`_i`$ be chosen. The factor set $`\mathrm{f}^{_1,\sigma _1}+\mathrm{f}^{_2,\sigma _2}`$ is a factor set of $`[_1]+[_2]`$. On the other hand, a factor set for $`+^M[(_1_2)\mathrm{\Delta }_A]`$ can be computed from the factor set $`\mathrm{f}^{_1_2,\sigma }`$ computed in theorem 8.3. It is easy to see that these factor sets are equal. ∎
## 9. Abelian Extensions as a Bifunctor
We want to be able to treat the $`𝐖`$-algebra of abelian extensions of $`A`$ by $`M`$ as a bifunctor. On objects, we define the bifunctor
$$𝐄_𝐕:(\text{Ov}[A,𝐕])^{\text{op}}\times \text{Ab}[A,𝐕]\text{Ab}$$
by the formula
$$𝐄_𝐕(Q,M)=𝐄_𝐕(AQ,{}_{\pi _Q}{}^{}\mathrm{Res}M),$$
where we recall that $`\pi _Q:AQA`$ is defined by $`\pi _Q:a,qa`$.
On arrows, it suffices to define an abelian group homomorphism $`𝐄_𝐕(r,1_M)`$, which we will write as a composition $`[][]r`$, and an abelian group homomorphism $`𝐄_𝐕(1_Q,\dot{g})`$, which we will write as a composition on the other side, $`[]\dot{g}[]`$, and to prove the following:
1. $`([]r)s=[](rs)`$;
2. $`\dot{h}(\dot{g}([]))=(\dot{h}\dot{g})[]`$; and
3. $`\dot{g}([]g)=(\dot{g}[])g`$,
when those compositions are defined.
If $`𝐄_𝐕(Q,M)`$, $`Q^{}\text{Ov}[A,𝐕]`$, and $`r:Q^{}Q`$ is a homomorphism, then we define $`[]r=[](Ar)`$. $`[]r`$ is an element of
$`𝐄_𝐕(AQ^{},{}_{Ar}{}^{}\mathrm{Res}({}_{\pi _Q}{}^{}\mathrm{Res}M))`$ $`=𝐄_𝐕(AQ^{},{}_{\pi _Q^{}}{}^{}\mathrm{Res}M)`$
$`=𝐄_𝐕(Q^{},M).`$
From corollary 7.4, this mapping is an abelian group homomorphism, and it is easy to see that property (1) holds, from the corresponding fact (corollary 7.3) for composition of extensions with homomorphisms of $`𝐕`$.
On the other hand, if $`M^{}`$ is another object of $`\text{Ab}[A,𝐕]`$ and $`\dot{g}:MM^{}`$ is a homomorphism, we define
$$\dot{g}[]=({}_{\pi _Q}{}^{}\mathrm{Res}\dot{g})[].$$
This is an element of $`𝐄_𝐕(Q,M^{})`$ and it is clear from theorem 7.8 that the mapping is an abelian group homomorphism, and from theorem 7.7 that property (2) holds.
Finally, if $``$, $`Q^{}`$, $`r`$, $`M^{}`$, and $`\dot{g}`$ are all given, we must prove that $`(\dot{g}[])r=\dot{g}([]r)`$, and this will complete the proof that $`𝐄_𝐕`$ is a bifunctor as we have defined it. Using theorem 7.9, we have
$`(\dot{g}[])r`$ $`=(({}_{\pi _Q}{}^{}\mathrm{Res}\dot{g})[])(Ar)`$
$`={}_{Ar}{}^{}\mathrm{Res}({}_{\pi _Q}{}^{}\mathrm{Res}\dot{g})([](Ar))`$
$`={}_{\pi _Q^{}}{}^{}\mathrm{Res}\dot{g}([]r)`$
$`=\dot{g}([]r).`$
## 10. Module Extensions
### $`R`$-modules as a variety of algebras
If $`R`$ is a ring, then the left $`R`$-modules can be treated as algebras having an underlying abelian group structure and one unary operation for each element of $`R`$. The reader can easily supply a list of identities defining the variety of left $`R`$-modules. This is a congruence-modular variety, as can easily be proved, because of the underlying abelian group structures.
It is important to note for what follows that if $`v\mathrm{Clo}_nR`$-Mod, then $`v(𝐦)`$ has the form $`\mathrm{\Sigma }_ir_im_i`$.
### Abelian group overalgebras in the variety of left $`R`$-modules
We need to prove some facts about abelian group overalgebras totally in $`R`$-Mod:
###### Theorem 10.1.
Let $`A`$ and $`B`$ be left $`R`$-modules, and $`f:AB`$ an onto homomorphism. Then every abelian group $`A`$-overalgebra totally in $`R`$-Mod is isomorphic to the restriction, along $`f`$, of an abelian group $`B`$-overalgebra.
###### Proof.
Let $`M`$ be an abelian group $`A`$-overalgebra totally in $`R`$-Mod, and let $`M^{},\eta `$ be a universal arrow to the functor $`{}_{f}{}^{}\mathrm{Res}`$. That is, $`M^{}`$ is an induced abelian group $`B`$-overalgebra in $`R`$-Mod of $`M`$ along $`f`$. By \[14, theorem 10.6\], since $`f`$ is onto, we can equally well consider $`M^{}`$ to be an induced pointed $`B`$-overalgebra of $`M`$ along $`f`$. Thus, \[14, theorem C.6.5\] applies, and the diagram
$$\begin{array}{ccc}AM& \underset{\iota _M}{}& A\\ \mu \left(A\eta _M\right)& & f& & \\ BM^{}& \underset{\iota _M^{}}{}& B\end{array}$$
is a pushout diagram, where $`\mu `$ is the *mashing homomorphism* defined by $`\mu :a,xf(a),x`$. However, $`\pi _M\iota _M=1_A`$, whence $`AMA{}_{0}{}^{}M`$. Similarly, $`BM^{}B{}_{0}{}^{}M_{}^{}`$. It follows easily that $`{}_{0}{}^{}M{}_{0}{}^{}M_{}^{}`$, and that $`M{}_{f}{}^{}\mathrm{Res}M^{}`$. ∎
###### Corollary 10.2.
If $`M`$ is an abelian group $`A`$-overalgebra totally in $`R`$-Mod, then $`M`$ is isomorphic to a restriction, along the unique homomorphism $`\pi _0:A0`$, of an abelian group $`0`$-overalgebra.
###### Corollary 10.3.
If $`M`$, $`M^{}`$ are abelian group $`A`$-overalgebras totally in $`R`$-Mod, and $`\chi :MM^{}`$ is a homomorphism, then $`\chi `$ is determined by $`{}_{0}{}^{}\chi `$.
###### Proof.
By \[14, theorem C.6.4\] restriction along $`\pi `$ is a full functor from $`\text{Pnt}[0,R`$-$`\text{Mod}]`$ to $`\text{Pnt}[A,R`$-$`\text{Mod}]`$. Since $`M`$ and $`M^{}`$ are isomorphic to restrictions along $`\pi `$, $`\chi =\varphi ^1({}_{\pi }{}^{}\mathrm{Res}\widehat{\chi })\varphi ^{}`$, where $`\varphi :M{}_{\pi }{}^{}\mathrm{Res}\widehat{M}`$, $`\varphi ^{}:M^{}{}_{\pi }{}^{}\mathrm{Res}\widehat{M}^{}`$, and $`\widehat{\chi }:\widehat{M}\widehat{M}^{}`$. It follows that $`\chi `$ is determined by any of its components $`{}_{a}{}^{}\chi `$. ∎
### Abelian extensions and module extensions
###### Theorem 10.4.
Let $`Q\text{Ov}[0,R`$-$`\text{Mod}]`$, and let $`M\text{Ab}[0,R`$-$`\text{Mod}]`$. Then $`𝐄_𝐕(Q,M)\mathrm{Ext}(0Q,{}_{0}{}^{}M)`$.
###### Proof.
Given an extension $`\chi ,E,\pi `$ of $`0Q`$ by $`{}_{\pi _Q}{}^{}\mathrm{Res}M`$, let $`\iota ={}_{0}{}^{}\chi `$, and view $`\iota `$ as a homomorphism from $`{}_{0}{}^{}M`$ to $`E`$. Then we have a module extension
$$0{}_{0}{}^{}M\stackrel{𝜄}{}E\stackrel{𝜋}{}0Q0$$
of $`0Q`$ by $`{}_{0}{}^{}M`$.
On the other hand, given $`\iota `$, $`E`$, and $`\pi `$, we define $`{}_{e}{}^{}\chi (m)=e+\iota (m)`$. This gives an $`E`$-function which is clearly one-one and onto, and is an $`E`$-overalgebra homomorphism from $`{}_{\pi _Q\pi }{}^{}\mathrm{Res}M`$ to $`\kappa ^{}`$, where $`\kappa =\mathrm{ker}\pi `$, because
$`v_𝐞^\kappa ^{}({}_{𝐞}{}^{}\chi (𝐦))`$ $`=v_𝐞^\kappa ^{}(e_1+\iota (m_1),\mathrm{},e_n+\iota (m_n))`$
$`=v(e_1+\iota (m_1),\mathrm{},e_n+\iota (m_n))`$
$`=v(𝐞)+v(\iota (𝐦))`$
$`=v(𝐞)+\iota (v^M(𝐦))`$
$`=v(𝐞)+\iota (v_{0,\mathrm{},0}^M(𝐦))`$
$`={}_{v(𝐞)}{}^{}\chi (v_{0,\mathrm{},0}^M(𝐦))`$
$`={}_{v(𝐞)}{}^{}\chi (v_𝐞^{{}_{\pi _Q\pi }{}^{}\mathrm{Res}M}(𝐦));`$
thus, $`\chi `$ is an isomorphism, and $`\chi ,E,\pi `$ is an extension in $`R`$-Mod of $`0Q`$ by $`{}_{\pi _Q}{}^{}\mathrm{Res}M`$, i.e., an element of $`𝐄_𝐕(Q,M)`$.
The first construction, of $`\iota `$ from $`\chi `$, is one-one by corollary 10.3, but it is also onto, because if we start with $`\iota `$ and construct $`\chi `$, then we get $`\iota `$ back as $`\iota ={}_{0}{}^{}\chi `$. Thus, the two mappings are inverses to each other.
If two extensions $`=\chi ,E,\pi `$ and $`\stackrel{~}{}=\stackrel{~}{\chi },\stackrel{~}{E},\stackrel{~}{\pi }`$ in $`R`$-Mod of $`0Q`$ by $`M`$ are given, and $`\gamma :\stackrel{~}{}`$, then from $`\gamma ^{}\chi ={}_{\gamma }{}^{}\mathrm{Res}\stackrel{~}{\chi }`$ we obtain that $`\gamma {}_{0}{}^{}\chi ={}_{0}{}^{}\gamma _{}^{}{}_{0}{}^{}\chi ={}_{0}{}^{}\stackrel{~}{\chi }`$, whence we have a commutative diagram
$$\begin{array}{ccccc}{}_{0}{}^{}M& \stackrel{{}_{0}{}^{}\chi }{}& E& \underset{\pi }{}& 0Q\\ & & \gamma & & & & \\ {}_{0}{}^{}M& \underset{{}_{0}{}^{}\stackrel{~}{\chi }}{}& \stackrel{~}{E}& \underset{\stackrel{~}{\pi }}{}& 0Q\end{array}$$
showing that the constructed module extensions remain equivalent, by the customary definition. On the other hand, given a commutative diagram
$$\begin{array}{ccccc}{}_{0}{}^{}M& \stackrel{\iota }{}& E& \underset{\pi }{}& 0Q\\ & & \gamma & & & & \\ {}_{0}{}^{}M& \underset{\stackrel{~}{\iota }}{}& \stackrel{~}{E}& \underset{\stackrel{~}{\pi }}{}& 0Q\end{array}$$
then let $`\chi `$ and $`\stackrel{~}{\chi }`$ be constructed from $`\iota `$ and $`\stackrel{~}{\iota }`$: we have $`\stackrel{~}{\pi }\gamma =\pi `$ and for all $`eE`$, and $`m{}_{0}{}^{}M`$,
$`{}_{e}{}^{}\gamma _{}^{}{}_{e}{}^{}\chi (m)`$ $`={}_{e}{}^{}\gamma _{}^{}(e+\iota (m))`$
$`=\gamma (e)+\gamma \iota (m)`$
$`=\gamma (e)+\stackrel{~}{\iota }(m)`$
$`={}_{\gamma (e)}{}^{}\stackrel{~}{\chi }(m)`$
$`={}_{e}{}^{}({}_{\gamma }{}^{}\mathrm{Res}\stackrel{~}{\chi })(m),`$
whence $`\gamma ^{}\chi ={}_{\gamma }{}^{}\mathrm{Res}\stackrel{~}{\chi }`$. Thus, the constructed elements of $`𝐄_𝐕(Q,M)`$ are equivalent by our definition.
We omit the verifications that this isomorphism between the bifunctors $`𝐄_𝐕(Q,M)`$ and $`\mathrm{Ext}(0Q,{}_{0}{}^{}M)`$ of $`M`$ and $`Q`$ is natural in both arguments, and that it preserves the abelian group operations. ∎
## 11. Clone Cohomology
The partial cochain complex which we introduced in §6, leading to the cohomology group we have identified as the set of extensions $`𝐄_𝐕(A,M)`$, can be extended to a full positive cochain complex (i.e., with objects $`C^i\text{Ab}[A,𝐕]`$ for all $`i0`$) and we will do so in this section. As we do, we want to make some small changes in our development. We want to develop the theory in such a way that, as we did in defining $`𝐄_𝐕(Q,M)`$ in §9, the resulting cohomology objects are bifunctors contravariant in $`\text{Ov}[A,𝐕]`$ and covariant in $`\text{Ab}[A,𝐕]`$. In addition, we want to express the cochain complex as a result of applying a hom functor to a chain complex we will define. Finally, we want to make explicit the simplicial aspect of the definition. We will call the resulting cohomology theory *clone cohomology*, because of the role played in the definition by the clone of the variety $`𝐕`$.
### The $`A`$-sets $`X_𝐕^i(Q)`$
As a first step, we will define functors $`X_𝐕^i:\text{Ov}[A,𝐕]A\text{-}\text{Set}`$. The $`A`$-sets $`X_𝐕^i(A)`$, for $`i=0`$, $`1`$, and $`2`$, can be seen as the special cases of $`X_𝐕^i(Q)`$ where $`Q`$ is the $`A`$-set $`[[A,1_A]]`$.
If $`Q`$ is an object of $`\text{Ov}[A,𝐕]`$, we will define $`X_𝐕^0(Q)`$ to be the underlying $`A`$-set of $`[[AQ,\pi _Q]]`$, with the element $`a,q{}_{a}{}^{}X_{𝐕}^{0}(Q)`$ being written as $`[q]_a`$.
We define $`X_𝐕^1(Q)`$ to be the $`A`$-set given by letting $`{}_{a}{}^{}X_{𝐕}^{1}(Q)`$ be the set of triples, written $`[v;𝐪]_𝐚`$, where $`v\mathrm{Clo}_n𝐕`$ for some $`n`$, $`𝐚A^n`$, and $`𝐪{}_{𝐚}{}^{}Q_{}^{n}`$, and such that $`v(𝐚)=a`$.
If $`i>1`$, we will define $`{}_{a}{}^{}X_{𝐕}^{i}(Q)`$ to be the set of $`(i+2)`$-tuples, written
$$[v_0,𝐯_1,\mathrm{},𝐯_{i1};𝐪]_𝐚,$$
where $`v_0`$ is an element of $`\mathrm{Clo}_{n_0}𝐕`$, $`𝐯_j`$ for $`0<j<i`$ is an $`n_{j1}`$-tuple of elements of $`\mathrm{Clo}_{n_j}(𝐕)`$, $`𝐚A^{n_{i1}}`$, and $`𝐪{}_{𝐚}{}^{}Q_{}^{n_i}`$, all for some natural numbers $`n_0`$, $`\mathrm{}`$, $`n_{i1}`$, and such that $`v_0𝐯_1\mathrm{}𝐯_{i1}(𝐚)=a`$.
Note that $`v_0`$, $`𝐯_1`$, $`\mathrm{}`$, $`𝐯_{i1}`$ can be considered as composable arrows of the theory constructed from $`\mathrm{Clo}𝐕`$.
### $``$
First, we will define $``$ directly, making $`A`$-$`\text{Set}(X_𝐕^{}(Q),M)`$ into a chain complex. Given an $`A`$-function $`\mathrm{f}:X_𝐕^i(Q)M`$, we define $`\mathrm{f}`$ by
$$(\mathrm{f})[v;𝐪]_𝐚=v_𝐚^M\mathrm{f}[𝐪]_𝐚\mathrm{f}[v_𝐚^Q(𝐪)]_{v(𝐚)}$$
if $`i=0`$, where $`[𝐪]_𝐚`$ stands for $`[q_1]_{a_1},\mathrm{},[q_n]_{a_n}`$, and otherwise, by
$`(\mathrm{f})[v_0,𝐯_1,\mathrm{},𝐯_{i1};𝐪]_𝐚`$ $`=(v_0)_{𝐯_1\mathrm{}𝐯_{i1}(𝐚)}^M\mathrm{f}[𝐯_1,\mathrm{},𝐯_{i1};𝐪]_𝐚`$
$`\mathrm{f}[v_0𝐯_1,𝐯_2,\mathrm{},𝐯_{i1};𝐪]_𝐚`$
$`+\mathrm{f}[v_0,𝐯_1𝐯_2,\mathrm{},𝐯_{i1};𝐪]_𝐚`$
$`\mathrm{}`$
$`+(1)^j\mathrm{f}[v_0,𝐯_1,\mathrm{},𝐯_{j1}𝐯_j,\mathrm{},𝐯_{i1};𝐪]_𝐚`$
$`\mathrm{}`$
$`+(1)^i\mathrm{f}[v_0,𝐯_1,\mathrm{},𝐯_{i2};(𝐯_{i1})_𝐚^Q(𝐪)]_{𝐯_{i1}(𝐚)},`$
where $`\mathrm{f}[𝐯_1,\mathrm{},𝐯_{i1};𝐪]_𝐚`$ stands for $`\mathrm{f}[v_{11},𝐯_2,\mathrm{},𝐯_{i1};𝐪]_𝐚,\mathrm{},\mathrm{f}[v_{1n_0},𝐯_2,\mathrm{},𝐯_{i1};𝐪]_𝐚`$.
A straightforward computation shows that $`=0`$, making $`A`$-$`\text{Set}(X_𝐕^{}(Q),M)`$ into a cochain complex, which we denote by $`C_𝐕^{}(Q,M)`$..
Another way to arrive at the same cochain complex is to apply the free functor from $`A\text{-}\text{Set}`$ to $`\text{Ab}[A,𝐕]`$ to the $`A`$-sets $`X_𝐕^i(Q)`$, yielding objects $`C_i(Q,𝐕)\text{Ab}[A,𝐕]`$. We then consider the simplicial complex of objects of $`\text{Ab}[A,𝐕]`$ given by the $`C_i(Q,𝐕)`$ and face maps $`_j^i`$, for $`i0`$ and $`0ji`$, defined on generators by
$$_0^0[q]_a=0,$$
by
$$_0^1[v;𝐪]_𝐚=v_𝐚^M[𝐪]_𝐚\text{ and }_1^1[v;𝐪]_𝐚=[v_𝐚^Q(𝐪)]_{v(𝐚)},$$
and by
$`_j^i[v_0,𝐯_1`$ $`,\mathrm{},𝐯_{i1};𝐪]_𝐚`$
$`=\{\begin{array}{cc}(v_0)_{𝐯_1\mathrm{}𝐯_{i1}(𝐚)}^M[𝐯_1,\mathrm{},𝐯_{i1};𝐪]_𝐚,\hfill & \text{for }j=0\text{,}\hfill \\ [v_0,𝐯_1,\mathrm{},𝐯_{j1}𝐯_j,\mathrm{},𝐯_{i1};𝐪]_𝐚,\hfill & \text{for }0<j<i\text{, and}\hfill \\ [v_0,𝐯_1,\mathrm{},𝐯_{i2};(𝐯_{i1})_𝐚^Q(𝐪)]_{𝐯_{i1}(𝐚)},\hfill & \text{for }j=i\text{,}\hfill \end{array}`$
for $`i2`$. (Edge maps can also be defined, using projection elements of $`\mathrm{Clo}𝐕`$, but we have no need to do so.) We then form $``$ as the alternating sum of the face maps $`_j^i`$ over $`j`$, and obtain a chain complex $`C_{}(Q,𝐕)`$ such that $`\text{Ab}[A,𝐕](C_{}(Q,𝐕),M)=C_𝐕^{}(Q,M)`$, the same cochain complex we described previously.
### Clone cohomology
We define the *clone cohomology objects for $`Q`$, with coefficients in $`M`$*, to be the cohomology groups of the cochain complex $`C_𝐕^{}(Q,M)`$, and denote them by $`H_𝐕^i(Q,M)`$. Clearly, they are bifunctors $`H^i:\text{Ov}[A,𝐕]^{\text{op}}\times \text{Ab}[A,𝐕]\text{Ab}`$.
Note that the clone cohomology objects are defined whether or not the variety $`𝐕`$ is congruence-modular.
### Factor sets in terms of $`Q`$
We previously defined factor sets of extensions of $`A`$ by $`M`$, and then defined extensions of an $`A`$-overalgebra $`Q`$ by $`M`$. We have not defined factor sets in terms of $`Q`$ yet, and will not do so in detail. The key fact which allows us to relate our previous work with factor sets, and the objects $`X_𝐕^{}(Q)`$ and $`C_𝐕^{}(Q,M)`$, is as follows:
###### Theorem 11.1.
For $`i=0`$, $`1`$, and $`2`$, the three functors to $`\text{Ab}[A,𝐕]`$,
$$A\text{-}\text{Set}(X_𝐕^i(AQ),{}_{\pi _Q}{}^{}\mathrm{Res}M),$$
$$A\text{-}\text{Set}(X_𝐕^i(Q),M),\text{ and}$$
$$C_𝐕^i(Q,M)$$
are naturally isomorphic as bifunctors in $`Q`$ and $`M`$.
###### Proof.
The first and second functors are naturally isomorphic, because of an adjunction between the functor of restriction of $`A`$-sets along $`\pi _Q`$ and the functor from $`(AQ)`$-sets to $`A`$-sets given by sending an $`(AQ)`$-set $`[[B,\pi ]]`$ to the $`A`$-set $`[[B,\pi _Q\pi ]]`$.
The second and third functors are naturally isomorphic, because of the adjunction between the free functor from $`A`$-Set to $`\text{Ab}[A,𝐕]`$, and the corresponding forgetful functor. ∎
## 12. Relative Clone Cohomology $`𝐕`$
The clone cohomology functors $`H_𝐕^i(Q,M)`$ are specific to a given variety of algebras $`𝐕`$, such that $`Q`$ and $`M`$ are totally in $`𝐕`$. However, the construction of the chain complex used in computing them can use a smaller variety $`𝐕^{}`$, if $`M`$ is totally in $`𝐕^{}`$. That is, in that case, we can form the free object in $`\text{Ab}[A,𝐕^{}]`$ on $`A`$-set of generators $`X_𝐕^i(Q)`$, which we denote by $`C_{i,𝐕^{}}(Q,𝐕)`$, rather than the free object in $`\text{Ab}[A,𝐕]`$ which we denoted by $`C_i(Q,𝐕)`$. Because $`M`$ is totally in $`𝐕^{}`$, the hom functor $`\text{Ab}[A,𝐕](,M)`$ takes these objects to the same abelian group. The definition of $``$ also makes sense. As a result,
###### Theorem 12.1.
If $`M`$ is totally in $`𝐕^{}`$, then the objects $`C_𝐕^i(Q,M)`$ and $`H_𝐕^i(Q,M)`$ are also totally in $`𝐕^{}`$.
If $`Q`$ is also totally in $`𝐕^{}`$, then the definition of $``$ for the chain complex $`C_{}(Q,𝐕^{})`$ makes sense, and there is an obvious homomorphism $`\pi _i`$ of $`C_i(Q,𝐕)`$ onto $`C_i(Q,𝐕^{})`$ for each $`i`$. As a result, we have a short exact sequence of complexes in $`\text{Ab}[A,𝐕]`$,
$$\begin{array}{ccccccc}K_0(Q,𝐕^{},𝐕)& & K_1(Q,𝐕^{},𝐕)& & K_2(Q,𝐕^{},𝐕)& & \mathrm{}\\ & & & & & & \\ C_{0,𝐕^{}}(Q,𝐕)& & C_{1,𝐕^{}}(Q,𝐕)& & C_{2,𝐕^{}}(Q,𝐕)& & \mathrm{}\\ \pi _0& & \pi _1& & \pi _2& & \\ C_0(Q,𝐕^{})& & C_1(Q,𝐕^{})& & C_2(Q,𝐕^{})& & \mathrm{}\end{array}$$
where the $`K_i(Q,𝐕^{},𝐕)`$ are the kernels of the onto homomorphisms $`\pi _i`$. Note that the vertical short exact sequences split. For, each generator of $`C_i(Q,𝐕^{})`$ can be mapped to a preimage in $`C_{i,𝐕^{}}(Q,𝐕)`$, leading to a splitting of $`\pi _i`$.
Applying the hom functor $`\text{Ab}[A,𝐕](,M)`$ and denoting $`\text{Ab}[A,𝐕](K_i(Q,𝐕^{},𝐕),M)`$ by $`C_{𝐕^{},𝐕}^i(Q,M)`$, we obtain a short exact sequence of complexes of objects of $`\text{Ab}[A,𝐕^{}]`$:
$$\begin{array}{ccccccc}C_{𝐕^{},𝐕}^0(Q,M)& & C_{𝐕^{},𝐕}^1(Q,M)& & C_{𝐕^{},𝐕}^2(Q,M)& & \mathrm{}\\ & & & & & & \\ C_𝐕^0(Q,M)& & C_𝐕^1(Q,M)& & C_𝐕^2(Q,M)& & \mathrm{}\\ & & & & & & \\ C_𝐕^{}^0(Q,M)& & C_𝐕^{}^1(Q,M)& & C_𝐕^{}^2(Q,M)& & \mathrm{}\end{array}$$
where, again, the vertical sequences are exact.
We define the *relative clone cohomology of $`Q`$, with coefficients in $`M`$, with respect to the inclusion $`𝐕^{}𝐕`$*, to be the cohomology objects of the complex $`C_{𝐕^{},𝐕}^{}(Q,M)`$, and denote these objects by $`H_{𝐕^{},𝐕}^i(Q,M)`$. It is clear that $`C_𝐕^0(Q,M)`$ and $`C_𝐕^{}^0(Q,M)`$ are isomorphic; thus, $`C_{𝐕^{},𝐕}^0(Q,M)=H_{𝐕^{},𝐕}^0(Q,M)=0`$. We then see that there is a long exact sequence of cohomology groups
$$0H_𝐕^{}^1(Q,M)H_𝐕^1(Q,M)H_{𝐕^{},𝐕}^1(Q,M)H_𝐕^{}^2(Q,M)\mathrm{}$$
relating the three sets of cohomology objects $`H_𝐕^{}^{}(Q,M)`$, $`H_𝐕^{}(Q,M)`$, and $`H_{𝐕^{},𝐕}^{}(Q,M)`$.
If we have three varieties $`𝐕^{\prime \prime }𝐕^{}𝐕`$ such that $`Q`$ and $`M`$ are totally in $`𝐕^{\prime \prime }`$, then similar methods, and standard methods from homological algebra, yield a long exact sequence
$$0H_{𝐕^{\prime \prime },𝐕^{}}^1(Q,M)H_{𝐕^{\prime \prime },𝐕}^1(Q,M)H_{𝐕^{},𝐕}^1(Q,M)H_{𝐕^{\prime \prime },𝐕^{}}^2(Q,M)\mathrm{}$$
relating the relative cohomology objects.
## Discussion
The study of abelian extensions, and the recognition that they form a cohomology group, has a considerable history, as does the investigation of cohomology theories in general. We have focussed our attention on the abelian group of extensions, and have defined a cohomology theory that has this algebra as cohomology group in dimension one. Many questions remain to be answered about this new cohomology theory, and about its relationship to previously-studied theories of cohomology of algebras. We will raise some of those questions in this section.
Categorical algebraists have invented a cohomology theory called comonadic cohomology. The theory can be applied to any comonad, but usually, the comonad being used when this theory is mentioned is a comonad derived from the free and forgetful functors for the variety in question. See , , and . As in clone cohomology, the theory gives the cohomology group of an $`A`$-overalgebra (actually, an algebra over $`A`$, as the theory is usually developed) totally in a variety $`𝐕`$ to which $`A`$ belongs, with coefficients in an $`A`$-module $`M`$ totally in $`𝐕`$ (or, as usually expressed, in a Beck module over $`A`$). As in clone cohomology, the resolution that gives rise to the cohomology objects comes from taking alternating sums of face maps of a simplicial complex.
We are led to ask, what is the relationship of the clone cohomology objects to the comonadic cohomology objects? In , the comonadic cohomology groups in dimension 0 and 1 were studied and interpreted. This was done in the generality of $`𝐕`$ an arbitrary (i.e., not necessarily congruence-modular) variety of algebras of some type. The interpretation of the group in dimension 1 is different from our interpretation of the clone cohomology group in dimension 1, but, we have proved directly that, for $`𝐕`$ a congruence-modular variety of algebras, the groups are naturally isomorphic. We have not discussed dimension 0 in this paper, but the groups in dimension 0 are also isomorphic. (We have not included the proofs of these results in this paper.)
This raises the question, are the groups isomorphic in all dimensions, when $`𝐕`$ is congruence-modular? More generally, are they isomorphic when $`𝐕`$ is not congruence-modular? It is not hard to show that they are isomorphic in dimension 0, but the answer is not known for higher dimensions, or even for dimensions higher than 1 in the congruence-modular case.
For clone cohomology, and for comonadic cohomology for that matter, there is the question, when are the resolutions used in the derivation of the cohomology groups exact? Note that in both cases, the resolutions are free. Thus, when the resolutions are exact, the cohomology functors will be derived functors in the standard sense. Otherwise, for comonadic cohomology, there is at least a uniqueness theorem () which characterizes the cohomology functors. For clone cohomology, there is not yet such a theorem.
For some well-known cases, such as the variety of groups, the comonadic and clone cohomology groups coincide with the standard cohomology groups which can be defined as derived functors in the usual sense. Thus, we can say that we have characterized these groups using two different universal properties. What is the significance of the fact that the two different universal properties arrive at the same answer?
Finally, we should mention that there is a formal, at least, resemblence between the derivation of the clone cohomology objects, and the construction of the cohomology groups of a category. What is the significance of this resemblence?
To our knowledge, none of the questions have yet been answered.
### Acknowledgement
It is a pleasure to acknowledge helpful discussions and correspondence about this topic with Saunders Mac Lane.
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# Pairing Gaps from Nuclear Mean–Field Models
## 1 Introduction
Pair correlations, which play a crucial role in superconducting solids Bar57a , also constitute an important complement of nuclear shell structure Boh58a ; Bel58a . Most often, pairing is treated in the so–called BCS approximation EGIII ; Ringbook ; Nilssonbook , which was introduced in the original paper of Bardeen, Cooper and Schrieffer Bar57a . The nuclear applications are particular in two respects: first, nuclei are finite, in fact rather small, objects, and second, we do not yet have a sufficiently reliable microscopic nuclear many–body theory from which we could deduce the appropriate pairing interaction and its strength. The second problem causes two further questions. First, one needs to develop a reliable and manageable form for the pairing energy functional, and second, one has to determine an appropriate pairing strength for a given functional.
There exist various prescriptions for the pairing energy functional. Schematic pairing forces basically consist of defining a small band of pairing–active states and parameterizing one typical pairing matrix element in dependence on the system size, i.e. neutron and proton number. They are convenient and successful in many respects. But they are plagued by serious disadvantages: the coupling to continuum states is much exaggerated and the parameterization in terms of system size becomes questionable for deformed systems along the fission path. To avoid these problems local two–body pairing forces are increasingly used which is particularly satisfying in connection with self–consistent mean–field calculations Ton79a ; Kri90a ; Dob95a ; Dob96c . There are even some mean–field models like the Gogny forces Dob96c ; Gogrev or the particular Skyrme force SkP Dob84a which aim at a simultaneous description of the particle–particle and particle–hole channels of the effective interaction using the same force. That is not a necessary condition. The simple local pairing energy functionals also provide a very good description of pairing properties throughout the whole chart of isotopes using only two universal strength parameters, one for protons and one for neutrons. This makes these forms for the pairing energy functional more reliable for calculations of deformed nuclei and of nuclei far off the valley of stability than the widely used schematic pairing forces. In the following we will concentrate on the class of local pairing interactions.
The small particle number of nuclei often interferes with the fact that the BCS ground state mixes particle numbers with a relative spread of order $`1/\sqrt{N}`$. In principle, one has to perform a projection of the BCS state before variation, but this exact particle number projection can be very cumbersome, see, e.g., Mang ; MONSTER . The Lipkin–Nogami (LN) scheme offers a reliable approximate projection method Lip60a ; Nog64a ; Pra73a ; Flo . It has the technical advantage that it is formulated completely in terms of BCS expectation values which makes its numerical implementation very simple, for detailed discussion see Pra73a ; Flo ; Zhe92a ; Dob92a . We will discuss the LN scheme side by side with the BCS treatment.
There remains as the last and crucial problem the determination of an appropriate pairing strength. Insufficient microscopic information requires that one recurs to a phenomenological assessment. This line of development has been followed with increasing accuracy since the introduction of pairing in nuclei, for a comprehensive compilation of pairing with schematic forces see Mol92a . Also the local two–body pairing interactions leave the overall strength as a free parameter to be fixed phenomenologically. This has to be done with respect to an observable sensitive to pairing correlations. Such an observable would ideally be provided by the pairing gap which, however, is not directly accessible in experiments. The observable quantity with probably closest relation to the pairing gap is the pronounced odd–even mass staggering of nuclei, which is usually used for fitting the pairing strength. Even here, though, there remains a choice of several recipes on how to extract the pairing gap or odd–even staggering, respectively, from a combination of neighboring mass values. From the theoretical side, of course, one has more direct access to a gap. But even there ambiguities emerge and one is left with several possibilities to define measures for the pairing correlations.
It is the aim of this paper to compare and discuss the various definitions of a pairing gap in models of local two-body pairing interactions in connection with self–consistent mean–field approaches, both at the level of the conventional BCS approach as well as of the LN method. Furthermore, we will compare two variants of the local zero–range pairing forces, a delta force and a density–dependent form Dob95a ; Dob96c . Last but not least, the discussion will extend to exotic nuclei where differences between the various definitions for the pairing gap and options for the pairing method become particularly obvious.
The paper is outlined as follows: In Section 2 we present the basic equations of the Lipkin–Nogami model employing local interactions in the framework of the Skyrme–Hartree–Fock model needed for our discussion. In Section 3 various approximations for the pairing gap are discussed, in Section 4 the results are presented, Section 5 summarizes our findings.
## 2 The Theoretical Framework
We investigate the pairing gap in the framework of self–consistent mean–field models. Pairing correlations are treated on the HF+BCS level where the equations of motion are derived by independent variation with respect to single–particle wave functions and occupation amplitudes. This is a widely used approximation to the more involved HFB approach where wave functions and occupation amplitudes are varied simultaneously Ringbook . The BCS approximation is applicable for all well–bound nuclei, i.e. for most of the nuclei discussed throughout this paper. It becomes critical only for nuclei close to the neutron or proton drip–line, see e.g. Dob84a , and such nuclei are not considered here. All conclusions about pairing gaps drawn in this paper can thus safely be derived in the BCS approximation.
### 2.1 The Mean Field
Presently the most widely used self–consistent mean–field models are the non–relativistic Hartree–Fock approach with either the Skyrme (SHF) SHFrev or the Gogny force Gogrev , and the relativistic mean–field model RMFrev . They all provide a well-adjusted effective energy functional for nuclear mean–field calculations. We choose the SHF model for the present investigation. The description starts from an energy functional
$$=_{\mathrm{mf}}+_{\mathrm{pair}},$$
(1)
whose mean–field part $`_{\mathrm{mf}}`$ is formulated in terms of the local distributions of density $`\rho `$, kinetic density $`\tau `$, and spin–orbit current $`𝐉`$. At this point it is not necessary to unfold all details of this rather elaborate functional, we abbreviate the dependence with the most general case, the full one–body density matrix
$$\widehat{\rho }\rho (\stackrel{}{x},\stackrel{}{x}^{})=\widehat{\psi }^{}(\stackrel{}{x}^{})\widehat{\psi }(\stackrel{}{x})$$
(2)
from which all local densities and currents can be derived, see appendix A.2 for details. $`\widehat{\psi }^{}(\stackrel{}{x})`$ creates a particle with spin projection $`\sigma /2`$ at the space point $`\stackrel{}{r}`$. Throughout this paper $`\mathrm{}`$ denotes BCS expectation values. For the Skyrme energy functional we choose the rather recent parameterization SkI4 Rei95a . The pairing energy functional depends additionally on the local pair density $`\chi `$ as introduced in Section 2.2. Variation of the energy functional $``$ with respect to the single–particle wave functions $`\varphi _k`$ yields the mean–field equations
$$\widehat{h}\varphi _k=\epsilon _k\varphi _k\text{with}\widehat{h}=\frac{\delta }{\delta \widehat{\rho }},$$
(3)
where we have neglected the contributions from the variation of the pairing density $`\chi `$ in the energy functional to the equations–of–motion of the single–particle states $`\varphi _k`$. This constitutes the BCS approximation to pairing (see Rei97a for the discussion of the full HFB equations in the representation in natural orbitals that is used here). The mean–field equations are solved on a grid in coordinate space with the damped gradient iteration method and a Fourier representation of the derivatives. The numerical techniques are summarized in Blum .
### 2.2 The Pairing Energy Functional
We parameterize the effective pairing interaction in terms of a local pairing energy functional of the form
$$_{\mathrm{pair}}=\frac{1}{4}\underset{q\{p,n\}}{}\mathrm{d}^3r\chi _q^{}(\text{r})\chi _q(\text{r})G_q(\text{r}),$$
(4)
which allows for a spatial modulation of the strength $`G(\text{r})`$. $`\chi (\text{r})`$ is the local part of the pair density matrix
$`\chi _q(\stackrel{}{r})`$ $`=`$ $`{\displaystyle \underset{\sigma =\pm }{}}\chi _q(\stackrel{}{r},\sigma ;\stackrel{}{r},\sigma )={\displaystyle \underset{\sigma =\pm }{}}\sigma \widehat{\psi }_q(\stackrel{}{r},\sigma )\widehat{\psi }_q(\stackrel{}{r},\sigma )`$ (5)
$`=`$ $`2{\displaystyle \underset{\genfrac{}{}{0pt}{}{k\mathrm{\Omega }_q}{k>0}}{}}u_kv_k|\varphi _k(\text{r})|^2,`$
with $`q\{p,n\}`$. The $`\varphi _k`$ are the single–particle wave functions and $`v_k`$, $`u_k=\sqrt{1v_k^2}`$ the pairing amplitudes. We restrict ourselves to stationary states of time–reversal invariant systems and pairing between like particles only. This sorts the single–particle states into conjugate pairs $`k\overline{k}k`$ with the same spatial properties but opposite projection of the total angular momentum and allows to restrict the summation in the pair density to $`k>0`$. The time–reversal symmetry also renders the pair density real, i.e. $`\chi ^{}(\text{r})=\chi (\text{r})`$.
Two models for the spatial modulation of the pairing strength are considered here
$$G_q(\text{r})=\{\begin{array}{ccc}V_{0,q}\hfill & & \text{DF,}\hfill \\ V_{0,q}\left[1\left(\frac{\rho (\text{r})}{\rho _0}\right)^\gamma \right]\hfill & & \text{DDDI.}\hfill \end{array}$$
(6)
The simpler case (DF) can be deduced from a *delta force* for the pairing interaction Ton79a ; Kri90a ; Dob95a , $`V_{\mathrm{pair}}(\text{r},\text{r}^{})=V_{0,q}\delta (\text{r}\text{r}^{})`$, while the other (DDDI) corresponds to a *density–dependent delta interaction* Dob95a ; Taj93b ; Fay94a . The additional parameters of the DDDI force are set here to $`\gamma =1`$ and $`\rho _0=0.16\mathrm{fm}^3`$ (i.e. the saturation density of symmetric nuclear matter). More general choices are conceivable Dob84a ; Fay96a but very hard to adjust phenomenologically. Thus we keep to the simplest choice above. Note that a separate pairing strength $`V_{0,p}`$ or $`V_{0,n}`$ is associated to each nucleon sort. This explicit breaking of the isospin symmetry in the pairing energy functional is standard in nearly all pairing forces and schematic models, see e.g. Mol92a ; Mad88a .
Although the pairing matrix elements deduced from the pairing functional (4) suppress the contribution from unbound states located outside the nucleus considerably (as compared to the schematic pairing force), the implicit zero–range nature of the pairing force still tends to overestimate the coupling to continuum states. This defect can be cured to some extent using finite–range forces like the Gogny force Dob96c ; Gogrev which are, however, cumbersome to handle. We prefer to simulate the effect of finite range by introducing smooth energy–dependent cutoff weights Kri90a
$$f_k=\frac{1}{1+\mathrm{exp}[(ϵ_k\lambda _q\mathrm{\Delta }E_q)/\mu _q]}$$
(7)
in the evaluation of the local pair density
$$\chi _q(\text{r})2\underset{\genfrac{}{}{0pt}{}{k\mathrm{\Omega }_q}{k>0}}{}f_ku_kv_k|\varphi _k(\text{r})|^2.$$
(8)
The cutoff parameters $`\mathrm{\Delta }E_q`$ and $`\mu _q=\mathrm{\Delta }E_q/10`$ are chosen self–adjusting to the actual level density in the vicinity of the Fermi energy. $`\mathrm{\Delta }E_q`$ is fixed from the condition that the sum of the cutoff weights includes approximately one additional shell of single-particle states above the Fermi surface
$$\underset{k\mathrm{\Omega }_q}{}f_k=N_q+1.65N_q^{2/3}.$$
(9)
### 2.3 The Lipkin–Nogami Equations
The LN scheme serves as an approximation to particle–number projected BCS. It can be derived by a momentum expansion of the projected BCS equations Lip60a ; Nog64a ; Pra73a ; Flo . At the end this boils down to adding overlaps with the variance of the particle number $`(\mathrm{\Delta }\widehat{N}_q)^2`$ at various places.
In most cases, the LN method is used with a simple schematic pairing interaction in the framework of macroscopic–microscopic models and self–consistent models for ground states and potential energy surfaces Mol92a ; Ben89a ; Naz90a ; Que90a ; Taj92a ; Hee93a ; Ska93 ; Mag95a ; Rei96a as well as high–spin states Mag93a ; Sat94a ; Sat94b ; Gal94a ; Wys95a ; Hee95a ; Ter95 . Only recently, the LN scheme was employed for a local delta pairing force Cwi96a and the Gogny force Val96a ; Val97a . Usually only the correction of the pairing energy is calculated; but in self–consistent models the contribution of the mean field to the total binding energy is calculated from the BCS state as well, so that the correction of the pairing energy has to be complemented by a correction of the mean–field energy as considered in Rei96a ; Val96a ; Val97a . In this paper, we present and employ the LN equations in the context of self–consistent mean–field models and for the case of local pairing energy functionals. As pairing gaps are the theme of this paper, particular emphasis is laid on the properties of the LN scheme relevant for the discussion of pairing gaps.
The LN equations are derived by variation of
$$𝒦=\underset{qp,n}{}\left(\lambda _{1,q}\widehat{N}_q+\lambda _{2,q}\widehat{N}_q^2\right).$$
(10)
Variation of (10) with respect to the occupation amplitudes $`v_k`$ leads to
$$v_k^2=\frac{1}{2}\left[1\frac{ϵ_k^{}\lambda _q}{\sqrt{(ϵ_k^{}\lambda _q)^2+f_k^2\mathrm{\Delta }_k^2}}\right].$$
(11)
This is the standard expression for the occupation number $`v_k^2`$ in the BCS model EGIII ; Ringbook ; Nilssonbook , here containing a state–dependent single–particle gap $`\mathrm{\Delta }_k`$ times the cutoff factor $`f_k`$ and a generalized Fermi energy
$$\lambda _q=\lambda _{1,q}+4\lambda _{2,q}(N_q+1),$$
(12)
which is determined from a constraint on particle number. The quantity $`ϵ_k^{}`$ is a renormalized single–particle energy
$$ϵ_k^{}=ϵ_k+4\lambda _{2,q}v_k^2.$$
(13)
In case of time-reversal invariance, the state–dependent single–particle gaps are given by
$$\mathrm{\Delta }_k=\mathrm{d}^3r\varphi _k^{}(\text{r})\mathrm{\Delta }_q(\text{r})\varphi _k(\text{r}),$$
(14)
i.e. they are matrix elements of the local pair potential
$$\mathrm{\Delta }_q(\text{r})=\frac{\delta _{\mathrm{pair}}}{\delta \chi _q(\text{r})}=\frac{1}{2}\chi _q(\text{r})G_q(\text{r}).$$
(15)
Note that $`\lambda _2`$ is not a Lagrange parameter Flo ; Que90a . It is held fixed during the variation and is determined after variation from the additional condition Pra73a
$$\lambda _{2,q}=\frac{(\widehat{H}_{\mathrm{mf}}+\widehat{H}_{\mathrm{pair}})(\mathrm{\Delta }\widehat{N}_{2,q})^2}{\widehat{N}_q(\mathrm{\Delta }\widehat{N}_{2,q})^2}.$$
(16)
For simplicity of the presentation, Eq. (16) is written for the case of an underlying many–body Hamiltonian $`\widehat{H}`$. The discussion of the more general case of an energy functional used in this paper is presented in Appendix A.1. $`\widehat{N}_{2,q}`$ is the part of the particle–number operator that projects onto two–quasiparticle states
$$\widehat{N}_{2,q}=\underset{k\mathrm{\Omega }_q}{}u_kv_k(\widehat{\alpha }_k^{}\widehat{\alpha }_{\overline{k}}^{}+\widehat{\alpha }_{\overline{k}}\widehat{\alpha }_k),$$
(17)
while $`\mathrm{\Delta }\widehat{N}_{2,q}=\widehat{N}_{2,q}\widehat{N}_{2,q}`$. The numerator of Eq. (16) contains, besides the familiar contribution from the pairing functional, an additional one from the linear response of the mean–field to the particle–number projection, see Rei96a . The total binding energy after approximate particle–number projection is given by
$$E^{\mathrm{LN}}=\underset{q\mathrm{p},\mathrm{n}}{}\lambda _{2,q}(\mathrm{\Delta }\widehat{N}_{2,q})^2.$$
(18)
For arbitrary one–body operators the LN expectation value can be calculated introducing effective LN occupation numbers and local densities, see Que90a ; Rei96a .
Thus far the presentation applies to the more involved LN method. The BCS approximation is recovered simply by setting $`\lambda _{2,q}=0`$ in the above equations.
### 2.4 Blocking
The evaluation of the odd–even staggering involves also nuclei with odd mass number where one pair of conjugate states has to be blocked, i.e. taken out of the pairing scheme. One of the blocked states has the occupation $`v_k=1`$, the other $`v_{\overline{k}}=0`$. In the standard textbook approach the blocked many–body state is constructed non–self–consistently from the BCS ground state as a one–quasiparticle excitation, see e.g. EGIII ; Ringbook ; Nilssonbook and Section 3.3. The generalization of this approach for energy functionals is outlined in Appendix B.
In a self–consistent approach to the blocked many–body state the single–particle wave functions and occupation amplitudes have to be determined from a variational principle. Blocking a pair of states changes the density matrix, Eq. (2) and with that the single–particle Hamiltonian, Eq. (3). Besides a rearrangement of the single–particle wave functions the unpaired nucleon causes a polarization of the core by breaking rotational and time–reversal invariance in the intrinsic frame (we assume spherical BCS ground states only), see e.g. ugpap . This requires a deformed calculation of the odd–mass nucleus considering also time–odd contributions to the single–particle Hamiltonian (3) as discussed in ugpap ; Dob95c .
The change in binding energy due to the core polarization depends on the properties of the time–odd spin and spin–isospin channels of the effective interaction which are not yet well adjusted for current mean–field models, see e.g. Eng99a and references therein. Calculations with the currently available models suggest that the polarization is non–negligible for the description of the odd–even staggering ugpap ; Sat98a ; Sat98b ; Xu99a ; uggap , but effective interactions with properly adjusted spin and spin–isospin channels are needed before the effect can be treated quantitatively. In view of these uncertainties we restrict ourselves here to the simpler spherical blocking approximation, where one replaces the blocked single–particle state by an average over the degenerate states in its $`j`$ shell, restoring rotational and time–reversal invariance of the many–body system in the intrinsic frame. In practice this means that the weight of the blocked $`j`$ shell is given by $`(2j1)uv`$ when calculating the pair density and $`(2j1)v^2+1`$ when calculating the local densities and currents. All other states enter with their full degeneracy $`2j+1`$. This approximation includes the large part of the rearrangement effects from monopole polarization, but omits the polarization effects from multipole deformations and time–odd currents.
Owing to the rearrangement effects blocking of the single–particle state with smallest quasiparticle energy (as defined in Sect. 3.3) does not necessarily lead to the largest possible total binding energy. Therefore one has to perform calculations for a number of blocked single–particle states around the Fermi energy and search for the configuration giving the largest total binding energy.
## 3 The Pairing Gap
### 3.1 Nuclear Masses and Odd–Even Staggering
The key feature of pairing correlations is the occurrence of an energy gap in the excitation spectrum. This gap manifests itself in two different kinds of energetic observables: First, there is a gap in the quasiparticle excitation spectra of even–even nuclei, which does not appear in the spectra of odd–mass number or odd–odd nuclei, and second, there occurs a shift between the interpolating curves of the ground–state binding energies of even–even as compared to odd–mass nuclei, which is called the odd–even mass staggering. Usually, the second phenomenon is exploited to define the experimental pairing gaps assuming Mad88a
$`E_{\mathrm{even}\mathrm{even}}(Z,N)`$ $`=`$ $`E_0(Z,N),`$
$`E_{\mathrm{odd}Z}(Z,N)`$ $`=`$ $`E_0(Z,N)+\mathrm{\Delta }_\mathrm{p}(Z,N),`$ (19)
$`E_{\mathrm{odd}N}(Z,N)`$ $`=`$ $`E_0(Z,N)+\mathrm{\Delta }_\mathrm{n}(Z,N).`$
In odd–odd nuclei there is additionally the residual interaction between the unpaired proton and neutron, but this case will not be considered in the present discussion. In a self–consistent mean–field approach $`E_0`$ is the (negative) energy of the fully paired many–particle wave function, i.e. the BCS ground state, while $`\mathrm{\Delta }_q`$ is the energy lost due to the blocking of a pair of conjugate states by the odd nucleon. The gap introduced with (3.1) has to be interpreted carefully. The gap as defined through the separation (3.1) contains more than pure pairing correlations. It includes unavoidably all polarization effects from the mean field as outlined in Sect. 2.4. Despite of these uncertainties (3.1) provides the starting point for the definition of experimentally accessible pairing gaps which will be discussed in Sect. 3.2. The gap defined with (3.1) serves than as a point of reference, as it can be calculated directly within the mean–field model
$$\mathrm{\Delta }^{(\mathrm{b})}(Z,N):=E_{\mathrm{block}}(Z,N)E_0(Z,N)$$
(20)
as the difference in binding energy between the one blocked state of an odd–mass number nucleus with the largest binding energy and its fully paired (fictitious) BCS vacuum. Both calculations have, of course, to be performed self–consistently. We will call $`\mathrm{\Delta }^{(\mathrm{b})}`$ the “blocking gap” in the following.
### 3.2 Finite–Difference Formulae
The gap from Eq. (20) is a purely theoretical construct. The problem is that $`E_0`$ is not measurable for odd–mass nuclei. We need a definition which is also experimentally accessible. It should fulfill two requirements which are useful for the phenomenological adjustment of the pairing strength: first, it should be easy to calculate both theoretically and from experimental data, and second, it should be influenced as little as possible by the properties of the underlying mean field in order to decouple mean–field and pairing properties. The odd–even staggering is related to differences of nuclear masses and as such easy to measure as well as to compute, As outlined above, the odd–even staggering of experimental masses is not a pure measure of pairing correlations, but also has non–negligible contributions from the response of the underlying mean field to the blocking of a single–particle state. Here we can take advantage of the fact that various difference formulae are conceivable and take the recipe which best decouples mean–field and pairing properties.
The odd–even staggering needs to be deduced from energy systematics. To that end, there are several finite–difference formulae in the literature to calculate $`\mathrm{\Delta }_q`$ from binding energies of adjacent nuclei. All available finite–difference formulae for $`\mathrm{\Delta }_q`$ are derived from the Taylor expansion of the nuclear mass in nucleon–number differences Mad88a ; Jen84a
$`E(N)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle \frac{^nE_0}{N^n}}|_{N_0}(NN_0)^n+D(N_0)`$ (21)
where $`E_0`$ is defined in (3.1) and the Gap $`D`$ given by
$$D=\{\begin{array}{cc}0& \text{even proton and neutron number},\hfill \\ \mathrm{\Delta }_\mathrm{n}& \text{odd neutron number},\hfill \\ \mathrm{\Delta }_\mathrm{p}& \text{odd proton number}.\hfill \end{array}$$
(22)
The number of the other kind of nucleons is assumed to be even and the same for all terms. Combining the expansion (21) of several adjacent nuclei leads to energy–difference formulae which can be used to approximate the gap $`\mathrm{\Delta }^{(\mathrm{b})}`$. The two–point (first–order) formula leads to the one–nucleon separation energy, which mixes mean–field, single–particle and pairing effects strongly and should better not be used to fit the pairing strength. The next higher–order is the three–point difference
$`E(N_0+1)2E(N_0)+E(N_01)`$
$`=`$ $`{\displaystyle \frac{^2E_0}{N^2}}|_{N_0}+{\displaystyle \frac{1}{12}}{\displaystyle \frac{^4E_0}{N^4}}|_{N_0}`$
$`+\mathrm{}+D(N_0+1)2D(N_0)+D(N_01),`$
Assuming that the gap $`D`$ varies only slowly with nucleon number and that the remaining contribution from the second derivative of $`E_0`$ is negligible (which is not so well fulfilled in some cases, see Sect. 4.2 and Sat98a ; Sat98b ) this equation can be resolved into an approximative expression for the pairing gap
$`\mathrm{\Delta }_q^{(3)}(N_0)`$ $`:=`$ $`{\displaystyle \frac{\pi _{N_0}^{}}{2}}\left[E(N_01)2E(N_0)+E(N_0+1)\right]`$ (24)
where $`\pi _{N_0}=(1)^{N_0}`$ is the number parity. $`\mathrm{\Delta }^{(3)}`$ calculated from pure HF states without pairing but considering polarization of the mean field is discussed in Refs. Sat98a ; Sat98b in great detail.
The next order corresponds to a four–point difference formula, but this order, employing an even number of nuclei, gives an expression which is asymmetric around the nucleus with $`N_0`$ and therefore offers two choices. With the same assumptions used going from (3.2) to (24), one possibility for the four–point gap is
$`\mathrm{\Delta }_q^{(4)}(N_0)`$ $`:=`$ $`{\displaystyle \frac{\pi _{N_0}}{4}}[E(N_02)3E(N_01)`$ (25)
$`+3E(N_0)E(N_0+1)].`$
This is an approximation for the gap at $`N_01/2`$. The other possible four–point formula gives the gap at $`N_0+1/2`$. The lowest–order derivative of $`E_0`$ hidden in the four–point formula is now of third order.
Although this four–point definition is widely used in the literature Kri90a ; Cwi96a ; BM we prefer the five–point formula
$`\mathrm{\Delta }_q^{(5)}(N_0)`$ $`:=`$ $`{\displaystyle \frac{\pi _{N_0}}{8}}[E(N_0+2)4E(N_0+1)+6E(N_0)`$ (26)
$`4E(N_01)+E(N_02)],`$
which is symmetric and yields the best decoupling from mean–field effects as we will see. The smooth contributions from the mean field to the gap are further suppressed compared to the lower–order formulae, the remaining derivative of $`E_0`$ entering $`\mathrm{\Delta }^{(5)}`$ is of fourth order.
### 3.3 Quasiparticle Energies
Another widely used approximation for the pairing gap is to calculate the energy difference (20) by constructing the blocked many–body wave function of an odd–mass number nucleus non–self–consistently from its BCS ground state as a so–called single–quasiparticle excitation EGIII ; Ringbook ; Nilssonbook . The lowest single–quasiparticle energy
$$E_{\mathrm{quasi}}=\text{min}(E_k)$$
(27)
– which we will simply denote as “quasiparticle energy” $`E_{\mathrm{quasi}}`$ in the following – is then another approximation for the odd–even staggering. In the LN scheme the $`E_k`$ are given in first–order approximation by
$$E_k\sqrt{(ϵ_k^{}\lambda _q)^2+f_k^2\mathrm{\Delta }_k^2}+\lambda _{2,q},$$
(28)
see Appendix B for details. The important difference between $`E_{\mathrm{quasi}}`$ and $`\mathrm{\Delta }^{(\mathrm{b})}`$ is that for $`E_{\mathrm{quasi}}`$ the blocked many–body wave function is not calculated self–consistently.
### 3.4 Spectral Gaps
The calculation of $`\mathrm{\Delta }^{(5)}`$ from mean–field models is a bit cumbersome since it requires information on five nuclei including nuclei with odd mass number. Moreover, the definition becomes inapplicable to describe the variation of pairing correlations with deformation for a given nucleus. At this point, we could recur to the purely theoretical definition (20) involving a blocked and an unblocked BCS calculation. This still requires involved calculations and can become unwieldy in deformed calculations.
Therefore, a commonly used approach is to estimate the pairing gap from spectral properties of a nucleus. In schematic pairing models using the same pairing matrix element for all states the single–particle gaps (14) turn out to be state–independent and are thus immediately a measure for the pairing correlations. With local forces as used here we obtain state–dependent single–particle gaps $`\mathrm{\Delta }_k`$ and have to define an average gap as representative for the strength of the pairing correlations. The authors of Dob84a have proposed to use the average of the single–particle gaps (14) weighted with the occupation probability $`v_k^2`$
$$v^2\mathrm{\Delta }_q=\frac{_{k\mathrm{\Omega }_q}f_kv_k^2\mathrm{\Delta }_k}{_{k\mathrm{\Omega }_q}f_kv_k^2}.$$
(29)
This definition, however, puts too much weight on deeply–bound states whereas pairing is a mechanism most active near the Fermi surface. We therefore propose an average with the same factor $`v_ku_k`$ as it appears in the accumulation of the pair density $`\chi (\text{r})`$, see Eq. (5), yielding the spectral gap as
$$uv\mathrm{\Delta }_q=\frac{_{k\mathrm{\Omega }_q}f_kv_ku_k\mathrm{\Delta }_k}{_{k\mathrm{\Omega }_q}f_ku_kv_k}.$$
(30)
Note that the spectral gaps are an estimate for the pairing gap and therefore the contribution of pairing correlations to the odd–even staggering (20). Assuming that the spectral gaps (29) and (30) are an approximation for the square–root term in (28), approximately particle–number projected spectral gaps are given by
$`v^2\mathrm{\Delta }_q^{(\mathrm{LN})}=v^2\mathrm{\Delta }_q+\lambda _{2,q},`$ (31)
$`uv\mathrm{\Delta }_q^{(\mathrm{LN})}=uv\mathrm{\Delta }_q+\lambda _{2,q},`$ (32)
which was proposed by the authors of Cwi96a for the average gap $`v^2\mathrm{\Delta }^{(\mathrm{LN})}`$. These spectral gaps will be discussed and compared with other alternatives for the calculated pairing gap in Section 4.
## 4 Results and Discussion
### 4.1 Fit of Pairing Strength
The first step is to determine appropriate pairing strengths $`V_{0,q}`$ for the DF and the DDDI functionals. We do that on the grounds of the five–point gap $`\mathrm{\Delta }^{(5)}`$ and adjust the pairing strength by fitting calculated values for $`\mathrm{\Delta }^{(5)}`$ to experimental ones for a large set of semi–magic nuclei, i.e. the isotope chains $`{}_{22}{}^{44}\text{Ca}`$, $`{}_{\mathrm{\hspace{0.33em}\hspace{0.25em}56}}{}^{106}\text{Sn}`$$`{}_{\mathrm{\hspace{0.33em}\hspace{0.25em}78}}{}^{128}\text{Sn}`$, and $`{}_{119}{}^{201}\text{Pb}`$$`{}_{124}{}^{206}\text{Pb}`$, for the neutrons and the isotone chains $`{}_{28}{}^{52}\text{Cr}`$, $`{}_{50}{}^{82}\text{Ge}_{32}^{}`$$`{}_{50}{}^{94}\text{Ru}_{44}^{}`$, $`{}_{\mathrm{\hspace{0.33em}\hspace{0.25em}82}}{}^{136}\text{Xe}_{54}^{}`$$`{}_{\mathrm{\hspace{0.33em}\hspace{0.25em}82}}{}^{147}\text{Tb}_{65}^{}`$, and $`{}_{126}{}^{210}\text{Po}_{84}^{}`$$`{}_{126}{}^{215}\text{Ac}_{89}^{}`$ for the protons. A seperate fit has been performed for each one of the pairing functionals, DF or DDDI, and for each appraoch, BCS or LN. Each nucleon sort, proton or neutron, aquires its own pairing strength adjusted to isotonic chains, or isotopic chains respectively.
The experimental data are taken from Wapstra . The resulting values for the pairing strength are listed in Table 1. We obtain a reasonable fit of the pairing gaps for all pairing schemes, see also Figures 4 (compare “Expt. with $`\mathrm{\Delta }^{(5)}`$) and 8 in what follows. We have checked that the pairing strengths are not significantly changed (i.e. on the order of $`1\%`$) when fitting instead theoretical $`\mathrm{\Delta }^{(3)}`$ to experimental $`\mathrm{\Delta }^{(3)}`$ or similarly the $`\mathrm{\Delta }^{(4)}`$.
### 4.2 Comparison of the Finite–Difference Formulae
The various finite–difference gaps were introduced as experimentally accessible approximations to $`\mathrm{\Delta }^{(\mathrm{b})}`$, i.e. the energy differences of blocked and unblocked calculations of odd–mass nuclei. Figure 1 compares directly the performance of the $`\mathrm{\Delta }^{(i)}`$, $`i=3,4,5`$, in this respect for calculations within the BCS approach. The three–point gaps $`\mathrm{\Delta }^{(3)}`$ show a pronounced odd–even staggering for all $`N=82`$ isotones and the tin isotopes with neutron numbers below the $`N=82`$ shell closure. Note that we have disentangled that by drawing a separate line for even–even and odd–mass nuclei. These two lines for $`\mathrm{\Delta }^{(3)}`$ embrace the gap $`\mathrm{\Delta }^{(\mathrm{b})}`$. The staggering is caused by non–vanishing mean–field contributions to $`\mathrm{\Delta }^{(3)}`$ (mainly the second derivative term in Eq. (3.2)), which enter with a different sign for even–even and odd–mass nuclei. The four–point gaps $`\mathrm{\Delta }^{(4)}`$ from Eq. (25) give a smoother approximation for $`\mathrm{\Delta }^{(\mathrm{b})}`$. But as expected the $`\mathrm{\Delta }^{(4)}`$ are slightly shifted versus the $`\mathrm{\Delta }^{(\mathrm{b})}`$ which becomes rather obvious where the gaps change rapidly, e.g. around $`N=63`$, $`N=70`$, $`Z=57`$, $`Z=65`$. The five–point gaps $`\mathrm{\Delta }^{(5)}`$ give the best overall agreement with the $`\mathrm{\Delta }^{(\mathrm{b})}`$.
The oscillation of the $`\mathrm{\Delta }^{(3)}`$ around the values for $`\mathrm{\Delta }^{(\mathrm{b})}`$ has a simple geometrical reason, as can be seen from Fig. 2. $`\mathrm{\Delta }^{(\mathrm{b})}`$ is per definition (20) the shift between the smooth curve connecting the (unblocked) BCS ground–state energies $`E_0`$ of all nuclei (thick line through circles) and the smooth curve that connects the (blocked) ground–state energies $`E_\mathrm{b}`$ of odd–mass nuclei (dotted line through full boxes). The three–point gap of an even–even nucleus $`\mathrm{\Delta }_{\mathrm{even}}^{(3)}(N_0)=\frac{1}{2}[E_\mathrm{b}(N_01)+E_\mathrm{b}(N_0+1)]E_0(N_0)`$ is the shift between the smooth curve connecting the $`E_0`$ and the average energy of the two adjacent odd–mass nuclei, which lies outside the band given by $`\mathrm{\Delta }^{(\mathrm{b})}`$. From this follows immediately $`\mathrm{\Delta }_{\mathrm{even}}^{(3)}>\mathrm{\Delta }^{(\mathrm{b})}`$ for bound nuclei. A similar construction leads to $`\mathrm{\Delta }_{\mathrm{odd}}^{(3)}<\mathrm{\Delta }^{(\mathrm{b})}`$, see Fig. 2. For $`\mathrm{\Delta }^{(4)}`$ and $`\mathrm{\Delta }^{(5)}`$ the deviation from $`\mathrm{\Delta }^{(\mathrm{b})}`$ becomes of course much smaller because the higher–order finite–difference formulae give a better approximation of the smooth curves connecting the $`E_0`$ and $`E_\mathrm{b}`$ respectively. This qualitative result does not depend on the level of sophistication for the calculation of the odd–mass–number nuclei. Considering polarization effects will shift $`E_\mathrm{b}`$ with respect to $`E_0`$ (which has to be counterweighted by a refit of the pairing strength Xu99a ; uggap ) and may distort the surface of the $`E_\mathrm{b}`$, but will not change the sign of its curvature.
There remains a significant difference between all $`\mathrm{\Delta }^{(i)}`$, $`i=3,4,5`$ and $`\mathrm{\Delta }^{(\mathrm{b})}`$ around shell closures where finite–difference formulae (except $`\mathrm{\Delta }^{(3)}`$ for odd nuclei) produce a peak which becomes broader with increasing order of the difference formula. We want to discuss the origin of this peak for the example of the five–point gap $`\mathrm{\Delta }^{(5)}`$. The five–point approximation (26) for $`\mathrm{\Delta }^{(\mathrm{b})}`$ makes three assumptions: (i) the binding energy is an analytical function of the nucleon numbers and therefore can be expanded in a Taylor series in the range of two mass units around the considered nucleus; (ii) derivatives of $`E_0`$ of higher than third order are negligible; and (iii) the gap varies only slowly with nucleon number. The first two assumptions are strongly violated at shell closures, where the kink in the systematics of binding energies does not allow the Taylor expansion (21) and leads to a spurious contribution from the mean–field functional to the finite–difference gaps.
To visualize this effect we compare in Fig. 3 $`\mathrm{\Delta }^{(5)}`$ with five–point gaps $`\mathrm{\Delta }^{(5)}^{(\mathrm{nb})}`$ using the binding energies of not blocked calculations of the odd–mass–number nuclei, which carries only the the spurious mean–field contribution to $`\mathrm{\Delta }^{(5)}`$. Subtracting $`\mathrm{\Delta }^{(5)}^{(\mathrm{nb})}`$ from $`\mathrm{\Delta }^{(5)}`$ one gets the contribution from the blocking of the odd particle to $`\mathrm{\Delta }^{(5)}`$. The peaks at closed shells disappear yielding a smooth curve for $`\mathrm{\Delta }^{(5)}\mathrm{\Delta }^{(5)}^{(\mathrm{nb})}`$ which now follows the values of $`\mathrm{\Delta }^{(\mathrm{b})}`$ everywhere.
Besides the immediate vicinity of shell closures, we find that $`\mathrm{\Delta }^{(5)}`$ is a very good approximation for the staggering gap $`\mathrm{\Delta }^{(\mathrm{b})}`$. For a few additional nuclei there remains a small difference between $`\mathrm{\Delta }^{(5)}`$ and $`\mathrm{\Delta }^{(\mathrm{b})}`$, see Fig. 3. This occurs for example around $`N=68`$ and $`N=72`$ for the tin isotopes. For those nuclei the assumption of a slow variation of the gap with nucleon number which enters the five–point formula (26) is not valid: the gap changes in these regions by about $`40\%`$ (due to a sudden change of the density of single–particle levels at the Fermi surface in these nuclei). But even then the deviation between $`\mathrm{\Delta }^{(5)}`$ and $`\mathrm{\Delta }^{(\mathrm{b})}`$ remains acceptably small. We thus prefer $`\mathrm{\Delta }^{(5)}`$ for the fit of pairing strengths. The small difference between $`\mathrm{\Delta }^{(5)}`$ and $`\mathrm{\Delta }^{(\mathrm{b})}`$ even suggests a simplified fitting procedure where calculated $`\mathrm{\Delta }^{(\mathrm{b})}`$ are adjusted to experimental values for $`\mathrm{\Delta }^{(5)}`$ in chains of spherical semi–magic nuclei.
As explained in Sect. 3.1, the $`\mathrm{\Delta }^{(\mathrm{b})}`$ unavoidably contain a contribution from the polarization of the mean field in odd nuclei. It is therefore somewhat unlucky and confusing that $`\mathrm{\Delta }^{(\mathrm{b})}`$ is usually denoted as “pairing gap” in the literature, see e.g. Mol92a ; Mad88a ; BM . The commonly used fist formulae like $`\mathrm{\Delta }^{(\mathrm{b})}12/\sqrt{A}`$ are intended to represent the average trend of the odd–even staggering, not the pairing gap. The discrepancies between the two quantities can be expected to decrease with increasing mass number and to be very small for heavy nuclei.
An analysis of gaps from a complementary point of view is given in Sat98a for the case of $`\mathrm{\Delta }^{(3)}`$. This study omits pairing altogether and concentrates exclusively on polarization effects for $`\mathrm{\Delta }^{(3)}`$ using pure deformed HF calculations. The focus is on small nuclei because these have most pronounced deformation effects. In this framework, it is found that 3–point gaps show a pronounced odd–even staggering with the $`\mathrm{\Delta }_{\mathrm{odd}}^{(3)}`$ being close to zero while most of the $`\mathrm{\Delta }_{\mathrm{even}}^{(3)}`$ have large positive values in most cases. This finding is explained in terms of the macroscopic–microscopic model in that a large contribution from the symmetry energy is counterweighted by the difference of single–particle energies when calculating $`\mathrm{\Delta }_{\mathrm{odd}}^{(3)}`$. This observation led the authors of Sat98a to the conclusions that $`\mathrm{\Delta }_{\mathrm{odd}}^{(3)}`$ is very close to the pure pairing gap while $`\mathrm{\Delta }_{\mathrm{even}}^{(3)}`$ contains a contribution from the mean field. The higher–order gaps $`\mathrm{\Delta }^{(4)}`$ and $`\mathrm{\Delta }^{(5)}`$ turn out to be rather useless in that environment. We want to point out that the odd–even staggering of $`\mathrm{\Delta }^{(3)}`$ observed in Sat98a is not the staggering of $`\mathrm{\Delta }^{(3)}`$ around the values for $`\mathrm{\Delta }^{(\mathrm{b})}`$ as we discuss it here, see e.g. Fig. 1. The deformation staggering is a phenomenon much similar to the spurious peak of finite–difference gaps at major shell closures discussed above. Mind that pure deformed HF calculations in small nuclei produce a subshell closure for each even–even nucleus by virtue of the Jahn–Teller effect. This, in turn, leads to a spurious contribution to $`\mathrm{\Delta }_{\mathrm{even}}^{(3)}`$ for nearly all even–even nuclei which is half of the energy difference between single–particle levels Sat98a , and it has devastating consequences for the systematics of 4–point and 5–point gaps. The picture changes dramatically when pairing is included, as we do here. Pairing induces a drive to spherical shapes and thus reduces deformation effects dramatically while rendering blocking the the dominant contribution to the gaps. This smoothes the systematics of binding energies and thus of any $`\mathrm{\Delta }^{(i)}`$. Moreover, we are discussing here semi–magic medium and heavy nuclei which have spherical BCS ground states. We are thus considering a sample with minimal mean–field effects, just appropriate to concentrate on pairing features. The interplay of polarization effects which we neglect here and pairing correlations will probably play a role for small non–magic nuclei. It deserves further inspection in the future.
### 4.3 Spectral Gaps in BCS Pairing
Having discussed the various finite difference gaps, we now take the five–point gap as reference value and study the relation to the spectral gaps and quasiparticle energies. Figure 4 compares calculated results for $`\mathrm{\Delta }^{(5)}`$, $`v^2\mathrm{\Delta }`$, $`uv\mathrm{\Delta }`$ and $`E_{\mathrm{quasi}}`$ computed within the BCS approach.
Let us start the discussion by looking at the neutron gaps in the chain of tin isotopes calculated with the DF pairing interaction in the BCS approach (upper left panel). The spectral gaps $`uv\mathrm{\Delta }`$ and $`v^2\mathrm{\Delta }`$ show a pronounced odd–even staggering where the gaps of even–even nuclei have larger values than those of odd–mass nuclei. This is caused by the blocking of one state with a large weight $`uv`$ in the odd–mass nucleus. The blocked state does not contribute to the pairing potential, leading to overall smaller single–particle gaps (14). The amplitude of the odd–even staggering decreases, of course, with increasing neutron number $`N`$ because the relative contribution of a particular state to the pair density becomes smaller with increasing level density in the heavier isotopes.
Both $`v^2\mathrm{\Delta }`$ and $`uv\mathrm{\Delta }`$ are exactly zero for closed shell nuclei, i.e. $`N=50`$, $`N=82`$, $`N=112`$, and the adjacent odd mass–number nuclei. In these nuclei the BCS scheme breaks down. This is a deficiency of the BCS scheme which is related to the particle–number uncertainty of the BCS state Ringbook .
The quasiparticle energies $`E_{\mathrm{quasi}}`$ show a similar dependence on the neutron number as the $`\mathrm{\Delta }^{(5)}`$, but for most nuclei they are 100–250 keV larger (and thus the same amount larger with respect to the $`\mathrm{\Delta }^{(\mathrm{b})}`$). This is caused by calculating the blocked many–body wave function entering $`E_{\mathrm{quasi}}`$ not self–consistently. The variational principle behind the self–consistent calculation of the odd–mass nuclei entering the $`\mathrm{\Delta }^{(5)}`$ leads always to larger a binding energy of the odd–mass nuclei, lowering the calculated gap.
The quasiparticle energies show the same peak at shell closures as the finite–difference gaps which is again related to a spurious contribution from the mean field. For non–magic nuclei the lowest single–quasiparticle state usually corresponds to a single–particle state with $`ϵ_k^{}\lambda _q`$ leading to the single–quasiparticle energy $`E_{\mathrm{quasi}}\mathrm{\Delta }_k+\lambda _{2,q}`$. In magic nuclei one has an additional contribution from the first term in the square root in Eq. (28) since the Fermi energy is approximately in the middle of the gap in the single–particle spectrum (In the BCS scheme, where the pairing breaks down for closed–shell nuclei the derivation of Eq. (28) is not valid. Then one has different Fermi energies for the removal and addition of a particle, which are the single–particle energies of the last occupied and first unoccupied single–particle state respectively). The large jump in the Fermi energy is the reason for the kink in the systematics of binding energies at shell closures, which in turn causes the peak in the finite–difference gaps discussed above.
While all definitions of the gap give similar values in the valley of stability, there appear large differences between the $`\mathrm{\Delta }^{(5)}`$ and $`E_{\mathrm{quasi}}`$ on one hand and the $`v^2\mathrm{\Delta }`$ and $`uv\mathrm{\Delta }`$ on the other hand for neutron–rich nuclei beyond the $`N=82`$ shell closure. The spectral gaps overestimate the $`\mathrm{\Delta }^{(5)}`$ (which closely follow the $`\mathrm{\Delta }^{(\mathrm{b})}`$ as discussed earlier). The $`uv\mathrm{\Delta }`$ are smaller than the $`v^2\mathrm{\Delta }`$ for all tin isotopes and in most systems the $`uv\mathrm{\Delta }`$ are closer to the $`\mathrm{\Delta }^{(5)}`$ than the $`v^2\mathrm{\Delta }`$.
Now we want to look at the changes when employing the DDDI pairing functional instead of the simpler delta force, see the lower left panel of Fig. 4. There are two significant differences to the results obtained with the DF functional: (i) all gaps except $`v^2\mathrm{\Delta }`$ are larger by roughly $`20\%`$ for neutron–rich nuclei with $`N>82`$, and (ii) the $`v^2\mathrm{\Delta }`$ are always much smaller than the $`uv\mathrm{\Delta }`$. This trend is just the opposite from the one for the DF functional. This is caused by the different choice of weights in the definition of $`v^2\mathrm{\Delta }`$ and $`uv\mathrm{\Delta }`$ in combination with the spectral distribution of the single–particle gaps $`\mathrm{\Delta }_k`$. While the DF pairing potential follows roughly the nuclear density distribution, the DDDI functional gives a pairing potential which is sharply peaked at the nuclear surface, see Fig. 5. This leads in case of the DF interaction to single–particle gaps $`\mathrm{\Delta }_k`$ of comparable size for all bound states while in case of the DDDI interaction the $`\mathrm{\Delta }_k`$ of deeply bound single–particle states are rather small, see the middle panel of Fig. 6. Together with the weight factors used to calculate $`uv\mathrm{\Delta }`$ and $`v^2\mathrm{\Delta }`$ – see the upper panel in Fig. 6 – this gives the observed pattern for the spectral gaps: when calculated with the DF interaction they are rather insensitive to the choice for the weight factors while there is a huge difference in case of the DDDI interaction.
This explains also why the $`\mathrm{\Delta }^{(5)}`$ extrapolate quite differently when comparing DF and DDDI pairing for neutron–rich nuclei. In these nuclei the states at the Fermi surface are only loosely bound and therefore have a large spatial extension but only small overlap with the volume–like DF pairing potential. An extreme example is the tin isotope with $`N=112`$ where DF pairing breaks down but DDDI pairing is still fully active. Experimental data on the excitation spectra of neutron–rich nuclei will give in the future valuable information to distinguish between volume–like (DF) and surface–peaked pairing interactions.
This different behavior of $`v^2\mathrm{\Delta }`$ on one hand and $`uv\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{(5)}`$ on the other hand hints that is possibly dangerous to use different definitions of the gap for experimental and calculated gaps when fitting the pairing strength and comparing calculated and experimental values.
The right panels of Fig. 4 show the pairing gaps of the protons in the chain of $`N=82`$ isotones. Qualitatively the results are similar to those for the neutron gaps in the tin isotopes. But the overall reproduction of the experimental values is much better than in the case of tin isotopes and the differences between the various gaps is much smaller when going towards the drip–line. The Coulomb potential stabilizes even loosely–bound protons. Therefore the difference between the gaps comparing volume–like and surface–like pairing potentials is quite small. Only for the DDDI force remains the large difference between the spectral gaps $`v^2\mathrm{\Delta }`$ and $`uv\mathrm{\Delta }`$ explained above.
### 4.4 Lipkin–Nogami Pairing
Figure 7 shows the gaps of the neutrons in the chain of tin isotopes, now calculated with the LN scheme. The spectral gaps $`v^2\mathrm{\Delta }`$ are omitted here. Instead we compare the “bare” $`uv\mathrm{\Delta }`$ (30) with the particle–number corrected values $`uv\mathrm{\Delta }+\lambda _2`$ (32).
The global pattern of the various gaps looks very similar to the one obtained with the BCS scheme, see Fig. 4. The most obvious difference between the BCS and LN methods is that the LN scheme does not break down for closed–shell nuclei. Therefore the “bare” spectral gap $`uv\mathrm{\Delta }`$ has a finite – but still somewhat too small – value around magic nuclei. Adding $`\lambda _2`$ gives better results around shell closures, but away from shell closures the difference between $`uv\mathrm{\Delta }`$ and $`uv\mathrm{\Delta }+\lambda _2`$ for the tin isotopes is too small to decide on one preferred definition.
From the difference between $`uv\mathrm{\Delta }`$ and $`uv\mathrm{\Delta }+\lambda _2`$ it can be seen that $`\lambda _2`$ is largest around shell closures. But this indicates also that the LN approximation might not be sufficient for magic nuclei, a variational calculation of $`\lambda _2`$ or even full projection of the many–body wave function is needed Dob92a there.
The single–quasiparticle energies $`E_{\mathrm{quasi}}`$ follow closely the particle–number projected $`uv\mathrm{\Delta }_\mathrm{p}+\lambda _{2,\mathrm{p}}`$, which is easily understood remembering that $`\lambda _2`$ is added to both quantities. At shell closures, however, the single–quasiparticle energies overestimate the experimental gaps. Like in the case of the BCS scheme the $`E_{\mathrm{quasi}}`$ are nearly always larger than the calculated five–point gaps $`\mathrm{\Delta }^{(5)}`$.
### 4.5 Comparison of all Models
Finally, we want to compare the four pairing models, i.e. any combination of BCS or LN and DF or DDDI pairing. The comparison is done with respect to the five–point gap $`\mathrm{\Delta }^{(5)}`$, which we prefer as the most robust empirical definition and which has turned out to be the most useful definition of the pairing gap. In Fig. 8 the five–point gaps calculated from various pairing schemes are compared with experimental values. The differences between BCS and LN are generally very small. It is to be remembered, however, that the effective pairing strength is readjusted for the LN scheme. The (approximate) particle–number projection increases the total binding energy, but this effect is renormalized by virtue of the fit delivering a slightly smaller pairing strength in the LN scheme. There is one detail where BCS and LN differ: the peak of $`\mathrm{\Delta }^{(5)}`$ in the vicinity of closed shells is more spread out in case of LN which is probably due to the softening of the shell closure by LN.
The differences between the pairing forces (DF versus DDDI) are much larger. For the neutron gaps of tin isotopes close to the valley of stability and the proton gaps in the $`N=82`$ isotones all schemes and forces still give similar results, but large differences between the DF and DDDI force occur for neutron gaps around $`N=60`$ and very neutron–rich nuclei with $`N>82`$. Only the DDDI model can describe the gaps in both the light and heavy known tin isotopes, while the DF interaction overestimates the gaps in the light ones by up to 15–20$`\%`$. The particle–number projection has only a small effect on the gaps when the strength of the pairing interaction is readjusted. The better description of the tin isotopes around $`N=60`$ with the DDDI force gives a hint that this pairing interaction may be more realistic than a delta force. This, however, has to be taken with a grain of salt: it may also be a spurious effect due to a deficiency of the underlying mean–field. The disagreement between calculated and experimental $`\mathrm{\Delta }^{(5)}`$ around $`N=70`$ is rather robust in that all forces and schemes give very similar results. For a profound decision which pairing interaction gives the most realistic results throughout the chart of nuclei more and other observables have to be investigated using various forces for the underlying mean–field. Research in that direction is underway.
We have checked that one obtains very similar results when comparing experimental and calculated three–point and four–point gaps. However, it is important that the experimental and calculated values are computed from the same formula. As already shown in Fig. 1, the various finite–difference formulae may give results which differ by $`25\%`$.
## 5 Summary
We have compared various approximations for the calculation of the pairing gap: three–point, four–point and five–point finite difference formulae where the gap is estimated from total binding energies of adjacent nuclei, spectral gaps with different weights which put bias on well–bound states or levels at the Fermi surface, and the single–quasiparticle energy. Predictions of four different pairing models for the gaps were compared, namely the BCS and LN pairing schemes employing the DF or DDDI interactions. Experimental values for the pairing gaps are usually calculated from a finite–difference formula. For the calculated gaps we find that apart from shell closures there are non–negligible deviations up to $`25\%`$ between the various definitions. Some definitions of the pairing gap cannot be used for closed–shell nuclei. Therefore, it is the safest choice to compute the pairing gaps from mean–field calculations in the same way as the experimental values.
The natural definition for the calculated gap is the difference in binding energy $`\mathrm{\Delta }^{(\mathrm{b})}`$ between the fully paired BCS ground state and the blocked one of odd–mass nuclei. The five–point gap $`\mathrm{\Delta }^{(5)}`$ provides a reliable quantity to fit the effective pairing strength, among all approximations for the pairing gap it is closest to $`\mathrm{\Delta }^{(\mathrm{b})}`$. Therefore we use it to calculate the experimental pairing gaps and take it as point of reference for all other approximations for the pairing gap. Like the other finite–difference formulae $`\mathrm{\Delta }^{(5)}`$ contains a spurious contribution from the mean field around magic numbers, which is related to the jump of the Fermi energy at shell closures.
The four–point gap $`\mathrm{\Delta }^{(4)}`$ has nearly the same overall quality as the five–point gap as compared to $`\mathrm{\Delta }^{(\mathrm{b})}`$, but its definition has an ambiguity, so that we prefer the five–point gap. The three–point gap $`\mathrm{\Delta }^{(3)}`$ has a large contribution from the mean–field, which becomes rather obvious looking at proton gaps, consequently this quantity should not be used for the fit of the pairing strength since deficiencies of the underlying mean–field (especially concerning the symmetry energy) may be visible in the $`\mathrm{\Delta }^{(3)}`$. When comparing calculated and experimental $`\mathrm{\Delta }^{(3)}`$ respectively $`\mathrm{\Delta }^{(4)}`$ for a well–adjusted interaction, however, both gaps show the same quality like the five–point gaps $`\mathrm{\Delta }^{(5)}`$.
In situations where it is not possible to calculate $`\mathrm{\Delta }^{(5)}`$, e.g. looking at potential energy surfaces, a reasonable approximation for the pairing gap is provided by $`uv\mathrm{\Delta }+\lambda _2`$. By adding $`\lambda _2`$ to the spectral gap the approximate particle–number projection gives an improved reproduction of the calculated $`\mathrm{\Delta }^{(5)}`$ in most cases. The spectral gap $`uv\mathrm{\Delta }+\lambda _2`$ is in much better agreement with the calculated $`\mathrm{\Delta }^{(5)}`$ than the averaged gap $`v^2\mathrm{\Delta }+\lambda _2`$, which sets too much bias on deeply bound states and therefore should not be used in models in which the pairing interaction acts mainly at the nuclear surface. The single–quasiparticle energies always overestimate the $`\mathrm{\Delta }^{(5)}`$, but this is not too surprising because they are calculated from a non–self–consistent wave function of the odd nucleus.
The results show the danger of comparing calculated and experimental gaps computed from different definitions. The deviations between the various definitions depend on the actual nucleus and become generally larger when going towards the drip–lines. This is important especially when adjusting the free parameters of the pairing interaction. Experimental and calculated gaps should be calculated in the same manner, the five–point gap provides a useful tool for that.
The choice of the test cases (heavy semi–magic nuclei) and restriction to spherical nuclei have minimized the impact of polarization effects on the gaps. Recent explorations of dynamical polarization effects Xu99a ; uggap hint that these may be non–negligible at a quantitative level. Conclusive answers are yet inhibited due to uncertainties of present mean–field models in the time–odd channel. This point deserves attention in future work.
Another open question is the functional form of the pairing interaction. We find in this paper that the DDDI functional gives a slightly better description of pairing gaps compared with a delta pairing force, but to give a definitive answer one has to look at more nuclei and other data as well. Work in this direction is underway.
## Acknowledgments
The authors would like to thank W. Nazarewicz for stimulating discussions and challenging comments. This work was supported by Bundesministerium für Bildung und Forschung (BMBF), Project No. 06 ER 808, by Gesellschaft für Schwerionenforschung (GSI), by Graduiertenkolleg Schwerionenphysik and by the U.S. Department of Energy under Contract No. DE–FG02–97ER41019 with the University of North Carolina and Contract No. DE–FG02–96ER40963 with the University of Tennessee and by the NATO grant SA.5–2–05 (CRG.971541). The Joint Institute for Heavy Ion Research has as member institutions the University of Tennessee, Vanderbilt University, and the Oak Ridge National Laboratory; it is supported by the members and by the Department of Energy through Contract No. DE–FG05–87ER40361 with the University of Tennessee.
## Appendix A The calculation of $`\lambda _2`$
### A.1 The General Expression
In case of an underlying many–body Hamiltonian $`\widehat{H}`$ the parameter $`\lambda _2`$ in the variational equation is fixed by the condition
$$\widehat{K}\widehat{N}_2^2=0,$$
(33)
where $`\widehat{K}`$ is given by
$$\widehat{K}=\widehat{H}\underset{qp,n}{}\left(\lambda _{1,q}\widehat{N}_q+\lambda _{2,q}\widehat{N}_q^2\right)$$
(34)
and $`\widehat{N}_2`$ is the two–quasiparticle part of the particle–number operator (17). Because we look at like-particle pairing only the equations for protons and neutrons separate. Therefore the index $`q`$ for the isospin of the particle–number operator, the single–particle states etc can suppressed to get a compact notation. Introducing the shifted many–body state Rei96a
$$|\xi =\mathrm{e}^{\mathrm{i}\xi \widehat{N}_2}|0\text{with}|0=|\xi |_{\xi =0},$$
(35)
the condition (33) is equivalent to
$$_\xi ^20|\widehat{K}|\xi |_{\xi =0}=0.$$
(36)
This can be resolved into an expression for $`\lambda _2`$
$$\lambda _2=\frac{_\xi ^20|\widehat{H}|\xi |_{\xi =0}}{_\xi ^20|\widehat{N}^2|\xi |_{\xi =0}}.$$
(37)
The denominator of (37) can be calculated easily using Wick’s theorem
$`_\xi ^20|\widehat{N}^2|\xi |_{\xi =0}`$ $`=`$ $`8\left[{\displaystyle \underset{k0}{}}u_k^2v_k^2\right]^216{\displaystyle \underset{k0}{}}u_k^4v_k^4.`$ (38)
So far we have formulated the numerator in terms of an underlying many–particle Hamiltonian, but the formulation (36) has the advantage that it can be translated into the formal framework of energy functionals Rei96a
$$0|\widehat{H}|\xi ^{(\xi )}=[\widehat{\rho }^{(\xi )},\widehat{\chi }^{(\xi )},\widehat{\chi }^{(\xi )}].$$
(39)
$`\widehat{\rho }^{(\xi )}`$ and $`\widehat{\chi }^{(\xi )}`$ are the shifted density matrix and pair density matrix respectively, which have to be calculated now as non-diagonal matrix elements
$`\widehat{\rho }^{(\xi )}`$ $``$ $`\rho ^{(\xi )}(\stackrel{}{x},\stackrel{}{x}^{})=0|\widehat{\psi }^{}(\stackrel{}{x}^{})\widehat{\psi }(\stackrel{}{x})|\xi `$ (40a)
$`\widehat{\chi }^{(\xi )}`$ $``$ $`\chi ^{(\xi )}(\stackrel{}{x},\stackrel{}{x}^{})=\sigma ^{}0|\widehat{\psi }(\stackrel{}{r}^{},\sigma ^{})\widehat{\psi }(\stackrel{}{x})|\xi `$ (40b)
$`\widehat{\chi }^{(\xi )}`$ $``$ $`\chi ^{(\xi )}(\stackrel{}{x},\stackrel{}{x}^{})=\sigma 0|\widehat{\psi }^{}(\stackrel{}{x}^{})\widehat{\psi }^{}(\stackrel{}{r},\sigma )|\xi `$ (40c)
All local densities and currents the Skyrme and pairing energy functionals depend on can be derived from these density matrices. The second derivative of the energy functional is given by
$`_\xi ^2^{(\xi )}`$ $`=`$ $`\mathrm{tr}\left\{{\displaystyle \frac{\delta }{\delta \widehat{\rho }}}_\xi ^2\widehat{\rho }+{\displaystyle \frac{\delta }{\delta \widehat{\chi }}}_\xi ^2\widehat{\chi }+{\displaystyle \frac{\delta }{\delta \widehat{\chi }^{}}}_\xi ^2\widehat{\chi }^{}\right\}`$ (41)
$`+\mathrm{tr}\mathrm{tr}\{{\displaystyle \frac{\delta ^2}{\delta \widehat{\rho }_1\delta \widehat{\rho }_2}}_\xi \widehat{\rho }_1_\xi \widehat{\rho }_2+2{\displaystyle \frac{\delta ^2}{\delta \widehat{\chi }\delta \widehat{\chi }^{}}}_\xi \widehat{\chi }_\xi \widehat{\chi }^{}`$
$`+2{\displaystyle \frac{\delta ^2}{\delta \widehat{\rho }\delta \widehat{\chi }}}_\xi \widehat{\rho }_\xi \widehat{\chi }+2{\displaystyle \frac{\delta ^2}{\delta \widehat{\rho }\delta \widehat{\chi }^{}}}_\xi \widehat{\rho }_\xi \widehat{\chi }^{}\}.`$
$`\delta `$ stands for a functional derivative. The trace is a shorthand notation for the integration and summation over all coordinates
$$\mathrm{tr}\left\{\frac{\delta }{\delta \widehat{\rho }}\widehat{\rho }\right\}dxdx^{}\frac{\delta }{\delta \rho (\stackrel{}{x},\stackrel{}{x}^{})}\rho (\stackrel{}{x},\stackrel{}{x}^{}).$$
(42)
The terms with a single trace in (41) vanish in case of a BCS ground state, while the terms with double traces can be simplified defining the response density matrices $`\stackrel{~}{\rho }`$, $`\stackrel{~}{\chi }`$ and $`\stackrel{~}{\chi }^{}`$
$$\stackrel{~}{\rho }=\mathrm{i}_\xi \widehat{\rho }^{(\xi )}|_{\xi =0},\stackrel{~}{\chi }=\mathrm{i}_\xi \widehat{\chi }^{(\xi )}|_{\xi =0}.$$
(43)
leading to the final expression
$`_\xi ^2^{(\xi )}|_{\xi =0}`$ $`=`$ $`\mathrm{tr}\mathrm{tr}\{{\displaystyle \frac{\delta ^2}{\delta \widehat{\rho }_1\delta \widehat{\rho }_2}}\stackrel{~}{\rho }_1\stackrel{~}{\rho }_2+2{\displaystyle \frac{\delta ^2}{\delta \widehat{\chi }\delta \widehat{\chi }^{}}}\stackrel{~}{\chi }\stackrel{~}{\chi }^{}`$ (44)
$`+2{\displaystyle \frac{\delta ^2}{\delta \widehat{\rho }\delta \widehat{\chi }}}\stackrel{~}{\rho }\stackrel{~}{\chi }+2{\displaystyle \frac{\delta ^2}{\delta \widehat{\rho }\delta \widehat{\chi }^{}}}\stackrel{~}{\rho }\stackrel{~}{\chi }^{}\}`$
which has to be inserted into (37). Since we are considering pairing between like particles only there are no mixed terms with derivatives with respect to proton and neutron densities. This expression has to be evaluated now for the pairing and the mean–field energy functional. This density–functional approach to the calculation of $`\lambda _2`$ incorporates in a natural way the additional contributions to the LN equations arising for density–dependent interactions discussed in Val96a ; Val97a .
### A.2 The Linear Response of a Skyrme Energy Functional
The Skyrme energy functionals are constructed to be effective interactions for nuclear mean–field calculations. For even–even nuclei, the Skyrme energy functional used in this paper
$$=_{\mathrm{kin}}+_{\mathrm{Sk}}+_\mathrm{C},$$
(45)
is the sum of the functional of the kinetic energy $`_{\mathrm{kin}}`$, the effective functional for the strong interaction $`_{\mathrm{Sk}}`$ and the Coulomb interaction $`_\mathrm{C}`$ including the exchange term in Slater approximation. The actual functionals are given by
$`_{\mathrm{kin}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \mathrm{d}^3r\tau },`$
$`_\mathrm{C}`$ $`=`$ $`{\displaystyle \frac{e^2}{2}}{\displaystyle \mathrm{d}^3r\mathrm{d}^3r^{}\frac{\rho _\mathrm{p}(\text{r})\rho _\mathrm{p}(\text{r}^{})}{|\text{r}\text{r}^{}|}}{\displaystyle \frac{3e^2}{4}}\left({\displaystyle \frac{3}{\pi }}\right)^{1/3}{\displaystyle \mathrm{d}^3r\rho _\mathrm{p}^{4/3}}`$
$`_{\mathrm{Sk}}`$ $`=`$ $`{\displaystyle }\mathrm{d}^3r[{\displaystyle \frac{b_0}{2}}\rho ^2+b_1\rho \tau {\displaystyle \frac{b_2}{2}}\rho \mathrm{\Delta }\rho +{\displaystyle \frac{b_3}{3}}\rho ^{\alpha +2}b_4\rho 𝐉`$ (46)
$`{\displaystyle \underset{q}{}}({\displaystyle \frac{b_0^{}}{2}}\rho _q^2+b_1^{}\rho _q\tau _q{\displaystyle \frac{b_2^{}}{2}}\rho _q\mathrm{\Delta }\rho _q`$
$`+{\displaystyle \frac{b_3^{}}{3}}\rho ^\alpha \rho _q^2+b_4^{}\rho _q𝐉_q)].`$
The local density $`\rho _q`$, kinetic density $`\tau _q`$ and spin–orbit current $`𝐉_q`$ entering the functional are given by
$`\rho _q(\stackrel{}{r})`$ $`=`$ $`{\displaystyle \underset{\sigma =\pm }{}}\rho _q(\stackrel{}{r},\sigma ;\stackrel{}{r},\sigma )`$
$`=`$ $`{\displaystyle \underset{k\mathrm{\Omega }_q}{}}v_k^2|\varphi _k(\stackrel{}{r})|^2,`$
$`\tau _q(\stackrel{}{r})`$ $`=`$ $`{\displaystyle \underset{\sigma =\pm }{}}{}_{}{}^{}\rho _{q}^{}(\stackrel{}{r},\sigma ;\stackrel{}{r}^{},\sigma )|_{\stackrel{}{r}=\stackrel{}{r}^{}}`$
$`=`$ $`{\displaystyle \underset{k\mathrm{\Omega }_q}{}}v_k^2|\varphi _k(\stackrel{}{r})|^2,`$
$`𝐉_q(\stackrel{}{r})`$ $`=`$ $`\frac{\mathrm{i}}{2}(^{})\times {\displaystyle \underset{\sigma ,\sigma ^{}=\pm }{}}\rho _q(\stackrel{}{r},\sigma ;\stackrel{}{r}^{},\sigma ^{})\stackrel{}{\sigma }_{\sigma ^{}\sigma }|_{\stackrel{}{r}=\stackrel{}{r}^{}}`$ (47)
$`=`$ $`\frac{\mathrm{i}}{2}{\displaystyle \underset{k\mathrm{\Omega }_q}{}}v_k^2\left[\varphi _k^{}(\stackrel{}{r})\times \widehat{\sigma }\varphi _k(\stackrel{}{r})\text{ h.c. }\right],`$
with $`q\{\mathrm{p},\mathrm{n}\}`$. $`\stackrel{}{\sigma }_{\sigma ^{}\sigma }`$ is the matrix element of the vector of the Pauli spin matrices between the unit spinors with spin projection $`\sigma ^{}/2`$ and $`\sigma /2`$. Densities without index in (A.2) denote total densities, e.g. $`\rho =\rho _\mathrm{p}+\rho _\mathrm{n}`$. The $`\varphi _k`$ are the spinors of the single–particle wave functions, the $`v^2`$ occupation probabilities. The parameters $`b_i`$ and $`b_i^{}`$ used in the above definition are chosen to give a most compact formulation of the energy functional, the corresponding mean–field Hamiltonian and residual interaction Rei92 . The Skyrme energy functional contains an extended spin–orbit coupling with an explicit isovector degree–of–freedom as used in the parameterization SkI4 Rei95a .
The response densities needed for the evaluation of (44) are given by
$`\stackrel{~}{\rho }_q`$ $`=`$ $`2{\displaystyle \underset{k\mathrm{\Omega }_q}{}}u_k^2v_k^2|\varphi _k|^2,`$
$`\stackrel{~}{\tau }_q`$ $`=`$ $`2{\displaystyle \underset{k\mathrm{\Omega }_q}{}}u_k^2v_k^2|\varphi _k|^2,`$
$`\stackrel{~}{𝐉}_q`$ $`=`$ $`\mathrm{i}{\displaystyle \underset{k\mathrm{\Omega }_q}{}}u_k^2v_k^2\left[\varphi _k^{}\times \widehat{\sigma }\varphi _k(\times \widehat{\sigma }\varphi _k)^{}\varphi _k\right].`$
Evaluating Eq. (44) for the Skyrme energy functional (A.2) leads to
$`_\xi ^2_{\mathrm{Sk}}^{(\xi )}|_{\xi =0}`$
$`=`$ $`\mathrm{tr}\mathrm{tr}\{\stackrel{~}{\rho }_{q,1}{\displaystyle \frac{\delta ^2_{\mathrm{Sk}}}{\delta \rho _{q,1}\delta \rho _{q,2}}}\stackrel{~}{\rho }_{q,2}+2\stackrel{~}{\rho }_q{\displaystyle \frac{\delta ^2_{\mathrm{Sk}}}{\delta \rho _q\delta \tau _q}}\stackrel{~}{\tau }_q`$
$`+2\stackrel{~}{\rho }_q{\displaystyle \frac{\delta ^2_{\mathrm{Sk}}}{\delta \rho _q\delta 𝐉_q}}\stackrel{~}{𝐉}_q\}`$
$`=`$ $`{\displaystyle }\mathrm{d}^3r\{(b_0b_0^{})\stackrel{~}{\rho }_q^2+2(b_1b_1^{})\stackrel{~}{\rho }_q\stackrel{~}{\tau }_q`$
$`\left(b_2b_2^{}\right)\stackrel{~}{\rho }_q\mathrm{\Delta }\stackrel{~}{\rho }_q2\left(b_4+b_4^{}\right)\stackrel{~}{\rho }_q\stackrel{~}{𝐉}_q`$
$`+\frac{1}{3}\left[(\alpha +2)(\alpha +1)b_32b_3^{}\right]\rho ^\alpha \stackrel{~}{\rho }_q^2`$
$`b_3^{}\frac{1}{3}[4\alpha \rho ^{\alpha 1}\rho _q+\alpha (\alpha 1)\rho ^{\alpha 2}{\displaystyle \underset{q^{}}{}}\rho _q^{}^2]\stackrel{~}{\rho }_q^2\}.`$
For the local densities appearing in (A.2) the trace reduces to a spatial integral. For the protons one has an additional contribution from the Coulomb interaction
$`_\xi ^2_\mathrm{C}^{(\xi )}|_{\xi =0}`$ (49)
$`=`$ $`\mathrm{tr}\mathrm{tr}\left\{\stackrel{~}{\rho }_\text{p}{\displaystyle \frac{\delta ^2_\text{C}}{\delta \rho _\text{p}\delta \rho _\text{p}}}\stackrel{~}{\rho }_\mathrm{p}\right\}`$
$`=`$ $`e^2{\displaystyle \mathrm{d}^3r\mathrm{d}^3r^{}\frac{\stackrel{~}{\rho }_\text{p}(𝐫)\stackrel{~}{\rho }_\text{p}(𝐫^{})}{|𝐫𝐫^{}|}}{\displaystyle \frac{e^2}{3}}\left({\displaystyle \frac{3}{\pi }}\right)^{1/3}{\displaystyle \mathrm{d}^3r\frac{\stackrel{~}{\rho }_\text{p}^2}{\rho _\text{p}^{2/3}}}`$
$`=`$ $`{\displaystyle \mathrm{d}^3r\stackrel{~}{\rho }_\text{p}\stackrel{~}{V}_{\text{coul}}}{\displaystyle \frac{e^2}{3}}\left({\displaystyle \frac{3}{\pi }}\right)^{1/3}{\displaystyle \mathrm{d}^3r\frac{\stackrel{~}{\rho }_\text{p}^2}{\rho _\text{p}^{2/3}}}.`$
Poisson’s equation for the response Coulomb potential
$$\mathrm{\Delta }\stackrel{~}{V}_{\text{coul}}=4\pi \stackrel{~}{\rho }_\text{p}.$$
(50)
is solved numerically using the techniques explained in Coul . The kinetic energy gives no contribution to $`\lambda _2`$. We omit the contribution from the center–of–mass correction (and therefore the approximate particle–number correction of this term).
### A.3 The Linear Response of the Pairing Energy Functional
For the calculation of the contribution of a pairing energy functional of type (4) to Eq. (44) the local response pair density $`\stackrel{~}{\chi }_q`$ is needed
$$\stackrel{~}{\chi }_q=\mathrm{i}_\xi \chi _q^{(\xi )}|_{\xi =0}=4\underset{\genfrac{}{}{0pt}{}{k\mathrm{\Omega }_q}{k>0}}{}f_k^2u_k^3v_k|\varphi _k|^2.$$
(51)
The response pair density is not hermitian, the adjoint response pair density $`\stackrel{~}{\chi }_q^{}`$ reads
$$\stackrel{~}{\chi }_q^{}=\mathrm{i}_\xi \chi _q^{(\xi )}|_{\xi =0}=4\underset{\genfrac{}{}{0pt}{}{k\mathrm{\Omega }_q}{k>0}}{}f_k^2u_kv_k^3|\varphi _k|^2.$$
(52)
In case of the DF pairing energy functional there is only one contribution from the derivatives with respect to the pair density (44)
$`_\xi ^2_{\mathrm{DF}}^{(\xi )}|_{\xi =0}`$ $`=`$ $`2\mathrm{tr}\mathrm{tr}\left\{\stackrel{~}{\chi }_q{\displaystyle \frac{\delta ^2_{\mathrm{DF}}}{\delta \chi _q\delta \chi _q^{}}}\stackrel{~}{\chi }_q^{}\right\}`$ (53)
$`=`$ $`{\displaystyle \frac{V_q}{2}}{\displaystyle \mathrm{d}^3r\stackrel{~}{\chi }_q^{}\stackrel{~}{\chi }_q}.`$
In case of the DDDI pairing energy functional there are additional contributions from derivatives with respect to the local density
$`_\xi ^2_{\mathrm{DDDI}}^{(\xi )}|_{\xi =0}`$ (54)
$`=`$ $`\mathrm{tr}\mathrm{tr}\{2\stackrel{~}{\chi }_q{\displaystyle \frac{\delta ^2_{\text{DDDI}}}{\delta \chi _q\delta \chi _q^{}}}\stackrel{~}{\chi }_q^{}+2\stackrel{~}{\rho }_q{\displaystyle \frac{\delta ^2_{\text{DDDI}}}{\delta \rho _q\delta \chi _q}}\stackrel{~}{\chi }_q`$
$`+2\stackrel{~}{\rho }_q{\displaystyle \frac{\delta ^2_{\text{DDDI}}}{\delta \rho _q\delta \chi _q^{}}}\stackrel{~}{\chi }_q^{}+\stackrel{~}{\rho }_{q,1}{\displaystyle \frac{\delta ^2_{\text{DDDI}}}{\delta \rho _{q,1}\delta \rho _{q,2}}}\stackrel{~}{\rho }_{q,2}\}`$
$`=`$ $`{\displaystyle \frac{V_q}{2}}{\displaystyle }\mathrm{d}^3r\{\stackrel{~}{\chi }_q^{}\stackrel{~}{\chi }_q[1\left(\frac{\rho }{\rho _0}\right)^\gamma ]`$
$`{\displaystyle \frac{\gamma }{\rho _0^\gamma }}\stackrel{~}{\rho }_q\rho ^{\gamma 1}\left(\stackrel{~}{\chi }_q^{}\chi _q+\chi _q^{}\stackrel{~}{\chi }_q\right)`$
$`{\displaystyle \frac{\gamma (\gamma 1)}{2\rho _0^\gamma }}\stackrel{~}{\rho }_q^2\rho ^{\gamma 2}\chi _q^{}\chi _q\}.`$
## Appendix B Single–Quasiparticle Energies
The single–quasiparticle energies are deduced from the non–self–consistent *ansatz*
$$|k=\widehat{\alpha }_k^{}|\mathrm{BCS}=v_k\widehat{a}_k^{}\underset{\genfrac{}{}{0pt}{}{m>0}{mk}}{}(u_m+v_m\widehat{a}_m^{}\widehat{a}_{\overline{m}}^{})|0$$
(55)
for the blocked many–body wave function of the odd mass number nucleus EGIII ; Ringbook ; Nilssonbook . The single–particle wave functions of all states and the occupation probabilities of all unblocked states $`mk`$ are taken from the self–consistent calculation of the BCS ground state. The normalization of $`|k`$ requires $`v_k=1`$, $`u_k=v_k=v_{\overline{k}}=0`$. Note that $`|k`$ does not yield the proper particle number of the excited state because the occupation of the blocked pair of states $`(k,\overline{k})`$ in the BCS ground state in general will differ from one.
The excitation energy $`E_k`$ of the lowest one–quasiparticle state is an approximation for the odd–even mass difference. To take the difference in particle number between the fully paired BCS state and $`|k`$ into account, $`E_k`$ has to be calculated from the difference of $`𝒦`$ given by (10) calculated for $`|k`$ and the BCS ground state. This leads to
$$E_k=^{(k)}+k|\widehat{N}^{}|k^{(0)}\widehat{N}^{}$$
(56)
where we have introduced the abbreviation
$$\widehat{N}^{}=\underset{q\{\mathrm{p},\mathrm{n}\}}{}\left(\lambda _{1,q}\widehat{N}_q+\lambda _{2,q}\widehat{N}_q^2\right).$$
(57)
$`^{(0)}`$ is the energy functional of the BCS ground state, while $`^{(k)}`$ is the energy functional evaluated for the one–quasiparticle state $`|k`$.
The expectation values of $`\widehat{N}^{}`$ are conveniently calculated with standard operator techniques EGIII ; Ringbook ; Nilssonbook
$$k|\widehat{N}^{}|k=\widehat{\alpha }_k^{}\widehat{N}^{}\widehat{\alpha }_k^{}=\widehat{N}_{00}^{}+\widehat{\alpha }_k^{}\widehat{N}_{11}^{}\widehat{\alpha }_k^{},$$
(58)
where $`\widehat{N}_{00}^{}`$ is the BCS expectation value of $`\widehat{N}^{}`$ while $`\widehat{N}_{11}^{}`$ is its (11) component. The Bogolyubov transformation of the particle–number operator reads
$`\widehat{N}`$ $`=`$ $`{\displaystyle \underset{k0}{}}\widehat{a}_k^{}\widehat{a}_k^{}`$ (59)
$`=`$ $`{\displaystyle \underset{k0}{}}\left[v_k^2+(u_k^2v_k^2)\widehat{\alpha }_k^{}\widehat{\alpha }_k^{}+u_kv_k(\widehat{\alpha }_k^{}\widehat{\alpha }_{\overline{k}}^{}+\widehat{\alpha }_{\overline{k}}^{}\widehat{\alpha }_k^{})\right]`$
$`=`$ $`\widehat{N}_{00}+\widehat{N}_{11}+\widehat{N}_{20}+\widehat{N}_{02}.`$
$`(\widehat{N}^2)_{11}`$ can be calculated from (59) without repetition of the quasiparticle transformation. One obtains
$$(\widehat{N}^2)_{11}=\underset{k0}{}\left\{(u_k^2v_k^2)[2(N+1)4v_k^2]1\right\}\widehat{\alpha }_k^{}\widehat{\alpha }_k^{}.$$
(60)
The calculation of the contribution of a many–body Hamiltonian operator to $`E_k`$ can be found in many textbooks, see e.g. EGIII ; Ringbook ; Nilssonbook . Here, however, we want to calculate it in the framework of effective energy functionals. The energy functional $`^{(k)}=[\widehat{\rho }^{(k)},\chi ^{(k)},\chi ^{(k)}]`$ of the one–quasiparticle state is to be calculated from the density matrix and pair density matrix evaluated for the one–quasiparticle state $`|k`$
$`\rho ^{(k)}(\stackrel{}{x};\stackrel{}{x}^{})=k|\widehat{\psi }^{}(\stackrel{}{x}^{})\widehat{\psi }(\stackrel{}{x})|k`$ (61a)
$`=`$ $`\varphi _k^{}(\stackrel{}{x}^{})\varphi _k(\stackrel{}{x})+{\displaystyle \underset{\genfrac{}{}{0pt}{}{m>0}{mk}}{}}v_m^2\varphi _m^{}(\stackrel{}{x}^{})\varphi _m(\stackrel{}{x}),`$
$`\chi ^{(k)}(\stackrel{}{r})={\displaystyle \underset{\sigma =\pm }{}}\sigma k|\widehat{\psi }(\stackrel{}{r},\sigma )\widehat{\psi }(\stackrel{}{x})|k`$ (61b)
$`=`$ $`2{\displaystyle \underset{\genfrac{}{}{0pt}{}{m>0}{mk}}{}}f_mu_mv_m|\varphi (\stackrel{}{r})|^2.`$
$`\rho ^{(k)}(\stackrel{}{x};\stackrel{}{x}^{})`$ and $`\chi ^{(k)}(\stackrel{}{r})`$ are inserted into the definition of the energy functionals (45) and (4). This leads to
$$E_k=^{(k)}^{(0)}(\lambda _q4\lambda _{2,q}v_k^2)(u_k^2v_k^2)+\lambda _{2,q}$$
(62)
where $`\lambda _q=\lambda _{1,q}+4\lambda _{2,q}(N_q+1)`$ was used.
The usual approximation for the one–quasiparticle energy is obtained taking only the linear change in the densities into account
$$^{(k)}^{(0)}+\mathrm{tr}\left\{\frac{\delta }{\delta \widehat{\rho }}\delta \widehat{\rho }+\frac{1}{2}\frac{\delta }{\delta \chi }\delta \chi +\frac{1}{2}\frac{\delta }{\delta \chi ^{}}\delta \chi ^{}\right\}$$
(63)
with
$`\delta \widehat{\rho }`$ $`=`$ $`\widehat{\rho }^{(k)}\widehat{\rho }=(12v_k^2)\varphi _k^{}(\stackrel{}{x}^{})\varphi _k(\stackrel{}{x}),`$ (64a)
$`\delta \chi `$ $`=`$ $`\delta \chi ^{}=\chi ^{(k)}\chi =2u_kv_k|\varphi _k|^2.`$ (64b)
This yields
$$^{(k)}^{(0)}ϵ_k(u_k^2v_k^2)+2f_k\mathrm{\Delta }_ku_kv_k.$$
(65)
Together with the contributions from the particle–number operators, Eqns. (59) and (60), one obtains
$`E_k`$ $`=`$ $`(ϵ_k^{}\lambda _q)(u_k^2v_k^2)+2f_k\mathrm{\Delta }_ku_kv_k+\lambda _2`$ (66)
$`=`$ $`\sqrt{(ϵ_k^{}\lambda _q)^2+f_k^2\mathrm{\Delta }_k^2}+\lambda _{2,q},`$
where the definitions of the Fermi energy (12), the renormalized single–particle energy (13) and the expressions for $`v_k^2`$ and $`u_k^2=1v_k^2`$ have been inserted.
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